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# Untitled Document
Algebra and Logic, Vol. 38, No. 4, 1999, 259–276.
A QUANTUM ANALOG OF THE
POINCARE–BIRKHOFF–WITT THEOREM
V. K. Kharchenko <sup>1</sup><sup>1</sup>1Supported by the National Society of Researchers, México (SNI, exp. 18740, 1997-2000).
Translated from Algebra i Logika, Vol. 38, No. 4, pp. 476-507, July-August, 1999. Original article submitted June 29, 1998.
> We reduce the basis construction problem for Hopf algebras generated by skew-primitive semi-invariants to a study of special elements, called “super-letters,” which are defined by Shirshov standard words. In this way we show that above Hopf algebras always have sets of PBW-generators (“hard” super-letters). It is shown also that these Hopf algebras having not more than finitely many “hard” super-letters share some of the properties of universal enveloping algebras of finite-dimensional Lie algebras. The background for the proofs is the construction of a filtration such that the associated graded algebra is obtained by iterating the skew polynomials construction, possibly followed with factorization.
INTRODUCTION
In this article we deal with the basis construction problem for character Hopf algebras, i.e., for the Hopf algebras generated by skew primitive semi-invariants and by an Abelian group of all group-like elements. These algebras constitute an important class actively studied within the frames of the quantum group theory. The class includes all known to the date quantizations with the coalgebra structure of a Lie algebra, and probably we may think of it as an abstractly defined class of all “quantum” universal enveloping algebras. In line with this approach, “quantum” Lie algebras are couched in terms of spaces of all skew primitive elements of the character Hopf algebras endowed with the natural structure of an Yetter–Drinfeld module and equipped with partial quantum operations; see .
In the present article the basis construction problem will be reduced to treating special elements defined by Shirshov standard words, which we call “super-letters.” The main result, Theorem 2, states that the set of all monotonic restricted words in “hard” super-letters constitute a basis for a Hopf algebra. If the Hopf algebra is generated by ordinary primitive elements, the set of all “hard” super-letters constitute a basis for the Lie algebra of all primitive elements. By this token, Theorem 2 may be conceived of as one of the possible quantum analogs for the Poincare–Birkhoff–Witt theorem.
The proof and the statement of the main theorem are based on Shirshov’s combinatorial method, originally developed for Lie algebras. Using this method, Shirshov solved a number of important problems in combinatorial theory of Lie algebras. Among them are the characterization problem for a free Lie algebra over an arbitrary operator ring , the basis construction problem for a free Lie algebra over a field (which is an outgrowth of M. Hall’s ideas in ), and the equality problem for Lie algebras with one defining relation or with finitely many homogeneous relations . Independently, most fundamental concepts of that method were pronounced in where the basis construction problem was dealt with for dual groups of the lower central series of a finitely generated free group.
A weak point in our modification of Shirshov’s method is that essential use will be made of the so-called “through” ordering of words, standard words, and super-words, for which the set of all standard words (super-letters) may be not completely ordered. For this reason, the main theorem is proved only for finitely generated Hopf algebras. In this connection, it is worth mentioning that Shirshov’s original method does not presume the use of a “through” ordering only. What it calls for is a weak restriction on the order: the end of a standard word should be less than the word itself. An example is M. Hall’s ordering from , which is in fact also used in the present article. However, we opt to not bring in both of the orders, to avoid (or at least minimize) misunderstanding. The reason why we do not use the M. Hall’s ordering as the main to our reasoning is because Lemma 8 becomes almost uninformative in this case.
In order to generalize the main theorem to the case of infinitely generated algebras, instead of searching suitable orders, one may apply the famous local method by Mal’tsev (cf. ), whereby the proof of the theorem reduces to a logical analysis of its formulation.
The main theorem can also be used to construct bases for known quantizations of Lie algebras. For the Drinfeld–Jimbo quantizations, such were constructed in Rosso , Yamane , Lusztig , and Kashiwara . Curiously, no one of the methods by Rosso, Yamane, Lusztig, or Kashiwara presupposes that active use be made of the coalgebraic structure — instead — they all presume a detailed treatment of the algebraic. At the same time, the coproduct in Drinfeld–Jimbo quantizations, and also in every pointed Hopf algebra (cf. ), consists only of a “skew primitive” leading part and a linear tensor combination of a lesser degree. — This opens up unbounded prospects for inductive proofs.
An approach attempted here is aimed at a study of effects brought about by the existence of a coproduct. The Poincare–Birkhoff–Witt theorem (PBW-theorem) can also be proved in terms of a coproduct, provided that a given Lie algebra is presupposed to be embedded in a (cocommutative connected) Hopf algebra. This was in fact done in Milnor and Moore \[15, Secs. 5 and 6\]. The PBW-theorem in the Milnor-Moore form carries no information about primitive elements (the given Lie algebra) but gives a complete solution to the basis construction problem for a Hopf algebra modulo its solution for a given Lie algebra. The mentioned above decreasing process (unlike detailed algebraic accounts) has sharply delineated boundaries of application — it cannot give any information about the structure of a set of skew primitive elements (that is of the structure of the quantum Lie algebra itself). Therefore, it might be interesting to investigate these sets in known quantizations as quantum Lie algebras, that is together with all partial quantum operations over them .
In Sec. 1, we introduce basic notions and give a formulation of Shirshov’s theorem \[3, Lemma 1\] needed for our further constructions. All statements under this section were proved by Shirshov sometimes in a more general form. In a slightly different guise, some of them were discovered independently of Shirshov in .
In Sec. 2, we replace the classical commutator with a skew commutator whose “curvature” depends on the parameters of specified elements in approximately the same way as it does in color Lie super-algebras. In our case, however, the bicharacter is not assumed symmetric. Still, identity (8), which is analogous to the Jacobi identity, is valid. And so are derivative identities (9) and (11), which link the skew commutator and the basic product. The bulk of the information needed is given in Lemmas 6 and 8, in which two decreasing processes are described. One is an analog of the Hall–Shirshov construction for nonassociative words and the other is concerned with a coproduct in the way mentioned above.
In Sec. 3, we pass from quantum variables to arbitrary skew primitive generators, and using the two above-mentioned decreasing processes, prove the main result, Theorem 2. On this theorem, each character Hopf algebra has the same basis as the universal enveloping algebra of a (restricted) Lie algebra. The role of a basis for the Lie algebra is played by special elements, which we call hard super-letters. The main lemmas are stated in such a way as to fit in dealing with skew primitive elements.
In Sec. 4, we derive some immediate consequences of the main theorem. In particular, it is shown that character Hopf algebras having not more than finitely many hard super-letters share some of the properties of universal enveloping algebras of finite-dimensional Lie algebras. The background for our proofs is the construction of a filtration such that the associated graded algebra is obtained by iterating the skew polynomials construction, possibly followed with factorization. Note also that the main theorem, as well as its corollaries, remain true for $`(G,\lambda )`$-graded Hopf algebras and for braided bigraded Hopf algebras. In this event a group $`G`$ merely defines a grading, but the algebra in question does not itself contain the $`G`$. Therefore, additional restrictions on a group are unnecessary.
Finally, the quantum Serre relations can be expressed in terms of some super-letters being equal to zero. If, in these super-letters, we replace the skew commutator operation with the classical one then the original Serre relations will appear; see \[1, Thm. 6.1\]. Therefore, it seems absolutely realistic that all hard super-letters of the Drinfeld–Jimbo quantized enveloping algebras arise from a suitable basis for the Lie algebra by merely replacing the commutator with the skew commutator. This is likely to be true not only for the case of Drinfeld–Jimbo quantizations.
1. SHIRSHOV STANDARD WORDS
Let $`x_1,\mathrm{},x_n`$ be a set of variables. Consider this set as an alphabet. On a set of all words in this alphabet, define the lexicographical order such that $`x_1>x_2>\mathrm{}>x_n`$. This means that two words $`v`$ and $`w`$ are compared by moving from left to right until the first distinct letter is encountered. If not, i.e., one of the words is the beginning of the other, then a shorter word is assumed to be greater than the longer (as is common practice in dictionaries). For example, all words of length at most two in two variables respect the following order:
$$x_1>x_1^2>x_1x_2>x_2>x_2x_1>x_2^2.$$
(1)
This order is stable under left multiplication and unstable under right. Nevertheless, if $`u>v`$ and $`u`$ is not the beginning of $`v`$, then the inequality is preserved under right multiplication, even by different words: $`uw>vt`$.
Every noncommutative polynomial $`f`$ in $`x_1,\mathrm{},x_n`$ is a linear combination of words $`f=\alpha _iu_i`$. By $`\overline{f}`$ we denote a leading word which occurs in this decomposition with a nonzero coefficient. In the general case the leading word of a product does not equal the product of leading words of the factors. For example, if $`f=x_1+x_1x_2`$ and $`g=x_3+x_3x_2`$ then $`\overline{fg}=x_1x_2x_3x_1x_3=\overline{f}\overline{g}`$. But if the leading word of $`f`$ is not the beginning of any other word in $`f`$, then
$$\overline{fg}=\overline{f}\overline{g}.$$
(2)
Indeed, the inequalities $`\overline{f}>u_j`$ can be multiplied from the right by (possibly distinct) elements $`\overline{f}v_k>u_jv_s`$. In particular, if $`f`$ is an homogeneous polynomial, i.e., all words $`u_i`$ have the same length, then formula (2) is true.
The set of all words is not completely ordered since there exist infinite decreasing chains — for instance,
$$x_1>x_1^2>x_1^3>\mathrm{}>x_1^n>\mathrm{}.$$
(3)
Yet, all of its finite subsets are completely ordered. This will allow us to use induction on the leading word, provided that bounds are set on the lengths of words, $`l(v)`$, or on degrees of the polynomials under consideration.
Definition 1. A word $`u`$ is called standard (in the sense of Shirshov) if, for each representation $`u=u_1u_2`$, where $`u_1`$ and $`u_2`$ are nonempty words, the inequality $`u>u_2u_1`$ holds. For example, in (3), there is only one standard word, namely, $`x_1`$, and in (1), there are three: $`x_1`$, $`x_1x_2`$, and $`x_2`$.
LEMMA 1. If $`u=sv`$ is a standard word and $`s`$ is nonempty then $`v`$ is not the beginning of $`u`$.
Proof. Suppose, to the contrary, that $`u=vs^{}`$. By the definition of a standard word, we then have $`sv=vs^{}>s^{}v`$, i.e., $`s>s^{}`$. Similarly, $`vs^{}=sv>vs`$, whence $`s^{}>s`$, a contradiction.
LEMMA 2. A word $`u`$ is standard if and only if it is greater than any one of its endings.
Proof. If the word $`u`$ is standard and $`u=vv_1`$ then $`vv_1>v_1v`$ by definition. By the previous lemma, $`v_1`$ is not the beginning of $`vv_1`$, and hence $`u=vv_1`$ and $`v_1v`$ differ already in their first $`l(v_1)`$ letters. Therefore, $`u>v_1`$. Conversely, if $`u=u_1u_2`$ and $`u>u_2`$ then $`u`$ is not the beginning of $`u_2`$, and so the above inequality remains true under right multiplication of the right hand side by $`u_1`$.
LEMMA 3. If $`u`$ and $`v`$, $`u>v`$, are standard words then $`u^h>v`$.
Proof. If $`u`$ is not the beginning of $`v`$ then $`u>v`$ can be multiplied from the right by different words. Suppose that $`v=u^kv^{}`$ and that $`v^{}`$ does not begin with $`u`$. If $`kh`$ then $`u^h>v`$ as a beginning. If $`k<h`$ then $`v^{}`$ is nonempty and $`v^{}<v<u`$. It follows that $`v=u^kv^{}<u^kuu^{hk1}=u^h`$.
LEMMA 4. Let $`u`$ and $`u_1`$ be standard words such that $`u=u_3u_2`$ and $`u_2>u_1`$. Then
$$uu_1>u_3u_1,uu_1>u_2u_1.$$
(4)
Proof. First we show that $`u_2u_1>u_1`$. If $`u_1`$ does not begin with $`u_2`$, the inequality follows immediately from $`u_2>u_1`$. Assume that $`u_1=u_2^ku_1^{}`$ and $`u_2`$ is not the beginning of $`u_1^{}`$. Since $`u_1`$ is standard, we have $`u_2^ku_1^{}>u_2^{k1}u_1^{}`$, i.e., $`u_2u_1^{}>u_1^{}`$. Hence $`u_2u_1=u_2^ku_2u_1^{}>u_2^ku_1^{}=u_1`$. Multiplying this inequality from the left by $`u_3`$ yields the first inequality required. Consider the second. Since $`u`$ is a standard word, $`u_3u_2>u_2`$ by Lemma 2, and $`u_3u_2`$ is of course not the beginning of $`u_2`$. We can therefore multiply the latter inequality from the right by $`u_1`$.
Recall that a nonassociative word is one where $`[,]`$ are somehow arranged to show how multiplication applies. The set of nonassociative words can be defined inductively by the following axioms:
(1) all letters are nonassociative words;
(2) if $`[v]`$ and $`[w]`$ are nonassociative words then $`\left[[v][w]\right]`$ is a nonassociative word;
(3) there are no other nonassociative words.
Definition 2. A nonassociative word $`[u]`$ is said to be standard (in the sense of Shirshov) if:
(1) an (associative) word $`u`$ obtained from this word by removing the brackets is standard;
(2) if $`[u]=\left[[v][w]\right]`$ then $`[v]`$ and $`[w]`$ are standard nonassociative words;
(3) if $`[u]=\left[\left[[v_1][v_2]\right][w]\right]`$ then $`v_2w`$.
The Shirshov theorem (cf. \[3, Lemma 1\]). Each standard word can be uniquely bracketed so that the resulting nonassociative word is standard.
This theorem, combined with the inductive definition of a set of all nonassociative words, immediately implies that every standard (associative) word $`u`$ has a decomposition $`u=vw`$, where $`v>w`$ and $`v`$ and $`w`$ are standard. Yet, for the associative decomposition (as distinct from nonassociative one), the words $`v`$ and $`w`$ are not defined uniquely. The factors $`v`$ and $`w`$ in the nonassociative decomposition $`[u]=[[v][w]]`$, we note, can be defined to be standard words such that $`u=vw`$, where $`v`$ has a least possible length; see .
2. DECOMPOSITION OF QUANTUM POLYNOMIALS INTO LINEAR COMBINATIONS OF MONOTONIC SUPER-WORDS
Let $`x_1,\mathrm{},x_n`$ be quantum variables, i.e., associated with each letter $`x_i`$ are an element $`g_i`$ of a fixed Abelian group $`G`$ and a character $`\chi ^i:G𝐤^{}`$. For every word $`u`$, denote by $`g_u`$ an element of the group $`G`$ which results from $`u`$ by replacing each occurrence of the letter $`x_i`$ with $`g_i`$. This group-like element is denoted also by $`G(u)`$, provided that $`u`$ is an unwieldy expression. Likewise, by $`\chi ^u`$ we denote a character which results from $`u`$ by replacing all $`x_i`$ with $`\chi ^i`$. For a pair of words $`u`$ and $`v`$, put
$$p_{uv}=\chi ^u(g_v).$$
(5)
Obviously, the following equalities hold:
$$p_{uu_1v}=p_{uv}p_{u_1v},p_{uvv_1}=p_{uv}p_{uv_1},$$
(6)
that is the operator $`p`$ is a bicharacter defined on a semigroup of all words. Sometimes we denote this operator by $`p(u,v)`$. Define a bilinear termal operation, a skew commutator, on a set of all quantum polynomials by setting
$$[u,v]=uvp_{uv}vu.$$
(7)
This satisfies the following identity:
$$[[u,v],w]=[u,[v,w]]+p_{wv}^1[[u,w],v]+(p_{vw}p_{wv}^1)[u,w]v,$$
(8)
which is similar to the Jacobi identity, where $`cdot`$ stands for usual multiplication in a free algebra, and which can be easily verified by direct computations using (7). Likewise, the following formulas of skew derivations, by which the skew commutator is linked to multiplication, are valid:
$$[u,vw]=[uv]w+p_{uv}v[uw],$$
(9)
$$[uv,w]=p_{vw}[uw]v+u[vw].$$
(10)
Definition 3. A super-letter is a polynomial equal to a standard nonassociative word with brackets defined as in operation (7).
By the Shirshov theorem, every standard word $`u`$ is associated with a super-letter $`[u]`$. If we remove the brackets in $`[u]`$ as is done in definition (7) we obtain an homogeneous polynomial whose leading word is equal to $`u,`$ and this leading word occurs in the decomposition of $`[u]`$ with coefficient 1. This is easily verified by induction on the degree. Indeed, if $`[u]=[[v][w]]`$ then the super-letter $`[u]`$ is equal to $`[v][w]p_{uw}[w][v]`$. By the induction hypothesis, $`[v]`$ and $`[w]`$ are homogeneous polynomials with the leading words $`v`$ and $`w`$, respectively. Therefore, the leading word of the first summand equals $`vw`$ and has coefficient 1; the leading word of the second equals $`wv`$ and is less than $`vw`$ by definition.
Thus, in correspondence with distinct standard words $`u`$ and $`v`$ are distinct super-letters $`[u]`$ and $`[v]`$, and the order on a set of super-letters can be defined as follows:
$$[u]>[v]u>v.$$
(11)
Definition 4. A word in super-letters is called a super-word. A super-word is said to be monotonic if it has the form
$$W=[u_1]^{k_1}[u_2]^{k_2}\mathrm{}[u_m]^{k_m},$$
(12)
where $`u_1<u_2<\mathrm{}<u_m`$.
We recall that the constitution of $`u`$ is a sequence of integers $`(m_1,m_2,\mathrm{},m_n)`$ such that $`u`$ has degree $`m_1`$ in $`x_1`$, degree $`m_2`$ in $`x_2`$, etc. Since super-letters and super-words are homogeneous in each of the variables, their constitutions can be defined in the obvious manner. Because $`G`$ is commutative, the elements $`g_u`$ and the characters $`\chi ^u`$ are the same for all words of a same constitution. For super-letters and super-words, therefore, $`G(W)=g_w`$ and $`p(U,V)=p_{uv}`$ are defined uniquely.
On the set of all super-words, consider a lexicographic order defined by the ordering of super-letters in (11).
LEMMA 5. A monotonic super-word $`W=[w_1]^{k_1}[w_2]^{k_2}\mathrm{}[w_m]^{k_m}`$ is greater than a monotonic super-word $`V=[v_1]^{m_1}[v_2]^{m_2}\mathrm{}[v_k]^{m_k}`$ if and only if the word $`w=w_1^{k_1}w_2^{k_2}\mathrm{}w_m^{k_m}`$ is greater than the word $`v=v_1^{m_1}v_2^{m_2}\mathrm{}v_k^{m_k}`$. Moreover, the leading word of the polynomial $`W`$, when decomposed into a sum of monomials, equals $`w`$ and has coefficient 1.
Proof. Let $`W>V`$. Then $`w_1v_1`$ in view of the ordering of super-letters. If $`w_1=v_1`$, we can remove one factor from the left of both $`V`$ and $`W`$, and then proceed by induction. Therefore, we will put $`w_1>v_1`$. If $`w_1`$ is not the beginning of $`v_1`$, then the latter inequality can be multiplied from the right by suitable distinct elements, which yields $`w>v`$, as required. Let $`v_1=(w_1^{k_1}w_2^{k_2}\mathrm{}w_{s1}^{k_{s1}})w_s^lv_1^{}`$, where $`0l<k_s`$. Note that the term between the parentheses may be missing (in which case $`s=1`$, $`l>0`$), and $`w_s`$ is not the beginning of $`v_1^{}`$. If $`v_1^{}`$ is a nonempty word, then $`v_1^{}<v_1<w_1w_s`$, since $`v_1`$ is standard. To obtain $`v<w`$, the inequality $`v_1^{}<w_s`$ will be multiplied from the left by one element $`(w_1^{k_1}w_2^{k_2}\mathrm{}w_{s1}^{k_{s1}})w_s^l`$, and from the right by (possibly) different elements. If $`v_1^{}`$ is the empty word, again we arrive at a contradiction with $`v_1`$ being standard. Indeed, if $`l>0`$, then the word $`v_1`$ should be greater than its end $`w_s`$; therefore, $`w_1>v_1>w_s`$, which contradicts the fact that $`w_1w_s`$ is valid for all $`s1`$. If $`l=0`$, then $`s>1`$, since $`v_1`$ begins with $`w_1`$. It follows that $`v_1`$ is greater than its end $`w_{s1}`$, which is again a contradiction with $`w_1>v_1>w_{s1}`$.
The second part of the lemma follows from the fact that the leading word of a product of homogeneous polynomials equals the product of leading words of the factors.
The lemma cannot be extended to the case of nonmonotonic super-words, for example, $`[x_1][x_3]>[x_1x_2]`$ and $`x_1x_3<x_1x_2`$.
LEMMA 6. Let $`u`$ and $`u_1`$ be standard words and $`u>u_1`$. Then the polynomial $`[[u][u_1]]`$ is a linear combination of super-words in the super-letters $`[w]`$ which lie properly between $`[u]`$ and $`[u_1]`$ and are such that $`wuu_1`$. In this case the degree of every summand in each of the variables $`x_1,\mathrm{},x_n`$ is equal to a respective degree of $`uu_1`$.
Proof. If the nonassociative word $`[[u][u_1]]`$ is standard then it defines a super-letter $`[w]`$. In this case $`u>w`$ since $`u`$ is the beginning of $`w`$, and $`w>u_1`$ by Lemma 2. In particular, the lemma is valid if the degrees of $`u`$ and $`u_1`$ are equal to 1. And we can therefore proceed by induction on the length of $`uu_1`$.
Suppose that our lemma is true if the length of $`uu_1`$ is less than $`m`$. Choose a pair $`u,u_1`$ with a greatest word $`[u]`$, so that the polynomial $`[[u][u_1]]`$ does not enjoy the required decomposition and the length of $`uu_1`$ equals $`m`$. Then the word $`[[u][u_1]]`$ is not standard, i.e., $`[u]=[[u_3][u_2]]`$ with $`u_2>u_1`$. We introduce the notation for super-letters $`U_i=[u_i]`$, $`i=1,2,3`$. By Jacobi identity (8), we can write
$$[[U_3U_2]]U_1]=[U_3[U_2U_1]]+p_{u_1u_2}^1[[U_3U_1]U_2]+(p_{u_2u_1}p_{u_1u_2}^1)[U_3U_1]U_2.$$
(13)
It follows that $`u_3>u>u_2>u_1`$. By the induction hypothesis, $`[U_3U_1]`$ can be represented as $`\underset{i}{}\alpha _i\underset{k}{}[w_{ik}]`$, where $`u_3>u_3u_1w_{ik}>u_1`$. Using Lemma 4, we obtain $`u>uu_1>u_3u_1w_{ik}`$, i.e., all super-letters $`[w_{ik}]`$ satisfy the requirements of the present lemma. Furthermore, the word $`u`$ cannot be the beginning of $`u_2`$, and so $`u>u_2`$ implies $`uu_1>u_2`$. Thus the super-letter $`U_2`$, too, satisfies the requirements. Consequently, the second \[in view of (7)\] and third summands of (13) have the required decomposition.
Using the induction hypothesis, for the first summand we obtain
$$[U_2U_1]=\underset{i}{}\beta _i\underset{k}{}[v_{ik}],$$
(14)
where $`u_2>u_2u_1v_{ik}>u_1`$. By Lemma 4, $`uu_1>u_2u_1v_{ik}`$, i.e., the super-letters $`[v_{ik}]`$ satisfy the conditions of the lemma. Rewrite the first summand using skew-derivation formula (9), with the first factor replaced by (14). With this, the first summand turns into a linear combination of words in the super-letters $`[v_{ik}]`$ and skew commutators $`[[u_3][v_{ik}]]`$. Since $`u_3>u>u_2>v_{ik}`$ and the length of $`v_{ik}`$ does not exceed that of $`u_2u_1`$, the induction hypothesis applies to yield
$$[[u_3][v_{ik}]]=\underset{j}{}\gamma _j\underset{t}{}[w_{jt}],$$
(15)
where $`u_3>u_3v_{ik}w_{jt}>v_{ik}`$. In this case $`u_2u_1v_{ik}`$ implies $`uu_1=u_3u_2u_1u_3v_{ik}w_{jt}`$; in addition, $`w_{jt}>v_{ik}>u_1`$, i.e., the super-letters $`[w_{jt}]`$ also satisfy the conditions.
LEMMA 7. Every nonmonotonic super-word is a linear combination of lesser monotonic super-words of a same constitution, whose super-letters all lie (not strictly) between the greatest and the least super-letters of a super-word given.
Proof. We proceed by induction on the degree. Whenever super-letters of a given super-word are rearranged, the degree of a polynomial remains fixed; therefore, the least super-word of degree $`m`$ will be monotonic. Assume that the lemma is true for super-words of degree $`<m`$, letting $`W=UU_1\mathrm{}U_t`$ be a least super-word of degree $`m`$ for which our lemma fails. If the super-word $`U_1\mathrm{}U_t`$ is not monotonic, by the induction hypothesis, then, it is a linear combination of lesser monotonic super-words $`W_i`$. And we can now apply the induction hypothesis to $`UW_i`$. Let
$$W=UU_1^{k_1}\mathrm{}U_t^{k_t},U_1<U_2<\mathrm{}<U_t.$$
(16)
If $`UU_1`$ then $`W`$ is monotonic, and there is nothing to prove. Let $`U>U_1`$. Then
$$W=[UU_1]U_1^{k_11}\mathrm{}U_t^{k_t}+p_{uu_1}U_1UU_1^{k_11}\mathrm{}U_t^{k_t}.$$
(17)
The second summand being a super-word is less than $`W`$, and so we can write it in the required form. By Lemma 6, the factor $`[UU_1]`$ in the first term can be represented as $`\underset{i}{}\alpha _i\underset{j}{}[w_{ij}]`$, where the super-letters $`[w_{ij}]`$ are less than $`U`$. Consequently, the super-letters $`\underset{j}{}[w_{ij}]U_1^{k_11}\mathrm{}U_t^{k_t}`$ are less than $`W`$, i.e., the first term, and hence also $`W`$, will have the required representation.
THEOREM 1. The set of all monotonic super-words constitute a basis for a free algebra $`𝐤x_1,\mathrm{},x_n`$.
Proof. Since the letters $`x_1,\mathrm{},x_n`$ are super-letters, every polynomial is a linear combination of monotonic super-words by Lemma 7. Our present goal is to prove that the set of all monotonic super-words is linearly independent. Let
$$\underset{i}{}\alpha _iW_i=0$$
(18)
and assume that $`W=[w_1]^{k_1}[w_2]^{k_2}\mathrm{}[w_m]^{k_m}`$ is a leading super-word in (18). By Lemma 5, the leading word of $`W`$ equals $`w=w_1^{k_1}w_2^{k_2}\mathrm{}w_m^{k_m}`$. Note that this word occurs exactly once in (18). Suppose, to the contrary, that $`W`$ does also occur in the decomposition $`V=[v_1]^{m_1}[v_2]^{m_2}\mathrm{}[v_k]^{m_k}`$. Then the word $`w`$ is less than or equal to the leading word $`v=v_1^{m_1}v_2^{m_2}\mathrm{}v_k^{m_k}`$ in the decomposition of $`V`$, which contradicts the fact that $`W>V`$ by Lemma 5.
Consider a free enveloping algebra in a given set of quantum variables $`Hx_1\mathrm{},x_n=G𝐤x_1,\mathrm{},x_n`$, on which the coproduct is defined by setting
$$\mathrm{\Delta }(x_i)=x_i1+g_{x_i}x_i,\mathrm{\Delta }(g)=gg,$$
(19)
and group-like elements commute with variables via $`xg=\chi ^x(g)gx`$. It follows that $`G𝐤x_1,\mathrm{},x_n`$ turns into a Hopf algebra; for details, see \[1, Sec. 3\].
LEMMA 8. The coproduct at a super-letter $`W=[w]`$ is represented thus:
$$\mathrm{\Delta }([w])=[w]1+g_w[w]+\underset{i}{}\alpha _iG(W_i^{\prime \prime })W_i^{}W_i^{\prime \prime },$$
(20)
where $`W_i^{}`$ are nonempty words in less super-letters than is $`[w]`$ . Moreover, the sum of degrees of super-words $`W_i^{}`$ and $`W_i^{\prime \prime }`$ in each variable $`x_j`$ equals the degree of $`W`$ in that variable, i.e., the sum of structural elements of $`W_i^{}`$ and $`W_i^{\prime \prime }`$ is equal to the constitution of $`W`$.
Proof. We use induction on the length of a word $`w`$. For letters, by (19), there is nothing to prove. Let $`W=[U,V]`$, $`U=[u]`$, and $`V=[v]`$. Assume that the decompositions
$$\mathrm{\Delta }(U)=U1+g_uU+\underset{i}{}\alpha _iG(U_i^{\prime \prime })U_i^{}U_i^{\prime \prime },$$
(21)
and
$$\mathrm{\Delta }(V)=V1+g_vV+\underset{j}{}\beta _jG(V_j^{\prime \prime })V_j^{}V_j^{\prime \prime }$$
(22)
satisfy the requirements of the lemma. Using (7) and properties of a bicharacter $`p`$, we can write
$$\mathrm{\Delta }(W)=\mathrm{\Delta }(U)\mathrm{\Delta }(V)p_{uv}\mathrm{\Delta }(V)\mathrm{\Delta }(U)=W1+g_wW+$$
$$(1p_{uv}p_{vu})g_uVU+\beta _jp(U,V_j^{\prime \prime })G(V_j^{\prime \prime })[UV_j^{}]V_j^{\prime \prime }+$$
$$\beta _jg_uG(V_j^{\prime \prime })V_j^{}(UV_j^{\prime \prime }p_{uv}p(V_j^{},U)V_j^{\prime \prime }U)+$$
(23)
$$\alpha _iG(U_i^{\prime \prime })(U_i^{}Vp_{uv}p(V,U_i^{\prime \prime })VU_i^{})U_i^{\prime \prime }+\alpha _ip(U_i^{},V)g_vG(U_i^{\prime \prime })U_i^{}[U_i^{\prime \prime }V]+$$
$$\alpha _i\beta _jG(U_i^{\prime \prime }V_j^{\prime \prime })(p(U_i^{},V_j^{\prime \prime })U_i^{}V_j^{}U_i^{\prime \prime }V_j^{\prime \prime }p_{uv}p(V_j^{},U_i^{\prime \prime })V_j^{}U_i^{}V_j^{\prime \prime }U_i^{\prime \prime }).$$
Collecting similar terms in this formula will result in the canceling of terms of the form $`g_vUV`$ only. We claim that all left parts of the remaining tensors in (23) admit the required decomposition. First, in view of the induction hypothesis, all super-letters of all super-words $`V_j^{}`$ are less than $`V`$, which are in turn less than $`W`$ because $`v`$ is the end of a standard word $`w`$. Moreover, by the induction hypothesis again, $`u`$ cannot be the beginning of any word $`u^{}`$ such that the super-letter $`[u^{}]`$ would occur in super-words $`U_i^{}`$. Therefore, $`u>u^{}`$ implies $`uv>u^{}`$ or $`W>[u^{}]`$. Thus all but the first and fourth super-words on the left-hand sides of all tensors depend only on super-letters which are less than $`W`$.
We want to apply Lemma 6 to the fourth tensor. Let $`V_j^{}=\underset{k}{}V_{ik}`$, where $`V_{ik}=[v_{ik}]`$ are less than $`V`$. By formula (9) the polynomial $`[U,V_j^{}]`$ is a linear combination of words in the super-letters $`V_{ik}`$ and skew commutators $`[U,V_{ik}]`$. By Lemma 6, each of these commutators is a linear combination of words in the super-letters $`[v^{}]`$ such that $`v^{}uv_{ik}`$. In view of $`v_{ik}<v`$, we obtain $`v^{}<uv=w`$.
The statement concerning the constitution follows immediately from formula (23) and the induction hypothesis.
LEMMA 9. The coproduct at a super-word $`W`$ is represented thus:
$$\mathrm{\Delta }(W)=W1+G(W)W+\underset{i}{}\alpha _iG(W_i^{\prime \prime })W_i^{}W_i^{\prime \prime },$$
(24)
where the sum of constitutions of $`W_i^{}`$ and $`W_i^{\prime \prime }`$ equals the constitution of $`W`$.
Proof. It suffices to observe that $`\mathrm{\Delta }`$ is an homomorphism of algebras. Here, we can no longer assert that $`W_i^{}<W`$.
3. BASIS FOR A CHARACTER HOPF ALGEBRA
Consider a Hopf algebra $`H`$ generated by a set of skew primitive semi-invariants $`a_1,\mathrm{},a_n`$ and by an Abelian group $`G`$ of all group-like elements. Denote by $`H_a`$ a subalgebra generated by $`a_1,\mathrm{},a_n`$. Then $`H=GH_a`$ since by definition, semi-invariants obey the following commutation rule:
$$ag=\chi ^a(g)ga.$$
(25)
Let $`x_1,\mathrm{},x_n`$ be quantum variables with the same parameters as $`a_1,\mathrm{},a_n`$, respectively, that is $`\chi ^{x_i}=\chi ^{a_i}`$ and $`g_{x_i}=g_{a_i}`$. Then there exists an homomorphism
$$\phi :𝐤x_1,\mathrm{},x_nH_a$$
(26)
which maps $`x_i`$ to $`a_i`$. This allows us to extend all the combinatorial notions applied to the words in $`x_1,\mathrm{},x_n`$ in the above sections to the words in $`a_1,\mathrm{},a_n`$.
With $`a_1,\mathrm{},a_n`$ we associate the respective natural degrees $`d_1,\mathrm{},d_n`$.<sup>2</sup><sup>2</sup>2Note that further argument will remain true for the case where $`d_1,\mathrm{},d_n`$ are arbitrary positive elements of a linearly ordered additive Abelian group. In this way, every word, super-letter, and super-word of a constitution $`(m_1,\mathrm{},m_n)`$ have degree $`m_1d_1+\mathrm{}+m_nd_n`$.
Definition 5. A $`G`$-super-word is a product of the form $`gW`$, where $`gG`$ and $`W`$ is a super-word. The degree, constitution, length, and other concepts which apply with $`G`$-super-words are defined by the super-word $`W`$. Alternatively, we assume that the degree and the constitution of $`gG`$ are equal to zero. In view of (25), every product of super-letters and group-like elements equals a linear combination of $`G`$-super-words of the same constitution.
Definition 6. A super-letter $`[u]`$ is said to be hard if it is not a linear combination of words of the same degree in less super-letters than is $`[u]`$ and of $`G`$-super-words of a lesser degree.
Definition 7. We say that the height of a super-letter $`[u]`$ of degree $`d`$ equals a natural number $`h`$ if $`h`$ is least with the following properties:
(1) $`p_{uu}`$ is a primitive root of unity of degree $`t1`$, and either $`h=t`$ or $`h=tl^k`$, where $`l`$ is the characteristic of the ground field;
(2) a super-word $`[u]^h`$ is a linear combination of super-words of degree $`hd`$ in less super-letters than is $`[u]`$ and of $`G`$-super-words of a lesser degree.
If, for the super-letter $`[u]`$, the number $`h`$ with the above properties does not exist then we say that the height of $`[u]`$ is infinite.
Definition 8. The monotonic $`G`$-super-word
$$g[u_1]^{n_1}[u_2]^{n_2}\mathrm{}[u_k]^{n_k}$$
is said to be restricted if each of the numbers $`n_i`$ is less than the height of the super-letter $`u_i`$.
THEOREM 2. If a Hopf algebra $`H`$ is generated by a set skew-primitive semi-invariants $`a_1,\mathrm{},a_n`$ and by an Abelian group $`G`$ of all group-like elements, then the set of all monotonic restricted $`G`$-super-words in hard super-letters constitute a basis for $`H`$.
The proof will proceed through a number of lemmas. For brevity, we call a super-word (a $`G`$-super-word) admissible if it is monotonic restricted and is a word in hard super-letters only.
LEMMA 10. Every nonadmissible super-word of degree $`d`$ is a linear combination of lesser admissible super-words of degree $`d`$ and of admissible $`G`$-super-words of a lesser degree. Also, all super-letters occurring in super-words of degree $`d`$ of this linear combination are less than or equal to a greatest super-letter of the super-word given.
The proof is by induction on the degree. Assume that the lemma is valid for super-words of degree $`<m`$. Let $`W`$ be a least super-word of degree $`m`$ for which the required representation fails. By Lemma 7, the super-word $`W`$ is monotonic. If it has a nonhard super-letter, by definition, we can replace it with a linear combination of $`G`$-super-words of a lesser degree and of words in less super-letters of the same degree. Removing the parentheses turns $`W`$ into a linear combination of $`G`$-super-words of a lesser degree and of lesser super-words of the same degree, a contradiction with the choice of $`W`$. If $`W`$ contains a subword $`[u]^k`$, where $`k`$ equals the height of $`[u]`$, then we can replace it as is specified above, which gives us a contradiction again. Thus the $`W`$ is itself monotonic restricted and is a word in hard super-letters only.
In order to prove Theorem 2, it remains to show that admissible $`G`$-super-words are linearly independent. Consider an arbitrary linear combination $`𝐓`$ of admissible $`G`$-super-words and let $`U=U_1^{n_1}U_2^{n_2}\mathrm{}U_k^{n_k}`$ be its leading super-word of degree $`m`$. Multiplying, if necessary, that combination by a group-like element, we can assume that $`U`$ occurs once without a group-like element:
$$𝐓=U+\underset{j=1}{\overset{r}{}}\alpha _jg_jU+\underset{i}{}\alpha _ig_iV_{i1}^{n_{i1}}V_{i2}^{n_{i2}}\mathrm{}V_{is}^{n_{is}}.$$
(27)
In the next three lemmas, we accept the following inductive assumption on $`m`$ and on $`r`$:
$$()\begin{array}{c}\text{the set of all admissible}G\text{-super-words of degree}m\text{which are less than}U,\\ \text{of admissible}G\text{-super-words of degree}<m,\text{and of}G\text{-super-words}\\ g_jU,1jr,\text{is linearly independent [we can assume}r=0\text{ in (27)].}\end{array}$$
In view of this assumption and Lemma 10, every super-word of degree $`m`$ which is less than $`U`$, and every super-word of degree $`<m`$, can be uniquely decomposed into a linear combination of admissible $`G`$-super-words. For brevity, such will be referred to as a basis decomposition.
LEMMA 11. If $`𝐓`$ is a skew primitive element then $`r=0`$ and all $`G`$-super-words of degree $`m`$ in (27) are super-words.
Proof. Rewrite the linear combination $`𝐓`$ as follows:
$$𝐓=U+\underset{i=1}{\overset{k}{}}\gamma _ig_iW_i+W^{},$$
(28)
where $`g_iW_i`$ are distinct $`G`$-super-words of degree $`m`$ in (27) and $`W^{}`$ a linear combination of $`G`$-super-words of degree $`<m`$. In the expression
$$\mathrm{\Delta }(𝐓)𝐓1g_t𝐓,$$
(29)
consider all tensors of the form $`gW\mathrm{}`$, where $`W`$ is of degree $`m`$. By Lemma 9, the sum of all such tensors equals
$$\underset{i=1}{\overset{r}{}}\gamma _ig_iW_ig_i\underset{i=1}{\overset{r}{}}\gamma _ig_iW_i1=\underset{i=1}{\overset{r}{}}\gamma _ig_iW_i(g_i1).$$
(30)
By inductive assumption $`()`$, the elements $`g_iW_i`$ are linearly independent modulo all left parts of tensors of degree $`<m`$ in (29). Therefore, if (29) vanishes then either $`\gamma _i=0`$ or $`g_i=1`$ for every $`i`$, as required.
LEMMA 12. If $`𝐓`$ is a skew primitive element then $`U=U_1^{n_1}`$ and all super-words of degree $`m`$ except $`U`$ are words in less super-letters than is $`U_1`$.
Proof. By the preceding lemma, we can assume that
$$𝐓=\underset{i}{}\alpha _ig_iV_{i1}^{n_{i1}}V_{i2}^{n_{i2}}\mathrm{}V_{is}^{n_{is}},$$
(31)
where $`V_{ij}=[v_{ij}]`$ are hard super-letters, $`\alpha _i`$ are nonzero coefficients, and $`g_i=1`$ if $`V_i`$ is of degree $`m`$. We apply coproduct to (31). By (8), then, the right-hand side assumes the form
$$\underset{i}{}\alpha _i(g_ig_i)\underset{j=1}{\overset{s}{}}(V_{ij}1+g_{ij}V_{ij}+\underset{\theta }{}g_{ij\theta }V_{ij\theta }^{}V_{ij\theta }^{\prime \prime })^{n_{ij}},$$
(32)
where $`V_{ij\theta }^{}<V_{ij}`$ and $`\mathrm{deg}V_{ij\theta }^{}+\mathrm{deg}V_{ij\theta }^{\prime \prime }=\mathrm{deg}V_{ij}`$.
Let $`[v]`$ be the largest super-letter occurring in super-words of degree $`m`$ in (31). Since all super-words of (31) are monotonic, this super-letter stands at the end of some super-words $`V_i`$, i.e., $`[v]=V_{is}`$. If one of these super-words depends only on $`[v]`$, i.e., $`V_i=[v]^k`$, then $`V_i`$ is a leading term, as required. Therefore, we assume that every super-word of degree $`m`$ ending with $`[v]`$ is a word in more than one super-letter.
Let $`k`$ be the largest exponent $`n_{is}`$ of $`[v]`$ in $`𝐓`$. Consider all tensors of the form $`g[v]^k\mathrm{}`$ obtained in (32) by removing the parentheses and applying the basis decomposition to all left parts of tensors in all terms except $`𝐓1`$ (all of these terms are of degree $`<m`$).
All left parts of tensors which appear in
$$\mathrm{\Delta }(V_i)=(g_ig_i)\underset{j=1}{\overset{s}{}}(V_{ij}1+g_{ij}V_{ij}+\underset{\theta }{}g_{ij\theta }V_{ij\theta }^{}V_{ij\theta }^{\prime \prime })^{n_{ij}}$$
by removing the parentheses arise from the word $`V_i=\alpha _ig_iV_{i1}^{n_{i1}}V_{i2}^{n_{i2}}\mathrm{}V_{is}^{n_{is}}`$ by replacing some of the super-letters $`V_{ij}`$ either with group-like elements $`g_{ij}`$ or with $`G`$-super-words $`g_{ij\theta }V_{ij\theta }^{}`$ of a lesser degree in less super-letters. The right parts are, respectively, products obtained by replacing super-letters $`V_{ij}`$ or super-words $`V_{ij\theta }^{\prime \prime }`$ multiplied from the left by $`g_i`$.
If, under the replacements above, a new super-word is greater in degree than $`[v]^k`$, then its basis decomposition will give rise to terms of the form $`g[v]^k\mathrm{}`$. In this case, however, the right parts of those terms are of degree less than $`mk\mathrm{deg}([v])`$ since the sum of degrees of both parts of the tensors either remains equal to $`m`$ or decreases.
If a new super-word is of degree less than the degree of $`[v]^k`$, or the super-word is itself less than $`[v]^k`$ then its basis decomposition will be freed of terms of the form $`g[v]^k\mathrm{}`$; see Lemma 10.
If a new super-word is of degree equal to that of $`[v]^k`$ and $`V_i`$ is of degree less than $`m`$ then the new super-word can be greater than or equal to $`[v]^k`$. In this case the right-hand sides of the new tensors are of degree less than $`mk\mathrm{deg}([v])`$ because the sum of degrees of the left- and right-hand sides of the tensors is less than $`m`$.
If a new super-word is of degree equal to the degree of $`[v]^k`$, but $`V_i`$ does not end with $`[v]^k`$, i.e., $`V_i=W_i[v]^s`$, $`0s<k`$, then the new super-word is less than $`[v]^k`$ since its first super-letter is less than $`[v]`$. (All super-letters of $`W_i`$ cannot be replaced with group-like elements, since otherwise the new word would be of degree less than or equal to the degree of $`[v]^s`$.)
Finally, if $`V_i=W_i[v]^k`$ then a super-word of degree $`k\mathrm{deg}([v])`$, which is greater than or equal to $`[v]^k`$, may appear only if all super-letters of the super-words $`W_i`$ are replaced with group-like elements, but $`[v]`$ is not. Here, the resulting tensor is of the form $`g(W_i)[v]^k\alpha _iW_i`$.
We fix an index $`t`$ such that $`V_t`$ ends with $`[v]^k`$, letting $`t=1`$. Then the sum of all tensors of the form $`G(W_1)[v]^k\mathrm{}`$ in $`\mathrm{\Delta }(𝐓)𝐓1`$ is equal to
$$G(W_1)[v]^k(\underset{j}{}\alpha _jW_j+𝐖^{}),$$
(33)
where $`𝐖^{}`$ is a linear combination of basis elements of degree less than $`mk\mathrm{deg}([v])`$, and $`j`$ runs through the set of all indices $`i`$ such that $`V_i=W_i[v]^k`$, $`G(W_i)=G(W_1)`$, and the degree of $`W_i`$ equals $`mk\mathrm{deg}([v])`$. Since $`W_i`$ are distinct nonempty basis super-words of degree less than $`m`$, tensor (33) is nonzero.
LEMMA 13. Under the conditions of Lemma $`12,`$ either $`n_1=1`$ or $`p(U_1,U_1)`$ is a primitive root of unity of degree $`t1`$, in which case $`n_1=t`$ or the characteristic of a base field equals $`l>0`$, and $`n_1=tl^k`$.
Proof. By the previous lemma, the linear combination $`𝐓`$ can be written in the form
$$𝐓=U^k+\underset{i}{}\alpha _ig_iV_{i1}^{n_{i1}}V_{i2}^{n_{i2}}\mathrm{}V_{is}^{n_{is}},$$
(34)
where $`U=[u]`$ is greater than all super-letters $`V_{ij}`$ for $`V_i`$ of degree $`m`$. First let $`\xi =1+p_{uu}+p_{uu}^2+\mathrm{}+p_{uu}^{k1}0`$ and assume $`k>1`$.
In the basis decomposition of $`\mathrm{\Delta }(𝐓)𝐓1`$, consider tensors of the form $`U^{k1}\mathrm{}`$. All super-letters $`V_{ij}`$ in super-words of degree $`m`$ are less than $`[u]`$; therefore, tensors of this form may appear under the basis decomposition of a tensor of $`\mathrm{\Delta }(V_i)V_i1`$ only if either the left part of that tensor is of degree greater than $`(k1)\mathrm{deg}([u])`$ or $`V_i`$ is of degree less than $`m`$. In either case the right part is of less degree than is $`[u]`$. As above, if we remove the parentheses in
$$\mathrm{\Delta }(U^k)=(U1+g_uU+\underset{\tau }{}U_\tau ^{}U_\tau ^{\prime \prime })^k,$$
(35)
we see that the left parts of the resulting tensors arise from the super-word $`U^k`$ by replacing some super-letters $`U`$ either with $`g_u`$ or with super-words $`U_\tau ^{}`$ of a lesser degree in less super-letters than is $`U`$. It follows that a super-word of degree $`(k1)\mathrm{deg}(U)`$ which is greater than or equal to $`U^{k1}`$ appears only if exactly one super-letter is replaced with a group element. Using the commutation rule $`U^sg_u=p_{uu}^sg_uU^s`$, we see that the sum of all tensors of the form $`g_uU^{k1}\mathrm{}`$ equals
$$g_uU^{k1}(\xi U+𝐖),$$
where $`𝐖`$ is a linear combination of basis $`G`$-super-words of degree less than $`\mathrm{deg}(U)`$. Consequently, (29) is nonzero for $`k1`$.
Now let $`\xi =0`$. Then $`p_{uu}^k=1`$. Therefore, $`p_{uu}`$ is a primitive root of unity of some degree $`t`$ (we put $`t=1`$ if $`p=1`$) and the number $`k`$ is divisible by $`t`$. We can write $`k`$ in one of the forms $`tq`$ or $`tl^kq`$, where $`l`$ is the characteristic of a base field, in which $`q10`$. Put $`h=t`$ or $`h=tl^k`$, respectively. Since $`(U1)(g_uU)=p_{uu}(g_uU)(U1)`$, use will be made of the quantum Newton binomial formula
$$(U1+g_uU)^h=U^h1+g_{uu}^hU^h.$$
(36)
This implies that if we remove the parentheses in
$$\mathrm{\Delta }(U^h)=((U1+g_{uu}U)+\underset{i}{}G(U_i^{\prime \prime })U_i^{}U_i^{\prime \prime })^h,$$
(37)
then Lemma 8 gives
$$\mathrm{\Delta }(U^h)=U^h1+g_{uu}^hU^h+\underset{\theta }{}G(U_\theta ^{\prime \prime })U_\theta ^{}U_\theta ^{\prime \prime },$$
(38)
where all super-words $`U_\theta ^{}`$ are less than $`U^h`$ and are of less degree than is $`U^h`$. In this formula, we note, all terms $`U^r\mathrm{}`$, $`r<h`$, whose left parts are greater than $`U^h`$, are banished. This allows us to treat $`U^h`$ in (34) as a single block, or as a new formal super-letter $`\{U^h\}`$ such that $`\{U^h\}<U`$, and $`\{U^h\}>[v_{ij}]`$ if $`u^h>v_{ij}`$ (which is equivalent to $`u>v_{ij}`$ by Lemma 3), i.e.,
$$𝐓=\{U^h\}^q+\underset{i}{}\alpha _ig_iV_{i1}^{n_{i1}}V_{i2}^{n_{i2}}\mathrm{}V_{is}^{n_{is}}.$$
(39)
Since $`p(U^h,U^h)=p^{hh}=1`$, we have
$$\xi _1=1+p(U^h,U^h)+\mathrm{}+p(U^h,U^h)^{q1}=q0.$$
As in the case above, assuming that $`\{U^h\}`$ is a single block, we can compute the sum of all tensors of the form $`g_u^h\{U^h\}^{q1}\mathrm{}`$ in the basis decomposition of $`\mathrm{\Delta }(𝐓)𝐓1`$ (provided that $`q>1`$):
$$g_{uu}^h\{U^h\}^{q1}(q\{U^h\}+𝐖),$$
(40)
where $`𝐖`$ is a linear combination of basis $`G`$-super-words of less degree than is $`U^h`$. By the induction hypothesis, tensor (40) is nonzero, and so therefore is (29).
The equality $`𝐓=0`$ does not hold. Indeed, if it did, then $`𝐓`$ would be a skew primitive element, which is nonzero in view of Lemma 13 and definitions of hard super-letters and their heights. Inductive assumption $`()`$ for $`r=0`$ is obviously valid if $`U`$ is smallest among generators $`a_i`$, since group-like elements, i.e., $`G`$-super-words of degree zero, are always linearly independent. Theorem 2 is proved.
4. SOME COROLLARIES
In this section, again we write $`H`$ for a Hopf algebra generated by an Abelian group $`G`$ of all group-like elements and by skew primitive semi-invariants $`a_1,\mathrm{},a_n`$ with which degrees $`d_1,\mathrm{},d_n`$ are associated.
COROLLARY 1. The set of all $`G`$-words in $`a_1,\mathrm{},a_n`$, obtained by dropping all brackets from monotonic restricted $`G`$-super-words in hard super-letters, constitute a basis for $`H`$.
Proof. Decompose an arbitrary word $`v`$ in $`a_1,\mathrm{},a_n`$ as is specified in Theorem 1, namely, $`v=\underset{j}{}\alpha _jV_j`$, where $`V_j=[v_{j1}]^{n_1}\mathrm{}[v_{jk}]^{n_k}`$ are monotonic super-words of the same constitution. By Lemma 5, the leading word appearing in $`\underset{j}{}\alpha _jV_j`$ under decomposition (7) equals $`v_s=v_{s1}^{n_1}\mathrm{}v_{sk}^{n_k}`$, where $`V_s`$ is the leading super-word among all $`V_j`$. Therefore, $`v=v_s`$, $`\alpha _s=1`$ — this is still a decomposition in the free algebra.
We use induction on the degree. Let $`w`$ be a minimal word of degree $`d`$ which is not a linear combination of the $`G`$-words specified in the statement. The word, as in the preceding paragraph, is decomposed thus: $`w=\underset{j}{}\alpha _jW_j`$. If the leading super-word $`W_s`$ is admissible, then $`w`$ arises from $`W_s`$ by dropping the brackets, and so there is nothing to prove. If $`W_s`$ is not admissible, $`W_s`$ is the required linear combination by Lemma 10 and inductive assumption $`()`$. We have $`w=(\underset{js}{}\alpha _jW_j)+W_s`$, where the first summand is a linear combination of words which are less than $`w`$, and again the inductive assumption applies.
We argue for linear independence. Let $`\underset{it}{}\beta _{it}g_{it}w_i=0`$. Then $`w_i=\alpha _{ij}W_{ij}`$, where $`w_i`$ is obtained by dropping all brackets from the leading super-word $`W_{is}`$, and $`\alpha _{is}=1`$. Therefore, $`W_{is}`$ are admissible super-words. Now the equality $`\underset{ijt}{}\beta _{it}\alpha _{ij}g_{it}W_{ij}=0`$ leads us to a contradiction. Indeed, by Lemma 10, the nonleading super-words $`W_{ij}`$ decrease under the basis decomposition, either in degree or in ordering.
Sometimes we find it useful to apply the following criterion which allows us to forget about skew commutators in computing hard super-letters.
COROLLARY 2. A super-letter $`[u]`$ is hard if and only if the standard word $`u`$ is not a linear combination of lesser words of degree $`\mathrm{deg}(u)`$ and of $`G`$-words of a lesser degree.
Proof. Let $`u=\underset{i}{}\alpha _iw_i+u_0`$, where $`w_i<u`$ and $`\mathrm{deg}(u_0)<\mathrm{deg}(u)`$. Decompose the words $`u`$ and $`w_i`$ as was done at the beginning of Corollary 1. We obtain $`u=[u]+\underset{j}{}\beta _jU_j`$ and $`w_i=\underset{t}{}\beta _tW_{it}`$, where the super-words $`U_j`$ are less than $`[u]`$, and $`w_i`$ equals the leading word of a polynomial defined by the leading super-word $`W_{is}`$. Since $`u>w_i`$, we have $`[u]>W_{is}`$ by Lemma 5, and hence $`[u]`$ is greater than all $`W_{it}`$. Consequently, the basis decompositions of $`U_i`$ and $`W_{it}`$ have only super-words which either are less than $`[u]`$ or of a lesser degree. For hard super-letters $`[u]`$, therefore, the equality $`[u]+\underset{j}{}\beta _jU_j\underset{i}{}\alpha _i\underset{t}{}\beta _tW_{it}u_0=0`$ is an impossibility.
Conversely, if $`[u]=\alpha _iW_i+U_0`$, where $`W_i`$ depends on super-letters less than $`[u]`$, then
$$u=[u]+(u[u])=\alpha _iW_i+U_0+(u[u]),$$
and the polynomial in the right part has no monomials whose degree equals the degree of $`u`$ and which are greater than or equal to $`u`$.
COROLLARY 3. Lemmas $`11`$-$`1`$3 are valid without assumption $`()`$.
COROLLARY 4. A Hopf algebra $`H`$ is finite-dimensional if and only if the group $`G`$ and the set of all hard super-letters are finite, and each hard super-letter has finite height.
COROLLARY 5. If $`H`$ has only a finite number of hard super-letters and $`G`$ is finitely generated, then $`H`$ is (left and right) Noetherian.
COROLLARY 6. Let the group algebra $`𝐤[G]`$ has no zero divisors. If $`H`$ has only a finite number of hard super-letters, of which each has infinite height, then $`H`$ has no zero divisors and has a classical skew field of quotients.
As in the case of classical Lie algebras, in order to prove the last two corollaries, we need only construct on $`H`$ a filtration $`H_0H_1\mathrm{}H_k\mathrm{}`$ such that the associated graded algebra $`D(H)`$ satisfies the required properties; see, e.g., \[17, Ch. V, Sec. 3, Thms. 4, 5\].
Construction of the filtration. Assume that $`H`$ has finitely many hard super-letters. Consider a set $`𝐑`$ of all words in $`a_1,\mathrm{},a_n`$, whose degree does not exceed the maximal degree of a hard super-letter multiplied by a maximal finite height, or by 2 if all heights are infinite. In this case $`𝐑`$ is composed of all standard words defining hard super-letters and of all products $`uv`$, $`u^h`$, where $`[u]`$ and $`[v]`$ are hard super-letters and $`h`$ is the height of $`u`$. Let words of $`𝐑`$ all respect the lexicographical ordering described at the beginning of Sec. 1 and $`n(u)`$ be the number of words in $`𝐑`$ which are less than or equal to $`u`$. The largest word $`a_1`$ is defined by the number $`L=n(a_1)`$. Denote by $`M`$ an arbitrary natural number which is greater than the length of any word in $`𝐑`$. Define the filtration degree on hard super-letters using the formula
$$\mathrm{Deg}([u])=M^{L+1}\mathrm{deg}(u)+M^{n(u)},$$
(41)
where $`\mathrm{deg}(u)`$ is specified by the constitution of $`u`$, and $`\mathrm{deg}(u)=d_1m_1+\mathrm{}+d_nm_n`$. The filtration degree of a basis element $`gW`$ equals the sum of filtration degrees of all of its super-letters. The filtration degree of an arbitrary element $`𝐓H`$ equals the maximal filtration degree of the basis elements occurring in its basis decomposition.
LEMMA 14. The function $`\mathrm{Deg}`$ defines a filtration on $`H`$, so that
$$H_0=𝐤[G],H_k=\{𝐓H|\mathrm{Deg}(𝐓)k\}.$$
Proof. We have to show that $`H_kH_sH_{k+s}`$, i.e., $`\mathrm{Deg}(𝐓_1𝐓_2)\mathrm{Deg}(𝐓_1)+\mathrm{Deg}(𝐓_2)`$. To do this, we construct an additional degree function $`D^{}`$ on a set of all linear combinations of (not necessarily admissible) super-words in the super-letters defined by all standard words of the vocabulary $`𝐑`$ The $`D^{}`$-degree of a super-letter is defined by formula (41). The $`D^{}`$-degree of a product of super-letters equals the sum of degrees of its factors. Accordingly, the $`D^{}`$-degree of a linear combination equals the maximum of $`D^{}`$-degrees of its summands. Of course we do not claim that the various linear combinations defining equal elements of $`H`$ have the same $`D^{}`$-degrees.
If we assume that $`𝐓_1𝐓_2`$ is obtained from $`𝐓_1`$ and $`𝐓_2`$ by merely removing the parentheses, then $`\mathrm{Deg}(𝐓_1)+\mathrm{Deg}(𝐓_2)=D^{}(𝐓_1𝐓_2)`$. Therefore, it suffices to specify how a basis decomposition of super-words proceeds in a way that $`D^{}`$ is kept unincreased. Our plan is as follows. First, we replace nonhard super-letters via Definition 6, next replace all subwords $`[u]^h`$, where $`h`$ is the height of a hard super-letter $`[u]`$, then apply the decreasing decomposition of Lemma 7, and again replace nonhard super-letters, etc.
Let $`[u]`$ be a nonhard super-letter defined by $`u`$ in $`𝐑`$ as follows:
$$[u]=\underset{i}{}\alpha _i\underset{j}{}[w_{ij}]+\underset{s}{}\alpha _sg_s\underset{t}{}[v_{st}],$$
(42)
where $`[w_{ij}]`$ are less than $`[u]`$, $`n(w_{ij})n(u)1`$, of the combined ordinary degree $`\mathrm{deg}(u)`$, and $`\underset{t}{}\mathrm{deg}(v_{st})\mathrm{deg}(u)1`$. By the definition of $`R`$, all words $`w_{ij}`$ and $`v_{st}`$, as well as the words which result from the summands of (42) by dropping all brackets and multiplication signs, belong to the vocabulary $`𝐑`$. In particular, their lengths are less than $`M`$. Therefore, the number of factors in every summand of (42) is less than $`M`$. Thus we may write
$$D^{}(\underset{j}{}[w_{ij}])=\underset{j}{}M^{L+1}\mathrm{deg}(w_{ij})+\underset{j}{}M^{n(w_{ij})}<M^{L+1}\mathrm{deg}(u)+MM^{n(u)1}=D^{}(u).$$
(43)
Similarly,
$$\begin{array}{c}D^{}(_t[v_{st}])=_tM^{L+1}\mathrm{deg}(v_{st})+_tM^{n(v_{st})}<\\ M^{L+1}(\mathrm{deg}(u)1)+MM^L=D^{}(u)M^{n(u)}M^{L+1}+M^{L+1}<D^{}(u).\end{array}$$
(44)
The argument is the same if $`[u]^h`$ is changed with $`\underset{i}{}\underset{j}{}[w_{ij}]+\underset{s}{}g_s\underset{t}{}[v_{st}]`$, where $`[u]`$ is a hard super-letter of height $`h`$. We have
$$D^{}(\underset{i}{}\underset{j}{}[w_{ij}]+\underset{s}{}g_s\underset{t}{}[v_{st}])<D^{}([u]^h).$$
(45)
We proceed to the decreasing process of Lemma 7. The second summand in (17) has the same $`D^{}`$-degree as $`W`$. By Lemma 6, the factor $`[UU_1]`$ in the first summand can be represented as $`\underset{i}{}\alpha _i\underset{j}{}[w_{ij}]`$, where the super-letters $`[w_{ij}]`$ are less than $`U`$, and they all are in $`𝐑`$ since their ordinary degrees are less than or equal to the ordinary degree of $`UU_1`$. Therefore, $`n(w_{ij})n(u)1`$. And the constitutions being equal indicates that the number of factors in $`\underset{j}{}`$ does not exceed $`M`$. We have
$$\begin{array}{c}D^{}(_j[w_{ij}])=_jM^{L+1}\mathrm{deg}(w_{ij})+M^{n(w_{ij})}<\\ M^{L+1}\mathrm{deg}(UU_1)+MM^{n(u)1}<M^{L+1}\mathrm{deg}(UU_1)+M^{n(u)}+M^{n(u_1)}=D^{}(UU_1).\end{array}$$
(46)
The lemma is proved.
Note that the strict inequality signs in (46) show that the inequality
$$\mathrm{Deg}([u][v]p_{uv}[v][u])<\mathrm{Deg}(u)+\mathrm{Deg}(v)$$
(47)
is valid for all hard super-letters $`U=[u]`$, $`U_1=[v]`$, $`u>v`$. Similarly, (45) gives us
$$\mathrm{Deg}([u]^h)<h\mathrm{Deg}(u),$$
(48)
where $`h`$ is the height of a hard super-letter $`[u]`$.
Associated graded algebra. With each hard super-letter $`[u]`$ we associate a new variable $`x_u`$. Denote by $`H^{}`$ an algebra generated by $`x_u`$ and defined by the relations $`x_ux_v=p_{uv}x_vx_u`$, where $`u>v`$. This algebra can be constructed by iterating the skew polynomials construction. Namely, let $`[u_1]<[u_2]<\mathrm{}<[u_s]`$ all be hard super-letters. Denote by $`H_k^{}`$ a subalgebra generated by $`x_{u_1},\mathrm{},x_{u_k}`$. Then $`H_1^{}`$ is isomorphic to an algebra of polynomials in one variable. The map $`\phi :x_vp_{uv}x_v`$, where $`u=u_{k+1}`$, $`v=u_i`$, $`ik`$, determines an automorphism of $`H_k^{}`$. The commutation rule $`x_ux_v=p_{uv}x_vx_u`$ can be written in the form $`x_ux_v=x_v^\phi x_u`$. Therefore, the algebra $`H_{k+1}^{}`$ is isomorphic to an algebra of skew polynomials over $`H_k^{}`$. In particular, the algebra $`H^{}`$ is Noetherian and has no zero divisors.
Define the action of $`G`$ on $`H^{}`$ by the formula $`x_u^g=\chi ^u(g)x_u`$. Let $`H^{}[G]`$ be a skew group algebra of $`G`$ with coefficients from $`H^{}`$.
THEOREM 3. The associated graded algebra $`D(H)`$ is isomorphic to the quotient algebra of $`H^{}[G]`$ with respect to relations $`x_u^h=0`$, where $`[u]`$ runs through the set of all hard super-letters of finite height and $`h`$ equals the height of $`[u]`$.
Proof. Denote by $`x_u`$ an element of $`D(H)`$ defined by the coset $`[u]+H_{m1}`$, where $`[u]`$ is a hard super-letter and $`m=\mathrm{Deg}(u)`$. Then the zero component $`𝐤[G]`$ and the elements $`x_u`$ generate $`D(H)`$. Formulas (47) and (48) show that $`x_ux_v=p_{uv}x_vx_u`$ and $`x_u^h=0`$ hold if $`u>v`$ and $`h`$ is the height of $`[u]`$. All the monotonic restricted $`G`$-words in $`x_u`$ are linearly independent in $`D(H)`$ since the filtration degree of every admissible word equals the sum of filtration degrees of all super-letters of that word.
Braided bigraded Hopf algebras. The above results can be easily extended to the case of $`(G,\lambda )`$-graded Hopf algebras (see, e.g., \[18, p. 206\]) and to the case of braided bigraded Hopf algebras (cf. ). These objects are not Hopf algebras in the ordinary sense, but still they are Hopf algebras in some categories. In view of Radford’s results, these algebras have embeddings in the ordinary Hopf algebras; see \[19, Thm. 1 and Prop. 2\]. The embeddings are obtained by adding the elements of $`G`$ treated as group-like elements. In this case primitive elements correspond to skew primitive ones. If a given $`(G,\lambda )`$-graded Hopf algebra $`H`$ is generated by primitive elements then the enveloping Hopf algebra is character and separated, i.e., $`H_a𝐤[G]=0`$.
All the concepts of this article can be easily extended to the case of $`(G,\lambda )`$-graded Hopf algebras by vanishing group-like elements and by replacing the bicharacter $`p(u,v)`$ with $`\lambda ^1(g_v,g_u)`$. We are now in a position to formulate the following corollarries.
COROLLARY 7. If a $`(G,\lambda )`$-graded Hopf algebra $`H`$ is generated by a set of primitive elements $`a_1,\mathrm{},a_n`$ then the set of all monotonic restricted words in hard super-letters constitutes a basis for $`H`$.
COROLLARY 8. Let $`H`$ be a $`(G,\lambda )`$-graded Hopf algebra generated by primitive elements. $`H`$ is finite-dimensional if and only if the set of all hard super-letters is finite and each hard super-letter has finite height.
COROLLARY 9. Let $`H`$ be a $`(G,\lambda )`$-graded Hopf algebra generated by primitive elements. If $`H`$ has only a finite number of hard super-letters then it is Noetherian.
COROLLARY 10. Let $`H`$ be a $`(G,\lambda )`$-graded Hopf algebra which is generated by primitive elements and has only a finite number of hard super-letters. If all these hard super-letters have infinite height then $`H`$ has no zero divisors and has a classical field of quotients.
In the corollaries above, we note, restrictions on a group are unnecessary since the group merely defines a bigrading, but the algebra in question would not contain it.
Acknowledgement. I want to thank Dr. Juan Antonio Montaraz, Director of FES-C, Dr. Suemi Rodriguez-Romo, and Virginia Lara Sagahon for providing beautiful facilities for my research work at FES-C UNAM, México. Thanks also are due to participants of Shirshov Seminar on Ring Theory (Institute of Mathematics of RAS) for interesting comments on the subject matter.
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# I. Introduction
## I. Introduction
After decades of investigations, string perturbation theory has now become a well-established old subject. It is yet a nontrivial task to uncover spacetime properties, given a first quantized worldsheet theory. Some of our previous endeavors on strings are regarded as a search for a formulation in which these spacetime properties come out in a more transparent way. String theory with constant $`B_{ij}`$ background offers an intriguing situation in which the emerging spacetime picture is given in terms of noncommutative geometry and is a focus of the recent intensive studies .
Several important steps have been taken in . In particular, the proper spacetime metric on the D$`p`$-brane (the open string metric) has been extracted from the worldsheet theory of an open string with its both ends on a D$`p`$-brane, (the $`p`$-$`p`$ open string system): the distances measured with respect to this metric are kept finite at all scales. The attendant noncommutative field theory in the zero slope limit lives on this metric together with the parameter representing noncommutativity of spacetime.
In the previous paper , we have examined the more complex $`p`$-$`p^{}`$ $`(p<p^{})`$ open string system where the both ends of the open string are on a D$`p`$-brane and on a D$`p^{}`$-brane respectively. We have obtained the open string metric and the noncommutativity parameter on the D$`p^{}`$-brane from the worldsheet two-point function. They in fact agree with those of the $`p`$-$`p`$ system. We have computed the spectrum in each case of $`(p,p^{})`$, uncovered the emergence of a large number of nearly massless states in some cases and clarified the connections among the GSO projection, branes at angles and supersymmetry. Yet a number of other properties, in particular, spacetime properties on the D$`p^{}`$-brane with D$`p`$-brane inside have remained elusive. Elucidating upon these is a major goal of the present paper.
It should be mentioned that, in the vanishing $`B_{ij}`$ background, several properties of the $`p`$-$`p^{}`$ open string system have been studied. For example, amplitudes of some scattering processes taking place on the D$`p`$-brane worldvolume (i.e. that of the lower dimensional D-brane) has been evaluated in and the conformal field theory correlation functions have been studied in . In the presence of $`B_{ij}`$ background the four point tachyon amplitudes have been given in for the $`p^{}=p+2`$ case.
In several properties of the $`0`$-$`4`$ system have been derived. The description of the system from the D$`0`$-brane has offered a new perspective to the moduli space of noncommutative instantons where the noncommutativity of the system is measured by the presence of the Fayet-Iliopoulos $`D`$ term. In contrast to this moduli space point of view, the thrust of the present paper is to uncover the spacetime properties of the system viewed from the higher dimensional D-brane, namely the D$`p^{}`$-brane. This can be accomplished by placing vertex operators on the worldvolume of the D$`p^{}`$-brane and by considering the scattering processes. This is an extension of the computation of scattering amplitudes in the $`p`$-$`p`$ open string system with and without $`B_{ij}`$ background . This line of reasonings has led us to carry out a systematic study which begins with evaluating superspace two-point functions in the relevant CFT with the spin and twist fields, proceeds to the computation of scattering amplitudes on the D$`p^{}`$-brane and ends with identifying a proper low energy noncommutative field theory in the zero slope limit.
In the next section, we begin with quantizing an open string ending on D$`p`$ and D$`p^{}`$ $`(p<p^{})`$ branes. We exploit superspace formulation, which we find extremely efficient in the calculation pursued in the subsequent sections. We evaluate two distinct two-point functions on superspace and observe the importance of the renormal ordering procedure from the $`SL(2,𝐑)`$ invariant vacuum to the oscillator vacuum. This procedure leads us to consider the difference of these two two-point functions as well. This third quantity plays an important role in section $`4`$ and will be referred to as subtracted two-point function. The conformal weights of the twist and spin fields are readily computed from the renormal ordering procedure.
In section $`3`$, we examine the tachyon vertex operator of the $`p`$-$`p^{}`$ open string and the vector vertex operator of the $`p^{}`$-$`p^{}`$ open string. We derive the on-shell conditions in terms of the open string metric, $`p+1`$ dimensional momenta of the tachyon, $`p^{}+1`$ dimensional momenta and polarizations of the massless vector and mass of the tachyon.
In section $`4`$, we consider the multiparticle tree scattering amplitudes consisting of external states of $`N2`$ vectors and two tachyons. We are able to derive an integral (Koba-Nielsen) representation of these quantities as integrals over $`N3`$ locations of the vertex operators as well as the $`N2`$ Grassmann counterparts and the $`N2`$ Grassmann sources conjugate to the polarization vectors. (See for the case of the $`p`$-$`p`$ string with vanishing $`B`$ field.) Several striking properties emerge from this representation. Among other things, we find a momentum dependent exponential factor to each external vector leg, which the subtracted two-point function is responsible for, as well as a new symplectic tensor multiplying the polarizations of the massless vector. We observe that some parts of the amplitudes are expressible in terms of the inner products of polarizations, momenta and the symplectic tensor with respect to the open string metric, while there are a host of other parts which do not permit such generic description by the inner product. We evaluate the amplitudes for the $`N=3,4`$ cases explicitly.
In section $`5`$, we examine the zero slope limit of the system. We find that the parts of the amplitudes expressible in terms of the inner product in this limit (the principal parts) can be summarized as a noncommutative field theory of a scalar field and a noncommutative $`U(1)`$ gauge field in $`p^{}+1`$ dimensions in which the scalar field decays exponentially in the $`x^{p+1},\mathrm{},x^p^{}`$-directions. They interact via the minimal coupling and a new interaction which consists of the field strength and the scalar bilinear contracted with the symplectic tensor. The contributions from the residual parts are consistent with the propagations in the $`t`$-channel of a large number of nearly massless states found in .
## II. Basic Properties of $`p`$-$`p^{}`$ System with $`B_{ij}`$ Field
In this section, we will provide the two-point functions, and the twist and the spin fields for a $`p`$-$`p^{}`$ open string in constant $`B`$ field background. This will also help us to establish our notations. We introduce two types of normal ordering: the one is taken with respect to the $`SL(2,𝐑)`$ invariant vacuum and the other is with respect to the oscillator vacuum. We will establish the relationship between these two, which will be important for our calculation in the subsequent sections.
### A. Action and boundary condition
The action of the NSR superstring in the constant $`B`$ background takes the form of
$$S=\frac{1}{2\pi }d^2\xi 𝑑\theta 𝑑\overline{\theta }\left(g_{\mu \nu }+2\pi \alpha ^{}B_{\mu \nu }\right)\overline{D}𝐗^\mu (𝐳,\overline{𝐳})D𝐗^\nu (𝐳,\overline{𝐳}),$$
(2.1)
where $`𝐳=(z,\theta )`$ and $`\overline{𝐳}=(\overline{z},\overline{\theta })`$ are the superspace coordinates on the worldsheet, $`D=\frac{}{\theta }+\theta \frac{}{z}`$ and $`\overline{D}=\frac{}{\overline{\theta }}+\overline{\theta }\frac{}{\overline{z}}`$ are the superspace covariant derivatives and $`g_{\mu \nu }`$ denotes the space-time metric which is taken to be flat. $`z=\xi ^1+i\xi ^2`$ and $`\overline{z}=\xi ^1i\xi ^2`$ are complex coordinates on the plane which are related to the strip coordinates $`(\tau ,\sigma )`$ by $`z=e^{\tau +i\sigma }`$ and $`\overline{z}=e^{\tau i\sigma }`$ respectively. The superfield $`𝐗^\mu (𝐳,\overline{𝐳})`$ is the string coordinate which is expressed in terms of component fields as
$$𝐗^\mu (𝐳,\overline{𝐳})=\sqrt{\frac{2}{\alpha ^{}}}X^\mu (z,\overline{z})+i\theta \psi ^\mu (z,\overline{z})+i\overline{\theta }\stackrel{~}{\psi }^\mu (z,\overline{z})+i\theta \overline{\theta }F^\mu (z,\overline{z}).$$
(2.2)
In terms of the component fields the action (2.1) is given by
$$S=\frac{1}{2\pi }d^2\xi \left(g_{\mu \nu }+2\pi \alpha ^{}B_{\mu \nu }\right)\left(\frac{2}{\alpha ^{}}\overline{}X^\mu X^\nu \overline{}\psi ^\mu \psi ^\nu +\stackrel{~}{\psi }^\mu \stackrel{~}{\psi }^\nu \right).$$
(2.3)
Here we have eliminated the auxiliary field $`F^\mu (z,\overline{z})`$ by using the equation of motion $`F^\mu =0`$, and the operators $``$ and $`\overline{}`$ denote $`\frac{}{z}`$ and $`\frac{}{\overline{z}}`$ respectively. Since the $`B`$ dependent terms do not couple to the worldsheet metric, the energy-momentum tensor of this system has the same form as that of the string without $`B`$ field background:
$`\mathrm{T}(𝐳)T_F(z)+\theta T_B(z){\displaystyle \frac{1}{2}}g_{\mu \nu }D𝐗^\mu D^2𝐗^\nu ,`$ (2.4)
$`\text{with}\{\begin{array}{c}T_F(z)={\displaystyle \frac{i}{2}}\sqrt{{\displaystyle \frac{2}{\alpha ^{}}}}g_{\mu \nu }\psi ^\mu X^\nu \hfill \\ T_B(z)={\displaystyle \frac{1}{\alpha ^{}}}g_{\mu \nu }X^\mu X^\nu {\displaystyle \frac{1}{2}}g_{\mu \nu }\psi ^\mu \psi ^\nu \hfill \end{array}.`$ (2.7)
Let us consider a $`p`$-$`p^{}`$ open string in the type IIA theory with $`p<p^{}`$ and $`p`$ and $`p^{}`$ being even integers. We concentrate on the situation in which a D$`p`$-brane extends in the $`(x^0,x^1,\mathrm{},x^p)`$-directions and a D$`p^{}`$-brane extends in the $`(x^0,x^1,\mathrm{},x^p^{})`$-directions with the D$`p`$-brane inside. The worldsheet of the open string corresponds to the upper half-plane: $`\mathrm{Im}z0`$ $`(0\sigma \pi )`$. The D$`p`$-brane worldvolume contains the boundary $`\sigma =0`$ while the D$`p^{}`$-brane worldvolume contains the boundary $`\sigma =\pi `$. As the space-time is flat with the metric
$$g_{\mu \nu }=\left(\begin{array}{ccccccc}1& & & & & & \\ & & & & & & \\ & \multicolumn{3}{c}{}& & & \\ & & g_{ij}& & & & \\ & & & & & & \\ & & & & & & \\ & & & & 1& & \\ & & & & & \mathrm{}& \\ & & & & & & 1\end{array}\right),g_{ij}=\epsilon \delta _{ij}(i,j=1,\mathrm{},p^{}),$$
(2.8)
we can bring $`B_{\mu \nu }`$ into a canonical form
$$B_{ij}=\frac{\epsilon }{2\pi \alpha ^{}}\left(\begin{array}{ccccc}0& b_1& & & \\ b_1& 0& & & \\ & & 0& b_2& \\ & & b_2& 0& \\ & & & & \mathrm{}\end{array}\right)(i,j=1,\mathrm{},p^{}),\text{otherwise }B_{\mu \nu }=0.$$
(2.9)
In what follows we will investigate the system on this background.
The boundary conditions for the string coordinates in the NS sector are
$`D𝐗^0\overline{D}𝐗^0|_{\sigma =0,\pi \theta =\overline{\theta }}=0,D𝐗^{p^{}+1,\mathrm{},9}+\overline{D}𝐗^{p^{}+1,\mathrm{},9}|_{\sigma =0,\pi \theta =\overline{\theta }}=0,`$
$`g_{kl}(D𝐗^l\overline{D}𝐗^l)+2\pi \alpha ^{}B_{kl}(D𝐗^l+\overline{D}𝐗^l)|_{\sigma =0,\pi \theta =\overline{\theta }}=0(k,l=1,\mathrm{},p)`$
$`D𝐗^i+\overline{D}𝐗^i|_{\sigma =0\theta =\overline{\theta }}=g_{ij}(D𝐗^j\overline{D}𝐗^j)+2\pi \alpha ^{}B_{ij}(D𝐗^j+\overline{D}𝐗^j)|_{\sigma =\pi \theta =\overline{\theta }}=0`$
$`(i,j=p+1,\mathrm{},p^{}).`$ (2.10)
For the bosonic components these conditions read
$`(\overline{})X^0|_{\sigma =0,\pi }=0,(+\overline{})X^{p^{}+1,\mathrm{},9}|_{\sigma =0,\pi }=0,`$
$`g_{kl}(\overline{})X^l+2\pi \alpha ^{}B_{kl}(+\overline{})X^l|_{\sigma =0,\pi }=0(k,l=1,\mathrm{},p),`$
$`(+\overline{})X^i|_{\sigma =0}=g_{ij}(\overline{})X^j+2\pi \alpha ^{}B_{ij}(+\overline{})X^j|_{\sigma =\pi }=0(i,j=p+1,\mathrm{},p^{}),`$ (2.11)
and for the fermionic components
$`\psi ^0\stackrel{~}{\psi }^0|_{\sigma =0,\pi }=0,\psi ^{p^{}+1,\mathrm{},9}+\stackrel{~}{\psi }^{p^{}+1,\mathrm{},9}|_{\sigma =0,\pi }=0,`$
$`g_{kl}(\psi ^l\stackrel{~}{\psi }^l)+2\pi \alpha ^{}B_{kl}(\psi ^l+\stackrel{~}{\psi }^l)|_{\sigma =0,\pi }=0(k,l=1,\mathrm{},p)`$
$`\psi ^i+\stackrel{~}{\psi }^i|_{\sigma =0}=g_{ij}(\psi ^j\stackrel{~}{\psi }^j)+2\pi \alpha ^{}B_{ij}(\psi ^j+\stackrel{~}{\psi }^j)|_{\sigma =\pi }=0(i,j=p+1,\mathrm{},p^{}).`$ (2.12)
The boundary conditions for the R fermions are
$`\psi ^0\stackrel{~}{\psi }^0|_{\sigma =0}=\psi ^0+\stackrel{~}{\psi }^0|_{\sigma =\pi }=0,\psi ^{p^{}+1,\mathrm{},9}+\stackrel{~}{\psi }^{p^{}+1,\mathrm{},9}|_{\sigma =0}=\psi ^{p^{}+1,\mathrm{},9}\stackrel{~}{\psi }^{p^{}+1,\mathrm{},9}|_{\sigma =\pi }=0,`$
$`g_{kl}(\psi ^l\stackrel{~}{\psi }^l)+2\pi \alpha ^{}B_{kl}(\psi ^l+\stackrel{~}{\psi }^l)|_{\sigma =0}=g_{kl}(\psi ^l+\stackrel{~}{\psi }^l)+2\pi \alpha ^{}B_{kl}(\psi ^l\stackrel{~}{\psi }^l)|_{\sigma =\pi }=0`$
$`(k,l=1,\mathrm{},p),`$
$`\psi ^i+\stackrel{~}{\psi }^i|_{\sigma =0}=g_{ij}(\psi ^j+\stackrel{~}{\psi }^j)+2\pi \alpha ^{}B_{ij}(\psi ^j\stackrel{~}{\psi }^j)|_{\sigma =\pi }=0(i,j=p+1,\mathrm{},p^{}).`$ (2.13)
It should be noted that these boundary conditions are written on the complex plane, while the boundary conditions in are written on the strip.
### B. Quantization of a $`p`$-$`p^{}`$ string
In this subsection we will give the mode expansions of the string coordinates of a $`p`$-$`p^{}`$ open string and the commutation relations among their modes.
Let us first consider the $`x^0`$-direction. In this direction the string coordinate obeys the Neumann boundary condition on both ends. The coordinate $`X^0(z,\overline{z})`$ is expanded as
$$X^0(z,\overline{z})=x^0i\alpha ^{}p^0\mathrm{ln}(z\overline{z})+i\sqrt{\frac{\alpha ^{}}{2}}\underset{m0}{}\frac{\alpha _m^0}{m}\left(z^m+\overline{z}^m\right).$$
(2.14)
The modes satisfy the following commutation relations,
$$[x^0,p^0]=ig^{00},[x^0,x^0]=[p^0,p^0]=0;[\alpha _m^0,\alpha _n^0]=g^{00}m\delta _{m+n}.$$
(2.15)
The mode expansions of the NS fermions $`\psi ^0`$ and $`\stackrel{~}{\psi }^0`$ become
$$\psi ^0(z)=\underset{r𝐙+1/2}{}b_r^0z^{r\frac{1}{2}},\stackrel{~}{\psi }^0(\overline{z})=\underset{r𝐙+1/2}{}b_r^0\overline{z}^{r\frac{1}{2}}.$$
(2.16)
The oscillators satisfy
$$\{b_r^0,b_s^0\}=g^{00}\delta _{r+s}.$$
(2.17)
In the R sector, we have
$$\psi ^0(z)=\underset{m𝐙}{}d_m^0z^{m\frac{1}{2}},\stackrel{~}{\psi }^0(\overline{z})=\underset{m𝐙}{}d_m^0\overline{z}^{m\frac{1}{2}},$$
(2.18)
and
$$\{d_m^0,d_n^0\}=g^{00}\delta _{m+n}.$$
(2.19)
Next we consider the $`x^i`$-directions, $`i=1,\mathrm{},p`$. In these directions, the boundary conditions on the string coordinates are the same as those on the string coordinates along the D$`p`$-branes in the D$`p`$-D$`p`$ system with $`B`$ field. This implies that the mode expansions and the commutation relations among the oscillators take the same form as those in the D$`p`$-D$`p`$ system in the $`B`$ field background. The mode expansions of the bosonic coordinates $`X^i(z,\overline{z})`$ are
$`X^i(z,\overline{z})=x^ii\alpha ^{}\left[g^1(g2\pi \alpha ^{}B)\right]_{}^{i}{}_{j}{}^{}p^j\mathrm{ln}zi\alpha ^{}\left[g^1(g+2\pi \alpha ^{}B)\right]_{}^{i}{}_{j}{}^{}p^j\mathrm{ln}\overline{z}`$ (2.20)
$`+i\sqrt{{\displaystyle \frac{\alpha ^{}}{2}}}{\displaystyle \underset{m0}{}}\left[\left(g^1(g2\pi \alpha ^{}B)\right)_{}^{i}{}_{j}{}^{}z^m+\left(g^1(g+2\pi \alpha ^{}B)\right)_{}^{i}{}_{j}{}^{}\overline{z}^m\right]{\displaystyle \frac{\alpha _m^j}{m}}.`$
Under this expansion, the commutation relations among the oscillators are given by
$$[x^i,x^j]=i\theta ^{ij},[p^i,p^j]=0,[x^i,p^j]=iG^{ij};[\alpha _m^i,\alpha _n^j]=G^{ij}m\delta _{m+n},$$
(2.21)
where $`\theta ^{ij}`$ is the noncommutativity parameter and $`G^{ij}`$ is the inverse of the open string metric $`G_{ij}`$ defined respectively as
$$\theta ^{ij}=(2\pi \alpha ^{})^2\left(\frac{1}{g+2\pi \alpha ^{}B}B\frac{1}{g2\pi \alpha ^{}B}\right)^{ij},G^{ij}=\left(\frac{1}{g+2\pi \alpha ^{}B}g\frac{1}{g2\pi \alpha ^{}B}\right)^{ij}.$$
(2.22)
We will further generalize $`G^{ij}`$ to include the time direction. This is denoted by $`G^{\mu \nu }`$ and will simplify our formulas in the subsequent sections. Taking eqs. (2.8) and (2.9) into account, we obtain
$$G^{\sigma \rho }=\left(\begin{array}{cccc}g^{00}& & & \\ & & & \\ & & & \\ & & G^{ij}& \end{array}\right)=\left(\begin{array}{cccccc}1& & & & & \\ & & & & & \\ & \frac{1}{\epsilon (1+b_1^2)}& 0& & & \\ & 0& \frac{1}{\epsilon (1+b_1^2)}& & & \\ & & & \frac{1}{\epsilon (1+b_2^2)}& 0& \\ & & & 0& \frac{1}{\epsilon (1+b_2^2)}& \\ & & & & & \mathrm{}\end{array}\right).$$
(2.23)
We obtain for the NS fermions
$$\{\begin{array}{c}\psi ^i(z)=\underset{r𝐙+1/2}{}\left[g^1(g2\pi \alpha ^{}B)\right]_{}^{i}{}_{j}{}^{}b_r^jz^{r\frac{1}{2}}\hfill \\ \stackrel{~}{\psi }^i(\overline{z})=\underset{r𝐙+1/2}{}\left[g^1(g+2\pi \alpha ^{}B)\right]_{}^{i}{}_{j}{}^{}b_r^j\overline{z}^{r\frac{1}{2}}\hfill \end{array},\text{with }\{b_r^i,b_s^j\}=G^{ij}\delta _{r+s},$$
(2.24)
and for the R fermions
$$\{\begin{array}{c}\psi ^i(z)=\underset{m𝐙}{}\left[g^1(g2\pi \alpha ^{}B)\right]_{}^{i}{}_{j}{}^{}d_m^jz^{m\frac{1}{2}}\hfill \\ \stackrel{~}{\psi }^i(\overline{z})=\underset{m𝐙}{}\left[g^1(g+2\pi \alpha ^{}B)\right]_{}^{i}{}_{j}{}^{}d_m^j\overline{z}^{m\frac{1}{2}}\hfill \end{array},\text{with }\{d_m^i,d_n^j\}=G^{ij}\delta _{m+n}.$$
(2.25)
Finally we investigate the $`x^i`$-directions $`(i=p+1,\mathrm{},p^{})`$. We complexify the string coordinates $`𝐗^i(𝐳,\overline{𝐳})`$ in these directions as
$`𝐙^I(𝐳,\overline{𝐳})`$ $`=`$ $`𝐗^{2I1}(𝐳,\overline{𝐳})+i𝐗^{2I}(𝐳,\overline{𝐳})=\sqrt{{\displaystyle \frac{2}{\alpha ^{}}}}Z^I(z,\overline{z})+i\theta \mathrm{\Psi }^I(z)+i\overline{\theta }\stackrel{~}{\mathrm{\Psi }}^I(\overline{z}),`$
$`\overline{𝐙}^{\overline{I}}(𝐳,\overline{𝐳})`$ $`=`$ $`𝐗^{2I1}(𝐳,\overline{𝐳})i𝐗^{2I}(𝐳,\overline{𝐳})=\sqrt{{\displaystyle \frac{2}{\alpha ^{}}}}\overline{Z}^{\overline{I}}(z,\overline{z})+i\theta \overline{\mathrm{\Psi }}^{\overline{I}}(z)+i\overline{\theta }\stackrel{~}{\overline{\mathrm{\Psi }}}^{\overline{I}}(\overline{z}),`$ (2.26)
where $`I,\overline{I}=\frac{p+2}{2},\mathrm{},\frac{p^{}}{2}`$ and we have eliminated the auxiliary field $`F^i`$. From the boundary conditions eq. (2.11) and equations of motion, we find that the mode expansions of $`Z^I`$ and $`\overline{Z}^{\overline{I}}`$ are given by
$`Z^I(z,\overline{z})=i\sqrt{{\displaystyle \frac{\alpha ^{}}{2}}}{\displaystyle \underset{n𝐙}{}}{\displaystyle \frac{\alpha _{n\nu _I}^I}{n\nu _I}}\left(z^{(n\nu _I)}\overline{z}^{(n\nu _I)}\right),`$
$`\overline{Z}^{\overline{I}}(z,\overline{z})=i\sqrt{{\displaystyle \frac{\alpha ^{}}{2}}}{\displaystyle \underset{m𝐙}{}}{\displaystyle \frac{\overline{\alpha }_{m+\nu _I}^{\overline{I}}}{m+\nu _I}}\left(z^{(m+\nu _I)}\overline{z}^{(m+\nu _I)}\right),`$ (2.27)
where $`\nu _I`$ are defined by
$$e^{2\pi i\nu _I}=\frac{1+ib_I}{1ib_I},0<\nu _I<1.$$
(2.28)
Now we can introduce the open string metric $`G^{IJ}`$, $`G^{\overline{I}\overline{J}}`$, $`G^{I\overline{J}}`$ and $`G^{\overline{I}J}`$ concerning the $`x^{p+1},\mathrm{},x^p^{}`$ directions,
$$G^{IJ}=G^{\overline{I}\overline{J}}=0,G^{I\overline{J}}=G^{\overline{J}I}=\frac{2}{\epsilon (1+b_I^2)}\delta ^{I\overline{J}}.$$
(2.29)
Similarly, the boundary conditions eq. (2.12), eq. (2.13) and equations of motion lead us to the mode expansions of the NS-fermions,
$`\mathrm{\Psi }^I(z)={\displaystyle \underset{r𝐙+1/2}{}}b_{r\nu _I}^Iz^{(r\nu _I)\frac{1}{2}},`$ $`\stackrel{~}{\mathrm{\Psi }}^I(\overline{z})={\displaystyle \underset{r𝐙+1/2}{}}b_{r\nu _I}^I\overline{z}^{(r\nu _I)\frac{1}{2}},`$
$`\overline{\mathrm{\Psi }}^{\overline{I}}(z)={\displaystyle \underset{s𝐙+1/2}{}}\overline{b}_{s+\nu _I}^{\overline{I}}z^{(s+\nu _I)\frac{1}{2}},`$ $`\stackrel{~}{\overline{\mathrm{\Psi }}}^{\overline{I}}(\overline{z})={\displaystyle \underset{s𝐙+1/2}{}}\overline{b}_{s+\nu _I}^{\overline{I}}\overline{z}^{(s+\nu _I)\frac{1}{2}},`$ (2.30)
and that of R-fermions,
$`\mathrm{\Psi }^I(z)={\displaystyle \underset{n𝐙}{}}d_{n\nu _I}^Iz^{(n\nu _I)\frac{1}{2}},`$ $`\stackrel{~}{\mathrm{\Psi }}^I(\overline{z})={\displaystyle \underset{n𝐙}{}}d_{n\nu _I}^I\overline{z}^{(n\nu _I)\frac{1}{2}},`$
$`\overline{\mathrm{\Psi }}^{\overline{I}}(z)={\displaystyle \underset{m𝐙}{}}\overline{d}_{m+\nu _I}^{\overline{I}}z^{(m+\nu _I)\frac{1}{2}},`$ $`\stackrel{~}{\overline{\mathrm{\Psi }}}^{\overline{I}}(\overline{z})={\displaystyle \underset{m𝐙}{}}\overline{d}_{m+\nu _I}^{\overline{I}}\overline{z}^{(m+\nu _I)\frac{1}{2}}.`$ (2.31)
The commutation relations are
$`[\alpha _{n\nu _I}^I,\overline{\alpha }_{m+\nu _J}^{\overline{J}}]={\displaystyle \frac{2}{\epsilon }}\delta ^{I\overline{J}}(n\nu _I)\delta _{n+m},`$
$`\{b_{r\nu _I}^I,\overline{b}_{s+\nu _J}^{\overline{J}}\}={\displaystyle \frac{2}{\epsilon }}\delta ^{I\overline{J}}\delta _{r+s},\{d_{n\nu _I}^I,\overline{d}_{m+\nu _J}^{\overline{J}}\}={\displaystyle \frac{2}{\epsilon }}\delta ^{I\overline{J}}\delta _{n+m}.`$ (2.32)
As for the $`x^{p^{}+1},\mathrm{},x^9`$-directions, the string coordinates obey the Dirichlet boundary condition. Our analysis in the remaining part of this paper does not involve these directions.
### C. Two-point functions on superspace
In this subsection we construct two-point functions of the $`p`$-$`p^{}`$ open string coordinates on superspace. For this purpose, we begin by defining the oscillator vacuum of the system.
As shown in the last subsection, the mode expansions in the $`x^0`$ and the $`x^i`$-directions $`(i=1,\mathrm{},p)`$ are similar to those of the usual open strings obeying Neumann boundary conditions in the sense that the bosons and the R fermions have integral moding oscillators and the NS fermions have half-integral moding ones. Therefore we can define the vacuum in these directions in the same way as the usual open string. In the NS sector the vacuum $`|0`$ is defined by
$$\{\begin{array}{ccc}\alpha _m^0|0=0,& \alpha _m^i|0=0,& \text{for }m0\hfill \\ b_r^0|0=0,& b_r^i|0=0,& \text{for }r\frac{1}{2}\hfill \end{array},$$
(2.33)
where $`\alpha _0^\mu =\sqrt{2\alpha ^{}}p^\mu `$ ($`\mu =0,1,\mathrm{},p`$). In the R sector the vacuum $`|S^\alpha `$ belongs to the spinor representation of the $`SO(p,1)`$ group. Now that the commutation relations and the vacuum $`|0`$ are determined, we can evaluate the two-point functions of the string coordinates in these directions ,
$`𝐆^{00}(𝐳_1,\overline{𝐳}_1|𝐳_2,\overline{𝐳}_2)0|𝐗^0(𝐳_1,\overline{𝐳}_1)𝐗^0(𝐳_2,\overline{𝐳}_2)|0`$
$`=g^{00}\left[\mathrm{ln}(z_1z_2\theta _1\theta _2)(\overline{z}_1\overline{z}_2\overline{\theta }_1\overline{\theta }_2)+\mathrm{ln}(z_1\overline{z}_2\theta _1\overline{\theta }_2)(\overline{z}_1z_2\overline{\theta }_1\theta _2)\right],`$
$`𝐆^{ij}(𝐳_1,\overline{𝐳}_1|𝐳_2,\overline{𝐳}_2)0|𝐗^i(𝐳_1,\overline{𝐳}_1)𝐗^j(𝐳_2,\overline{𝐳}_2)|0`$ (2.34)
$`=g^{ij}\mathrm{ln}(z_1z_2\theta _1\theta _2)(\overline{z}_1\overline{z}_2\overline{\theta }_1\overline{\theta }_2)+(g^{ij}2G^{ij})\mathrm{ln}(z_1\overline{z}_2\theta _1\overline{\theta }_2)(\overline{z}_1z_2\overline{\theta }_1\theta _2)`$
$`2{\displaystyle \frac{\theta ^{ij}}{2\pi \alpha ^{}}}\mathrm{ln}{\displaystyle \frac{z_1\overline{z}_2\theta _1\overline{\theta }_2}{\overline{z}_1z_2\overline{\theta }_1\theta _2}}2D^{ij},`$
where $``$ stands for the radial ordering. $`D^{ij}`$ are the contributions from the zero modes $`x^i`$ and these will be fixed conveniently as is done in . When we restrict these two-point functions onto the D$`p^{}`$-brane worldvolume, i.e. the worldsheet boundary characterized by $`z=e^{\tau +i\pi }=e^\tau `$ and $`\theta =\overline{\theta }`$, they become
$`𝐆^{00}(e^{\tau _1},\theta _1|e^{\tau _2},\theta _2)`$ $``$ $`𝐆^{00}(𝐳_1,\overline{𝐳}_1|𝐳_2,\overline{𝐳}_2)|_{\sigma =\pi ,\theta =\overline{\theta }}=2g^{00}\mathrm{ln}(e^{\tau _1}e^{\tau _2}+\theta _1\theta _2)^2,`$
$`𝐆^{ij}(e^{\tau _1},\theta _1|e^{\tau _2},\theta _2)`$ $``$ $`𝐆^{ij}(𝐳_1,\overline{𝐳}_1|𝐳_2,\overline{𝐳}_2)|_{\sigma =\pi ,\theta =\overline{\theta }}`$ (2.35)
$`=`$ $`2G^{ij}\mathrm{ln}(e^{\tau _1}e^{\tau _2}+\theta _1\theta _2)^2{\displaystyle \frac{i}{\alpha ^{}}}\theta ^{ij}ϵ(\tau _1\tau _2),`$
where $`ϵ(x)`$ is the sign function.
In the $`x^i`$-directions $`(i=p+1,\mathrm{},p^{})`$, the situation is more complex as the string coordinates are expanded in non-integer power of $`z`$ and $`\overline{z}`$. We define the oscillator vacuum $`|\sigma `$ for the bosonic sector so that this should be annihilated by the negative energy modes:
$$|\sigma =\underset{I}{}|\sigma _I\text{with}\{\begin{array}{cc}\alpha _{n\nu _I}^I|\sigma _I=0\hfill & n>\nu _I\hfill \\ \overline{\alpha }_{m+\nu _I}^{\overline{I}}|\sigma _I=0\hfill & m>\nu _I\hfill \end{array}.$$
(2.36)
For the fermions in the NS sector we define the oscillator vacuum $`|s`$ by<sup>1</sup><sup>1</sup>1The vacuum $`|s`$ is defined in order that the negative energy and the positive energy modes for $`0<\nu _I<1/2`$ should be the annihilation and the creation modes respectively. The energy carried by the lowest creation mode $`\overline{b}_{\frac{1}{2}+\nu _I}^{\overline{I}}`$ becomes negative when $`\nu _I`$ becomes greater than $`1/2`$. We could define another oscillator vacuum $`|\stackrel{~}{s}`$ by
$$|\stackrel{~}{s}=\underset{I}{}|\stackrel{~}{s}_I\text{with}\{\begin{array}{cc}b_{r\nu _I}^I|\stackrel{~}{s}_I=0\hfill & r\frac{3}{2}\hfill \\ \overline{b}_{s+\nu _I}^{\overline{I}}|\stackrel{~}{s}_I=0\hfill & s\frac{1}{2}\hfill \end{array},$$
which makes the negative energy and the positive energy modes for $`1/2<\nu _I<1`$ to be the annihilation and the creation modes respectively.
$$|s=\underset{I}{}|s_I\text{with}\{\begin{array}{cc}b_{r\nu _I}^I|s_I=0\hfill & r\frac{1}{2}\hfill \\ \overline{b}_{s+\nu _I}^{\overline{I}}|s_I=0\hfill & s\frac{1}{2}\hfill \end{array},$$
(2.37)
and for the R sector we define the oscillator vacuum $`|S`$ by
$$|S=\underset{I}{}|S_I\text{with}\{\begin{array}{cc}d_{n\nu _I}^I|S_I=0\hfill & n>\nu _I\hfill \\ \overline{d}_{m+\nu _I}^{\overline{I}}|S_I=0\hfill & m>\nu _I\hfill \end{array}.$$
(2.38)
By using the commutation relations and the defining relations of the vacua, we can calculate the two-point functions<sup>2</sup><sup>2</sup>2If we adopted $`|\stackrel{~}{s}`$ as the vacuum of the NS fermions instead of $`|s`$, the two-point function of the supercoordinates $`𝐙^I`$ and $`\overline{𝐙}^{\overline{I}}`$ would be $`\stackrel{~}{𝓖}^{I\overline{J}}(𝐳_1,\overline{𝐳}_1|𝐳_2,\overline{𝐳}_2)\sigma ,\stackrel{~}{s}\left|𝐙^I(𝐳_1,\overline{𝐳}_1)\overline{𝐙}^{\overline{J}}(𝐳_2,\overline{𝐳}_2)\right|\sigma ,\stackrel{~}{s}`$ $`=\mathrm{\Theta }(|z_1||z_2|){\displaystyle \frac{2\delta ^{I\overline{J}}}{\epsilon }}[(1\nu _I;{\displaystyle \frac{z_2}{z_1\theta _1\theta _2}})+(1\nu _I;{\displaystyle \frac{\overline{z}_2}{\overline{z}_1\overline{\theta }_1\overline{\theta }_2}})`$ $`(1\nu _I;{\displaystyle \frac{\overline{z}_2}{z_1\theta _1\overline{\theta }_2}})(1\nu _I;{\displaystyle \frac{z_2}{\overline{z}_1\overline{\theta }_1\theta _2}})]`$ $`+\mathrm{\Theta }(|z_2||z_1|){\displaystyle \frac{2\delta ^{I\overline{J}}}{\epsilon }}\left[(\nu _I;{\displaystyle \frac{z_1\theta _1\theta _2}{z_2}})+(\nu _I;{\displaystyle \frac{\overline{z}_1\overline{\theta }_1\overline{\theta }_2}{\overline{z}_2}})(\nu _I;{\displaystyle \frac{z_1\theta _1\overline{\theta }_2}{\overline{z}_2}})(\nu _I;{\displaystyle \frac{\overline{z}_1\overline{\theta }_1\theta _2}{z_2}})\right].`$ :
$`𝓖^{I\overline{J}}(𝐳_1,\overline{𝐳}_1|𝐳_2,\overline{𝐳}_2)\sigma ,s\left|𝐙^I(𝐳_1,\overline{𝐳}_1)\overline{𝐙}^{\overline{J}}(𝐳_2,\overline{𝐳}_2)\right|\sigma ,s`$
$`=\mathrm{\Theta }(|z_1||z_2|){\displaystyle \frac{2\delta ^{I\overline{J}}}{\epsilon }}[(1\nu _I;{\displaystyle \frac{z_2+\theta _1\theta _2}{z_1}})+(1\nu _I;{\displaystyle \frac{\overline{z}_2+\overline{\theta }_1\overline{\theta _2}}{\overline{z_1}}})`$
$`(1\nu _I;{\displaystyle \frac{\overline{z}_2+\theta _1\overline{\theta _2}}{z_1}})(1\nu _I;{\displaystyle \frac{z_2+\overline{\theta }_1\theta _2}{\overline{z_1}}})]`$
$`+\mathrm{\Theta }(|z_2||z_1|){\displaystyle \frac{2\delta ^{I\overline{J}}}{\epsilon }}[(\nu _I;{\displaystyle \frac{z_1}{z_2+\theta _1\theta _2}})+(\nu _I;{\displaystyle \frac{\overline{z_1}}{\overline{z}_2+\overline{\theta }_1\overline{\theta _2}}})`$
$`(\nu _I;{\displaystyle \frac{z_1}{\overline{z}_2+\theta _1\overline{\theta _2}}})(\nu _I;{\displaystyle \frac{\overline{z_1}}{z_2+\overline{\theta }_1\theta _2}})],`$ (2.39)
where $`\mathrm{\Theta }(x)`$ is the step function, $`(\nu ;z)`$ is defined as
$$(\nu ;z)=\frac{z^\nu }{\nu }F(1,\nu ;1+\nu ;z)=\underset{n=0}{\overset{\mathrm{}}{}}\frac{1}{n+\nu }z^{n+\nu },$$
(2.40)
and $`F(a,b;c;z)`$ is the hypergeometric function. When we restrict this two-point function onto the worldsheet boundary on the D$`p^{}`$-brane worldvolume, this becomes
$`𝓖^{I\overline{J}}(e^{\tau _1},\theta _1|e^{\tau _2},\theta _2)𝓖^{I\overline{J}}(𝐳_1,\overline{𝐳}_1|𝐳_2,\overline{𝐳}_2)|_{\sigma =\pi ,\theta =\overline{\theta }}`$ (2.41)
$`=4G^{I\overline{J}}\left[\mathrm{\Theta }(\tau _1\tau _2)(1\nu _I;{\displaystyle \frac{e^{\tau _2}\theta _1\theta _2}{e^{\tau _1}}})+\mathrm{\Theta }(\tau _2\tau _1)(\nu _I;{\displaystyle \frac{e^{\tau _1}}{e^{\tau _2}\theta _1\theta _2}})\right].`$
Now we would like to study the two-point function eq. (2.41) more closely. Let us recast the right hand side of eq. (2.41) into
$`4G^{I\overline{J}}\left[{\displaystyle \frac{1}{2}}\left\{(1\nu _I;{\displaystyle \frac{e^{\tau _2}}{e^{\tau _1}}})+(\nu _I;{\displaystyle \frac{e^{\tau _1}}{e^{\tau _2}}})\right\}\theta _1\theta _2{\displaystyle \frac{\left(\frac{e^{\tau _1}}{e^{\tau _2}}\right)^{\nu _I}}{e^{\tau _1}e^{\tau _2}}}\right]`$
$`+ϵ(\tau _1\tau _2){\displaystyle \frac{4}{\epsilon }}{\displaystyle \frac{\delta ^{I\overline{J}}}{1+b_I^2}}\left\{(1\nu _I;{\displaystyle \frac{e^{\tau _2}}{e^{\tau _1}}})(\nu _I;{\displaystyle \frac{e^{\tau _1}}{e^{\tau _2}}})\right\}.`$ (2.42)
As noncommutativity of the D$`p^{}`$-brane worldvolume originates from the term proportional to the sign function $`ϵ(\tau _1\tau _2)`$ in the above equation , we will also refer to this term as noncommutativity term in what follows. Here it should be noted that by using the relations,
$`{\displaystyle \frac{d}{dz}}\left[z^{c1}F(a,b;c;z)\right]`$ $`=`$ $`(c1)z^{c2}F(a,b;c1;z),`$
$`F(a,b;b;z)`$ $`=`$ $`(1z)^a.`$ (2.43)
we can obtain
$$\frac{d}{dz}\left[(1\nu _I;\frac{1}{z})(\nu _I;z)\right]=0.$$
(2.44)
This implies that the noncommutativity term in eq. (2.42) is constant. The value of this constant can be fixed by evaluating the noncommutativity term at a certain point on the real axis, such as $`\frac{e^{\tau _1}}{e^{\tau _2}}=1`$. By using the hypergeometric series, we find that
$$(1\nu _I;1)(\nu _I;1)=\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}\frac{1}{n+\nu _I}=\pi \mathrm{cot}\left(\pi \nu _I\right)=\pi b_I.$$
(2.45)
Thus we find that the noncommutativity term becomes
$$\frac{4}{\epsilon }\frac{\delta ^{I\overline{J}}}{1+b_I^2}\left\{(1\nu _I;\frac{e^{\tau _2}}{e^{\tau _1}})(\nu _I;\frac{e^{\tau _1}}{e^{\tau _2}})\right\}=\frac{4}{\epsilon }\frac{\delta ^{I\overline{J}}}{1+b_I^2}\pi b_I.$$
(2.46)
When we rewrite the complex string coordinates $`𝐙^I`$ and $`\overline{𝐙}^{\overline{I}}`$ into the real one $`𝐗^i`$ $`(i=p+1,\mathrm{},p^{})`$, this noncommutativity term takes the same form as that in eq. (2.35). This means that, as is pointed out in , the noncommutativity on the D-brane worldvolume in the $`p`$-$`p^{}`$ system is the same as that in the $`p`$-$`p`$ system . From eq. (2.46), we conclude that
$$𝓖^{I\overline{J}}(e^{\tau _1},\theta _1|e^{\tau _2},\theta _2)=4G^{I\overline{J}}(\nu _I;\frac{e^{\tau _1}}{e^{\tau _2}\theta _1\theta _2})+ϵ(\tau _1\tau _2)\frac{4}{\epsilon }\frac{\delta ^{I\overline{J}}}{1+b_I^2}\pi b_I,$$
(2.47)
where $`(\nu ;z)`$ is defined by using the hypergeometric series as
$$(\nu ;z)=\{\begin{array}{cc}(1\nu _I;\frac{1}{z})\frac{\pi }{2}b_I=\underset{n=0}{\overset{\mathrm{}}{}}\frac{z^{n1+\nu _I}}{n+1\nu _I}\frac{\pi }{2}b_I\hfill & \text{ for }|z|>1\hfill \\ (\nu _I;z)+\frac{\pi }{2}b_I=\underset{n=0}{\overset{\mathrm{}}{}}\frac{z^{n+\nu _I}}{n+\nu _I}+\frac{\pi }{2}b_I\hfill & \text{ for }|z|<1\hfill \end{array}.$$
(2.48)
The two infinite series in the above defining relation should be analytically continued to each other.
In Appendix we give another derivation of the two-point function (2.47) from eq. (2.41).
### D. Twist field and spin field
Here we would like to make further consideration on the oscillator vacuum consisting of the bosonic sector $`|\sigma `$ and the fermionic sector $`|s`$ . As is explained in the last subsection, the primary fields $`D𝐙^I`$, $`\overline{D}𝐙^I`$, $`D\overline{𝐙}^{\overline{I}}`$ and $`\overline{D}\overline{𝐙}^{\overline{I}}`$ defined on the upper half plane are expanded in non-integer powers of $`z`$ and $`\overline{z}`$. It follows that when we extend the defining region of these fields to the whole complex plane through the doubling trick these fields become multi-valued functions on the whole plane. For example, when the primary fields $`Z^I(z)`$, $`\overline{}Z^I(\overline{z})`$, $`\overline{Z}^{\overline{I}}(z)`$ and $`\overline{}\overline{Z}^{\overline{I}}(\overline{z})`$ on the whole plane are transported once around the origin, they gain phase factors:
$`Z^I(e^{2\pi i}z)=e^{2\pi i\nu _I}Z^I(z),`$ $`\overline{}Z^I(e^{2\pi i}\overline{z})=e^{2\pi i\nu _I}\overline{}Z^I(\overline{z})`$
$`\overline{Z}^{\overline{I}}(e^{2\pi i}z)=e^{2\pi i\nu _I}\overline{Z}^{\overline{I}}(z),`$ $`\overline{}\overline{Z}^{\overline{I}}(e^{2\pi i}\overline{z})=e^{2\pi i\nu _I}\overline{}\overline{Z}^{\overline{I}}(\overline{z}).`$ (2.49)
This implies that a twist field $`\sigma _I^+(\xi ^1)`$ and an anti-twist field $`\sigma _I^{}(\xi ^1)`$, both of which are mutually non-local with respect to $`Z^I`$ and $`\overline{Z}^I`$, are located at the origin and at infinity on the plane respectively. They create a branch cut between themselves. The twist field $`\sigma ^+`$ serves as a boundary changing operator from the $`p^{}`$-brane to the $`p`$-brane and the anti-twist field $`\sigma ^{}`$ acts in the opposite way . The incoming vacuum $`|\sigma _I`$ defined in eq. (2.36) should be interpreted as being excited from the $`SL(2,𝐑)`$-invariant vacuum $`|0`$ by the twist field $`\sigma _I^+`$:
$$|\sigma _I=\underset{\xi ^10}{lim}\sigma _I^+(\xi ^1)|0.$$
(2.50)
In the same way, the outgoing vacuum $`\sigma _I|`$ should be regarded as
$$\sigma _I|=\underset{\stackrel{~}{\xi }^10}{lim}\left(\frac{1}{\stackrel{~}{\xi }^1}\right)^{2h_{\sigma _I}}0|\sigma _I^{}\left(\frac{1}{\stackrel{~}{\xi }^1}\right),$$
(2.51)
where $`h_{\sigma _I}`$ denotes the weight of the (anti-) twist field. We will later explain that $`h_{\sigma _I}=\frac{1}{2}\nu _I(1\nu _I)`$ . We can read off the OPE’s of $`Z^I`$ and $`\overline{Z}^I`$ with $`\sigma _I^\pm `$ from eq. (2.36):
$`\{\begin{array}{cc}Z^I(z)\sigma _J^+(0)\delta _J^Iz^{(1\nu _I)}\tau _I^+(0),\hfill & \overline{}Z^I(\overline{z})\sigma _J^+(0)\delta _J^I\overline{z}^{(1\nu _I)}\stackrel{~}{\tau }_I^+(0),\hfill \\ \overline{Z}^{\overline{I}}(z)\sigma _J^+(0)\delta _J^{\overline{I}}z^{\nu _I}\tau _I^+(0),\hfill & \overline{}\overline{Z}^{\overline{I}}(\overline{z})\sigma _J^+(0)\delta _J^{\overline{I}}\overline{z}^{\nu _I}\stackrel{~}{\tau }_I^+(0),\hfill \end{array}`$ (2.54)
$`\{\begin{array}{cc}Z^I(z)\sigma _J^{}(0)\delta _J^Iz^{\nu _I}\tau _I^{}(0),\hfill & \overline{}Z^I(\overline{z})\sigma _J^{}(0)\delta _J^I\overline{z}^{\nu _I}\stackrel{~}{\tau }_I^{}(0),\hfill \\ \overline{Z}^{\overline{I}}(z)\sigma _J^{}(0)\delta _J^{\overline{I}}z^{(1\nu _I)}\tau _I^{}(0),\hfill & \overline{}\overline{Z}^{\overline{I}}(\overline{z})\sigma _J^{}(0)\delta _{\overline{J}}^I\overline{z}^{(1\nu _I)}\stackrel{~}{\tau }_I^{}(0),\hfill \end{array}`$ (2.57)
where $`\tau `$ ’s are excited twist fields.
Similar argument holds for the fermionic coordinates. In the NS sector, the spin fields $`s_I^+`$ and $`s_I^{}`$ are mutually non-local with respect to the fermions. They are located at the origin and at the infinity on the worldsheet respectively. They exchange the boundary conditions corresponding to the $`p`$-brane and those to the $`p^{}`$-brane by generating a branch cut between themselves. The incoming vacuum $`|s_I`$ and the outgoing vacuum $`s_I|`$ should be regarded as being excited from the $`SL(2,𝐑)`$-invariant vacuum by spin fields:
$$|s_I=\underset{\xi ^10}{lim}s_I^+(\xi ^1)|0,s_I|=\underset{\stackrel{~}{\xi }^10}{lim}\left(\frac{1}{\stackrel{~}{\xi }^1}\right)^{2h_{s_I}}0|s_I^{}\left(\frac{1}{\stackrel{~}{\xi }^1}\right),$$
(2.58)
where $`h_{s_I}`$ is the weight of the spin fields which will be found to be $`h_{s_I}=\frac{1}{2}\nu _{I}^{}{}_{}{}^{2}`$. The defining relation eq. (2.37) yields the OPE’s,
$`\{\begin{array}{cc}\mathrm{\Psi }^I(z)s_J^+(0)\delta _J^Iz^{+\nu _I}t_I^+(0),\hfill & \stackrel{~}{\mathrm{\Psi }}^I(\overline{z})s_J^+(0)\delta _J^I\overline{z}^{+\nu _I}\stackrel{~}{t}_I^+(0),\hfill \\ \overline{\mathrm{\Psi }}^{\overline{I}}(z)s_J^+(0)\delta _J^{\overline{I}}z^{\nu _I}t_I^+(0),\hfill & \stackrel{~}{\overline{\mathrm{\Psi }}}^{\overline{I}}(\overline{z})s_J^+(0)\delta _J^{\overline{I}}\overline{z}^{\nu _I}\stackrel{~}{t}^+(0),\hfill \end{array}`$ (2.61)
$`\{\begin{array}{cc}\mathrm{\Psi }^I(z)s_J^{}(0)\delta _J^Iz^{\nu _I}t_I^{}(0),\hfill & \stackrel{~}{\mathrm{\Psi }}^I(\overline{z})s_J^{}(0)\delta _J^I\overline{z}^{\nu _I}\stackrel{~}{t}_I^{}(0)\hfill \\ \overline{\mathrm{\Psi }}^{\overline{I}}(z)s_J^{}(0)\delta _J^{\overline{I}}z^{+\nu _I}t_I^{}(0),\hfill & \stackrel{~}{\overline{\mathrm{\Psi }}}^{\overline{I}}(\overline{z})s_J^{}(0)\delta _J^{\overline{I}}\overline{z}^{+\nu _I}\stackrel{~}{t}_I^{}(0),\hfill \end{array},`$ (2.64)
In the fermionic sector, the bosonization simplifies the treatment of the spin fields. While we will not invoke this treatment here, we include it here for the sake of completeness. We write
$$\mathrm{\Psi }^I(z)\sqrt{\frac{2}{\epsilon }}e^{iH^I(z)},\overline{\mathrm{\Psi }}^{\overline{I}}(z)\sqrt{\frac{2}{\epsilon }}e^{iH^I(z)},$$
(2.65)
where $`H^I(z)`$ are free bosons normalized as $`H^I(z)H^J(w)\delta ^{IJ}\mathrm{ln}(zw)`$. Then the OPE (2.64) tells us that the spin fields $`s_I^\pm (z)`$ should be bosonized as <sup>3</sup><sup>3</sup>3We are also able to perform the same analysis on the other incoming and the outgoing vacua, $`|\stackrel{~}{s}_I`$ and $`\stackrel{~}{s}|`$, in the NS sector defined in the footnote 1. These states are excited from the $`SL(2,𝐑)`$-invariant vacuum by $`\stackrel{~}{s}_I^+(z)`$ and $`\stackrel{~}{s}_I^{}(z)`$ respectively which are bosonized as
$$\stackrel{~}{s}_I^+(z)e^{+i(1+\nu _I)H^I(z)},\stackrel{~}{s}_I^{}(z)e^{i(1+\nu _I)H^I(z)}.$$
$$s_I^+(z)e^{+i\nu _IH^I(z)},s_I^{}(z)e^{i\nu _IH^I(z)}.$$
(2.66)
We can repeat the same analysis in the R sector. We find that the incoming vacuum eq. (2.38) and the outgoing vacuum are excited from the $`SL(2,𝐑)`$-invariant vacuum by the spin fields $`S_I^+(z)`$ and $`S^{}(z)`$ respectively. They are bosonized as
$$S_I^+(z)e^{i\left(\frac{1}{2}+\nu _I\right)H^I(z)},S_I^{}(z)e^{i\left(\frac{1}{2}+\nu _I\right)H^I(z)}.$$
(2.67)
### E. Subtracted two-point functions and weights of twist and spin fields
In the $`x^i`$-directions $`(i=p+1,\mathrm{},p^{})`$, we have two types of vacuum: the one is the $`SL(2,𝐑)`$-invariant vacuum and the other is the oscillator vacuum. We can define the normal ordering corresponding to each one. We will use the symbols $`::`$ and $`\stackrel{}{}`$$`\stackrel{}{}`$ to denote the normal orderings with respect to the $`SL(2,𝐑)`$-invariant vacuum and the oscillator vacuum respectively.
In the other directions we have a single type of vacuum, namely $`SL(2,𝐑)`$-invariant vacuum. We have $`::`$-normal ordered product only. For free bosons and free fermions in these directions, it is defined by a subtraction:
$$:𝐗^\mu (𝐳_1,\overline{𝐳}_1)𝐗^\nu (𝐳_2,\overline{𝐳}_2):=𝐗^\mu (𝐳_1,\overline{𝐳}_1)𝐗^\nu (𝐳_2,\overline{𝐳}_2)𝐆^{\mu \nu }(𝐳_1,\overline{𝐳}_1|𝐳_2,\overline{𝐳}_2),$$
(2.68)
for $`\mu ,\nu =0,1,\mathrm{},p`$. Here $`𝐆^{\mu \nu }(𝐳_1,\overline{𝐳}_1|𝐳_2,\overline{𝐳}_2)`$ is the two-point function defined in eq. (2.34).
In the same way we can define $`\stackrel{}{}`$$`\stackrel{}{}`$-normal ordered product for the free fields in the $`x^i`$-directions $`(i=p+1,\mathrm{},p^{})`$ as
$$\stackrel{}{}𝐙^I(𝐳_1,\overline{𝐳}_1)\overline{𝐙}^{\overline{J}}(𝐳_2,\overline{𝐳}_2)\stackrel{}{}=𝐙^I(𝐳_1,\overline{𝐳}_1)\overline{𝐙}^{\overline{J}}(𝐳_2,\overline{𝐳}_2)𝓖^{I\overline{J}}(𝐳_1,\overline{𝐳}_1|𝐳_2,\overline{𝐳}_2),$$
(2.69)
for $`I,\overline{J}=\frac{p+2}{2},\mathrm{},\frac{p^{}}{2}`$. Here $`𝓖^{I\overline{J}}(𝐳_1,\overline{𝐳}_1|𝐳_2,\overline{𝐳}_2)`$ is the two-point function defined in eq. (2.39). In addition to the $`\stackrel{}{}`$$`\stackrel{}{}`$-normal ordered product, we would like to define the $`::`$-normal ordered product for these free fields $`𝐙^I`$ and $`\overline{𝐙}^{\overline{I}}`$. In order to apply the definition (2.68) directly to these fields, we have to evaluate their two-point functions on the $`SL(2,𝐑)`$-invariant vacuum. As is pointed out, the twist and the spin fields play the role of boundary changing operators. This implies that if we delete the twist and the spin fields and thus remove the cut generated by them, the boundary conditions imposed on the positive real axis of the $`z`$-plane is expected to be identical to those imposed on the negative real axis. This leads us to conclude that the two point function of $`𝐙^I`$ and $`\overline{𝐙}^{\overline{I}}`$ evaluated on the $`SL(2,𝐑)`$-invariant vacuum takes the same form as that of the $`p^{}`$-$`p^{}`$ system with $`B`$ field:
$`𝐆^{I\overline{J}}(𝐳_1,\overline{𝐳}_1|𝐳_2,\overline{𝐳}_2)0|𝐙^I(𝐳_1,\overline{𝐳}_1)\overline{𝐙}^{\overline{J}}(𝐳_2,\overline{𝐳}_2)|0`$ (2.70)
$`={\displaystyle \frac{2\delta ^{I\overline{J}}}{\epsilon }}[\mathrm{ln}(z_1z_2\theta _1\theta _2)(\overline{z}_1\overline{z}_2\overline{\theta }_1\overline{\theta }_2)+\mathrm{ln}(z_1\overline{z}_2\theta _1\overline{\theta }_2)(\overline{z}_1z_2\overline{\theta }_1\theta _2)`$
$`{\displaystyle \frac{2}{1+b_I^2}}\mathrm{ln}(z_1\overline{z}_2\theta _1\overline{\theta }_2)(\overline{z}_1z_2\overline{\theta }_1\theta _2)2i{\displaystyle \frac{b_I}{1+b_I^2}}\mathrm{ln}{\displaystyle \frac{z_1\overline{z}_2\theta _1\overline{\theta }_2}{\overline{z}_1z_2\overline{\theta }_1\theta _2}}]+\text{D-term}`$
Computing this two point function at the boundary $`\sigma =\pi `$ and $`\theta =\overline{\theta }`$, we obtain
$`𝐆^{I\overline{J}}(e^{\tau _1},\theta _1|e^{\tau _2},\theta _2)𝐆^{I\overline{J}}(𝐳_1,\overline{𝐳}_1|𝐳_2,\overline{𝐳}_2)|_{\sigma =\pi ,\theta =\overline{\theta }}`$ (2.71)
$`={\displaystyle \frac{4\delta ^{I\overline{J}}}{\epsilon (1+b_I^2)}}\mathrm{ln}(e^{\tau _1}+e^{\tau _2}\theta _1\theta _2)^2+{\displaystyle \frac{4\pi \delta ^{I\overline{J}}b_I}{\epsilon (1+b_I^2)}}ϵ(\tau _1\tau _2).`$
From the fact that the normal ordered product is defined by a subtraction, we can readily find that for an arbitrary functional $`𝓞`$ of the free fields $`𝐙^I`$ and $`\overline{𝐙}^{\overline{I}}`$ the normal ordering is formally expressed as (see e.g. )
$`:𝓞:=\mathrm{exp}\left({\displaystyle d^2𝐳_1d^2𝐳_2𝐆^{I\overline{J}}(𝐳_1,\overline{𝐳}_1|𝐳_2,\overline{𝐳}_2)\frac{\delta }{\delta 𝐙^I(𝐳_1,\overline{𝐳}_1)}\frac{\delta }{\delta \overline{𝐙}^{\overline{J}}(𝐳_2,\overline{𝐳}_2)}}\right)𝓞,`$
$`\stackrel{}{}𝓞\stackrel{}{}=\mathrm{exp}({\displaystyle }d^2𝐳_\mathrm{𝟏}d^2𝐳_2𝓖^{I\overline{J}}(𝐳_1,\overline{𝐳}_1|𝐳_2,\overline{𝐳}_2){\displaystyle \frac{\delta }{\delta 𝐙^I(𝐳_1,\overline{𝐳}_1)}}{\displaystyle \frac{\delta }{\delta \overline{𝐙}^{\overline{J}}(𝐳_2,\overline{𝐳}_2)}})𝓞,`$ (2.72)
where $`d^2𝐳`$ is defined as $`d^2𝐳=d^2\xi d\theta d\overline{\theta }`$. From these general definitions of the normal orderings, we can read off the formula of the reordering between them:
$$:𝓞:=\mathrm{exp}\left(d^2𝐳_\mathrm{𝟏}d^2𝐳_2𝓖_{\mathrm{sub}}^{}{}_{}{}^{I\overline{J}}(𝐳_1,\overline{𝐳}_1|𝐳_2,\overline{𝐳}_2)\frac{\delta }{\delta 𝐙^I(𝐳_1,\overline{𝐳}_1)}\frac{\delta }{\delta \overline{𝐙}^{\overline{J}}(𝐳_2,\overline{𝐳}_2)}\right)\stackrel{}{}𝓞\stackrel{}{}.$$
(2.73)
Here $`𝓖_{\mathrm{sub}}^{}{}_{}{}^{I\overline{J}}(𝐳_1,\overline{𝐳}_1|𝐳_2,\overline{𝐳}_2)`$ is a subtracted two-point function defined as
$`𝓖_{\mathrm{sub}}^{}{}_{}{}^{I\overline{J}}(𝐳_1,\overline{𝐳}_1|𝐳_2,\overline{𝐳}_2)\sigma ,s|:𝐙^I(𝐳_1,\overline{𝐳}_1)\overline{𝐙}^{\overline{J}}(𝐳_2,\overline{𝐳}_2):|\sigma ,s`$
$`=𝓖^{I\overline{J}}(𝐳_1,\overline{𝐳}_1|𝐳_2,\overline{𝐳}_2)𝐆^{I\overline{J}}(𝐳_1,\overline{𝐳}_1|𝐳_2,\overline{𝐳}_2).`$ (2.74)
Let us compute the subtracted two-point function $`𝓖_{\mathrm{sub}}^{}{}_{}{}^{I\overline{J}}`$ at the worldsheet boundary $`\sigma =\pi `$ and $`\theta =\overline{\theta }`$. From eqs. (2.41) and (2.71), we find that
$`𝓖_{\mathrm{sub}}^{}{}_{}{}^{I\overline{J}}(e^{\tau _1},\theta _1|e^{\tau _2},\theta _2)\sigma ,s|:𝐙^I(e^{\tau _1},\theta _1)\overline{𝐙}^{\overline{J}}(e^{\tau _2},\theta _2):|\sigma ,s`$
$`={\displaystyle \frac{8\delta ^{I\overline{J}}}{\epsilon (1+b_I^2)}}[\gamma {\displaystyle \frac{𝝍(\nu _I)+𝝍(1\nu _I)}{2}}+{\displaystyle \frac{1}{2}}\mathrm{ln}e^{\tau _1+\tau _2}`$
$`+\mathrm{\Theta }(\tau _1\tau _2)\left\{{\displaystyle \frac{1}{2}}\mathrm{ln}{\displaystyle \frac{e^{\tau _1}}{e^{\tau _2}}}\left({\displaystyle \frac{e^{\tau _2}}{e^{\tau _1}}}\right)^{1\nu _I}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\left(1\nu _I\right)_n}{n!}}{\displaystyle \underset{m=0}{\overset{n1}{}}}{\displaystyle \frac{\nu _I}{(m+1)(m+1\nu _I)}}\left(1{\displaystyle \frac{e^{\tau _2}}{e^{\tau _1}}}\right)^n\right\}`$
$`+\mathrm{\Theta }(\tau _2\tau _1)\{{\displaystyle \frac{1}{2}}\mathrm{ln}{\displaystyle \frac{e^{\tau _2}}{e^{\tau _1}}}\left({\displaystyle \frac{e^{\tau _1}}{e^{\tau _2}}}\right)^{\nu _I}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\left(\nu _I\right)_n}{n!}}{\displaystyle \underset{m=0}{\overset{n1}{}}}{\displaystyle \frac{1\nu _I}{(m+1)(m+\nu _I)}}(1{\displaystyle \frac{e^{\tau _1}}{e^{\tau _2}}})^n\}]`$
$`+\theta _1\theta _2{\displaystyle \frac{8\delta ^{I\overline{J}}}{\epsilon (1+b_I^2)}}\left\{\left({\displaystyle \frac{e^{\tau _1}}{e^{\tau _2}}}\right)^{\nu _I}1\right\}{\displaystyle \frac{1}{e^{\tau _1}+e^{\tau _2}}},`$ (2.75)
where $`\gamma `$ is Euler’s constant, $`𝝍(w)`$ denotes the digamma function defined as $`𝝍(w)=\frac{d}{dw}\mathrm{ln}\mathrm{\Gamma }(w)`$ and $`(a)_n\frac{\mathrm{\Gamma }(a+n)}{\mathrm{\Gamma }(a)}`$. Here we have used the formulas (A.1) and (A.5).
Let us compute the conformal weights of the twist and spin fields. This helps us appreciate these subtracted two-point functions better. The relevant part of the energy-momentum tensor in this computation consists of the terms depending on the string coordinates in the $`x^i`$-directions, $`i=p+1,\mathrm{},p^{}`$: $`T_B^{(Z,\mathrm{\Psi })}(z)=T_B^{(Z,\overline{Z})}(z)+T_B^{(\mathrm{\Psi },\overline{\mathrm{\Psi }})}(z)`$ with
$$T_B^{(Z,\overline{Z})}(z)=\frac{\epsilon \delta _{I\overline{J}}}{\alpha ^{}}:Z^I\overline{Z}^{\overline{J}}(z):,T_B^{(\mathrm{\Psi },\overline{\mathrm{\Psi }})}(z)=\frac{\epsilon \delta _{I\overline{J}}}{4}[:\mathrm{\Psi }^I\overline{\mathrm{\Psi }}^{\overline{J}}(z)::\mathrm{\Psi }^I\overline{\mathrm{\Psi }}^{\overline{J}}(z):].$$
(2.76)
Using the subtracted two point function, we obtain
$`\sigma _I|T_B(z)|\sigma _I={\displaystyle \frac{\epsilon }{\alpha ^{}}}\underset{wz}{lim}\sigma _I|:Z^I(w)\overline{Z}^{\overline{I}}(z):|\sigma _I`$ (2.77)
$`=\underset{wz}{lim}\left[\left({\displaystyle \frac{w}{z}}\right)^{\nu _I1}{\displaystyle \frac{1}{wz}}\left({\displaystyle \frac{1\nu _I}{z}}+{\displaystyle \frac{1}{wz}}\right){\displaystyle \frac{1}{(wz)^2}}\right]={\displaystyle \frac{1}{z^2}}{\displaystyle \frac{\nu _I(1\nu _I)}{2}}.`$
This implies that the twist fields $`\sigma _I^\pm `$ have the weights $`h_{\sigma _I}=\frac{1}{2}\nu _I(1\nu _I)`$. In the same way, we obtain
$`s_I|T_B(z)|s_I={\displaystyle \frac{\epsilon }{4}}\underset{wz}{lim}s_I|:\left(\mathrm{\Psi }^I(w)\overline{\mathrm{\Psi }}^{\overline{I}}(z)\mathrm{\Psi }^I(w)\overline{\mathrm{\Psi }}^{\overline{I}}(z)\right):|s_I`$
$`=\underset{wz}{lim}\left[\left({\displaystyle \frac{w}{z}}\right)^{\nu _I}{\displaystyle \frac{1}{wz}}\left\{{\displaystyle \frac{\nu _I}{2}}\left({\displaystyle \frac{1}{w}}+{\displaystyle \frac{1}{z}}\right){\displaystyle \frac{1}{wz}}\right\}+{\displaystyle \frac{1}{(wz)^2}}\right]={\displaystyle \frac{1}{z^2}}{\displaystyle \frac{\nu _{I}^{}{}_{}{}^{2}}{2}}.`$ (2.78)
From this equation we can read <sup>4</sup><sup>4</sup>4This result can also be obtained from the fact that $`T_B^{(\mathrm{\Psi },\overline{\mathrm{\Psi }})}(z)`$ is bosonized as $`T_B^{(\mathrm{\Psi },\overline{\mathrm{\Psi }})}(z)\frac{1}{2}\delta _{IJ}:H^IH^J(z):`$ . that the spin fields $`s_I^\pm `$ have the weights $`h_{s_I}=\frac{1}{2}\nu _{I}^{}{}_{}{}^{2}`$.
## III. Vertex operators
Let us pay some attention to vertex operators of our system before we start calculating scattering amplitudes. We focus on two types of vertex operators. The one is the tachyon vertex operators, and the other is the massless vector vertex operators. These are the relevant ones in order for us to find out the spacetime processes occurring on the D$`p^{}`$-brane worldvolume.
Here we make a comment on the GSO projection. In this paper we take the GSO projection in the NS sector so that the oscillator vacuum corresponding to the tachyon vertex operator survives. It follows that the GSO projection adopted in this paper is not always the same as that in our previous work : in the cases of $`p^{}=p+2`$, $`p+6`$ they are opposite to each other, while in the cases of $`p^{}=p+4`$, $`p+8`$ they are the same. This is attributed to the sign convention of the D$`p`$-brane charge: in the cases of $`p^{}=p+2`$, $`p+6`$ the D$`p`$-branes in this paper should be referred to as anti-D$`p`$-branes in the convention of .
### A. Tachyon vertex operator of $`p`$-$`p^{}`$ string
First let us investigate vertex operators of the $`p`$-$`p^{}`$ open string which contain the twist and the spin fields. We will focus on the vertex operator that corresponds to the ground state in the NS sector of this open string. This vertex operator is seen for instance in :
$$𝐕_T^\pm (\xi ^1,\theta )=V_T^{\pm (1)}(\xi ^1)+\theta V_T^{\pm (0)}(\xi ^1)=𝓣^\pm (\xi ^1,\theta ):\mathrm{exp}\left(i\sqrt{\frac{\alpha ^{}}{2}}\underset{\mu =0}{\overset{p}{}}k_\mu 𝐗^\mu (\xi ^1,\theta )\right):,$$
(3.1)
where $`𝐗^\mu (\xi ^1,\theta )`$ is the boundary value of the superfield $`𝐗^\mu (𝐳,\overline{𝐳})`$: $`𝐗^\mu (\xi ^1,\theta )𝐗^\mu (𝐳,\overline{𝐳})|_{z=\overline{z}=\xi ^1,\theta =\overline{\theta }}`$, and $`𝓣^\pm (\xi ^1,\theta )`$ is the superfield whose lowest component $`𝒯_0^\pm (\xi ^1)`$ consists of the twist and the spin fields,
$$𝒯_0^\pm (\xi ^1)=\underset{I}{}\sigma _I^\pm (\xi ^1)s_I^\pm (\xi ^1).$$
(3.2)
The upper component is obtained by applying the supercurrent $`T_F(z)`$ to $`𝒯_0^\pm (\xi ^1)`$ . In eq. (3.1), $`V_T^{\pm (0)}(\xi ^1)`$ denotes the $`0`$-picture vertex and $`V_T^{\pm (1)}(\xi ^1)`$ is the matter field contribution to the $`(1)`$-picture vertex
$$𝒱_T^{\pm (1)}(\xi ^1)=e^\varphi (\xi ^1)V_T^{\pm (1)}(\xi ^1)=e^\varphi \underset{I}{}\sigma _I^\pm s_I^\pm :\mathrm{exp}\left(i\underset{\mu =0}{\overset{p}{}}k_\mu X^\mu \right):,$$
(3.3)
where $`e^\varphi `$ comes from the $`\beta \gamma `$-ghost sector.
It is worth noting that space-time momentum $`k_\mu `$ of the tachyon vertex operator eq. (3.1) is not $`(p^{}+1)`$ dimensional but $`(p+1)`$ dimensional. This is because the $`p`$-$`p^{}`$ string coordinates in $`x^{p+1},\mathrm{},x^p^{}`$-directions do not possess zero-modes and thus the momenta in these directions are not defined. This implies that an initial/final tachyon field corresponding to this vertex operator is frozen in these space-time directions.
Let us study the physical state conditions for this vertex operator. If the $`(1)`$-picture vertex satisfies the physical state condition, the $`0`$-picture vertex is automatically physical because of the worldsheet supersymmetry. We will therefore concentrate on the $`(1)`$-picture vertex. Through the operator-state mapping this vertex operator corresponds to the state<sup>5</sup><sup>5</sup>5Here we ignore the ghost sector.,
$$|V_T^{(1)}\underset{\xi ^10}{lim}V_T^{(1)}(\xi ^1)|0=|0;k_\mu |\sigma ,s,$$
(3.4)
where the state $`|0;k_\mu `$ is defined as $`|0;k_\mu =\mathrm{exp}\left(i{\displaystyle \underset{\mu =0}{\overset{p}{}}}k_\mu x^\mu \right)|0`$. Here $`x^\mu `$ are the zero modes of the bosonic coordinates. It is evident that $`|V_{T(1)}`$ is a primary state. We just require that this state should have weight $`\frac{1}{2}`$. From the calculation in the last subsection, we find that the state $`|\sigma ,s`$ has a weight
$$h\left[|\sigma ,s\right]=\underset{I}{}\left(\frac{\nu _I(1\nu _I)}{2}+\frac{\nu _{I}^{}{}_{}{}^{2}}{2}\right)=\underset{I}{}\frac{\nu _I}{2}.$$
(3.5)
Substituting the mode expansions eqs. (2.20) and (2.24) into the defining relation eq. (2.7), we see that the terms in $`T_B(z)`$ which depend on the string coordinates in $`x^\mu `$-directions $`(\mu =0,1,\mathrm{},p)`$ are
$`T_B^{(X,\psi )}(z){\displaystyle \underset{\mu =0}{\overset{p}{}}}\left({\displaystyle \frac{1}{\alpha ^{}}}g_{\mu \nu }X^\mu X^\nu (z)+{\displaystyle \frac{1}{2}}g_{\mu \nu }\psi ^\mu \psi ^\nu (z)\right){\displaystyle \underset{m𝐙}{}}L_mz^{m2},`$
$`\text{with}L_m={\displaystyle \frac{1}{2}}{\displaystyle \underset{n𝐙}{}}{\displaystyle \underset{\sigma ,\rho =0}{\overset{p}{}}}G_{\sigma \rho }:\alpha _{mn}^\sigma \alpha _n^\rho :+{\displaystyle \frac{1}{4}}{\displaystyle \underset{r𝐙+1/2}{}}(2rm){\displaystyle \underset{\sigma ,\rho =0}{\overset{p}{}}}G_{\sigma \rho }:b_{mr}^\sigma b_r^\rho :,`$ (3.6)
where $`\alpha _0^\mu =\sqrt{2\alpha ^{}}p^\mu `$. Here $`G_{\sigma \rho }`$ is the open string metric including time direction, its inverse is given in eq. (2.23). This yields
$$L_0|0;k_\mu =\alpha ^{}\underset{\sigma ,\rho =0}{\overset{p}{}}G^{\sigma \rho }k_\sigma k_\rho |0;k_\mu .$$
(3.7)
Gathering all results obtained above, we conclude that the weight of the state $`|V_T^{(1)}`$ is $`L_0={\displaystyle \underset{I}{}}{\displaystyle \frac{1}{2}}\nu _I+\alpha ^{}{\displaystyle \underset{\sigma ,\rho =0}{\overset{p}{}}}G^{\sigma \rho }k_\sigma k_\rho `$. Thus the on-shell condition $`L_0=\frac{1}{2}`$ requires that the mass squared of this state should be
$$\alpha ^{}m_T^2\alpha ^{}\underset{\sigma ,\rho =0}{\overset{p}{}}G^{\sigma \rho }k_\sigma k_\rho =\frac{1}{2}\left(1\underset{I}{}\nu _I\right).$$
(3.8)
### B. Massless vector vertex operator of $`p^{}`$-$`p^{}`$ string
The vector emission vertex operator takes the form of
$$𝐕_{\mathrm{vec}}(\xi ^1,\theta )V_{\mathrm{vec}}^{(1)}(\xi ^1)+\theta V_{\mathrm{vec}}^{(0)}(\xi ^1)\frac{i}{2}\underset{\mu =0}{\overset{p^{}}{}}:\zeta _\mu (k)\dot{𝐗}^\mu (\xi ^1,\theta )\mathrm{exp}\left(i\sqrt{\frac{\alpha ^{}}{2}}\underset{\rho =0}{\overset{p^{}}{}}k_\rho 𝐗^\rho (\xi ^1,\theta )\right):,$$
(3.9)
where $`\zeta _\mu (k)`$ denotes the polarization vector and $`\dot{𝐗}^\mu (\xi ^1,\theta )(D+\overline{D})𝐗^\mu (𝐳,\overline{𝐳})|_{z=\overline{z}=\xi ^1,\theta =\overline{\theta }}`$. The operator $`V_{\mathrm{vec}}^{(0)}(\xi ^1)`$ is the $`0`$-picture vertex and $`V_{\mathrm{vec}}^{(1)}(\xi ^1)`$ is the matter field contribution to the $`(1)`$-picture vertex $`𝒱_{\mathrm{vec}}^{(1)}(\xi ^1)=e^\varphi V_{\mathrm{vec}}^{(1)}(\xi ^1)`$. Their explicit forms are
$`V_{\mathrm{vec}}^{(1)}(\xi ^1)={\displaystyle \frac{1}{2}}{\displaystyle \underset{\mu =0}{\overset{p^{}}{}}}:\zeta _\mu (k)\left(\psi ^\mu +\stackrel{~}{\psi }^\mu \right)\mathrm{exp}\left(i{\displaystyle \underset{\rho =0}{\overset{p^{}}{}}}k_\rho X^\rho \right):,`$ (3.10)
$`V_{\mathrm{vec}}^{(0)}(\xi ^1)={\displaystyle \frac{1}{\sqrt{2\alpha ^{}}}}{\displaystyle \underset{\mu =0}{\overset{p^{}}{}}}:\zeta _\mu (k)\left[i\dot{X}^\mu +{\displaystyle \frac{\alpha ^{}}{2}}\left({\displaystyle \underset{\rho =0}{\overset{p^{}}{}}}k_\rho \left(\psi ^\rho +\stackrel{~}{\psi }^\rho \right)\right)\left(\psi ^\mu +\stackrel{~}{\psi }^\mu \right)\right]\mathrm{exp}\left(i{\displaystyle \underset{\lambda =0}{\overset{p^{}}{}}}k_\lambda X^\lambda \right):,`$
where $`\dot{X}^\mu (\xi ^1)`$ is defined as $`\dot{X}^\mu (\xi ^1)=(+\overline{})X^\mu (z,\overline{z})|_{z=\overline{z}=\xi ^1}`$.
The string coordinates of a $`p^{}`$-$`p^{}`$ open string in the $`x^i`$-directions $`(i=1,\mathrm{},p^{})`$ obey the same boundary conditions as the $`x^i`$-directions with $`i=1,\mathrm{},p`$ of the $`p`$-$`p^{}`$ string. They have the same mode expansions and the commutators among the oscillating modes. Therefore the vector emission vertex operator in each picture corresponds to the respective state
$`|V_{\mathrm{vec}}^{(1)}\underset{\xi ^10}{lim}V_{\mathrm{vec}}^{(1)}(\xi ^1)|0={\displaystyle \underset{\mu =0}{\overset{p^{}}{}}}\zeta _\mu (k)b_{\frac{1}{2}}^\mu |0;k_\rho ^{},`$
$`|V_{\mathrm{vec}}^{(0)}\underset{\xi ^10}{lim}V_{\mathrm{vec}}^{(0)}(\xi ^1)|0={\displaystyle \underset{\mu =0}{\overset{p^{}}{}}}\zeta _\mu (k)\left(\alpha _1^\mu +\sqrt{2\alpha ^{}}\left({\displaystyle \underset{\lambda =0}{\overset{p^{}}{}}}k_\lambda b_{\frac{1}{2}}^\lambda \right)b_{\frac{1}{2}}^\mu \right)|0;k_\rho ^{},`$ (3.11)
where $`|0;k_\rho ^{}`$ is defined as $`|0;k_\rho ^{}=\mathrm{exp}\left({\displaystyle \underset{\rho =0}{\overset{p^{}}{}}}k_\rho x^\rho \right)|0`$. Let us consider the physical state conditions on these states. The relevant part of the energy-momentum tensor for this analysis is eq. (3.6) with $`p`$ being replaced by $`p^{}`$. This yields
$`L_0|V_{\mathrm{vec}}^{(0)}=\left[\alpha ^{}{\displaystyle \underset{\sigma ,\rho =0}{\overset{p^{}}{}}}G^{\sigma \rho }k_\sigma k_\rho +1\right]|V_{\mathrm{vec}}^{(0)},`$
$`L_1|V_{\mathrm{vec}}^{(0)}=\sqrt{2\alpha ^{}}{\displaystyle \underset{\sigma ,\rho =0}{\overset{p^{}}{}}}G^{\sigma \rho }k_\sigma \zeta _\rho (k)|0;k_\rho ^{}.`$ (3.12)
From these relations we find that the physical state conditions, $`L_0=1`$ and $`L_1=0`$, require that
$$\alpha ^{}m_{\mathrm{vec}}^2\alpha ^{}\underset{\sigma ,\rho =0}{\overset{p^{}}{}}G^{\sigma \rho }k_\sigma k_\rho =0,\underset{\sigma ,\rho =0}{\overset{p^{}}{}}G^{\sigma \rho }k_\sigma \zeta _\rho (k)=0.$$
(3.13)
We will write the vector emission vertex operator $`𝐕_{\mathrm{vec}}(\xi ^1,\theta )`$ in an exponential form ,
$$𝐕_{\mathrm{vec}}(\xi ^1,\theta )=d\eta :\mathrm{exp}\left[i\underset{\mu =0}{\overset{p^{}}{}}\{\sqrt{\frac{\alpha ^{}}{2}}k_\mu +\eta \zeta _\mu (k)\frac{1}{2}(D+\overline{D})\}𝐗^\mu (𝐳,\overline{𝐳})\right]:|_{\genfrac{}{}{0pt}{}{z=\overline{z}=\xi ^1}{\theta =\overline{\theta }}},$$
(3.14)
by introducing a Grassmann parameter $`\eta `$.
We will write the part of vertex operator eq. (3.14) which depends on the string coordinates in the $`x^i`$-directions $`(i=p+1,\mathrm{},p^{})`$ as
$$:\mathrm{exp}[\underset{I=\frac{p+2}{2}}{\overset{p^{}/2}{}}𝐄_I(\zeta ,k,\eta )𝐙^I(𝐳,\overline{𝐳})+\underset{\overline{J}=\frac{p+2}{2}}{\overset{p^{}/2}{}}\overline{𝐄}_{\overline{J}}(\zeta ,k,\eta )\overline{𝐙}^{\overline{J}}(𝐳,\overline{𝐳})]:|_{z=\overline{z}=\xi ^1,\theta =\overline{\theta }},$$
(3.15)
using the complex variables. Here $`𝐄_I(\zeta ,k,\eta )`$ and $`\overline{𝐄}_{\overline{I}}(\zeta ,k,\eta )`$ are differential operators on the superspace defined as
$`𝐄_I(\zeta ,k,\eta )=i\sqrt{{\displaystyle \frac{\alpha ^{}}{2}}}\kappa _I+i\eta e_I(k){\displaystyle \frac{1}{2}}\left(D+\overline{D}\right),`$
$`\overline{𝐄}_{\overline{I}}(\zeta ,k,\eta )=i\sqrt{{\displaystyle \frac{\alpha ^{}}{2}}}\overline{\kappa }_{\overline{I}}+i\eta \overline{e}_{\overline{I}}(k){\displaystyle \frac{1}{2}}\left(D+\overline{D}\right),`$ (3.16)
with
$`\kappa _I={\displaystyle \frac{1}{2}}\left(k_{2I1}ik_{2I}\right),`$ $`\overline{\kappa }_{\overline{I}}={\displaystyle \frac{1}{2}}\left(k_{2I1}+ik_{2I}\right);`$
$`e_I(k)={\displaystyle \frac{1}{2}}\left(\zeta _{2I1}(k)i\zeta _{2I}(k)\right),`$ $`\overline{e}_{\overline{I}}(k)={\displaystyle \frac{1}{2}}\left(\zeta _{2I1}(k)+i\zeta _{2I}(k)\right).`$ (3.17)
## IV. Scattering Amplitudes
We would like to find out a proper description of the physical processes taking place on the worldvolume of the D$`p^{}`$-brane with the D$`p`$-brane inside ($`p<p^{}`$) in the case where the $`B`$ field is nonvanishing. In this section, we will consider multiparticle tree scattering amplitudes consisting of the external states of $`N2`$ vectors obtained from the mode of the $`p^{}`$-$`p^{}`$ open string and two tachyons from the mode of the $`p`$-$`p^{}`$ open string. The spacetime picture of this string scattering process is depicted in Fig.1. That this process is possible is easy to see once we draw a spacetime diagram and map the end points of the open strings onto a circle. See Fig.2.
Open string tree amplitudes in general are obtained by placing vertex operators on the boundary of the upper half plane, namely, the real axis, integrating over the positions of the vertex operators and dividing by the volume of the (super)conformal killing vectors. To obtain the amplitudes of our concern, we first locate each of the two kinds of the tachyon vertex operators $`𝐕_T^+(\xi ,\theta )`$, $`𝐕_T^{}(\xi ,\theta )`$ discussed in the last section at $`\xi =\xi _1`$ and at $`\xi =\xi _2`$ respectively. A cut is generated on the interval between these two locations as $`𝐕_T^+`$ and $`𝐕_T^{}`$ contain the twist field and the anti-twist field respectively. The worldvolume of the D$`p^{}`$-brane contains this interval on which we place the $`N2`$ vector emission vertex operators $`𝐕_{\mathrm{vec}}(\xi ,\theta )`$ of the $`p^{}`$-$`p^{}`$ open string. In what follows we will obtain the integral (Koba-Nielsen) representation of the amplitudes. The explicit expressions for the $`N=3,4`$ cases that we obtain will be exploited to determine the form of the low energy effective field theory in the subsequent section.
In manifestly supersymmetric formulation on superspace, the $`N`$ point tree amplitude in question reads
$`{\displaystyle \frac{c}{V_{\mathrm{SCKV}}}}{\displaystyle \underset{a=1}{\overset{N}{}}d\xi _ad\theta _a0|𝐕_T^+(\xi _1,\theta _1;k_{1\mu })𝐕_T^{}(\xi _2,\theta _2;k_{2\mu })\underset{c=3}{\overset{N}{}}𝐕_{\mathrm{vec}}(\xi _c,\theta _c;k_{c\mu },\zeta _{c\mu })|0},`$ (4.1)
where $`V_{\mathrm{SCKV}}`$ denotes the volume of the isometry group generated by the superconformal Killing vectors, namely, the graded extension of the $`SL(2,𝐑)`$ group. We have denoted by $`c`$ the overall constant which does not concern us in this paper. In eq. (4.1), the domain of $`\xi _a`$ integrations is not restricted except that $`\xi _3,\xi _4,\mathrm{},\xi _N`$ are located on the cut created on the interval between $`\xi _1`$ and $`\xi _2`$. This domain falls into a sum of the $`(N2)!`$ regions. In each region, an ordering among $`\xi _3,\xi _4,\mathrm{},\xi _N`$ is specified and integrals over each region give a contribution corresponding to a respective open string (dual) diagram<sup>6</sup><sup>6</sup>6 Recall that we have $`(N1)!`$ open string (dual) diagrams in the case of the $`N`$ point amplitude of a $`p`$-$`p`$ open string.. We will evaluate the contribution from the region $`\xi _2<\xi _3<\xi _4<\mathrm{}<\xi _N<\xi _1`$ and this is denoted by $`A_N`$. In most cases below, we will not write the region of integrations explicitly.
Eq. (4.1) is invariant under the graded $`SL(2,𝐑)`$ transformations after the physical state conditions are invoked at each vertex operator. For actual evaluation of the amplitude, we first set $`\theta _1=\theta _2=0`$ to fix the odd elements of the transformations. We then fix the even elements by giving fixed values to three of the locations of the vertex operators. These locations are chosen as $`\xi _1,\xi _2,`$ and $`\xi _3`$. This amounts to factoring out the following volume element from the integration,
$$d^3F(\xi _1,\xi _2,\xi _3)d\theta _1d\theta _2(\xi _1\xi _2),$$
(4.2)
where
$$d^3F(\xi _1,\xi _2,\xi _3)\frac{d\xi _1d\xi _2d\xi _3}{(\xi _1\xi _2)(\xi _2\xi _3)(\xi _3\xi _1)}.$$
(4.3)
Having done this, we obtain
$$A_N=c\frac{{\displaystyle \underset{a=1}{\overset{N}{}}}d\xi _a}{d^3F(\xi _1,\xi _2,\xi _3)}\frac{{\displaystyle \underset{c=3}{\overset{N}{}}}d\theta _c}{\xi _1\xi _2}0|V_T^{+(1)}(\xi _1;k_{1\mu })V_T^{(1)}(\xi _2;k_{2\mu })\underset{c^{}=3}{\overset{N}{}}𝐕_{\mathrm{vec}}(\xi _c^{},\theta _c^{};k_{c^{}\mu },\zeta _{c^{}\mu })|0.$$
(4.4)
The component corresponding to the $`(1)`$-picture has been selected at each of the tachyon vertex operators. Let us choose $`\xi _1=0`$, $`\xi _2=\mathrm{}`$ and $`\xi _3=1`$, so that the negative real axis becomes the worldsheet boundary ending on the D$`p^{}`$-brane. Introducing positive real variables $`x_a\xi _a\left(=e^{\tau _a}\right)>0`$ and adopting eq. (3.14) for the vector emission vertex operators, we find
$`A_N=c{\displaystyle \frac{{\displaystyle \underset{a=1}{\overset{N}{}}}dx_a}{d^3F(x_1,x_2,x_3)}\frac{{\displaystyle \underset{a^{}=3}{\overset{N}{}}}d\theta _a^{}d\eta _a^{}}{x_1x_2}\left(\frac{1}{x_2}\right)^{{\displaystyle \underset{I={\scriptscriptstyle \frac{p+2}{2}}}{\overset{p^{}/2}{}}}\nu _I}}`$
$`\times 0|{\displaystyle \underset{f=1}{\overset{N}{}}}:\mathrm{exp}\left[i{\displaystyle \underset{\mu =0}{\overset{p}{}}}\{\sqrt{{\displaystyle \frac{\alpha ^{}}{2}}}k_{f\mu }𝐗^\mu (x_f,\theta _f)+{\displaystyle \frac{1}{2}}\eta _f\zeta _{f\mu }\dot{𝐗}^\mu (x_f,\theta _f)\}\right]:|_{\genfrac{}{}{0pt}{}{\zeta _{1\mu }=\zeta _{2\mu }=0}{\theta _1=\theta _2=0}}|0`$
$`\times \sigma ,s|{\displaystyle \underset{c=3}{\overset{N}{}}}:\mathrm{exp}[{\displaystyle \underset{I=\frac{p+2}{2}}{\overset{p^{}/2}{}}}𝐄_I(\zeta _c,k_c,\eta _c)𝐙^I(x_c,\theta _c)+{\displaystyle \underset{\overline{J}=\frac{p+2}{2}}{\overset{p^{}/2}{}}}\overline{𝐄}_{\overline{J}}(\zeta _c,k_c,\eta _c)\overline{𝐙}^{\overline{J}}(x_c,\theta _c)]:|\sigma ,s,`$
where $`x_1=0`$, $`x_2=\mathrm{}`$ and $`x_3=1`$. Here the right hand side consists of two of the expectation values of the exponential operators: the one is obtained from the $`x^\mu `$-directions $`(\mu =0,\mathrm{},p)`$ and the other is from the $`x^i`$-directions $`(i=p+1,\mathrm{},p^{})`$, and we have used
$$\sigma _I,s_I|=\underset{x\mathrm{}}{lim}x^{\nu _I}0|\sigma _Is_I(x).$$
(4.6)
Let us examine the contribution from the $`x^i`$-directions $`(i=p+1,\mathrm{},p^{})`$. As is explained in section II., the operators inside $`\sigma ,s|\mathrm{}|\sigma ,s`$ in eq. (IV.) are normal ordered with respect to the $`SL(2,𝐑)`$ invariant vacuum and are not with respect to the oscillator vacuum. Applying the reordering formula eq. (2.73), we obtain
$`\sigma ,s|{\displaystyle \underset{a=3}{\overset{N}{}}}:\mathrm{exp}\left({\displaystyle \underset{I}{}}𝐄_{aI}𝐙^I(x_a,\theta _a)+{\displaystyle \underset{\overline{J}}{}}\overline{𝐄}_{a\overline{J}}\overline{𝐙}^{\overline{J}}(x_a,\theta _a)\right):|\sigma ,s`$ (4.7)
$`={\displaystyle \underset{a=3}{\overset{N}{}}}\mathrm{exp}\left[{\displaystyle \underset{I,\overline{J}}{}}𝐄_{aI}\overline{𝐄}_{a\overline{J}}𝓖_{\mathrm{sub}}^{}{}_{}{}^{I\overline{J}}((x_a,\theta _a))\right]`$
$`\times \sigma ,s\left|{\displaystyle \underset{c=3}{\overset{N}{}}}\stackrel{}{}\mathrm{exp}[{\displaystyle \underset{I}{}}𝐄_{cI}𝐙^I(x_c,\theta _c)+{\displaystyle \underset{\overline{J}}{}}\overline{𝐄}_{c\overline{J}}\overline{𝐙}^{\overline{J}}(x_c,\theta _c)]\stackrel{}{}\right|\sigma ,s,`$
$`={\displaystyle \underset{a=3}{\overset{N}{}}}x_a^{2\alpha ^{}_{I,\overline{J}}\kappa _{aI}\overline{\kappa }_{a\overline{J}}G^{I\overline{J}}}\mathrm{exp}\left[𝒞_a\left(\nu _I\right)+\sqrt{2\alpha ^{}}\eta _a{\displaystyle \underset{I\overline{J}}{}}e_{aI}\overline{\kappa }_{a\overline{J}}G^{I\overline{J}}{\displaystyle \frac{\theta _a}{x_a}}\right]`$
$`\times \sigma ,s\left|{\displaystyle \underset{c=3}{\overset{N}{}}}\stackrel{}{}\mathrm{exp}[{\displaystyle \underset{I}{}}𝐄_{cI}𝐙^I(x_c,\theta _c)+{\displaystyle \underset{\overline{J}}{}}\overline{𝐄}_{c\overline{J}}\overline{𝐙}^{\overline{J}}(x_c,\theta _c)]\stackrel{}{}\right|\sigma ,s,`$
where
$$𝒞_a(\nu _I)=\alpha ^{}\underset{I,\overline{J}}{}2\kappa _{aI}\overline{\kappa }_{a\overline{J}}G^{I\overline{J}}\left\{\gamma +\frac{1}{2}\left(𝝍(\nu _I)+𝝍(1\nu _I)\right)\right\},$$
(4.8)
and $`𝐄_{cI}`$ and $`\overline{𝐄}_{c\overline{J}}`$ stand for $`𝐄_I(\zeta _c,k_c,\eta _c)`$ and $`\overline{𝐄}_{\overline{J}}(\zeta _c,k_c,\eta _c)`$ respectively. Note that the part in eq. (4.7) that corresponds to self-contractions has been given by the subtracted Green function at the coincident point:
$`{\displaystyle \underset{I\overline{J}}{}}𝐄_{aI}\overline{𝐄}_{a\overline{J}}𝓖_{\mathrm{sub}}^{}{}_{}{}^{I\overline{J}}((x_a,\theta _a))\underset{\genfrac{}{}{0pt}{}{\theta _c\theta _a}{x_cx_a}}{lim}{\displaystyle \underset{I\overline{J}}{}}𝐄_{cI}\overline{𝐄}_{a\overline{J}}𝓖_{\mathrm{sub}}^{}{}_{}{}^{I\overline{J}}(x_c,\theta _c|x_a,\theta _a)`$
$`=𝒞_a(\nu _I)+{\displaystyle \underset{I,\overline{J}}{}}\left[2\alpha ^{}\kappa _{aI}\overline{\kappa }_{a\overline{J}}G^{I\overline{J}}\mathrm{ln}x_a+\sqrt{2\alpha ^{}}\eta _ae_{aI}\overline{\kappa }_{a\overline{J}}G^{I\overline{J}}{\displaystyle \frac{\theta _a}{x_a}}\right].`$ (4.9)
The last factor $`\sigma ,s|\mathrm{}|\sigma ,s`$ in eq. (4.7) is now calculated using the two-point function $`𝓖^{I\overline{J}}(x_c,\theta _c|x_c^{},\theta _c^{})`$:
$`\sigma ,s|{\displaystyle \underset{c=3}{\overset{N}{}}}\stackrel{}{}\mathrm{exp}\left[{\displaystyle \underset{I}{}}𝐄_{cI}𝐙^I(x_c,\theta _c)+{\displaystyle \underset{\overline{J}}{}}\overline{𝐄}_{c\overline{J}}\overline{𝐙}^{\overline{J}}(x_c,\theta _c)\right]\stackrel{}{}|\sigma ,s`$ (4.10)
$`={\displaystyle \underset{3c<c^{}N}{}}\mathrm{exp}\left[{\displaystyle \underset{I,\overline{J}}{}}\left\{𝐄_{cI}\overline{𝐄}_{c^{}\overline{J}}𝓖^{I\overline{J}}(x_c,\theta _c|x_c^{},\theta _c^{})+\overline{𝐄}_{c\overline{J}}𝐄_{c^{}I}𝓖^{I\overline{J}}(x_c^{},\theta _c^{}|x_c,\theta _c)\right\}\right]`$
$`={\displaystyle \underset{3c<c^{}N}{}}\mathrm{exp}[{\displaystyle \underset{I,\overline{J}}{}}G^{I\overline{J}}[2\alpha ^{}\kappa _{cI}\overline{\kappa }_{c^{}\overline{J}}(\nu _I;{\displaystyle \frac{x_c}{x_c^{}\theta _c\theta _c^{}}})2\alpha ^{}\overline{\kappa }_{c\overline{J}}\kappa _{c^{}I}(\nu _I;{\displaystyle \frac{x_c^{}}{x_c+\theta _c\theta _c^{}}})`$
$`+\eta _c\sqrt{2\alpha ^{}}\left\{e_{cI}\overline{\kappa }_{c^{}\overline{J}}{\displaystyle \frac{\left(\frac{x_c}{x_c^{}}\right)^{\nu _I}\theta _c^{}\left(\frac{x_c}{x_c^{}}\right)^{\nu _I1}\theta _c}{x_cx_c^{}}}+\overline{e}_{c\overline{J}}\kappa _{c^{}I}\left({\displaystyle \frac{x_c^{}}{x_c}}\right)^{\nu _I}{\displaystyle \frac{\theta _c^{}\theta _c}{x_cx_c^{}}}\right\}`$
$`+\eta _c^{}\sqrt{2\alpha ^{}}\left\{\kappa _{cI}\overline{e}_{c^{}\overline{J}}\left({\displaystyle \frac{x_c}{x_c^{}}}\right)^{\nu _I}{\displaystyle \frac{\theta _c^{}\theta _c}{x_cx_c^{}}}+\overline{\kappa }_{c\overline{J}}e_{c^{}I}{\displaystyle \frac{\left(\frac{x_c^{}}{x_c}\right)^{\nu _I1}\theta _c^{}\left(\frac{x_c^{}}{x_c}\right)^{\nu _I}\theta _c}{x_cx_c^{}}}\right\}`$
$`+\eta _c\eta _c^{}\{e_{cI}\overline{e}_{c^{}\overline{J}}{\displaystyle \frac{\left(\frac{x_c}{x_c^{}}\right)^{\nu _I}}{x_cx_c^{}}}(1\theta _c\theta _c^{}{\displaystyle \frac{(1\nu _I)+\nu _I\frac{x_c^{}}{x_c}}{x_cx_c^{}}})`$
$`+\overline{e}_{c\overline{J}}e_{c^{}I}{\displaystyle \frac{\left(\frac{x_c^{}}{x_c}\right)^{\nu _I}}{x_cx_c^{}}}(1\theta _c\theta _c^{}{\displaystyle \frac{(1\nu _I)+\nu _I\frac{x_c}{x_c^{}}}{x_cx_c^{}}})\}]`$
$`{\displaystyle \underset{I,\overline{J}}{}}ϵ(x_cx_c^{}){\displaystyle \frac{2\delta ^{I\overline{J}}\pi b_I}{\epsilon (1+b_I^2)}}\alpha ^{}(\kappa _{cI}\overline{\kappa }_{c^{\overline{J}}}\overline{\kappa }_{c\overline{J}}\kappa _{c^{}I})].`$
The factor coming from the $`x^\mu `$-directions $`(\mu =0,\mathrm{},p)`$ is handled by the two-point function $`𝐆^{\mu \mu ^{}}`$:
$`0|{\displaystyle \underset{f=1}{\overset{N}{}}}:\mathrm{exp}\left[{\displaystyle \underset{\mu =0}{\overset{p}{}}}_\mu (\zeta _f,k_f,\eta _f)𝐗^\mu (x_f,\theta _f)\right]:|0`$ (4.11)
$`=\mathrm{exp}\left[{\displaystyle \underset{1f<f^{}N}{}}{\displaystyle \underset{\mu ,\mu ^{}=0}{\overset{p}{}}}_\mu (\zeta _f,k_f,\eta _f)_\mu ^{}(\zeta _f^{},k_f^{},\eta _{})𝐆^{\mu \mu ^{}}(𝐳_f,\overline{𝐳}_f|𝐳_f^{},\overline{𝐳}_f^{})\right]|_{\genfrac{}{}{0pt}{}{z_f=\overline{z}_f=x_f,}{\theta _f=\overline{\theta }_f}}`$
$`\times 0|:\mathrm{exp}\left[{\displaystyle \underset{c=1}{\overset{N}{}}}{\displaystyle \underset{\rho =0}{\overset{p}{}}}_\rho (\zeta _c,k_c,\eta _c)𝐗^\rho (x_c,\theta _c)\right]:|0`$
$`=\mathrm{exp}[{\displaystyle \underset{1f<f^{}N}{}}\{{\displaystyle \underset{\sigma ,\rho =0}{\overset{p}{}}}\{\alpha ^{}G^{\sigma \rho }k_{f\sigma }k_{f^{}\rho }\mathrm{ln}(x_fx_f^{}+\theta _f\theta _f^{})^2`$
$`\sqrt{2\alpha ^{}}\left(\eta _fG^{\sigma \rho }\zeta _{f\sigma }k_{f^{}\rho }+\eta _f^{}G^{\sigma \rho }k_{f\sigma }\zeta _{f^{}\rho }\right){\displaystyle \frac{\theta _f\theta _f^{}}{x_fx_f^{}}}`$
$`+\eta _f\eta _f^{}G^{\sigma \rho }\zeta _{f\sigma }\zeta _{f^{}\rho }{\displaystyle \frac{1}{(x_fx_f^{}+\theta _f\theta _f^{})}}\}+{\displaystyle \underset{i,j=1}{\overset{p}{}}}{\displaystyle \frac{i}{2}}\theta ^{ij}k_{fi}k_{f^{}j}ϵ(x_fx_f^{})\}]`$
$`\times (2\pi )^{p+1}{\displaystyle \underset{\mu =0}{\overset{p}{}}}\delta \left(k_{1\mu }+\mathrm{}+k_{N\mu }\right),`$
where $`_\mu (\zeta _f,k_f,\eta _f)`$ is a differential operator defined as
$`_\mu (\zeta _f,k_f,\eta _f)=i\sqrt{{\displaystyle \frac{\alpha ^{}}{2}}}k_{f\mu }+i{\displaystyle \frac{1}{2}}\eta _f\zeta _{f\mu }(k)\left(D_f+\overline{D}_f\right),`$
$`_\mu (\zeta _f,k_f,\eta _f)𝐗^\mu (x_f,\theta _f)_\mu (\zeta _f,k_f,\eta _f)𝐗^\mu (𝐳_f,\overline{𝐳}_f)|_{\genfrac{}{}{0pt}{}{z_f=\overline{z}_f=x_f}{\theta _f=\overline{\theta }_f}}.`$ (4.12)
We have now evaluated the two of the expectation values in eq. (IV.) and are ready to present the integral representation of $`A_N`$. Let us first introduce several shorthand notations. We denote by $`{\displaystyle \genfrac{}{}{0pt}{}{}{\left(p\right)}}`$ the inner product of two $`(p+1)`$-dimensional vectors $`A_i`$ and $`B_j`$ lying on the $`Dp`$-brane worldvolume with respect to the open string metric. Namely
$$A\genfrac{}{}{0pt}{}{}{\left(p\right)}B=\underset{\sigma ,\rho =0}{\overset{p}{}}G^{\sigma \rho }A_\sigma B_\rho .$$
(4.13)
Similarly, we denote by $`{\displaystyle \genfrac{}{}{0pt}{}{}{\left(p^{}\right)}}`$ the inner product of two $`(p^{}+1)`$-dimensional vectors $`A_i`$ and $`B_j`$ lying on the $`Dp^{}`$-brane worldvolume with respect to the open string metric. We will also write $`A{\displaystyle \genfrac{}{}{0pt}{}{}{(p,p^{})}}B`$ to denote the inner product of the last $`(p^{}p)`$ components of the two vectors. For example, we have
$`k{\displaystyle \genfrac{}{}{0pt}{}{}{(p,p^{})}}\zeta ={\displaystyle \underset{I,\overline{J}=\frac{p+2}{2}}{\overset{\frac{p^{}}{2}}{}}}\left(G^{I\overline{J}}\kappa _I\overline{e}_{\overline{J}}+G^{\overline{J}I}\overline{\kappa }_{\overline{J}}e_I\right)={\displaystyle \underset{I,\overline{J}=\frac{p+2}{2}}{\overset{\frac{p^{}}{2}}{}}}G^{I\overline{J}}\left(\kappa _I\overline{e}_{\overline{J}}+\overline{\kappa }_{\overline{J}}e_I\right),`$
$`k{\displaystyle \genfrac{}{}{0pt}{}{}{(p,p^{})}}k={\displaystyle \underset{I,\overline{J}=\frac{p+2}{2}}{\overset{\frac{p^{}}{2}}{}}}2G^{I\overline{J}}\kappa _I\overline{\kappa }_{\overline{J}},\zeta {\displaystyle \genfrac{}{}{0pt}{}{}{(p,p^{})}}\zeta ={\displaystyle \underset{I,\overline{J}=\frac{p+2}{2}}{\overset{\frac{p^{}}{2}}{}}}2G^{I\overline{J}}e_I\overline{e}_{\overline{J}}.`$ (4.14)
We will use the notations $`{\displaystyle \genfrac{}{}{0pt}{}{}{(p,p^{})}}`$ and $`{\displaystyle \genfrac{}{}{0pt}{}{\times }{(p,p^{})}}`$ which denote
$$\left(k\genfrac{}{}{0pt}{}{}{(p,p^{})}\zeta \right)_I=\underset{\overline{J}}{}G^{I\overline{J}}(\kappa _I\overline{e}_{\overline{J}}+\overline{\kappa }_{\overline{J}}e_I),\left(k\genfrac{}{}{0pt}{}{\times }{(p,p^{})}\zeta \right)_I=\underset{\overline{J}}{}\frac{2\delta ^{I\overline{J}}(\kappa _I\overline{e}_{\overline{J}}\overline{\kappa }_{\overline{J}}e_I)}{\epsilon (1+b_I^2)},$$
(4.15)
etc. From these defining relations one can find that
$$\underset{I}{}\left(k\genfrac{}{}{0pt}{}{}{(p,p^{})}\zeta \right)_I=k\genfrac{}{}{0pt}{}{}{(p,p^{})}\zeta ,\underset{I}{}\left(k\genfrac{}{}{0pt}{}{\times }{(p,p^{})}\zeta \right)_I=ik\genfrac{}{}{0pt}{}{}{(p,p^{})}J\zeta ,$$
(4.16)
where $`J`$ is a $`(p^{}+1)\times (p^{}+1)`$ antisymmetric matrix defined as
$$J=\left(J_{\mu }^{}{}_{}{}^{\rho }\right)\begin{array}{cc}\hfill \left(\begin{array}{cccccccc}0& & & & & & & \\ & \mathrm{}& & & & & & \\ & & 0& & & & & \\ & & & & & & & \\ & & & 0& 1& & & \\ & & & 1& 0& & & \\ & & & & & \mathrm{}& & \\ & & & & & & 0& 1\\ & & & & & & 1& 0\end{array}\right)& \begin{array}{c}0\hfill \\ \mathrm{}\hfill \\ p\hfill \\ p+1\hfill \\ p+2\hfill \\ \mathrm{}\hfill \\ p^{}1\hfill \\ p^{}\hfill \end{array}\hfill \end{array}.$$
(4.17)
We will group the terms in the exponent by the number of $`\eta _a`$’s and by the number of $`\theta _a`$’s, using the notation $`[0,2]`$, $`[2,0]`$, $`[1,1]`$, $`[2,2]`$. The first number in the bracket indicates the number of $`\eta _a`$’s and the second number the number of $`\theta _a`$’s.
Having prepared these, we write eq. (IV.) as
$`A_N=c(2\pi )^{p+1}{\displaystyle \underset{\mu =0}{\overset{p}{}}}\delta \left({\displaystyle \underset{e=1}{\overset{N}{}}}k_{e\mu }\right){\displaystyle \underset{a=4}{\overset{N}{}}dx_a\underset{a^{}=3}{\overset{N}{}}d\theta _a^{}d\eta _a^{}\mathrm{exp}𝒞_a^{}(\nu _I)}`$ (4.18)
$`\times x_2^{_I\nu _I}(x_2x_3)(x_3x_1){\displaystyle \underset{c^{\prime \prime }=3}{\overset{N}{}}}x_{c^{\prime \prime }}^{\alpha ^{}k_{c^{\prime \prime }}{\displaystyle \genfrac{}{}{0pt}{}{}{(p,p^{})}}k_{c^{\prime \prime }}}{\displaystyle \underset{1c<c^{}N}{}}(x_c^{}x_c)^{2\alpha ^{}k_c^{}{\displaystyle \genfrac{}{}{0pt}{}{}{\left(p\right)}}k_c}`$
$`\times {\displaystyle \underset{3c<c^{}N}{}}\mathrm{exp}[2\alpha ^{}{\displaystyle \underset{I,\overline{J}}{}}G^{I\overline{J}}\{\kappa _{cI}\overline{\kappa }_{c^{}\overline{J}}(\nu _I;{\displaystyle \frac{x_c}{x_c^{}}})+\overline{\kappa }_{c\overline{J}}\kappa _{c^{}I}(\nu _I;{\displaystyle \frac{x_c^{}}{x_c}})\}]`$
$`\times \mathrm{exp}\left(\mathrm{NC}\right)\mathrm{exp}\left([0,2]+[2,0]+[1,1]+[2,2]\right).`$
Here
$`[0,2]=2\alpha ^{}{\displaystyle \underset{3c<c^{}N}{}}{\displaystyle \frac{\theta _c\theta _c^{}}{x_cx_c^{}}}[k_c{\displaystyle \genfrac{}{}{0pt}{}{}{\left(p\right)}}k_c^{}`$
$`+{\displaystyle \frac{1}{2}}{\displaystyle \underset{I}{}}\{\left(k_c{\displaystyle \genfrac{}{}{0pt}{}{}{(p,p^{})}}k_c^{}\right)_I[\left({\displaystyle \frac{x_c}{x_c^{}}}\right)^{\nu _I}+\left({\displaystyle \frac{x_c^{}}{x_c}}\right)^{\nu _I}]+\left(k_c{\displaystyle \genfrac{}{}{0pt}{}{\times }{(p,p^{})}}k_c^{}\right)_I[\left({\displaystyle \frac{x_c}{x_c^{}}}\right)^{\nu _I}\left({\displaystyle \frac{x_c^{}}{x_c}}\right)^{\nu _I}]\}],`$
$`[1,1]=\sqrt{2\alpha ^{}}{\displaystyle \underset{c=3}{\overset{N}{}}}\eta _c\theta _c\left[{\displaystyle \frac{1}{2}}\zeta _c{\displaystyle \genfrac{}{}{0pt}{}{}{(p,p^{})}}(1+iJ)k_c{\displaystyle \frac{1}{x_c}}+k_1{\displaystyle \genfrac{}{}{0pt}{}{}{\left(p\right)}}\zeta _c{\displaystyle \frac{1}{x_1x_c}}+k_2{\displaystyle \genfrac{}{}{0pt}{}{}{\left(p\right)}}\zeta _c{\displaystyle \frac{1}{x_2x_c}}\right]`$
$`\sqrt{2\alpha ^{}}{\displaystyle \underset{3c<c^{}N}{}}{\displaystyle \frac{1}{x_cx_c^{}}}[(\eta _c\zeta _c{\displaystyle \genfrac{}{}{0pt}{}{}{\left(p\right)}}k_c^{}+\eta _c^{}k_c{\displaystyle \genfrac{}{}{0pt}{}{}{\left(p\right)}}\zeta _c^{})(\theta _c\theta _c^{})`$
$`+{\displaystyle \frac{1}{2}}{\displaystyle \underset{I}{}}\{\eta _c\left(\zeta _c{\displaystyle \genfrac{}{}{0pt}{}{}{(p,p^{})}}k_c^{}\right)_I(\{\left({\displaystyle \frac{x_c}{x_c^{}}}\right)^{\nu _I1}+\left({\displaystyle \frac{x_c^{}}{x_c}}\right)^{\nu _I}\}\theta _c\{\left({\displaystyle \frac{x_c}{x_c^{}}}\right)^{\nu _I}+\left({\displaystyle \frac{x_c^{}}{x_c}}\right)^{\nu _I}\}\theta _c^{})`$
$`+\eta _c^{}\left(k_c{\displaystyle \genfrac{}{}{0pt}{}{}{(p,p^{})}}\zeta _c^{}\right)_I\left(\left\{\left({\displaystyle \frac{x_c}{x_c^{}}}\right)^{\nu _I}+\left({\displaystyle \frac{x_c^{}}{x_c}}\right)^{\nu _I}\right\}\theta _c\left\{\left({\displaystyle \frac{x_c}{x_c^{}}}\right)^{\nu _I}+\left({\displaystyle \frac{x_c^{}}{x_c}}\right)^{\nu _I1}\right\}\theta _c^{}\right)`$
$`+\eta _c\left(\zeta _c{\displaystyle \genfrac{}{}{0pt}{}{\times }{(p,p^{})}}k_c^{}\right)_I\left(\left\{\left({\displaystyle \frac{x_c}{x_c^{}}}\right)^{\nu _I1}\left({\displaystyle \frac{x_c^{}}{x_c}}\right)^{\nu _I}\right\}\theta _c\left\{\left({\displaystyle \frac{x_c}{x_c^{}}}\right)^{\nu _I}\left({\displaystyle \frac{x_c^{}}{x_c}}\right)^{\nu _I}\right\}\theta _c^{}\right)`$
$`+\eta _c^{}\left(k_c{\displaystyle \genfrac{}{}{0pt}{}{\times }{(p,p^{})}}\zeta _c^{}\right)_I(\{\left({\displaystyle \frac{x_c}{x_c^{}}}\right)^{\nu _I}\left({\displaystyle \frac{x_c^{}}{x_c}}\right)^{\nu _I}\}\theta _c\{\left({\displaystyle \frac{x_c}{x_c^{}}}\right)^{\nu _I}\left({\displaystyle \frac{x_c^{}}{x_c}}\right)^{\nu _I1}\}\theta _c^{})\}],`$
$`[2,0]={\displaystyle \underset{3c<c^{}N}{}}{\displaystyle \frac{\eta _c\eta _c^{}}{x_cx_c^{}}}[\zeta _c{\displaystyle \genfrac{}{}{0pt}{}{}{\left(p\right)}}\zeta _c^{}`$
$`+{\displaystyle \frac{1}{2}}{\displaystyle \underset{I}{}}\{\left(\zeta _c{\displaystyle \genfrac{}{}{0pt}{}{}{(p,p^{})}}\zeta _c^{}\right)_I[\left({\displaystyle \frac{x_c}{x_c^{}}}\right)^{\nu _I}+\left({\displaystyle \frac{x_c^{}}{x_c}}\right)^{\nu _I}]+\left(\zeta _c{\displaystyle \genfrac{}{}{0pt}{}{\times }{(p,p^{})}}\zeta _c^{}\right)_I[\left({\displaystyle \frac{x_c}{x_c^{}}}\right)^{\nu _I}\left({\displaystyle \frac{x_c^{}}{x_c}}\right)^{\nu _I}]\}],`$
$`[2,2]={\displaystyle \underset{3c<c^{}N}{}}{\displaystyle \frac{\eta _c\theta _c\eta _c^{}\theta _c^{}}{(x_cx_c^{})^2}}[\zeta _c{\displaystyle \genfrac{}{}{0pt}{}{}{\left(p\right)}}\zeta _c^{}`$
$`+{\displaystyle \frac{1}{2}}{\displaystyle \underset{I}{}}\{\left(\zeta _c{\displaystyle \genfrac{}{}{0pt}{}{}{(p,p^{})}}\zeta _c^{}\right)_I((1\nu _I)\{\left({\displaystyle \frac{x_c}{x_c^{}}}\right)^{\nu _I}+\left({\displaystyle \frac{x_c^{}}{x_c}}\right)^{\nu _I}\}+\nu _I\{\left({\displaystyle \frac{x_c}{x_c^{}}}\right)^{\nu _I1}+\left({\displaystyle \frac{x_c^{}}{x_c}}\right)^{\nu _I1}\})`$
$`+\left(\zeta _c{\displaystyle \genfrac{}{}{0pt}{}{\times }{(p,p^{})}}\zeta _c^{}\right)_I((1\nu _I)\{\left({\displaystyle \frac{x_c}{x_c^{}}}\right)^{\nu _I}\left({\displaystyle \frac{x_c^{}}{x_c}}\right)^{\nu _I}\}+\nu _I\{\left({\displaystyle \frac{x_c}{x_c^{}}}\right)^{\nu _I1}\left({\displaystyle \frac{x_c^{}}{x_c}}\right)^{\nu _I1}\})\}],`$
and $`(\mathrm{NC})`$ denotes the terms containing the sign function:
$`(\mathrm{NC})`$ $`=`$ $`{\displaystyle \underset{1a<a^{}N}{}}{\displaystyle \frac{i}{2}}ϵ(x_ax_a^{}){\displaystyle \underset{i,j=1}{\overset{p}{}}}\theta ^{ij}k_{ai}k_{a^{}j}`$ (4.20)
$`{\displaystyle \underset{3c<c^{}N}{}}ϵ(x_cx_c^{}){\displaystyle \underset{I,\overline{J}}{}}\alpha ^{}{\displaystyle \frac{2\delta ^{I\overline{J}}\pi b_I}{\epsilon (1+b_I^2)}}\left(\kappa _{cI}\overline{\kappa }_{c^{}\overline{J}}\overline{\kappa }_{c\overline{J}}\kappa _{c^{}I}\right)`$
$`=`$ $`{\displaystyle \underset{1a<a^{}N}{}}{\displaystyle \frac{i}{2}}ϵ(x_ax_a^{}){\displaystyle \underset{\mu ,\lambda =0}{\overset{p^{}}{}}}\theta ^{\mu \lambda }k_{a\mu }k_{a^{}\lambda },`$
with $`k_{1j}=k_{2j}=0`$ for $`(j=p+1,\mathrm{},p^{})`$. Here we have written the noncommutativity term in terms of the real variables and generalized the notation $`\theta ^{\mu \lambda }`$ to include the time components $`\theta ^{0i}=0`$. We also remind the readers that $`(\nu _I;\frac{x_c}{x_c^{}})`$ is given in eq. (2.48).
So far, we have not exploited that $`x_1=0,x_2=\mathrm{}`$ and $`x_3=1`$ except that the oscillator vacuum $`|\sigma ,s`$, $`\sigma ,s|`$ is obtained from the tachyon vertex operators. Firstly, by sending $`x_2=\mathrm{}`$, all factors in eq. (4.18) containing $`x_2`$ are removed. In fact, this is ensured by an equality
$$1\underset{I=\frac{p+2}{2}}{\overset{\frac{p^{}}{2}}{}}\nu _I+2\alpha ^{}\underset{c2}{}k_2\genfrac{}{}{0pt}{}{}{\left(p\right)}k_c=0,$$
(4.21)
which is obtained from the momentum conservation $`{\displaystyle \underset{\mu =0}{\overset{p}{}}}\delta \left({\displaystyle \underset{c=1}{\overset{N}{}}}k_{c\mu }\right)`$ and the on-shell condition (eq. (3.8)) for the tachyon. Secondly, setting $`x_1=0`$, we find
$`{\displaystyle \underset{\genfrac{}{}{0pt}{}{1c<c^{}N}{c,c^{}2}}{}}(x_c^{}x_c)^{2\alpha ^{}k_c^{}{\displaystyle \genfrac{}{}{0pt}{}{}{\left(p\right)}}k_c}{\displaystyle \underset{c^{\prime \prime }=3}{\overset{N}{}}}x_{c^{\prime \prime }}^{\alpha ^{}k_{c^{\prime \prime }}{\displaystyle \genfrac{}{}{0pt}{}{}{(p,p^{})}}k_{c^{\prime \prime }}}`$ (4.22)
$`={\displaystyle \underset{c=3}{\overset{N}{}}}x_c^{\alpha ^{}s_c+\alpha ^{}m_T^2}{\displaystyle \underset{3c<c^{}N}{}}(x_cx_c^{})^{2\alpha ^{}k_c^{}{\displaystyle \genfrac{}{}{0pt}{}{}{\left(p\right)}}k_c}.`$
Here $`s_c(k_c+k_1){\displaystyle \genfrac{}{}{0pt}{}{}{\left(p\right)}}(k_c+k_1)`$ and we have used the on-shell condition for the tachyon (eq. (3.8)) and that for the vector (eq. (3.13)). Finally we would like to convert the integrations at eq. (4.18) into those over a set of $`N3`$ $`SL(2,𝐑)`$ invariant cross ratios. We choose these cross ratios as
$$x^{(a+3)}\frac{(x_1x_{a+3})(x_2x_3)}{(x_1x_3)(x_2x_{a+3})}=\frac{x_{a+3}}{x_3},a=1,\mathrm{}N3.$$
(4.23)
We can therefore accomplish this conversion by rescaling $`x_{a+3}`$ by $`x_3`$ and setting $`x_3=1`$ in eq. (4.18) without changing the form of the integrand.
Putting all these considerations together, we obtain from eq. (4.18)
$`A_N=c(2\pi )^{p+1}{\displaystyle \underset{\mu =0}{\overset{p}{}}}\delta \left({\displaystyle \underset{e=1}{\overset{N}{}}}k_{e\mu }\right){\displaystyle \underset{a=4}{\overset{N}{}}dx_a\underset{a^{}=3}{\overset{N}{}}d\theta _a^{}d\eta _a^{}\mathrm{exp}𝒞_a^{}(\nu _I)}`$ (4.24)
$`\times {\displaystyle \underset{c=4}{\overset{N}{}}}\left[x_c^{\alpha ^{}s_c+\alpha ^{}m_T^2}(1x_c)^{2\alpha ^{}k_3{\displaystyle \genfrac{}{}{0pt}{}{}{\left(p\right)}}k_c}\right]{\displaystyle \underset{4c<c^{}N}{}}(x_cx_c^{})^{2\alpha ^{}k_c^{}{\displaystyle \genfrac{}{}{0pt}{}{}{\left(p\right)}}k_c}`$
$`\times {\displaystyle \underset{3c<c^{}N}{}}\mathrm{exp}[2\alpha ^{}{\displaystyle \underset{I,\overline{J}}{}}G^{I\overline{J}}\{\kappa _{cI}\overline{\kappa }_{c^{}\overline{J}}(\nu _I;{\displaystyle \frac{x_c}{x_c^{}}})+\overline{\kappa }_{c\overline{J}}\kappa _{c^{}I}(\nu _I;{\displaystyle \frac{x_c^{}}{x_c}})\}]`$
$`\times \mathrm{exp}\left(\mathrm{NC}\right)\mathrm{exp}\left([0,2]+[2,0]+[1,1]+[2,2]\right)|_{x_1=0,x_2=\mathrm{},x_3=1}.`$
This expression is regarded as an $`SL(2,𝐑)`$ invariant integral (Koba-Nielsen) representation for the amplitude of our concern. Let us list several features which are distinct from the corresponding formula in the case of a $`p`$-$`p`$ open string. (See ).
1. The term denoted by $`\mathrm{exp}\left(\mathrm{NC}\right)`$ which originated from the noncommutativity of the worldvolume extends into both the $`x^1,\mathrm{},x^p`$ directions and the remaining $`x^{p+1},\mathrm{},x^p^{}`$ directions.
2. To each external vector leg, we have a momentum dependent multiplicative factor $`\mathrm{exp}𝒞(\nu _I)`$.
3. A new tensor $`J`$ has appeared.
4. There are parts in the expression which are expressible in terms of the momenta of the tachyons, the momenta and the polarization tensors of the vectors and $`J`$ alone, using the inner product with respect to the open string metric. These parts come, however, with a host of other parts which do not permit such generic description in terms of the inner product.
Let us finally compute $`N=3,4`$ cases explicitly. For $`N=3`$ case, we need to pick up $`\theta _3`$ and $`\eta _3`$ from $`[1,1]`$. Using the transversality of the polarization vector (eq. (3.13)) and the $`(p+1)`$ dimensional momentum conservation, we find that
$$A_3=c(2\pi )^{p+1}\underset{\mu =0}{\overset{p}{}}\delta \left(\underset{a=1}{\overset{3}{}}k_{a\mu }\right)\sqrt{\frac{\alpha ^{}}{2}}\left\{(k_2k_1)\genfrac{}{}{0pt}{}{}{\left(p\right)}\zeta _3ik_3\genfrac{}{}{0pt}{}{}{(p,p^{})}J\zeta _3\right\}e^{𝒞_3(\nu _I)}e^{\frac{i}{2}\theta ^{ij}k_{1i}k_{2j}}.$$
For $`N=4`$ case, we have
$`[0,2]=2\alpha ^{}{\displaystyle \frac{\theta _3\theta _4}{1x}}\left[k_3{\displaystyle \genfrac{}{}{0pt}{}{}{\left(p\right)}}k_4+{\displaystyle \frac{1}{2}}{\displaystyle \underset{I}{}}\left\{\left(k_3{\displaystyle \genfrac{}{}{0pt}{}{}{(p,p^{})}}k_4\right)_I\left[x^{\nu _I}+x^{\nu _I}\right]+\left(k_3{\displaystyle \genfrac{}{}{0pt}{}{\times }{(p,p^{})}}k_4\right)_I\left[x^{\nu _I}x^{\nu _I}\right]\right\}\right],`$
$`[1,1]=\sqrt{2\alpha ^{}}[{\displaystyle \frac{\eta _3\theta _3}{2}}\{((k_2k_1){\displaystyle \genfrac{}{}{0pt}{}{}{\left(p\right)}}\zeta _3ik_3{\displaystyle \genfrac{}{}{0pt}{}{}{(p,p^{})}}J\zeta _3)+k_4{\displaystyle \genfrac{}{}{0pt}{}{}{\left(p\right)}}\zeta _3\}`$
$`+{\displaystyle \frac{\eta _4\theta _4}{2x}}\left\{\left((k_2k_1){\displaystyle \genfrac{}{}{0pt}{}{}{\left(p\right)}}\zeta _4ik_4{\displaystyle \genfrac{}{}{0pt}{}{}{(p,p^{})}}J\zeta _4\right)+k_3{\displaystyle \genfrac{}{}{0pt}{}{}{\left(p\right)}}\zeta _4\right\}`$
$`+{\displaystyle \frac{\theta _3\theta _4}{1x}}\left(\eta _3k_4{\displaystyle \genfrac{}{}{0pt}{}{}{\left(p\right)}}\zeta _3+\eta _4k_3{\displaystyle \genfrac{}{}{0pt}{}{}{\left(p\right)}}\zeta _4\right)`$
$`{\displaystyle \frac{1}{2}}{\displaystyle \frac{1}{1x}}{\displaystyle \underset{I}{}}\{\eta _3\left(\zeta _3{\displaystyle \genfrac{}{}{0pt}{}{}{(p,p^{})}}k_4\right)_I\{(x^{\nu _I+1}+x^{\nu _I})\theta _3(x^{\nu _I}+x^{\nu _I})\theta _4\}`$
$`+\eta _4\left(k_3{\displaystyle \genfrac{}{}{0pt}{}{}{(p,p^{})}}\zeta _4\right)_I\left\{\left(x^{\nu _I}+x^{\nu _I}\right)\theta _3\left(x^{\nu _I}+x^{\nu _I1}\right)\theta _4\right\}`$
$`+\eta _3\left(\zeta _3{\displaystyle \genfrac{}{}{0pt}{}{\times }{(p,p^{})}}k_4\right)_I\left\{\left(x^{\nu _I+1}x^{\nu _I}\right)\theta _3\left(x^{\nu _I}x^{\nu _I}\right)\theta _4\right\}`$
$`+\eta _4\left(k_3{\displaystyle \genfrac{}{}{0pt}{}{\times }{(p,p^{})}}\zeta _4\right)_I\{(x^{\nu _I}x^{\nu _I})\theta _3(x^{\nu _I}x^{\nu _I1})\theta _4\}\}],`$
$`[2,0]={\displaystyle \frac{\eta _3\eta _4}{1x}}\left[\zeta _3{\displaystyle \genfrac{}{}{0pt}{}{}{\left(p\right)}}\zeta _4+{\displaystyle \frac{1}{2}}{\displaystyle \underset{I}{}}\left\{\left(\zeta _3{\displaystyle \genfrac{}{}{0pt}{}{}{(p,p^{})}}\zeta _4\right)_I\left(x^{\nu _I}+x^{\nu _I}\right)+\left(\zeta _3{\displaystyle \genfrac{}{}{0pt}{}{\times }{(p,p^{})}}\zeta _4\right)_I\left(x^{\nu _I}x^{\nu _I}\right)\right\}\right],`$
$`[2,2]={\displaystyle \frac{\eta _3\theta _3\eta _4\theta _4}{(1x)^2}}[\zeta _3{\displaystyle \genfrac{}{}{0pt}{}{}{\left(p\right)}}\zeta _4+{\displaystyle \frac{1}{2}}{\displaystyle \underset{I}{}}\{\left(\zeta _3{\displaystyle \genfrac{}{}{0pt}{}{}{(p,p^{})}}\zeta _4\right)_I\{(1\nu _I)(x^{\nu _I}+x^{\nu _I})+\nu _I(x^{\nu _I+1}+x^{\nu _I1})\}`$ (4.25)
$`+\left(\zeta _3{\displaystyle \genfrac{}{}{0pt}{}{\times }{(p,p^{})}}\zeta _4\right)_I\{(1\nu _I)(x^{\nu _I}x^{\nu _I})+\nu _I(x^{\nu _I+1}x^{\nu _I1})\}\}],`$
where we have set $`x_3=1`$ and written $`x_4=x`$. By picking up terms from $`[2,2]+[0,2][2,0]+\frac{1}{2}[1,1]^2`$, we obtain
$`A_4=c(2\pi )^{p+1}{\displaystyle \underset{\mu =0}{\overset{p}{}}}\delta \left({\displaystyle \underset{a=1}{\overset{4}{}}}k_{a\mu }\right){\displaystyle _0^1}𝑑xx^{\alpha ^{}t+\alpha ^{}m_T^2}(1x)^{2\alpha ^{}k_3{\displaystyle \genfrac{}{}{0pt}{}{}{\left(p\right)}}k_4}\mathrm{exp}\left(𝒞_3(\nu _I)+𝒞_4(\nu _I)+(\mathrm{NC})\right)`$ (4.26)
$`\times \mathrm{exp}\left[\alpha ^{}{\displaystyle \underset{I}{}}\left\{\left(k_3{\displaystyle \genfrac{}{}{0pt}{}{}{(p,p^{})}}k_4+k_3{\displaystyle \genfrac{}{}{0pt}{}{\times }{(p,p^{})}}k_4\right)_I(\nu _I;{\displaystyle \frac{1}{x}})+\left(k_3{\displaystyle \genfrac{}{}{0pt}{}{}{(p,p^{})}}k_4k_3{\displaystyle \genfrac{}{}{0pt}{}{\times }{(p,p^{})}}k_4\right)_I(\nu _I;x)\right\}\right]`$
$`\times [{\displaystyle \frac{1}{(1x)^2}}\zeta _3{\displaystyle \genfrac{}{}{0pt}{}{}{\left(p\right)}}\zeta _4(12\alpha ^{}k_3{\displaystyle \genfrac{}{}{0pt}{}{}{\left(p\right)}}k_4)`$
$`+{\displaystyle \frac{\alpha ^{}}{2}}{\displaystyle \frac{1}{x}}\left\{\left[(k_2k_1){\displaystyle \genfrac{}{}{0pt}{}{}{\left(p\right)}}\zeta _3ik_3{\displaystyle \genfrac{}{}{0pt}{}{}{(p,p^{})}}J\zeta _3\right]k_4{\displaystyle \genfrac{}{}{0pt}{}{}{\left(p\right)}}\zeta _3\right\}\left\{\left[(k_2k_1){\displaystyle \genfrac{}{}{0pt}{}{}{\left(p\right)}}\zeta _4ik_4{\displaystyle \genfrac{}{}{0pt}{}{}{(p,p^{})}}J\zeta _4\right]+k_3{\displaystyle \genfrac{}{}{0pt}{}{}{\left(p\right)}}\zeta _4\right\}`$
$`+\alpha ^{}{\displaystyle \frac{1}{1x}}\left\{\left[(k_2k_1){\displaystyle \genfrac{}{}{0pt}{}{}{\left(p\right)}}\zeta _3ik_3{\displaystyle \genfrac{}{}{0pt}{}{}{(p,p^{})}}J\zeta _3\right]k_3{\displaystyle \genfrac{}{}{0pt}{}{}{\left(p\right)}}\zeta _4k_4{\displaystyle \genfrac{}{}{0pt}{}{}{\left(p\right)}}\zeta _3\left[(k_2k_1){\displaystyle \genfrac{}{}{0pt}{}{}{\left(p\right)}}\zeta _4ik_4{\displaystyle \genfrac{}{}{0pt}{}{}{(p,p^{})}}J\zeta _4\right]\right\}`$
$`+{\displaystyle \underset{I}{}}{\displaystyle \frac{x^{\nu _I}}{(1x)^2}}\{\alpha ^{}(k_3{\displaystyle \genfrac{}{}{0pt}{}{}{(p,p^{})}}k_4+k_3{\displaystyle \genfrac{}{}{0pt}{}{\times }{(p,p^{})}}k_4)_I\zeta _3{\displaystyle \genfrac{}{}{0pt}{}{}{\left(p\right)}}\zeta _4`$
$`+({\displaystyle \frac{1\nu _I}{2}}\alpha ^{}k_3{\displaystyle \genfrac{}{}{0pt}{}{}{\left(p\right)}}k_4)(\zeta _3{\displaystyle \genfrac{}{}{0pt}{}{}{(p,p^{})}}\zeta _4+\zeta _3{\displaystyle \genfrac{}{}{0pt}{}{\times }{(p,p^{})}}\zeta _4)_I\}`$
$`+{\displaystyle \underset{I}{}}{\displaystyle \frac{x^{\nu _I}}{(1x)^2}}\{\alpha ^{}(k_3{\displaystyle \genfrac{}{}{0pt}{}{}{(p,p^{})}}k_4k_3{\displaystyle \genfrac{}{}{0pt}{}{\times }{(p,p^{})}}k_4)_I\zeta _3{\displaystyle \genfrac{}{}{0pt}{}{}{\left(p\right)}}\zeta _4`$
$`+({\displaystyle \frac{1\nu _I}{2}}\alpha ^{}k_3{\displaystyle \genfrac{}{}{0pt}{}{}{\left(p\right)}}k_4)(\zeta _3{\displaystyle \genfrac{}{}{0pt}{}{}{(p,p^{})}}\zeta _4\zeta _3{\displaystyle \genfrac{}{}{0pt}{}{\times }{(p,p^{})}}\zeta _4)_I\}`$
$`+{\displaystyle \underset{I}{}}{\displaystyle \frac{\nu _I}{2}}{\displaystyle \frac{x^{\nu _I+1}}{(1x)^2}}\left(\zeta _3{\displaystyle \genfrac{}{}{0pt}{}{}{(p,p^{})}}\zeta _4+\zeta _3{\displaystyle \genfrac{}{}{0pt}{}{\times }{(p,p^{})}}\zeta _4\right)_I+{\displaystyle \underset{I}{}}{\displaystyle \frac{\nu _I}{2}}{\displaystyle \frac{x^{\nu _I1}}{(1x)^2}}\left(\zeta _3{\displaystyle \genfrac{}{}{0pt}{}{}{(p,p^{})}}\zeta _4\zeta _3{\displaystyle \genfrac{}{}{0pt}{}{\times }{(p,p^{})}}\zeta _4\right)_I`$
$`{\displaystyle \frac{\alpha ^{}}{2}}{\displaystyle \underset{I,L}{}}{\displaystyle \frac{x^{\nu _I\nu _L}}{(1x)^2}}\left(k_3{\displaystyle \genfrac{}{}{0pt}{}{}{(p,p^{})}}k_4+k_3{\displaystyle \genfrac{}{}{0pt}{}{\times }{(p,p^{})}}k_4\right)_I\left(\zeta _3{\displaystyle \genfrac{}{}{0pt}{}{}{(p,p^{})}}\zeta _4+\zeta _3{\displaystyle \genfrac{}{}{0pt}{}{\times }{(p,p^{})}}\zeta _4\right)_L`$
$`{\displaystyle \frac{\alpha ^{}}{2}}{\displaystyle \underset{I,L}{}}{\displaystyle \frac{x^{\nu _I+\nu _L}}{(1x)^2}}\left(k_3{\displaystyle \genfrac{}{}{0pt}{}{}{(p,p^{})}}k_4k_3{\displaystyle \genfrac{}{}{0pt}{}{\times }{(p,p^{})}}k_4\right)_I\left(\zeta _3{\displaystyle \genfrac{}{}{0pt}{}{}{(p,p^{})}}\zeta _4\zeta _3{\displaystyle \genfrac{}{}{0pt}{}{\times }{(p,p^{})}}\zeta _4\right)_L`$
$`{\displaystyle \frac{\alpha ^{}}{2}}{\displaystyle \underset{I,L}{}}{\displaystyle \frac{x^{\nu _I+\nu _L}}{(1x)^2}}\left(k_3{\displaystyle \genfrac{}{}{0pt}{}{}{(p,p^{})}}k_4+k_3{\displaystyle \genfrac{}{}{0pt}{}{\times }{(p,p^{})}}k_4\right)_I\left(\zeta _3{\displaystyle \genfrac{}{}{0pt}{}{}{(p,p^{})}}\zeta _4\zeta _3{\displaystyle \genfrac{}{}{0pt}{}{\times }{(p,p^{})}}\zeta _4\right)_L`$
$`{\displaystyle \frac{\alpha ^{}}{2}}{\displaystyle \underset{I,L}{}}{\displaystyle \frac{x^{\nu _I\nu _L}}{(1x)^2}}\left(k_3{\displaystyle \genfrac{}{}{0pt}{}{}{(p,p^{})}}k_4k_3{\displaystyle \genfrac{}{}{0pt}{}{\times }{(p,p^{})}}k_4\right)_I\left(\zeta _3{\displaystyle \genfrac{}{}{0pt}{}{}{(p,p^{})}}\zeta _4+\zeta _3{\displaystyle \genfrac{}{}{0pt}{}{\times }{(p,p^{})}}\zeta _4\right)_L`$
$`+{\displaystyle \frac{\alpha ^{}}{2}}{\displaystyle \underset{I,L}{}}{\displaystyle \frac{x^{\nu _I\nu _L}}{1x}}\left(k_4{\displaystyle \genfrac{}{}{0pt}{}{}{(p,p^{})}}\zeta _3k_4{\displaystyle \genfrac{}{}{0pt}{}{\times }{(p,p^{})}}\zeta _3\right)_I\left(k_3{\displaystyle \genfrac{}{}{0pt}{}{}{(p,p^{})}}\zeta _4+k_3{\displaystyle \genfrac{}{}{0pt}{}{\times }{(p,p^{})}}\zeta _4\right)_L`$
$`{\displaystyle \frac{\alpha ^{}}{2}}{\displaystyle \underset{I,L}{}}{\displaystyle \frac{x^{\nu _I+\nu _L1}}{1x}}\left(k_4{\displaystyle \genfrac{}{}{0pt}{}{}{(p,p^{})}}\zeta _3+k_4{\displaystyle \genfrac{}{}{0pt}{}{\times }{(p,p^{})}}\zeta _3\right)_I\left(k_3{\displaystyle \genfrac{}{}{0pt}{}{}{(p,p^{})}}\zeta _4k_3{\displaystyle \genfrac{}{}{0pt}{}{\times }{(p,p^{})}}\zeta _4\right)_L`$
$`+{\displaystyle \frac{\alpha ^{}}{2}}{\displaystyle \underset{I}{}}{\displaystyle \frac{x^{\nu _I1}}{1x}}\{([(k_2k_1){\displaystyle \genfrac{}{}{0pt}{}{}{\left(p\right)}}\zeta _3ik_3{\displaystyle \genfrac{}{}{0pt}{}{}{(p,p^{})}}J\zeta _3]k_4{\displaystyle \genfrac{}{}{0pt}{}{}{\left(p\right)}}\zeta _3)(k_3{\displaystyle \genfrac{}{}{0pt}{}{}{(p,p^{})}}\zeta _4k_3{\displaystyle \genfrac{}{}{0pt}{}{\times }{(p,p^{})}}\zeta _4)_I`$
$`([(k_2k_1){\displaystyle \genfrac{}{}{0pt}{}{}{\left(p\right)}}\zeta _4ik_4{\displaystyle \genfrac{}{}{0pt}{}{}{(p,p^{})}}J\zeta _4]+k_3{\displaystyle \genfrac{}{}{0pt}{}{}{\left(p\right)}}\zeta _4)(k_4{\displaystyle \genfrac{}{}{0pt}{}{}{(p,p^{})}}\zeta _3+k_4{\displaystyle \genfrac{}{}{0pt}{}{\times }{(p,p^{})}}\zeta _3)_I\}`$
$`+{\displaystyle \frac{\alpha ^{}}{2}}{\displaystyle \underset{I}{}}{\displaystyle \frac{x^{\nu _I}}{1x}}\{([(k_2k_1){\displaystyle \genfrac{}{}{0pt}{}{}{\left(p\right)}}\zeta _3ik_3{\displaystyle \genfrac{}{}{0pt}{}{}{(p,p^{})}}J\zeta _3]+k_4{\displaystyle \genfrac{}{}{0pt}{}{}{\left(p\right)}}\zeta _3)(k_3{\displaystyle \genfrac{}{}{0pt}{}{}{(p,p^{})}}\zeta _4+k_3{\displaystyle \genfrac{}{}{0pt}{}{\times }{(p,p^{})}}\zeta _4)_I`$
$`([(k_2k_1){\displaystyle \genfrac{}{}{0pt}{}{}{\left(p\right)}}\zeta _4ik_4{\displaystyle \genfrac{}{}{0pt}{}{}{(p,p^{})}}J\zeta _4]k_3{\displaystyle \genfrac{}{}{0pt}{}{}{\left(p\right)}}\zeta _4)(k_4{\displaystyle \genfrac{}{}{0pt}{}{}{(p,p^{})}}\zeta _3k_4{\displaystyle \genfrac{}{}{0pt}{}{\times }{(p,p^{})}}\zeta _3)_I\}],`$
where $`t`$ is defined as $`ts_4(k_4+k_1){\displaystyle \genfrac{}{}{0pt}{}{}{\left(p\right)}}(k_4+k_1)`$.
## V. The Zero Slope Limit and the Low Energy Effective Action
In the last section, we have evaluated the three and four point amplitudes with the initial and final tachyons and $`N2`$ vectors $`(N=3,4)`$ present. Let us try to extract physical significance from these.
The three point amplitude eq. (IV.) contains the two multiplicative factors $`e^{\frac{i}{2}\theta ^{ij}k_{1i}k_{2j}}`$ and $`e^{𝒞(\nu _I)}`$ both of which are listed in the last section as prominent features. The first factor represents the noncommutativity of the D$`p`$-brane worldvolume. The second factor will be discussed shortly. Aside from these factors and $`(p+1)`$-dimensional delta functions representing momentum conservation, eq. (IV.) is interpreted as coming from field theory vertex of
$$\mathrm{\Phi }i\stackrel{}{}\mathrm{\Phi }^{}\genfrac{}{}{0pt}{}{}{\left(p\right)}A+\frac{1}{2}\mathrm{\Phi }\mathrm{\Phi }^{}J^{MN}F_{MN},$$
(5.1)
where $`\mathrm{\Phi }`$ is a complex tachyon field and $`A_M`$ and $`F_{MN}`$ are the gauge field and its field strength respectively. The first term is the gauge-scalar derivative interaction while the second factor is a new interaction coming from our $`p`$-$`p^{}`$ open string system.
The four point amplitude eq. (4.26) is quite complex but one can still systematically investigate the singular behavior of the integrand around its end points $`x=0`$, $`1`$. This behavior is sufficient to tell us the zero slope limit of the amplitude and the content of the low energy field theory. We will focus upon this in the remainder of this section. To be more accurate we consider the sum of eq. (4.26) and the one obtained from this by $`k_3k_4`$, $`\zeta _3\zeta _4`$ in accordance with the two open string (dual) diagrams.
The nontrivial zero slope limit is given by sending $`\alpha ^{}0`$ while keeping the parameter $`\theta ^{ij}`$ of noncommutativity and the open string metric fixed. From eqs. (2.8) and (2.9) this means that
$$\begin{array}{ccc}\hfill \alpha ^{}& & \epsilon ^{1/2}0,\hfill \\ \hfill g& & \epsilon 0,\hfill \\ \hfill |b_I|& & \epsilon ^{1/2}\mathrm{}.\hfill \end{array}$$
(5.2)
In this limiting procedure, $`\alpha ^{}b_I`$ becomes finite: $`\alpha ^{}b_I\beta _I`$. Let us first look at the multiplicative factor $`𝒞(\nu _I)`$ defined in eq. (4.8). Using $`𝝍(1)=\gamma `$ and eq. (A.5), we obtain
$$𝒞(\nu _I)\pi \underset{I,\overline{J}}{}\left|\beta _I\right|\kappa _I\overline{\kappa }_{\overline{J}}G^{I\overline{J}}=\frac{\pi }{2}\underset{I}{}\left|\beta _I\right|\left(k\genfrac{}{}{0pt}{}{}{(p,p^{})}k\right)_I.$$
(5.3)
So this exponential multiplicative factor acts as a gaussian damping factor when vectors propagate into the $`x^{p+1},\mathrm{},x^p^{}`$ directions.
Let us come back to the four point amplitude (4.26). If the integrand is regular, one could take the $`\alpha ^{}0`$ limit inside the integral and this will not give us any nontrivial contribution. If the integrand is singular at some point, it will still not give us much as long as one can avoid such singularity by a contour deformation. The nontrivial contribution in the $`\alpha ^{}0`$ limit, therefore, is obtained only when we have end point singularities.
Let us focus on the behavior of the integrand near $`x=0`$ from which we can read off the mass of particles exchanged in the $`t`$-channel. Fig.3 indicates that the $`t`$-channel poles originate in the propagation of the $`p`$-$`p^{}`$ open string. Thus the complicated behavior of the integrand near $`x=0`$ should reflect the spectrum of the $`p`$-$`p^{}`$ open string . In order to identify the $`t`$-channel poles, we expand the integrand of the amplitude (4.26) around $`x=0`$:
$$A_4=_0^1𝑑x\underset{A}{}f_Ax^{\alpha ^{}t+K_A},$$
(5.4)
where the coefficients $`f_A`$ are functions of momenta and polarization tensors. The term $`f_Ax^{\alpha ^{}t+K_A}`$ in the integrand of the above equation yields the $`t`$-channel pole at $`\alpha ^{}t=K_A+1`$, when it is integrated near $`x=0`$: $`{\displaystyle _0^\delta }𝑑x\mathrm{}`$ . From explicit computation we find that the $`t`$-channel poles exist at
$`\alpha ^{}t=\alpha ^{}m_T^2+W+{\displaystyle \underset{I^{}}{}}M_I^{}(n+1\nu _I^{})+{\displaystyle \underset{L^{}}{}}M_L^{}^{}(n^{}+\nu _L^{})+N,`$ (5.5)
$`\begin{array}{ccc}\hfill \text{with}W& =& 0,1,1\nu _I,\nu _I,1+\nu _I,2\nu _I,\hfill \\ & & 1\nu _I\nu _L,\nu _I+\nu _L,1+\nu _I+\nu _L,\text{ or }1\nu _I+\nu _L,\hfill \end{array}`$ (5.8)
where $`n`$, $`n^{}`$, $`M_I^{}`$, $`M_L^{}^{}`$ and $`N`$ are non-negative integers. The terms proportional to $`M_I^{}`$ and $`M_L^{}^{}`$ come from the exponential of the hypergeometric function $``$,
$`\mathrm{exp}[\alpha ^{}{\displaystyle \underset{I}{}}\{(k_3{\displaystyle \genfrac{}{}{0pt}{}{}{(p,p^{})}}k_4+k_3{\displaystyle \genfrac{}{}{0pt}{}{\times }{(p,p^{})}}k_4)_I(\nu _I:{\displaystyle \frac{1}{x}})+(k_3{\displaystyle \genfrac{}{}{0pt}{}{}{(p,p^{})}}k_4k_3{\displaystyle \genfrac{}{}{0pt}{}{\times }{(p,p^{})}}k_4)_I(\nu _I:x)\}]`$
$`=\mathrm{exp}\left(\alpha ^{}{\displaystyle \underset{I}{}}\pi b_I\left(k_3{\displaystyle \genfrac{}{}{0pt}{}{\times }{(p,p^{})}}k_4\right)_I\right)`$
$`\times {\displaystyle \underset{M=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{M!}}(\alpha ^{}{\displaystyle \underset{I^{}}{}}(k_3{\displaystyle \genfrac{}{}{0pt}{}{}{(p,p^{})}}k_4+k_3{\displaystyle \genfrac{}{}{0pt}{}{\times }{(p,p^{})}}k_4)_I^{}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{x^{n+1\nu _I^{}}}{n+1\nu _I^{}}})^M`$
$`\times {\displaystyle \underset{M^{}}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{M^{}!}}(\alpha ^{}{\displaystyle \underset{L^{}}{}}(k_3{\displaystyle \genfrac{}{}{0pt}{}{}{(p,p^{})}}k_4k_3{\displaystyle \genfrac{}{}{0pt}{}{\times }{(p,p^{})}}k_4)_L^{}{\displaystyle \underset{n^{}=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{x^{n^{}+\nu _L^{}}}{n^{}+\nu _L^{}}})^M^{}.`$ (5.9)
The $`t`$-channel poles eq. (5.5) correspond to the spectrum of the $`p`$-$`p^{}`$ open string . In view of the analysis in , we expect that a large number of light states should be exchanged in the zero slope limit in which one of $`\nu _I`$ goes to unity and the others approach zero. In order to specify the situation, we assume, without loss of generality, that $`\nu \nu _{\frac{p+2}{2}}`$ goes to $`1`$ and $`\nu _{\stackrel{~}{I}}`$ $`(\stackrel{~}{I}\frac{p+2}{2})`$ go to $`0`$ in the zero slope limit. In this zero slope limit many light states are realized by the poles in eq. (5.5) with $`(n,n^{},M_{\stackrel{~}{I}^{}},M_{\frac{p+2}{2}}^{},N)=0`$,
$$W=0,1\nu ,\nu _{\stackrel{~}{I}},1\nu \nu _{\stackrel{~}{I}},1\nu +\nu _{\stackrel{~}{I}},$$
(5.10)
and $`M_{\frac{p+2}{2}}`$ and $`M_{\stackrel{~}{I}}^{}`$ being arbitrary non-negative integers. Aside from the multiplicative factors and the momentum conserving delta functions, the massless pole obtained in the zero slope limit turns out to be
$`\{{\displaystyle \frac{1}{tm_T^2}}\left[{\displaystyle \frac{1}{2}}\{(k_2(k_1+k_4)){\displaystyle \genfrac{}{}{0pt}{}{}{\left(p\right)}}\zeta _3ik_3{\displaystyle \genfrac{}{}{0pt}{}{}{(p,p^{})}}J\zeta _3\}\{((k_2+k_3)k_1){\displaystyle \genfrac{}{}{0pt}{}{}{\left(p\right)}}\zeta _4ik_4{\displaystyle \genfrac{}{}{0pt}{}{}{(p,p^{})}}J\zeta _4\}\right]`$
$`+{\displaystyle \frac{1}{t(m_T^2+\frac{1}{\pi |\beta |})}}[(k_3{\displaystyle \genfrac{}{}{0pt}{}{}{(p,p^{})}}k_4+k_3{\displaystyle \genfrac{}{}{0pt}{}{\times }{(p,p^{})}}k_4)_{\frac{p+2}{2}}\zeta _3{\displaystyle \genfrac{}{}{0pt}{}{}{\left(p\right)}}\zeta _4`$
$`+\left({\displaystyle \frac{1}{2\pi |\beta |}}k_3{\displaystyle \genfrac{}{}{0pt}{}{}{\left(p\right)}}k_4\right)(\zeta _3{\displaystyle \genfrac{}{}{0pt}{}{}{(p,p^{})}}\zeta _4+\zeta _3{\displaystyle \genfrac{}{}{0pt}{}{\times }{(p,p^{})}}\zeta _4)_{\frac{p+2}{2}}`$
$`+{\displaystyle \frac{1}{2}}\left\{(k_2k_1+k_4){\displaystyle \genfrac{}{}{0pt}{}{}{\left(p\right)}}\zeta _3ik_3{\displaystyle \genfrac{}{}{0pt}{}{}{(p,p^{})}}J\zeta _3\right\}(k_3{\displaystyle \genfrac{}{}{0pt}{}{}{(p,p^{})}}\zeta _4+k_3{\displaystyle \genfrac{}{}{0pt}{}{\times }{(p,p^{})}}\zeta _4)_{\frac{p+2}{2}}`$
$`{\displaystyle \frac{1}{2}}\{(k_2k_1k_3){\displaystyle \genfrac{}{}{0pt}{}{}{\left(p\right)}}\zeta _4ik_4{\displaystyle \genfrac{}{}{0pt}{}{}{(p,p^{})}}J\zeta _4\}(k_4{\displaystyle \genfrac{}{}{0pt}{}{}{(p,p^{})}}\zeta _3k_4{\displaystyle \genfrac{}{}{0pt}{}{\times }{(p,p^{})}}\zeta _3)_{\frac{p+2}{2}}]`$
$`+{\displaystyle \underset{\stackrel{~}{I}}{}}{\displaystyle \frac{1}{t(m_T^2+\frac{1}{\pi |\beta _{\stackrel{~}{I}}|})}}[{\displaystyle \frac{1}{2\pi |\beta _{\stackrel{~}{I}}|}}(\zeta _3{\displaystyle \genfrac{}{}{0pt}{}{}{(p,p^{})}}\zeta _4\zeta _3{\displaystyle \genfrac{}{}{0pt}{}{\times }{(p,p^{})}}\zeta _4)_{\stackrel{~}{I}}`$
$`+{\displaystyle \frac{1}{2}}\left\{\left(k_2(k_1+k_4)\right){\displaystyle \genfrac{}{}{0pt}{}{}{\left(p\right)}}\zeta _3ik_3{\displaystyle \genfrac{}{}{0pt}{}{}{(p,p^{})}}J\zeta _3\right\}\left(k_3{\displaystyle \genfrac{}{}{0pt}{}{}{(p,p^{})}}\zeta _4k_3{\displaystyle \genfrac{}{}{0pt}{}{\times }{(p,p^{})}}\zeta _4\right)_{\stackrel{~}{I}}`$
$`{\displaystyle \frac{1}{2}}\{((k_2+k_3)k_1){\displaystyle \genfrac{}{}{0pt}{}{}{\left(p\right)}}\zeta _4ik_4{\displaystyle \genfrac{}{}{0pt}{}{}{(p,p^{})}}J\zeta _4\}(k_4{\displaystyle \genfrac{}{}{0pt}{}{}{(p,p^{})}}\zeta _3+k_4{\displaystyle \genfrac{}{}{0pt}{}{\times }{(p,p^{})}}\zeta _3)_{\stackrel{~}{I}}]`$
$`{\displaystyle \underset{\stackrel{~}{I},\stackrel{~}{L}}{}}{\displaystyle \frac{1}{t(m_T^2+\frac{1}{\pi |\beta _{\stackrel{~}{I}}|}+\frac{1}{\pi |\beta _{\stackrel{~}{L}}|})}}{\displaystyle \frac{1}{2}}\left(k_4{\displaystyle \genfrac{}{}{0pt}{}{}{(p,p^{})}}\zeta _3+k_4{\displaystyle \genfrac{}{}{0pt}{}{\times }{(p,p^{})}}\zeta _3\right)_{\stackrel{~}{I}}\left(k_3{\displaystyle \genfrac{}{}{0pt}{}{}{(p,p^{})}}\zeta _4k_3{\displaystyle \genfrac{}{}{0pt}{}{\times }{(p,p^{})}}\zeta _4\right)_{\stackrel{~}{L}}`$
$`+{\displaystyle \underset{\stackrel{~}{I}}{}}{\displaystyle \frac{1}{t(m_T^2+\frac{1}{\pi |\beta |}\frac{1}{\pi |\beta _{\stackrel{~}{I}}|})}}\times [{\displaystyle \frac{1}{2}}(k_3{\displaystyle \genfrac{}{}{0pt}{}{}{(p,p^{})}}k_4+k_3{\displaystyle \genfrac{}{}{0pt}{}{\times }{(p,p^{})}}k_4)_{\frac{p+2}{2}}(\zeta _3{\displaystyle \genfrac{}{}{0pt}{}{}{(p,p^{})}}\zeta _4+\zeta _3{\displaystyle \genfrac{}{}{0pt}{}{\times }{(p,p^{})}}\zeta _4)_{\stackrel{~}{I}}`$
$`{\displaystyle \frac{1}{2}}\left(k_3{\displaystyle \genfrac{}{}{0pt}{}{}{(p,p^{})}}k_4+k_3{\displaystyle \genfrac{}{}{0pt}{}{\times }{(p,p^{})}}k_4\right)_{\stackrel{~}{I}}\left(\zeta _3{\displaystyle \genfrac{}{}{0pt}{}{}{(p,p^{})}}\zeta _4+\zeta _3{\displaystyle \genfrac{}{}{0pt}{}{\times }{(p,p^{})}}\zeta _4\right)_{\frac{p+2}{2}}`$
$`+{\displaystyle \frac{1}{2}}\left(k_4{\displaystyle \genfrac{}{}{0pt}{}{}{(p,p^{})}}\zeta _3k_4{\displaystyle \genfrac{}{}{0pt}{}{\times }{(p,p^{})}}\zeta _3\right)_{\frac{p+2}{2}}\left(k_3{\displaystyle \genfrac{}{}{0pt}{}{}{(p,p^{})}}\zeta _4+k_3{\displaystyle \genfrac{}{}{0pt}{}{\times }{(p,p^{})}}\zeta _4\right)_{\stackrel{~}{I}}`$
$`+{\displaystyle \frac{1}{2}}(k_4{\displaystyle \genfrac{}{}{0pt}{}{}{(p,p^{})}}\zeta _3k_4{\displaystyle \genfrac{}{}{0pt}{}{\times }{(p,p^{})}}\zeta _3)_{\stackrel{~}{I}}(k_3{\displaystyle \genfrac{}{}{0pt}{}{}{(p,p^{})}}\zeta _4+k_3{\displaystyle \genfrac{}{}{0pt}{}{\times }{(p,p^{})}}\zeta _4)_{\frac{p+2}{2}}]`$
$`{\displaystyle \underset{\stackrel{~}{I}}{}}{\displaystyle \frac{1}{t(m_T^2+\frac{1}{\pi |\beta |}+\frac{1}{\pi |\beta _{\stackrel{~}{I}}|})}}\times [{\displaystyle \frac{1}{2}}(k_3{\displaystyle \genfrac{}{}{0pt}{}{}{(p,p^{})}}k_4+k_3{\displaystyle \genfrac{}{}{0pt}{}{\times }{(p,p^{})}}k_4)_{\frac{p+2}{2}}(\zeta _3{\displaystyle \genfrac{}{}{0pt}{}{}{(p,p^{})}}\zeta _4\zeta _3{\displaystyle \genfrac{}{}{0pt}{}{\times }{(p,p^{})}}\zeta _4)_{\stackrel{~}{I}}`$
$`+{\displaystyle \frac{1}{2}}(k_3{\displaystyle \genfrac{}{}{0pt}{}{}{(p,p^{})}}k_4k_3{\displaystyle \genfrac{}{}{0pt}{}{\times }{(p,p^{})}}k_4)_{\stackrel{~}{I}}(\zeta _3{\displaystyle \genfrac{}{}{0pt}{}{}{(p,p^{})}}\zeta _4+\zeta _3{\displaystyle \genfrac{}{}{0pt}{}{\times }{(p,p^{})}}\zeta _4)_{\frac{p+2}{2}}]\}`$
$`\times \mathrm{exp}\left\{\pi {\displaystyle \underset{I}{}}\left|\beta _I\right|\left(k_3{\displaystyle \genfrac{}{}{0pt}{}{}{(p,p^{})}}k_4\right)_I\right\},`$ (5.11)
where we have used
$$\{\begin{array}{cccc}\nu 1\frac{1}{\pi b_{\frac{p+2}{2}}}\hfill & & \frac{\alpha ^{}}{1\nu }\pi \beta _{\frac{p+2}{2}}& (>0)\\ \nu _{\stackrel{~}{I}}\frac{1}{\pi b_{\stackrel{~}{I}}}\hfill & & \frac{\alpha ^{}}{\nu _{\stackrel{~}{I}}}\pi \beta _{\stackrel{~}{I}}& (>0)\end{array},$$
(5.12)
in the zero slope limit. The first term in eq. (5.11) comes from the tachyon state exchange. Here the vertices derived from three point amplitude emerge. From the $`t`$-channel diagram in Fig.4, we find that these vertices depend on momenta in a proper way. It is worth noting that combining the exponential factor in eq. (5.11) with the multiplicative factor $`\mathrm{exp}\left(𝒞_3(\nu _I)+𝒞_4(\nu _I)\right)`$, we obtain a gaussian damping factor in the $`x^{p+1},\mathrm{},x^p^{}`$ directions in the zero slope limit,
$$\mathrm{exp}\left[\frac{\pi }{2}\underset{I}{}\left|\beta _I\right|\left((k_3+k_4)\genfrac{}{}{0pt}{}{}{(p,p^{})}(k_3+k_4)\right)_I\right].$$
(5.13)
Next we focus on the behavior of the integrand near $`x=1`$ from which we can read off the $`s`$-channel poles. From Fig.3, the $`s`$-channel poles come from the propagation of the $`p^{}`$-$`p^{}`$ open string. In a similar way to the $`t`$-channel, by expanding the integrand around $`x=1`$ and integrating it near $`x=1`$: $`{\displaystyle _{1\delta }^1}𝑑x\mathrm{}`$ , we find that $`s`$-channel poles correspond to the $`p^{}`$-$`p^{}`$ open string spectrum. In particular, aside from the multiplicative factor and momentum conserving delta functions, the massless pole turns out to be
$`{\displaystyle \frac{1}{2s}}[\{(k_2k_1){\displaystyle \genfrac{}{}{0pt}{}{}{\left(p\right)}}\zeta _3i(k_3+k_4){\displaystyle \genfrac{}{}{0pt}{}{}{(p,p^{})}}J\zeta _3\}k_3{\displaystyle \genfrac{}{}{0pt}{}{}{\left(p^{}\right)}}\zeta _4`$
$`k_4{\displaystyle \genfrac{}{}{0pt}{}{}{\left(p^{}\right)}}\zeta _3\left\{(k_2k_1){\displaystyle \genfrac{}{}{0pt}{}{}{\left(p\right)}}\zeta _4i(k_3+k_4){\displaystyle \genfrac{}{}{0pt}{}{}{(p,p^{})}}J\zeta _4\right\}`$
$`2{\displaystyle \underset{I}{}}\nu _I\left(k_3{\displaystyle \genfrac{}{}{0pt}{}{\times }{(p,p^{})}}k_4\right)_I\zeta _3{\displaystyle \genfrac{}{}{0pt}{}{}{\left(p^{}\right)}}\zeta _42k_3{\displaystyle \genfrac{}{}{0pt}{}{}{\left(p^{}\right)}}k_4{\displaystyle \underset{I}{}}\nu _I\left(\zeta _3{\displaystyle \genfrac{}{}{0pt}{}{\times }{(p,p^{})}}\zeta _4\right)_I`$
$`+\{t+m_T^2+k_3{\displaystyle \genfrac{}{}{0pt}{}{}{(p,p^{})}}k_4{\displaystyle \underset{I}{}}(12\nu _I)\left(k_3{\displaystyle \genfrac{}{}{0pt}{}{\times }{(p,p^{})}}k_4\right)_I\}\zeta _3{\displaystyle \genfrac{}{}{0pt}{}{}{\left(p^{}\right)}}\zeta _4]`$
$`\times \mathrm{exp}\left[2\alpha ^{}{\displaystyle \underset{I}{}}\left(k_3{\displaystyle \genfrac{}{}{0pt}{}{}{(p,p^{})}}k_4\right)_I\left\{\gamma +{\displaystyle \frac{1}{2}}\left(𝝍(\nu _I)+𝝍(1\nu _I)\right)\right\}\right]`$
$`+(k_3k_4;\zeta _3\zeta _4),`$ (5.14)
where
$$s(k_3+k_4)\genfrac{}{}{0pt}{}{}{\left(p^{}\right)}(k_3+k_4)=2k_3\genfrac{}{}{0pt}{}{}{\left(p^{}\right)}k_4.$$
(5.15)
Here by using eqs. (2.43), (A.1) and (A.5) we have expanded the hypergeometric function $``$ around $`x=1`$ as
$`(\nu ;{\displaystyle \frac{1}{x}})={\displaystyle \frac{\pi }{2}}b\mathrm{ln}(1x)`$
$`+{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(1+\nu )_m}{m!}}(1x)^m{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(1\nu )_n}{n!}}\left\{\psi (n+1)\psi (n+1\nu )\right\}(1x)^n,`$
$`(\nu ;x)={\displaystyle \frac{\pi }{2}}b\mathrm{ln}(1x)`$
$`+{\displaystyle \underset{m=0}{}}{\displaystyle \frac{(\nu )_m}{m!}}(1x)^m{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(\nu )_n}{n!}}\left\{\psi (n+1)\psi (n+\nu )\right\}(1x)^n.`$ (5.16)
We find that in our result eq. (5.14) the first two terms are in accordance with the vertex seen at the three point amplitude (IV.). This vertex shows up again with proper momentum dependence (see Fig.4) as well as in the $`t`$-channel pole corresponding to the tachyon exchange. Combined with the multiplicative factor $`\mathrm{exp}\left(𝒞_3(\nu _I)+𝒞_4(\nu )\right)`$, the exponential factor in eq. (5.14) gives us a gaussian damping factor,
$$\mathrm{exp}\left[\alpha ^{}\underset{I}{}\left\{\gamma +\frac{1}{2}\left(𝝍(\nu _I)+𝝍(1\nu _I)\right)\right\}\left((k_3+k_4)\genfrac{}{}{0pt}{}{}{(p,p^{})}(k_3+k_4)\right)_I\right].$$
(5.17)
It is noteworthy that in the zero slope limit this gaussian damping factor turns out to be the same as that of $`t`$-channel (eq. (5.13)).
While we do not try to derive here the complete action of the low energy noncommutative field theory in $`p^{}+1`$ dimensions, it is still possible to exhibit the interactions which reproduce the parts of the amplitudes in the zero slope limit which are expressible in terms of the inner product with respect to the open string metric. We find that this part of the action is
$`S=S_0+S_1,`$
$`\text{with}S_0={\displaystyle \frac{1}{g_{YM}^{2}}}{\displaystyle d^{p^{}+1}x\sqrt{G}\left\{\left(D_\mu \mathrm{\Phi }\right)^{}\left(D^\mu \mathrm{\Phi }\right)m^2\mathrm{\Phi }^{}\mathrm{\Phi }\frac{1}{4}F_{MN}F^{MN}\right\}},`$
$`S_1={\displaystyle \frac{1}{2g_{YM}^{2}}}{\displaystyle d^{p^{}+1}x\sqrt{G}\mathrm{\Phi }^{}F_{MN}J^{MN}\mathrm{\Phi }},`$ (5.18)
where
$`D_\mu \mathrm{\Phi }=_\mu \mathrm{\Phi }iA_\mu \mathrm{\Phi },\left(D_\mu \mathrm{\Phi }\right)^{}=_\mu \mathrm{\Phi }^{}+i\mathrm{\Phi }^{}A_\mu ,`$
$`F_{MN}=_MA_N_NA_Mi[A_M,A_N]_{},[A_M,A_N]_{}=A_MA_NA_NA_M,`$ (5.19)
the $``$ product of two functions $`f`$ and $`g`$ is given by
$$f(x)g(x)=e^{\frac{i}{2}\theta ^{\mu \rho }\frac{}{y^\mu }\frac{}{z^\rho }}f(y)g(z)|_{y,zx},$$
(5.20)
and $`g_{YM}`$ is the effective Yang-Mills coupling defined by using the open sting coupling $`G_s`$ and that of the closed string $`g_s`$ as
$$\frac{1}{g_{YM}^{2}}=\frac{\left(\alpha ^{}\right)^{\frac{3p^{}}{2}}}{\left(2\pi \right)^{p^{}2}G_s}=\frac{\left(\alpha ^{}\right)^{\frac{3p^{}}{2}}}{\left(2\pi \right)^{p^{}2}g_s}\left(\frac{det(g+2\pi \alpha ^{}B)}{detG}\right)^{\frac{1}{2}}.$$
(5.21)
In eq. (5.18) we have determined the seagull interaction corresponding to the non-pole term by invoking the noncommutative $`U(1)`$ invariance.
Let us finally discuss the gaussian damping factor eq. (5.13) which have originated from the exponential multiplicative factor eq. (4.8) and the lowest modes in the hypergeometric function $``$. Recall that there is no momentum conservation for the $`x^{p+1},\mathrm{},x^p^{}`$-directions and that the tachyon momenta $`k_1`$ and $`k_2`$ are constrained to lie on the $`x^0,\mathrm{},x^p`$-directions. Without the gaussian damping factor, our picture would be that $`(N2)`$ incident noncommutative $`U(1)`$ photons travel freely in the $`x^{p+1},\mathrm{},x^p^{}`$-directions until they get stopped by the D$`p`$-brane. The actual spacetime picture which we have exhibited here is that the lowest mode of the $`p`$-$`p^{}`$ open string develops a physical scale $`\sqrt{|\beta _I|}`$ and that this mode creates a cloud around the D$`p`$-brane in the zero slope limit. The noncommutative $`U(1)`$ photons get decelerated by the presence of this cloud, which is reflected in our damping factor eq. (5.13). The mean free paths will be measured by $`\sqrt{|\beta _I|}`$.
In this situation, the tachyon field in these directions should be expanded by the coherent states $`\left\{x^{p+1},\mathrm{},x^p^{}|\nu _I\right\}`$ associated with the would-be zero modes $`\alpha _{1\nu _I}^I`$ (or $`\overline{\alpha }_{\nu _I}^{\overline{I}}`$). On this basis, the complete analysis of low-lying states and $`\nu _I`$ dependent interactions obtained from the residual parts of the amplitudes will lead to the full-fledged form of the tachyon-vector interactions in these directions. The appearance of the coherent states here suggests that the fields which have originated from the $`p`$-$`p^{}`$ open string should support noncommutative solitons on the D$`p^{}`$-brane worldvolume which has recently been found in .
## Acknowledgements:
We are grateful to Professor E. Date for helpful discussions on this subject.
## Appendix More on the two-point function at the worldsheet boundary
The radius of convergence of the hypergeometric series $`F(a,b;c;z)`$ is unity. Thus we have evaluated the hypergeometric series on its convergent circle in eq. (2.45) in deriving the noncommutativity term (2.46). In this appendix we will give another derivation of eq. (2.45) to verify the noncommutativity term (2.46) and the two-point function (2.47).
Let us focus on the relation,
$$(\nu ;z)=\mathrm{ln}(1z)+z^\nu \underset{n=0}{\overset{\mathrm{}}{}}\frac{(\nu )_n}{n!}\left\{𝝍(n+1)𝝍(n+\nu )\right\}(1z)^n.$$
(A.1)
One can obtain this relation by using eq. (2.43) and a formula for the hypergeometric function ,
$`F(a,b;a+b;z)={\displaystyle \frac{\mathrm{\Gamma }(a+b)}{\mathrm{\Gamma }(a)\mathrm{\Gamma }(b)}}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(a)_n(b)_n}{\left(n!\right)^2}}[2𝝍(n+1)𝝍(a+n)`$ (A.2)
$`𝝍(b+n)\mathrm{ln}(1z)](1z)^n,`$
which is derived from the following relation by putting $`c=a+b+\delta `$ and by taking the limit of $`\delta 0`$:
$`F(a,b;c;z)={\displaystyle \frac{\mathrm{\Gamma }(c)\mathrm{\Gamma }(cab)}{\mathrm{\Gamma }(ca)\mathrm{\Gamma }(cb)}}F(a,b;a+bc+1;1z)`$ (A.3)
$`+{\displaystyle \frac{\mathrm{\Gamma }(c)\mathrm{\Gamma }(a+bc)}{\mathrm{\Gamma }(a)\mathrm{\Gamma }(b)}}(1z)^{cab}F(ca,cb;cab+1;1z).`$
From eq. (A.1), one can find that $`\underset{z1}{lim}\left\{(1\nu _I;{\displaystyle \frac{1}{z}})(\nu _I;z)\right\}`$ is sensitive to the way of taking the limit. In the original expression (2.41) of the two point function, however, the way of sending $`z1`$ is fixed in a definite way on the real axis because of the step function in front of each hypergeometric function. The constant noncommutativity term (2.46) should be more precisely described as
$`{\displaystyle \frac{4}{\epsilon }}{\displaystyle \frac{\delta ^{I\overline{J}}}{1+b_I^2}}\left\{\underset{\tau _1\tau _2+0}{lim}(1\nu _I;{\displaystyle \frac{e^{\tau _2}}{e^{\tau _1}}})\underset{\tau _1\tau _20}{lim}(\nu _I;{\displaystyle \frac{e^{\tau _1}}{e^{\tau _2}}})\right\}`$ (A.4)
$`={\displaystyle \frac{4}{\epsilon }}{\displaystyle \frac{\delta ^{I\overline{J}}}{1+b_I^2}}\underset{\tau _1\tau _2}{lim}[\{\mathrm{ln}|1{\displaystyle \frac{e^{\tau _2}}{e^{\tau _1}}}|+\left({\displaystyle \frac{e^{\tau _2}}{e^{\tau _1}}}\right)^{1\nu _I}(𝝍(1)𝝍(1\nu _I))\}`$
$`\left\{\mathrm{ln}\left|1{\displaystyle \frac{e^{\tau _1}}{e^{\tau _2}}}\right|+\left({\displaystyle \frac{e^{\tau _1}}{e^{\tau _2}}}\right)^{\nu _I}\left(𝝍(1)𝝍(\nu _I)\right)\right\}`$
$`={\displaystyle \frac{4}{\epsilon }}{\displaystyle \frac{\delta ^{I\overline{J}}}{1+b_I^2}}\pi b_I,`$
where we have used a relation for the digamma function,
$$\pi b_I=\pi \mathrm{cot}(\pi \nu _I)=𝝍(\nu _I)𝝍(1\nu _I).$$
(A.5)
Thus we obtained the same noncommutativity term as eq. (2.46) through more careful treatment of the hypergeometric functions and this provides another verification to the two-point function (2.47).
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# 1 The difference of the critical temperature 𝑇_𝐶 and the sphaleron freezing out temperature 𝑇_{𝑜𝑢𝑡} determined by Γ_{𝑠𝑝ℎ}(𝑇_{𝑜𝑢𝑡})=𝐻(𝑇_{𝑜𝑢𝑡}) in dependence of 𝑇_𝐶.
Neutrino Majorana Mass and Baryon Number of the Universe below the Electroweak Symmetry breaking Scale
H.V. Klapdor-Kleingrothaus , St. Kolb and U. Sarkar<sup>†,‡</sup>
Max-Planck-Institut für Kernphysik, P.O. 10 39 80, D-69029 Heidelberg, Germany
Physical Research Laboratory, Ahmedabad, 380 009 India
## Abstract
If the neutrino is Majorana type and the electroweak phase transition is second or weak first order, neutrino-induced interactions together with sphaleron transitions have the potential to erase a previously generated baryon asymmetry of the universe. Taking correctly into account the evolution of the vacuum expectation of the Higgs field the effective light neutrino masses are constrained to be lighter than $`𝒪(10MeV)`$, while the effective heavy masses are constrained to be heavier than $`𝒪(10^7GeV)`$.
Recent experiments seem to indicate that the neutrino is massive . From a model-building point of view the most natural structure of the neutrino mass matrix contains both Lepton-number- ($`L`$-) conserving Dirac- and $`L`$-violating Majorana-type entries (see e.g. ). In principle observable consequences are $`L`$-violating processes such as neutrinoless double beta ($`0\nu \beta \beta `$) decay and $`L`$-violating lepton-gauge boson scattering (inverse $`0\nu \beta \beta `$ decay).
On the other hand, in the early universe $`L`$ violation together with sphaleron-mediated transitions has the potential to create the Baryon number ($`B`$) of the universe (BAU) (Leptogenesis) or erase an existing BAU. The latter may be the case both above or below the electroweak symmetry breaking scale. This consideration would limit the amount of $`L`$ violation and hence give a bound on the Majorana mass of the neutrinos. In the following we will reconsider the limits on neutrino masses and point out that the evolution of the vacuum expectation value ($`vev`$) of the Higgs responsible for electroweak symmetry breaking weakens the existing bound by as large as three orders of magnitude.
In an extension of the standard model, the Majorana mass of the neutrinos can come from an effective dimension-5 operator
$`{\displaystyle \frac{\alpha _{ij}}{M}}(L_i^TC^1\tau _2\stackrel{}{\tau }L_j)(H^T\tau _2\stackrel{}{\tau }H)`$ (1)
where $`i,j=1,2,3`$ are generation indices, $`M`$ is the $`L`$-violating mass-scale, $`\alpha `$ is an effective coupling and $`L`$ ($`H`$) are $`SU(2)_L`$ lepton (Higgs) doublets. After the electroweak symmetry breaking, when the higgs doublet scalar acquires a $`vev`$, the neutrinos get a Majorana mass of the order of
$`m_\nu {\displaystyle \frac{\alpha _{ij}v(T)^2}{M}}`$ (2)
where $`v(T)`$ is the $`vev`$ of $`H`$. The condition that the associated $`L`$ violation should not wash out the primordial BAU then gives an upper bound on the Majorana mass of the neutrinos of the order of a few keV . In most of these approaches the analysis was simplified by assuming the $`vev`$ to be constant and the rate of the $`L`$-violating interactions was considered to be less than the expansion rate of the universe. However, above the critical temperature $`T_C`$ of the electroweak phase transition (EWPT) the $`vev`$ of $`H`$ is zero. Below $`T_C`$ the $`vev`$ starts growing. On the other hand soon the sphalerons freeze out and they cannot wash out the BAU any longer. Thus, during the period when the sphalerons wash out the BAU, the $`vev`$ may still be quite small. This weakens the upper bound on the Majorana mass of the neutrinos. In the following this argument will be discussed in more detail.
Consider the see-saw mechanism of neutrino masses . The $`L`$-violating mass scale will have its origin from integrating out heavy $`SU(2)_L`$ singlet neutrinos. The neutrino mass-matrix has the general structure
$`^\nu =\left(\begin{array}{cc}0& (m^D)^T\\ m^D& M\end{array}\right)`$ (4)
where $`m^Dv(T)`$ is the Dirac mass matrix coupling $`SU(2)_L`$ doublet neutrinos $`\nu `$ to $`SU(2)_L`$ singlet neutrinos $`N`$, whereas $`M`$ is a Majorana mass matrix for the $`N`$’s. Below the EWS the Higgs field aquires a (temperature dependent) vacuum expectation value $`v(T)`$ and so that $`\nu `$ and $`N`$ mix with resulting masses
$`m_i={\displaystyle \underset{j}{}}U_{ij}_{ij},UU^{}=1`$ (5)
If the scale of $`M`$ is much larger than that of $`m^D`$ the diagonal mass-matrix consists of two blocks of light and heavy masses (see-saw mechanism)
$`m_{light}(m^D)^TM^1m^D,m_{heavy}M.`$ (6)
The off-diagonal blocks of the mixing matrix $`U`$ are approximately $`(m^D)^{}(M^1)^{}`$ and $`M^1m^D`$.
In the triplet higgs model one introduces a triplet higgs scalar $`\xi `$ with mass $`M`$. The couplings of the triplet higgs breaks lepton number explicitly at the scale $`M`$, but the $`vev`$ of the triplet higgs gets a see-saw contribution of amount
$`\xi {\displaystyle \frac{H^2}{M}}.`$ (7)
The direct coupling of the triplet higgs with the two neutrinos then give a Majorana mass to the neutrinos.
In the early universe neutrinos give rise to $`L`$-violating processes such as ($`i`$,$`j`$ are generation indices)
$`e_i^\pm e_j^\pm W^\pm W^\pm .`$ (8)
The masses $`m_{ij}^M`$ give rise to $`L`$-violating processes such as neutrinoless double beta decay ($`0\nu \beta \beta `$) (for an overview see e.g. ). For a linear collider this process has been studied in .
In the early universe $`m_{ij}^M`$ induced processes have the potential to erase the obeserved asymmetry in the baryon- ($`B`$) and antibaryon-number of the universe (BAU) (see e.g. ). This is due to the fact that as long as sphaleron transitions are in thermal equilibrium $`B`$ and $`L`$ are both proportional to $`(BL)`$ . Hence, if $`L`$ is erased by an $`L`$-violating process and sphalerons are still operative, $`B`$ is erased as well. If the electroweak phase transition (EWPT) with an associated critical temperature $`T_C`$ is strong first order sphalerons are never in thermal equilibrium and $`L`$-violating processes do not affect $`B`$ below $`T_C`$. On the other hand, if the EWPT is second or weak first order there is a period
$`T_{out}<T<T_C`$ (9)
($`T_{out}`$ denotes the sphaleron freezing-out temperature) during which $`L`$-violating processes have the potential to erase $`L`$ and consequently $`B`$. Hence, if the EWPT is second or weak first order the requirement that a preexisting BAU should not be washed out poses a limit on the amount of $`L`$ violation. In the case of a Majorana neutrino in previous works an estimated limit $`m_{ij}^M\text{ }\stackrel{<}{}\text{ }20keV`$ has been obtained. In this note it will be argued that this bound is in fact three orders of magnitude less stringent if the temperature dependence of the $`vev`$ is correctly taken into account as has been done recently for the case of $`L`$-violating sneutrinos in .
The temperature dependence of the $`vev`$ for a second or weak first order phase transition is given by
$`v(T)v(T=0)(1T^2/T_C^2)^{1/2},v(T=0)=246GeV`$ (10)
and the the sphaleron rate in the broken phase is
$`\mathrm{\Gamma }_{Sph}(T)2.810^5T^4\kappa \left({\displaystyle \frac{\alpha _W}{4\pi }}\right)^4\left({\displaystyle \frac{2m_W(T)}{\alpha _WT}}\right)^7\mathrm{exp}\left({\displaystyle \frac{E_{sp}(T)}{T}}\right)`$ (11)
where
$`m_W(T)={\displaystyle \frac{1}{2}}g_2v(T),`$ (12)
the free energy of the sphaleron configuration is given by
$`E_{Sph}(T)={\displaystyle \frac{2m_W(T)}{\alpha _W}}B\left({\displaystyle \frac{m_H}{m_W}}\right),`$ (13)
$`B(0)=1.52,B(\mathrm{})=2.72`$ and $`\kappa `$=exp(-3.6) . As usual $`T_{out}`$ is determined by the condition
$`\mathrm{\Gamma }_{sph}(T_{out})=H(T_{out})=1.7\sqrt{g_{}}{\displaystyle \frac{T_{out}^2}{M_{Pl}}}`$ (14)
where $`M_{Pl}10^{19}GeV`$ is the Planck scale and $`g_{}100`$ in the Standard Model. Lattice simulations suggest that for a Higgs mass of around $`m_H70GeV`$ $`T_C150GeV`$ and higher for larger values of $`m_H`$ . For our phenomenological purposes $`T_C`$ will be varied between $`50GeV`$ and $`250GeV`$. The temperature range eq. (9) is plotted in figure 1. It is smaller than $`1GeV`$ for $`T_C\text{ }\stackrel{<}{}\text{ }100GeV`$ but of order $`𝒪(10GeV)`$ for $`T_C\text{ }\stackrel{>}{}\text{ }200GeV`$.
Relevant processes for depleting a pre-existing $`L`$ number during the epoch (9) are $`L`$-violating $`22`$ scatterings $`W^\pm W^\pm e_i^\pm e_j^\pm `$, $`W^\pm e_i^{}W^{}e_j^\pm `$ and gauge boson decays. The depletion of an initial $`L`$ number $`L_i`$ is described by (see for example )
$`L(z)=L_i\mathrm{exp}[{\displaystyle \underset{z_c}{\overset{z_{out}}{}}}dz^{}z^{}[g_{}{\displaystyle \frac{n_{\stackrel{~}{\nu }}}{s}}\mathrm{\Gamma }_D(z^{})+n_\gamma \sigma |v|)]/H(T=m_{\stackrel{~}{\nu }})]`$ (15)
Working with the Boltzmann equation will then give us the amount of residual asymmetry after the spheleron transitions have frozen out. Unlike earlier works, where the interaction rate has been compared with the expansion rate of the universe, we consider the condition for erasure of the primordial BAU is that the asymmetry depletes by a factor of at least 10. In most of the cases when we get the bound, the depletion is more than two orders of magnitude.
For the see-saw mechanism case the thermally averaged contribution of the light neutrino states with masses $`m_kT_C`$ may be approximated by
$`\sigma |v|_{ij}{\displaystyle \frac{\alpha _W^2m_{ij}^2}{T^4}},m_{ij}={\displaystyle \underset{k}{}}U_{ik}U_{jk}m_k`$ (16)
and the contribution of the heavy states with masses $`M_nT_C`$ is
$`\sigma |v|{\displaystyle \frac{\alpha _W^2}{M_{mn}^2}},{\displaystyle \frac{1}{M_{ij}}}={\displaystyle \frac{_nU_{in}U_{jn}}{M_n}}`$ (17)
where $`\alpha _W`$ is the weak coupling constant. Compared to the zero-temperature case both cross-sections are suppressed by a factor
$`{\displaystyle \frac{v(T)^2}{v(T=0)^2}}=\left(1{\displaystyle \frac{T^2}{T_C^2}}\right)^2,`$ (18)
since for the light states $`m_kv(T)^2`$ and for the heavy states $`U_{in}U_{jn}v(T)^2`$, see above. The bounds on the quantities $`m_k`$ and $`M_n`$ are displayed in figure 2. For plausible values of the critical temperature $`T_C\text{ }\stackrel{>}{}\text{ }150GeV`$ the bound on the light states is of order $`m_{ij}\text{ }\stackrel{<}{}\text{ }𝒪(10MeV)`$ while the bound on the heavy states is of order $`M_{ij}\text{ }\stackrel{>}{}\text{ }10^7GeV`$.
For the triplet higgs mechanism, the bound on the mass of the neutrino comes out to be the same as the bound on the light neutrino state of the see-saw mechanism. In this case also the suppression is given by equation (18), since the neutrino mass is again proportional to $`m_\nu v(T)^2`$. Thus even in this case the bound comes out to be around 10 MeV. Given the generality of the dimension-5 operators for the Majorana neutrino mass, one may conclude that in all models of neutrino masses this bound is valid. Previous estimates of the light neutrino masses have been given in for the case eq. (1) and it has been argued that for every entry of the corresponding mass matrix the bound from the BAU is of order $`m_{ij}\text{ }\stackrel{<}{}\text{ }10keV`$, that is three orders of magnitude more stringent than if the evolution of the $`vev`$ is taken into account.
In summary, we included the effect of evolution of the higgs $`vev`$ and solved the Boltzmann equation in estimating the bound on the neutrino masses coming from the erasure of the baryon asymmetry of the universe. This makes the bounds three orders of magnitude weaker than the one obtained from earlier naive estimates.
Acknowledgement
U.S. wants to thank Max-Planck-Institut für Kernphysik for hospitality.
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# The exchange energy in the local Airy gas approximation
## Abstract
The Airy gas model of the edge electron gas is used to construct an exchange-energy functional which is an alternative to those obtained in the local density and generalized gradient approximations. Test calculations for rare gas atoms, molecules, solids and surfaces show that the Airy gas functional performs better than the local density approximation in all cases and better than the generalized gradient approximation for solids and surfaces.
Since the pioneering papers on density functional theory (DFT) there has been a constant search for exchange-correlation functionals of chemical accuracy. This includes the works on the generalized gradient approximation (GGA) which are dedicated efforts to construct local functionals for inhomogeneous systems ranging from atoms to solids based on the uniform electron gas, i.e., the local density approximation (LDA), and density gradient corrections, as well as the development of a number gradient level, semiempirical functionals . The GGA functionals have had a considerable impact upon the fields of quantum chemistry and solid state physics because they reduce the LDA overbinding and generally improve the calculated properties, relative to experiments, of molecules and bulk solids . However, they perform less well for the bulk properties of late transition metals and semiconductors , and the underestimate of the exchange energies of surfaces as well as the overestimate of the dissociation energies of the multiply bonded molecules indicate the necessity to go beyond the gradient level approximations and develop functionals that depend upon other inhomogeneity parameters, e.g., higher derivatives of the charge density or the Kohn-Sham kinetic energy density. One step in this direction is the meta generalized gradient approximation (meta-GGA) of Perdew, Kurth, Zupan, and Blaha which proves highly promising for both finite and extended systems .
In the present work we introduce and apply a new gradient level exchange energy functional based on the concept of the edge electron gas . Besides the formal interest in the development of density based, orbital independent functionals there are several reasons why in applications of DFT the focus is on the approximate, local exchange-correlation schemes. Within the Kohn-Sham approach to DFT the Kohn-Sham exchange energy may be determined exactly and as demostrated recently so may the corresponding local exchange potential. However, the exact Kohn-Sham exchange formalism is non-local and orbital-based, i.e. both the exchange energy and potential are highly complicated non-local functionals of the Kohn-Sham orbitals. In consequence, the application of exact exchange is computationally demanding. Furthermore, when exchange is treated exactly the error cancellations between the exchange and correlation energies on which all approximate schemes depend are lost owing to the poor description of correlation effects and, as a result, the total energies worsen . For these reasons the exact Kohn-Sham exchange energy has only been used in practice in connection with semi-empirical, hybrid approximations .
The concept of the edge electron gas was put forward by Kohn and Mattsson as an appropriate basis for the treatment of systems with edge surface outside of which all Kohn-Sham orbitals decay exponentially. Its simplest realization, the Airy gas model, is based on the linear potential approximation and may serve as the starting point for the construction of functionals which are alternative to the GGA. The Airy gas model has recently been used to construct an explicit kinetic energy functional for inhomogeneous systems which for atoms and surfaces has the accuracy of functionals based on a second order gradient expansion.
Here we have taken the exchange energy of the Airy gas model derived by Kohn and Mattsson and cast it in a form amenable to a simple, accurate parametrization. The procedure may be viewed as local mapping of the real system described by its density and scaled gradient onto the Airy gas model and represents one possible solution to the joining of the interior to the edge regions. The parametrized functional which we refer to as the local Airy gas (LAG) functional is tested in calculations of the exchange energies of rare gas atoms and of metallic surfaces within the jellium model where the exact results are known . In addition, we apply the LAG exchange functional in conjunction with the LDA for correlation in calculations of the molecular binding energies and bulk properties of solids.
The present LAG exchange functional has a number of advantages over previous GGA functionals: i) it explicitly includes the properties of the edge region where much interesting physics occurs, ii) its accuracy may be systematically improved by including higher order expansions of the effective potential of the model system, and iii) the resulting exchange-energy functional is as simple and well-defined as that of the standard LDA. i.e., it has no adjustable parameters.
The starting point for the Airy gas exchange energy functional is the potential
$`v_{eff}(z)=\{\begin{array}{ccc}\mathrm{}\hfill & \text{for}\hfill & zL\hfill \\ Fz\hfill & \text{for}\hfill & L<z<\mathrm{}\hfill \end{array},`$ (3)
which is linear in $`z`$, independent of $`x`$ and $`y`$, and has a hard wall at $`L`$ far from the electronic edge at $`z=0`$. The slope of the effective potential $`F=dv_{eff}/dz`$ leads to a characteristic length scale
$$l\left(\frac{\mathrm{}^2}{2mF}\right)^{1/3},$$
(4)
and the electron and exchange-energy densities are then given by
$$n(z)=l^3n(\zeta ),$$
(5)
and
$$\epsilon _x(z)=\frac{e^2}{2}l^4\epsilon _x(\zeta ),$$
(6)
where $`\zeta =z/l`$,
$$n(\zeta )=\frac{1}{2\pi }_0^{\mathrm{}}Ai^2(\zeta +\zeta ^{})\zeta ^{}𝑑\zeta ^{},$$
(7)
and
$`\epsilon _x(\zeta )`$ $`=`$ $`{\displaystyle \frac{1}{\pi }}{\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle _0^{\mathrm{}}}{\displaystyle _0^{\mathrm{}}}Ai(\zeta +ϵ)Ai(\zeta ^{}+ϵ)Ai(\zeta +ϵ^{})Ai(\zeta ^{}+ϵ^{})`$ (8)
$`\times `$ $`|\zeta ^{}\zeta |^3g(\sqrt{ϵ}|\zeta ^{}\zeta |,\sqrt{ϵ^{}}|\zeta ^{}\zeta |)d\zeta ^{}dϵdϵ^{}.`$ (9)
A contour plot of the universal function $`g(s,s^{})`$ may be found in Ref. . The exchange energy (6) may be written in the form
$$\epsilon _x(z)=\epsilon _x^{LDA}(z)F_x[s(z)],$$
(10)
where $`\epsilon _x^{LDA}(z)`$ is the exchange energy density of the uniform electron gas. The enhancement function
$$F_x(\zeta )\frac{2}{3}\left(\frac{\pi }{3}\right)^{1/3}\frac{\epsilon _x(\zeta )}{n^{4/3}(\zeta )},$$
(11)
is the unique function $`(F_x(\zeta ),s(\zeta ))`$ of the scaled gradient
$$s(\zeta )\frac{n^{}(\zeta )}{2(3\pi ^2)^{1/3}n^{4/3}(\zeta )},$$
(12)
plotted in Fig. 1. For comparison we also present results obtained by the GGA of Perdew, Burke, and Ernzerhof (PBE) as defined in Ref. , and the second order gradient expansion (GEA) . It follows from the figure that the exchange density (10) in the low gradient limit of the Airy gas model reduces to $`\epsilon _x^{LDA}(z)`$ as it should. In the large gradient limit $`\epsilon _x(\zeta )\frac{n(\zeta )}{2\zeta }`$, Ref. , and similar to the case of the kinetic energy density we use the properties of the Airy gas to find the following explicit asymptotic expression
$$\epsilon _x[n(z)]\frac{e^2}{2}\frac{n(z)}{2}\left[n(z)\frac{^3n(z)}{z^3}\frac{n(z)}{z}\frac{^2n(z)}{z^2}\right]\left[\frac{n(z)}{z}\right]^2,$$
(13)
in terms of the density and its derivatives.
The density of the exchange energy per electron of the Airy gas is plotted as a function of the distance $`z`$ from the electronic edge in Fig. 2. It is seen that the large gradient expression (13) is accurate for $`z/l>1.4`$ corresponding to $`s>0.5`$. It is also seen that neither the LDA nor the PBE GGA leads to the correct behaviour near and beyond the electronic edge at $`z=0`$.
The scaled gradient is conserved when going from the real electron gas to the Airy gas model and therefore the enhancement function $`F_x(s)`$ parametrized, for instance, in a modified Becke form
$$F_x^{LAG}(s)=1+\beta \frac{s^\alpha }{(1+\gamma s^\alpha )^\delta },$$
(14)
which includes the proper LDA limit, may be used to obtain the exchange energy density of the real electron gas from the local, scaled gradient $`s[n(z)]`$. For $`\alpha =2.626712,\beta =0.041106,\gamma =0.092070`$, and $`\delta =0.657946`$ we find that the local deviation between the exact result (11) and the parametrized form (14) integrated over the range $`0<s<20`$ is less than $`0.3\%`$. We note that the present parametrization, being an overall fit, does not reduce to the GEA in the low gradient limit. In contrast to the case of the kinetic energy we have not been able to find an explicit, analytical expression for the exchange energy for small $`s`$ values, and to establish the behaviour numerically has not been attempted because the $`s0`$ limit is reached only at $`z\mathrm{}`$ as seen in Fig. 2. The exact behaviour of the LAG exchange functional at $`s0`$ is therefore not know at present.
In the following we report the results of applying the LAG exchange functional to four test systems: i) rare gas atoms, ii) diatomic molecules, iii) jellium surfaces, and iv) solids. In all cases the total energy is calculated using self-consistent LDA densities. For molecules and solids the LAG exchange energy is combined with the LDA correlation energy , since correlation effects has not been worked out in the Airy gas model. The motivation of this combination is given in terms of the enhancement function over the local exchange energy , defined as $`F_{xc}(s)ϵ_{xc}[n]/ϵ^{LDA}(n)`$, where $`ϵ_{xc}[n]`$ denotes the exchange-correlation energy density.
Most of the currently applied approximate density functionals are based on error cancellations between the exchange and correlation energies . For physically interesting densities this cancellation leads to $`F_{xc}(s)`$ with negligible slope up to $`s1`$. Plots of the enhancement function over the local exchange energy for gradient level and meta-GGA approximations can be found in Refs. . In the present LAG exchange plus LDA correlation scheme this function becames $`F_{xc}^{LAG}(s)=F_x^{LAG}(s)+ϵ_c^{LDA}(n)/ϵ_x^{LDA}(n)`$, where $`ϵ_c^{LDA}(n)`$ is the correlation energy density of the uniform electron gas. Thus, the $`F_{xc}^{LAG}(s)`$ is determined only by the LAG enhancement function (14), which, for $`s<1`$, is a slowly increasing function of $`s`$. Therefore, we expect the present exchange-correlation scheme to preserve the excellent cancellation properties of the LDA and PBE GGA, and, at the same time, to bring the calculated properties in closer agreement with experiment than conventional LDA.
For the rare gas atoms included in Table I the GEA, PBE, and LAG functionals yield exchange energies which are, on the average, 6.4 %, 8.5 %, and 1.8 %, respectively, larger than those obtained in the LDA. The PBE values are in very good agreement with the exact Kohn-Sham results , which are given relative to the LDA energies in the last column of the table. The LAG approximation represents only a minor improvement relative to the LDA total atomic exchange energies.
The effect of the gradient correction to the LDA atomization energies for a few selected diatomic molecules is shown in Table II which also includes the relative difference between the LDA results and experimental data . Here, the LDA charge densities for the molecules have been generated using the full charge density (FCD) technique in conjunction with the exact muffin-tin orbital method (EMTO) . It is seen that the LAG approximation (i.e. LAG exchange and LDA correlation energy) and PBE GGA have comparable accuracy: Both functionals reduce the LDA overbinding, and yield atomization energies which are, on the average, 16.8 % (PBE) and 16.2 % (LAG) smaller than the LDA values.
In Fig. 3 we compare four exchange functionals applied to the jellium model of metallic surfaces . The fact that for a given $`r_s`$-value the exchange energies become increasingly negative in the order LDA, LAG, GEA, and PBE is a simple consequence of the enhancement functions shown in Fig. 1 and in agreement with the observation that the GGA significantly underestimate surface exchange energies . We note that LAG approximation represents an improvement over both the LDA and PBE and vary less with $`r_s`$ than either of the other two approximations.
As a final test of the LAG approximation we have calculated the atomic volumes and bulk moduli of several metals and semiconductors in their observed low temperature crystal structures by means of the FCD-EMTO method . The results for the equilibrium atomic radii are plotted in Fig. 4. For some selected metals and semiconductors, for which accurate LDA, PBE, and meta-GGA results have been published , the atomic radii and bulk moduli are presented in Table III. The comparison of our LDA atomic radii for the transition metal series with those obtained by the full-potential linear muffin-tin orbital and linear augmented plane wave methods using the same LDA gives mean deviations of 0.33 %, 0.43 %, and 0.49 % for the $`3d,4d`$, and $`5d`$ series, respectively. For Li and Na the present LDA results agree within 0.07 % with the full-potential values from Ref. . We therefore expect that the results of the present LAG and PBE calculations shown in Fig. 4 will deviate less than 0.5 % from full-potential calculations. The mean deviations between the present atomic radii and bulk moduli listed in Table III and those of Ref. obtained using the linear augmented plane wave method are 0.20 % and 3.28 % for the LDA and 0.27 % and 3.26 % for the PBE functionals.
The LDA atomic radii shown in Fig. 4 deviate, on average, by 2.26 % from the experimental values , while those calculated in the LAG model and the PBE deviate by $`0.83\%`$ and $`0.91\%`$, respectively. Among the energy functionals considered in Table III the LAG is found to give the lowest mean deviations for both atomic radii and bulk moduli. We note that for these solids the LAG approximation achieves the accuracy of the recently developed meta-GGA .
We have used the Airy gas model of the edge electron gas that is equivalent to the linear potential approximation to develop an exchange energy functional which may serve as an alternative to the functionals based on the generalized gradient appoximation, e.g., PBE GGA. Test calculations for finite and extended systems show that the LAG approximation is more accurate than the local density approximation in all cases. While the LAG results for atoms are very close to the LDA results and, hence inferior to the PBE GGA results, its accuracy for the atomization energies of diatomic molecules is similar to that of the PBE GGA. In bulk systems the LAG results are, on average, closer to the experimental values than those obtained in the PBE GGA. These results are very satisfactory in view of the fact that the LAG exchange functional is derived solely form the properties of the Airy gas, and, hence, with no a priory assumptions concerning the exchange enhancement factor. In this sense it is truely ab initio but for the correlation effects which needs to be worked out in the Airy gas model.
Acknowledgments: We gratefully acknowledge interesting and fruitful discussions with Professor Walter Kohn. L.V. acknowledges the valuable observations by Professor J. P. Perdew and Dr. Á. Nagy. L.V. and B.J. are grateful to the Swedish Natural Science Research Council and the Swedish Foundation for Strategic Research for financial support. Part of this work was supported by the research project OTKA 23390 of the Hungarian Scientific Research Fund. The Center for Atomic-scale Materials Physics is sponsored by the Danish National Research Foundation.
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# Rotational properties of 252,253,254No. Influence of pairing correlations.
## 1 Introduction
The ground-state rotational band of <sup>254</sup>No has been observed recently ?-?). Its properties show two interesting features. First, this nucleus has a large deformation. The deduced quadrupole deformation is $`\beta _2`$ = 0.27 $`\pm `$ 0.02, in good agreement with previous calculations ?-?). Secondly, a strong robustness of shell-effects against rotation in such heavy nuclei is expected from the calculations as the rotational band is observed up to a spin of 22$`\mathrm{}`$.
Rotational properties constitute a good opportunity to test theoretical models in an extreme mass region, far away from the ones where effective forces are usually adjusted. They are sensitive to both shell effects (related to the particle-hole channel) and pairing correlations (related to the particle-particle channel), which allows to probe all the components of the force. The ability of the models and forces to describe rotational bands in a mass region just below the super-heavy elements can be a discriminating test for them to judge their predictive power in the super-heavy region.
The aim of the present study is to calculate the ground-state rotational bands of <sup>252</sup>No and <sup>254</sup>No which are experimentally known and of <sup>253</sup>No, for which experiments are planned. Special attention will be paid to the influence of pairing correlations on rotational properties. A Skyrme-type interaction is used in the particle-hole channel and two parametrizations are analyzed for the particle-particle channel : a zero-range force with and without density dependence. This latter point is particularly relevant as pairing correlations are crucial in reproducing dynamical moment of inertia and quasi-particle Routhians.
The paper is organized as follows. Section 2 recalls the theoretical framework of the calculations and the different forces used in the particle-particle channel. In section 3, the relevance in the A $``$ 250 mass region of the delta density-dependent pairing force, previously fitted on the rotational bands of super-deformed nuclei in the A $``$ 150 mass region , is discussed. The consistency of this fit is checked on odd-even mass differences which are more specifically related to pairing correlations. Then, we present results on the quasi-particle Routhians, single-particle energies and dynamical moment of inertia of the three Nobelium isotopes considered here. The fission barriers of <sup>254</sup>No at spin 0$`\mathrm{}`$ and 20$`\mathrm{}`$ are also discussed in this section. Finally, we draw our conclusions in section 5.
## 2 The theoretical framework
The method used to calculate rotational bands in No isotopes has been presented in Ref. and applied to the study of super-deformed rotational bands in the A $``$ 150 and 190 regions ?-?). It is based on the self-consistent cranked Hartree-Fock-Bogolyubov method, with an approximate particle number projection by the Lipkin-Nogami prescription (HFBLN). This method is well suited to describe odd nuclei with self-consistent blocking.
In the particle-hole channel, we use a two-body force of the Skyrme-type, SLy4, which has been adjusted to reproduce also the characteristics of the infinite neutron matter and, consequently, should have good isospin properties . This force has been shown to describe satisfactorily nuclear properties for which it had not been adjusted such as super-deformed rotational bands and the structure and decay of super-heavy elements . It seems therefore an adequate choice for our study of No isotopes.
In the $`T`$ = 1 particle-particle channel, we use two types of force for n-n and p-p pairing. First, a density independent contact interaction (eq. 2.1) fitted to reproduce the average proton and neutron pairing gaps in <sup>254</sup>Fm ($`V_{0,n}^{vol}`$ = $``$ 250 MeV.fm<sup>3</sup>, $`V_{0,p}^{vol}`$ = $``$ 290 MeV.fm<sup>3</sup>):
$$\widehat{V_\tau }^{vol}=\frac{V_{0,\tau }^{vol}}{2}(1P_\sigma )\delta (\text{r}_1\text{r}_2),$$
(2.1)
where $`P_\sigma `$ is the spin exchange operator. In the second case, the pairing correlations have been treated with a surface-peaked delta force (eq. 2.2) adjusted on the low spin behavior of the moments of inertia of super-deformed bands in the A $``$ 150 region . It is given by
$$\widehat{V_\tau }^{surf}=\frac{V_{0,\tau }^{surf}}{2}(1P_\sigma )\delta (\text{r}_1\text{r}_2)(1\frac{\rho (\stackrel{}{R})}{\rho _c}),$$
(2.2)
where $`V_{0,n}^{surf}`$ = $``$ 1250 MeV.fm<sup>3</sup>, $`V_{0,p}^{surf}`$ = $``$ 1250 MeV.fm<sup>3</sup>, $`\rho (\stackrel{}{R})`$ is the local density of the nucleus and $`\rho _c`$ = 0.16 fm<sup>3</sup> the nuclear saturation density. In this way, this pairing interaction is mainly active in the surface region of the nucleus.
With a contact interaction, one has to define a cut-off procedure for the active pairing space; a smooth cut-off simulates the decay of coupling as it would occur with a finite-range force. For the volume pairing, the active pairing space includes all states up to 5 MeV above the Fermi energy. For the surface pairing, the active space includes roughly one major shell, from 5 MeV above to 5 MeV below the Fermi level.
## 3 A pairing exploration
The surface pairing interaction (see eq. 2.2) has been adjusted on moments of inertia of SD bands in the A $``$ 150 mass region. Systematic calculations have shown that the same parametrization leads also to a very good agreement with experimental data for even, odd and odd-odd nuclei in the A = 190 region . Before using that parametrization to describe rotational bands in No isotopes, it is worth testing its validity for the description of ground-state properties in the A $``$ 250 region.
For this purpose, we have calculated fifteen nuclei around <sup>252</sup>No and <sup>254</sup>No. The Nobelium isotopes from <sup>250</sup>No to <sup>256</sup>No, the <sup>252</sup>No isotones from <sup>250</sup>Fm to <sup>255</sup>Lr and the <sup>254</sup>No isotones from <sup>252</sup>Fm to <sup>256</sup>Rf. For each odd nucleus, the different 1-qp states have been self-consistently calculated and the lowest in energy has been selected. The experimental spins and theoretical spin projections on the symmetry axis (K-value) of the ground-state of these nuclei are compiled in table 1. The calculated K-value are consistent with experimentally known (or assigned) spins. This gives some confidence in the level scheme obtained with the SLy4 parametrization and also to the spin predictions made for nuclei for which no experimental data are available. This is also supported by self-consistent calculations performed with a delta density-independent pairing and a Woods-Saxon potential which found the <sup>253</sup>No and <sup>255</sup>Lr spins to be K<sup>π</sup> = 9/2<sup>-</sup> and 7/2<sup>-</sup> respectively.
The next step is to choose a relevant quantity to compare theoretical and experimental pairing correlations. One can use a pure theoretical definition of the “gap” such as the gap at the Fermi energy $`\mathrm{\Delta }_{ϵ_F}`$, a mean gap $`<\mathrm{\Delta }>=\frac{_ku_kv_k\mathrm{\Delta }_k}{_ku_kv_k}`$ or the lowest quasi-particle energy. These quantities may be compared to experimental “equivalents” defined by the excitation spectrum of even nuclei or by mass differences between odd and even isotopes. This is a practical procedure to deal with a systematic work on pairing because it only requires the theoretical calculation of a single nucleus at a time and avoids microscopic calculations of odd nuclei when considering even ones.
However, this procedure can be improved by imposing the following two requirements:
* The quantity used to extract pairing information must be derived in the same way both from experiment and theory,
* it should be mainly sensitive to pairing correlations and eliminate as much as possible mean-field contributions to total energies.
### 3.1 Odd-even mass formulas
Several well known quantities which aim at evaluating the neutron or proton “pairing gaps” are finite-difference mass formulas ?-?). They approximate the systematic additional binding energy due to pairing of an even-even nucleus with respect to an odd-even neighbor. We give here the two-point (first order) and three-point (second order) formulas
$`\mathrm{\Delta }_q^{(2)}(N)`$ $`=`$ $`(1)^N\left[E(N1)E(N)\right]`$ (3.3)
$`\mathrm{\Delta }_q^{(3)}(N)`$ $`=`$ $`{\displaystyle \frac{(1)^N}{2}}\left[E(N1)2E(N)+E(N+1)\right],`$ (3.4)
where the neighboring nuclei are taken along an isotopic (q = neutron) or isotonic chain (q = proton).
To fulfill the above requirement (1), the $`\mathrm{\Delta }_q^{(n)}`$’s have to be calculated theoretically with full consistency thanks to the correct treatment of odd-N (or Z) nuclei in the model.
However, these quantities are not a measure of the pairing gaps only. They include a contribution arising from the mean-field which depends upon the order (n) of the formula and remains to be evaluated. $`\mathrm{\Delta }_q^{(2)}(N)`$ can be roughly estimated as
$`\mathrm{\Delta }_q^{(2)}(N)`$ $`=`$ $`(1)^NS_q(N)`$
$``$ $`\mathrm{\Delta }_q^{(2)pairing}(N)+(1)^{N+1}\lambda _q,`$
where $`S_q(N)`$ is the “one-nucleon separation energy”, $`\lambda _q`$ is the neutron or proton chemical potential and $`\mathrm{\Delta }_q^{(2)pairing}(N)`$ defines the contribution to $`\mathrm{\Delta }_q^{(2)}(N)`$ coming from pairing correlations only.
As a consequence, $`\mathrm{\Delta }_q^{(2)}(N)`$ is dominated by the contribution from the mean-field since, apart for drip-line nuclei, we generally have $`\lambda \mathrm{\Delta }^{pairing}`$. For $`\mathrm{\Delta }_q^{(3)}(N)`$, however, the two contributions are of the same order of magnitude in light nuclei and the pairing dominates a priori in heavy ones. Actually, to isolate the influence of the pairing field, it is necessary to use high order formulas : the higher the order of the formula, the smaller the contribution coming from the particle-hole channel.
The use of $`\mathrm{\Delta }_q^{(3)}(N)`$ is sufficient as our aim is limited to test the validity of the pairing parametrizations that we use for No isotopes without performing a global fit in the whole nuclear chart. To fit a pairing force and fulfill requirement (2), one should take a higher order formula.
### 3.2 Pairing adjustment
Figure 1 shows the experimental and calculated values of $`\mathrm{\Delta }_q^{(3)}(N)`$ for neutrons along a chain of Nobelium isotopes and for protons along the neighboring isotones of <sup>254</sup>No. The surface pairing fitted to reproduce correctly the SD bands in the A $``$ 150 mass region gives satisfactory results, since all calculated values fall within the error bars. We will refer to this parametrization as pairing 1.
The experimental error bars are quite large and do not allow a very precise adjustment of the pairing interaction. Variations of the strengths $`V_{0,q}^{surf}`$ by 100 MeV still lead to $`\mathrm{\Delta }_q^{(3)}(N)`$ within the error bars. Small variations of the cut-off energies $`E_q`$ can also be explored. We call “pairing 2” the set of parameters $`V_{0,n}^{surf}`$ = $``$ 1250 MeV.fm<sup>3</sup>, $`V_{0,p}^{surf}`$ = $``$ 1350 MeV.fm<sup>3</sup> and an active window of 4.4 MeV and 5.4 MeV around the Fermi energy for neutron and proton respectively. This set leads to theoretical $`\mathrm{\Delta }_q^{(3)}(N)`$ the closest to the middle of the experimental error bars. This second fit allows us to explore a significant range of variations for the parameters with respect to known or estimated experimental errors.
Figure 2 displays the often used one neutron separation energy $`S_n(N)`$ from <sup>252</sup>No to <sup>256</sup>No. The mean deviation from experimental data is about 300 keV.
Surprisingly, there is no one to one visible relationship between the improvement of $`\mathrm{\Delta }_n^{pairing}`$ obtained with pairing 2 (see Fig.1) and the modification that this parametrization brings on the $`S_n(N)`$. This can be understood by analyzing the $`S_n(N)`$ formula (see eq. 3.1). Indeed, $`\mathrm{\Delta }_n^{pairing}`$ enters into $`S_n(N)`$ with a different sign for odd and even N. The change of $`\mathrm{\Delta }_n^{pairing}`$ being the main change in $`S_n(N)`$, the decrease of $`\mathrm{\Delta }_n^{pairing}`$ leads to an increase of $`S_n(N)`$ for odd N and to a decrease for even N. In this way, the odd $`S_n(N)`$ are improved contrary to the even one. Thus, there is an apparent paradox : while improving the pairing ($`\mathrm{\Delta }_q^{(3)}(N)`$), some $`S_n(N)`$ are spoiled because of this alternating sign. This result illustrates the interplay between the pairing and mean-field parts of the energy. The mean field is responsible for the global shift of the $`S_n`$ with respect to experiment. Thus, if it were corrected, the decrease of $`\mathrm{\Delta }_n^{pairing}`$ would also improve the even $`S_n`$, since they would be too high with pairing 1 in this case.
Consequently, the use of $`S_n(N)`$ to evaluate the influence of the pairing adjustment can reveal a weakness of the mean-field since this quantity is sensitive to both particle-particle and particle-hole channels (second order formula).
Figure 3 gives the quadrupole deformation $`\beta _2=\sqrt{\frac{\pi }{5}}\frac{<\widehat{Q}>}{A<r^2>}`$ along the isotopic chain using pairing 1. All these nuclei have axial quadrupole deformations in their ground state. The values obtained with pairing 2 differ by at most 0.5 $`\%`$ for all isotopes although the neutron gap are modified by up to 30 $`\%`$. The value obtained for <sup>254</sup>No is close to the one deduced from experiment ?-?) ($`\beta `$ = 0.27 $`\pm `$ 0.02). This result is supported by the ones obtained in HFB calculations with SLy4 and volume pairing , with the Gogny force , by the relativistic Hartree-Bogoliubov method and by a macroscopic-microscopic approach . Taking the corrections due to $`\beta _4`$ into account according to Ref , one obtains $`\beta _2`$ = 0.264.
The deformation along the whole line between <sup>250</sup>No and <sup>256</sup>No shows a very smooth global decrease with increasing $`A`$ toward shell closure at N = 184 . Besides this global trend the quadrupole deformation presents an odd-even staggering along the isotopic line. This very weak effect shows that even nuclei are comparatively less deformed than their odd neighbors. Figure 3 also gives the quadrupole deformation for the same isotopes except that the odd nuclei are calculated self-consistently as if they were even ones without any blocking, i.e. without breaking time invariance. In this case, no staggering occurs and only the smooth decrease can be seen. The additional deformation of an odd nucleus with respect to its virtual even partner is thus directly connected to the blocking which was, up to now, known to be responsible for the odd-even staggering of the binding energy. One can thus assess that the blocking procedure, through the weakening of the pairing correlations and the polarization of the core, is responsible for a second staggering effect on the deformation.
## 4 Rotational properties
### 4.1 Routhians
Figure 4 shows the proton and neutron qp Routhians of <sup>254</sup>No and Figure 5 the single-particle energies in the HF basis as a function of the frequency. The quasi-particles (qp) are labeled by their dominant Nilsson component in the HF basis and by the letter p (particle) or h (hole) indicating whether this component lies above or below the Fermi level.
The quasi-proton 7/2<sup>+</sup> orbitals originate from the intruder orbitals $`1i_{13/2}`$. As expected, they are strongly down-slopping and present a signature splitting of 0.2 MeV for $`\mathrm{}\omega `$=0.3 MeV. They should play an important role at high spins in the heavy isotones of <sup>254</sup>No, while the 7/2<sup>-</sup> qp should dominate in the light isotones.
The neutron quasi-particle orbitals corresponding to the 9/2<sup>-</sup> intruder states have the lowest excitation energies for hole type orbitals. The ground state of <sup>253</sup>No is based on this qp excitation. The other low lying qp excitations are of particle type. Consequently, several 1qp excitations should be close in energy in <sup>255</sup>No, leading to a large fragmentation of the population of the rotational bands.
Large proton and neutron gaps can be seen on Figure 5 at all frequencies. For the neutrons, they correspond to the lighter isotopes with N = 150 and 152 and for the protons to the heavier isotones with Z = 104 and Z = 108 although in this last case, the gap disappears above $`\mathrm{}\omega `$ = 0.2 MeV. We also obtain a gap for N = 162, which confirms the existence of such a deformed shell closure found in previous microscopic and macro-microscopic calculations .
### 4.2 Moments of inertia
We compare in Figure 6 the theoretical dynamical moment of inertia of <sup>252</sup>No and <sup>254</sup>No with the experimental data . For both nuclei, we have used the pairing 1 presented in section 3.2. In addition, we have tested three other pairing parametrizations for <sup>254</sup>No: the volume pairing defined in eq. 2.1), the pairing 2 discussed in section 3.2 and a third surface pairing in which the neutron and proton strengths have been arbitrarily decreased to $`V_{0,n}^{surf}`$= $``$ 1200 MeV.fm<sup>3</sup>. This last set has been introduced to determine whether the agreement with the data of the set 1 could be improved at low frequency staying within the error bars for the ground states.
The agreement with the data is quite good in all cases. The volume pairing slightly overestimates the moment of inertia of <sup>254</sup>No while all the surface pairing parametrizations slightly underestimate it at low frequencies. At higher frequencies, different behaviors for the two pairings are seen with a faster increase of $`𝒥^{(2)}`$ associated to the surface pairing. Reasonable changes of the pairing strength do not modify significantly the moments of inertia which seem to be more sensitive to the region of the space (surface or volume) where the pairing is active than to variations of its strength. An intermediate behavior between the extreme surface and volume pairings might improve the agreement with the data. It does not seem however appropriate to perform such a detailed fit of the pairing interaction by looking to properties only in this mass region.
There is no experimental data for the dynamical moment of inertia of <sup>253</sup>No. The predicted $`𝒥^{(2)}`$ for <sup>252</sup>No, <sup>253</sup>No and <sup>254</sup>No using pairing 1 are shown on Fig. 7. The calculated rotational band of <sup>253</sup>No is based on the intruder qp 9/2<sup>-</sup>. One observes an increase of the moment of inertia of <sup>253</sup>No compared to <sup>252</sup>No and <sup>254</sup>No at low frequency, due to the reduction of the pairing by blocking. At high frequency, the high $`j`$ blocked state still contributes strongly to the alignment but without contributing to the loss of pairing energy. This effect makes the dynamical moment of inertia of <sup>253</sup>No flatter than those of <sup>252</sup>No and <sup>254</sup>No at high spin.
The comparison between the moment of inertia of the two even nuclei is also fruitful. At low frequency, the two moments of inertia are identical, but $`𝒥^{(2)}(^{252}No)`$ increases faster than $`𝒥^{(2)}(^{254}No)`$ beyond 200 Kev. This feature is observed on Fig. 7. It is however more pronounced in the experimental data (see right hand side of Fig. 7). The present calculations are thus able to reproduce qualitatively this peculiar feature of Nobelium rotational properties, but more experimental data at higher spin would be of interest for quantitative comparisons.
### 4.3 Fission barrier
The rotational band of <sup>254</sup>No has been observed up to spin 22$`\mathrm{}`$ and an excitation energy larger than 6 MeV at that spin . To understand the stability of this band, we have calculated the fission barrier of <sup>254</sup>No and its dependence on angular momentum.
The result is shown on Figure 8 for a spin of 0$`\mathrm{}`$ and at spin $`20\mathrm{}`$. The calculation is performed with the same formalism used for the rotational band with an additional constraint on the quadrupole moment.
At spin 0$`\mathrm{}`$ two minima are observed; a first very deep one at $`𝒬_{20}`$ = 32.8 b and a second one, less pronounced, at $`𝒬_{20}`$ = 106 b with an excitation energy of 2.4 MeV. These two minima are separated by a barrier at $`𝒬_{20}`$ = 68 b whose energy is 12.6 MeV above the first minimum when constraining to axial deformation only. Allowing for triaxial deformations, the energy curve is only modified in the region of the first barrier which is lowered by 3 MeV. At the maximum of the barrier, the triaxial angle $`\gamma `$ reaches 10 . Apart for the first barrier, the nucleus is axially symmetric at all deformations.
The height of the second barrier is 2.9 MeV with respect to the second minimum. Consequently, this second minimum could hold a super-deformed state stable against fission. However, the inclusion of octupole deformation would strongly decrease its fission half-life ?-?). Moreover, this state is very likely unstable against neutron emission.
At $`20\mathrm{}`$, the structure of the barrier is similar, although the two minima are slightly closer in energy and the height of the barrier is decreased to 11.1 MeV and should be lowered by a few MeV by the inclusion of triaxiality. These calculations thus confirm what has been found experimentally, namely a barrier height $``$ 5 MeV for spin I = $`20\mathrm{}`$. This structure should remains the same up to spin $`3040\mathrm{}`$. A similar result result has been found with the Gogny force .
Figure 9 gives the <sup>254</sup>No ground-state density distribution for four different quadrupole axial deformations. These densities are given in a plane containing the z symmetry axis. A hollow in the center of the nucleus can be seen for $`𝒬_{20}`$ = 68 b (maximum of the barrier) and becomes deeper with deformation. At larger quadrupole deformations a strong increase of negative axial hexadecapole deformation takes place (see the right panels of Fig. 9). Beyond the second minimum, namely for values of $`𝒬_{20}`$ illustrated by the rightmost panel of Fig. 9, important octupole deformations are expected.
As mentioned above, triaxial deformations are important at the barrier. Figure 10 shows the density distribution for $`𝒬_{20}`$ = 68 b in the same plane (left panel) and in the transverse plane (right panel) for a triaxial calculation ($`\gamma `$ $``$ 10 ) corresponding to the saddle point between the two minima.
## 5 Conclusion
In this paper, we have calculated rotational properties of <sup>252,253,254</sup>No using the HFBLN formalism in the coordinate representation .
In the particle-particle channel, we have used a zero-range density-dependent pairing force adjusted on super-deformed bands in the A $``$ 150 mass region . We have checked the validity of this force in the present region of Nobelium using odd-even mass formulas. This previous fit gives results in agreement with the data, within the errors bars. The spins of even-odd nuclei obtained for nuclei around the Nobeliums are in good agreement with experiment. This gives credit to the single-particle spectra obtained with the SLy4 Skyrme parametrization in this mass region. Compiling ground-state quadrupole deformation for the Nobeliums isotopes between <sup>250</sup>No and <sup>256</sup>No, we have found a theoretically interesting odd-even staggering which originates, as the odd-even staggering of the binding energy, in the weakening of the pairing correlations and polarization of the core due to the blocking effect in odd nuclei.
We have calculated the dynamical moments of inertia corresponding to the two experimentally known rotational bands of <sup>252</sup>No and <sup>254</sup>No. We have also calculated the <sup>253</sup>No’s one for which experiments are planned. We have found a good agreement with experiment for two kinds of pairing force, namely the zero-range force and the zero-range force with a density dependence (active at the surface of the nucleus). We have used several fits for the zero-range density-dependent pairing force, however the dynamical moments of inertia revealed to be less sensitive to the fit used than to the region in space where the pairing is active. The above results are of particular interest because single-particle energies and pairing correlations are deeply involved in the description of rotational properties. Consequently, this gives a strong credit to the forces and fits used in the two channels, in a mass region just below the super-heavy elements.
Finally, we have calculated the fission barriers of <sup>254</sup>No at 0$`\mathrm{}`$ and 20$`\mathrm{}`$ to understand the observation of its rotational band up to spin 22$`\mathrm{}`$ and energy excitation around 6 MeV . The height of the first barrier is about 12.5 MeV at 0$`\mathrm{}`$, and still around 11 Mev at spin 20$`\mathrm{}`$. We expect it to remain important enough to hold a stable state against fission up to spin 30-40$`\mathrm{}`$.
Thus, in the language of macroscopic-microscopic models , we can assess that the shell-correction effects are sufficiently robust against rotation to explain the observation of the rotational band of <sup>254</sup>No up to such high spins. We have also calculated the densities for several quadrupole deformations of <sup>254</sup>No. These densities show a hollow in the center of the nucleus already for the ground-state and which increases with deformation. Such a hollow may affect the Coulomb part of the liquid drop energy in macroscopic-microscopic models.
## 6 Acknowledgment
We would like to thank J.L. Egido for fruitful discussions and T.L. Khoo for providing us results before publication.
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# A dual point description of mesoscopic superconductors
## I Introduction
The ability to detect and manipulate vortices with great sensivity in systems of small size such as mesoscopic superconductors or atomic condensates has generated an outgrow of interest in the mechanism of creation and annihilation of vortices and in the study of stable and metastable vortex configurations. In particular, recent advances in the technique of Hall magnetometry have allowed to measure the magnetization of small superconducting samples containing only a few vortices . These experiments are conducted on aluminium disks well below the superconducting transition temperature, whereas previous measurements were performed only in the vicinity of the normal/superconductor phase boundary . Besides, the magnetization measurements in are carried out on an individual disk and not on an ensemble of disks as in . The radius $`R`$ and the thickness $`d`$ of the sample used in the experiments are comparable to the superconducting characteristic lengths, i.e. the London penetration length ($`\lambda =`$ 70nm) and the coherence length ($`\xi =`$250nm). Such a sample can neither be considered to be macroscopic, nor microscopic. The system falls, rather, in a mesoscopic regime where surface effects are of the same order of magnitude as the bulk effects. Thus, the magnetic response of a mesoscopic superconducting disk to an applied field depends strongly on its size and is very different from that of a macroscopic superconductor. When the radius $`R`$ of the sample is much smaller than the coherence length $`\xi `$, no vortex can nucleate, the normal/superconductor phase transition is second order and the magnetization $`M`$, as a function of the external applied field $`H_e`$, has a non-linear behaviour (non-linear Meissner effect ). If $`R`$ is comparable to $`\xi `$, the superconducting phase transition is first order and a bistable hysteresis region appears in the $`MH_e`$ curve. For $`R`$ greater than $`\xi `$, the phase transition is again second order, and when the applied field exceeds a critical value $`H_1`$, the magnetization curve exhibits a series of discontinuous jumps corresponding to the successive entry of vortices into the sample. This qualitative interpretation is supported, at least for low applied magnetic fields, by the periodicity of the jumps which corresponds to the entrance of an additional superconducting quantum of flux into the disk. For larger fields, or equivalently for higher density of vortices, both the period and the height of the jumps become smaller, a behaviour related to the interactions between the vortices and to transitions between stable vortex configurations, with the same number of vortices.
The magnetization shows also a hysteretic behaviour depending on the direction of the field sweep, due to the presence of a confining energy barrier (the absence of remanent magnetization precludes pinning effects). In some metastable states, the sample may exhibit even a paramagnetic response whereas in thermodynamic equilibrium a superconductor is diamagnetic.
These experimental results have led to a renewed interest in the theory of mesoscopic superconductors. Numerical computations have shown that the phenomenological Ginzburg-Landau theory is well suited to describe a superconducting sample in the mesoscopic regime, even far from the critical temperature. These works have revealed physical phenomena that play an important role in such systems (for a review see ), like: the role of surface barriers for vortex nucleation and hysteresis ; the interplay between vortex/vortex and vortex/edge interactions that explains vortex structures in mesoscopic disks ; the transition between a giant multiple vortex state and a state with several vortices carrying a unit quantum of flux .
The Ginzburg-Landau free energy of a superconductor involves two fields, the (complex) order parameter $`\psi =|\psi |e^{i\chi }`$ and the vector potential $`\stackrel{}{A}`$. The minimization of this free energy leads to a set of two coupled non-linear partial differential equations for $`\psi `$ and $`\stackrel{}{A}`$, involving the two characteristic lengths $`\lambda `$ and $`\xi .`$ But the solutions depend only on one relevant number, the phenomenological Ginzburg-Landau parameter $`\kappa `$ defined by
$$\kappa =\frac{\lambda }{\xi }.$$
(1)
A macroscopic superconductor is said to be of type I if $`\kappa <\frac{1}{\sqrt{2}}`$ and of type II if $`\kappa >\frac{1}{\sqrt{2}}.`$ A macroscopic superconductor of Type II admits a stable Abrikosov vortex lattice phase when the applied field $`H_e`$ lies between the first penetration field and the upper critical field . For aluminium, $`\kappa `$ is smaller than $`\frac{1}{\sqrt{2}}`$, hence a macroscopic sample of Al is a type I superconductor.
Analytical studies of the Ginzburg-Landau equations in two-dimensional systems require the use of various approximations since, in general, exact solutions can not be found due to the non-linearity. One approach is to linearize the equations assuming $`|\psi |1,`$ and to decouple them by supposing that the magnetic field $`B`$ in the sample is equal to the applied field $`H_e`$. This approach describes correctly the superconducting/normal phase boundary , but fails to explain the behaviour of the sample deep inside the superconducting state. For example, in the linearized theory, all the vortices are at the center of the disk and therefore one can not study the role of surface barriers, the interaction between vortices, and the fragmentation of a giant vortex into unit vortices. Besides, the critical fields corresponding to the successive entrance of vortices into the sample do not scale correctly with the size of the system (e.g. experimentally, the entrance field $`H_1`$ of the first vortex scales as $`R^1`$ whereas the linear theory predicts a $`R^2`$ dependence). Of course, in the vicinity of the upper critical field the linearized theory agrees quantitatively with the experimental results.
A second approach is to use the London equation which can be derived from the Ginzburg-Landau equations by supposing that $`|\psi |=1`$ everywhere except on a finite number of isolated points, called vortices, where $`|\psi |=0.`$ London’s equation is valid rigorously when the parameter $`\kappa `$ goes to infinity, i.e. for extreme type II superconductors in which vortices are indeed point-like. Many theoretical results have been derived from the London equation, such as discrete nucleation of flux lines in a thin cylinder or in a thin disk , the existence of surface energy barriers , and the computation of polygonal ring configurations of vortices in finite samples . However, when $`\kappa \mathrm{},`$ the minimum energy is obtained for one flux quantum per vortex and vortices have a hard-core repulsive interaction impeding the formation of a giant vortex state. Moreover, the surface energy barriers calculated from London’s equation are quantitatively different from those obtained by numerically solving the Ginzburg-Landau equations . In fact, the experimental conditions are far off the London limit, although thin Al disks are likely to have an effective $`\kappa `$ greater than its measured value of 0.28 (in a thin disk, one can argue, following , that the effective London length is of the order of $`\lambda ^2/d`$, and this results in a higher value of $`\kappa `$).
We follow another approach, less explored in the literature, based on an exact result for the two-dimensional Ginzburg-Landau equations. In an infinite plane reduce to first order differential equations that can be decoupled when the parameter $`\kappa `$ takes the special value $`1/\sqrt{2},`$ called the dual point. At that point, the free energy is a topological invariant of the system. In , we generalized this method to a finite domain with boundaries; this enabled us to classify solutions with different number of vortices and to derive analytical expressions for the free energy and the magnetization of a mesoscopic disk as a function of the applied field. Our results agreed qualitatively with the experimental data, and even quantitatively when the number of vortices in the system is low. However, some important features such as the non-linear Meissner effect in a fractional fluxoid disk, the variation of the amplitude and the period of the jumps in the $`MH_e`$ curve could not be described. Moreover, in , we discussed only the case where $`R`$ is much larger than $`\xi `$ and did not obtain the different regimes of the magnetization curve when the ratio $`R/\xi `$ is varied.
In this paper, we study the Ginzburg-Landau free energy $``$ not only at the dual point $`\kappa =1/\sqrt{2}`$ but also in its vicinity where vortices start to interact weakly . Taking into account non-linear effects, our calculations describe the magnetic response of the sample as its size changes, providing an understanding of the non-linear Meissner effect and of the multi-vortex state. We shall also study non-equilibrium vortex configuration in order to determine the interaction between a vortex and edge currents.
The plan of this paper goes as follows: in section 2, some basic features of the Ginzburg-Landau theory of superconductivity are recalled. In section 3, after studying the case of an infinite system, we generalize the Bogomol’nyi’s approach to a finite size superconductor and calculate its free energy at the dual point. This result is applied to an infinite cylinder in section 4. The case of a mesoscopic disk is studied in section 5 and magnetization curves are obtained for systems of different sizes. In section 6, we obtain the free energy and the magnetization of a cylindrically symmetric system when $`\kappa `$ is close to the dual point. The surface energy barrier for a one vortex state out of thermodynamic equilibrium is calculated in section 7. In the last section we discuss our results and suggest some further generalizations. Some mathematical details are included in the two appendices.
## II The Ginzburg-Landau theory of superconductivity
We recall here some basic features of the Ginzburg-Landau theory and define our notations. The order parameter $`\psi =|\psi |e^{i\chi }`$ is a complex number and the potential vector $`\stackrel{}{A}`$ satisfies $`\stackrel{}{}\times \stackrel{}{A}=\stackrel{}{B}`$, where $`\stackrel{}{B}`$ is the local magnetic induction. The two characteristic lengths $`\lambda `$ and $`\xi `$ appear as phenomenological parameters. In this work, we measure lengths in units of $`\lambda \sqrt{2}`$, the magnetic field in units of $`\frac{\varphi _0}{4\pi \lambda ^2}`$ and the vector potential in units of $`\frac{\varphi _0}{2\sqrt{2}\pi \lambda }`$ where the flux quantum $`\varphi _0`$ is given by $`\varphi _0=\frac{hc}{2e}.`$ The Ginzburg-Landau free energy $``$, defined as the difference of the free energies $`=F_S(B)F_S(0)`$, is measured in units of $`\frac{H_c^2}{4\pi }`$ where $`H_c`$ the thermodynamic field satisfies $`H_c=\sqrt{2}\kappa \frac{\varphi _0}{4\pi \lambda ^2}`$. In these units, $``$ is given by
$$=_\mathrm{\Omega }\frac{1}{2}|B|^2+\kappa ^2|1|\psi |^2|^2+|(\stackrel{}{}i\stackrel{}{A})\psi |^2,$$
(2)
where the integration is over the superconducting domain $`\mathrm{\Omega }`$. The Ginzburg-Landau equations that minimize $``$, become
$`(\stackrel{}{}i\stackrel{}{A})^2\psi `$ $`=`$ $`2\kappa ^2\psi (1|\psi |^2)`$ (3)
$`\stackrel{}{}\times \stackrel{}{B}`$ $`=`$ $`2\stackrel{}{ȷ}`$ (4)
Equation (4) is the Maxwell-Ampère equation with a current density $`\stackrel{}{ȷ}=\text{Im}(\psi ^{}\stackrel{}{}\psi )|\psi |^2\stackrel{}{A}`$ related to the superfluid velocity $`\stackrel{}{v}_s`$ by
$$\stackrel{}{v}_s=\frac{\stackrel{}{ȷ}}{|\psi |^2}=\stackrel{}{}\chi \stackrel{}{A}$$
(5)
Outside the superconducting sample, $`\psi =0.`$ The boundary condition on the surface of the superconductor is obtained by requiring that the normal component of the current density vanishes (superconductor/insulator boundary condition ):
$$(\stackrel{}{}i\stackrel{}{A})\psi |_{\widehat{𝐧}}=0$$
(6)
here $`\widehat{𝐧}`$ is the unit vector normal at each point to the surface of the superconductor.
The London fluxoid is the quantity $`\left({\displaystyle \frac{\stackrel{}{ȷ}}{|\psi |^2}}+\stackrel{}{A}\right),`$ that is identical to $`\stackrel{}{}\chi `$. Since $`\chi `$ is the phase of the univalued function $`\psi ,`$ the circulation of the London fluxoid along a closed contour $`𝒞`$ is quantized :
$$_𝒞(\frac{\stackrel{}{ȷ}}{|\psi |^2}+\stackrel{}{A}).\stackrel{}{d}l=_𝒞\stackrel{}{}\chi .\stackrel{}{d}l=2\pi n$$
(7)
The integer $`n`$ is the winding number of the phase of the system along the contour $`𝒞`$ and is a topological characteristic of the system.
In this study, the superconducting sample is either an infinite cylinder or a thin disk, with cross-section of radius $`R`$, placed in an external magnetic field parallel to its axis. Since $`R`$ is an important parameter, we define the dimensionless quantity:
$$a=\frac{\lambda \sqrt{2}}{R}$$
(8)
$`a`$ is supposed to be small compared to 1 (typically $`a1/10`$ in the experiments) unless stated otherwise. The flux created by the external and uniform magnetic field $`H_e`$ (expressed in units of $`\frac{\varphi _0}{4\pi \lambda ^2}`$) through the cross section $`\pi R^2`$ of the sample is equal to $`\pi R^2H_e\frac{\varphi _0}{4\pi \lambda ^2}=\frac{H_e}{2a^2}\varphi _0.`$ The flux $`\varphi _e`$, in units of the flux quantum $`\varphi _0`$, is thus given by:
$$\varphi _e=\frac{H_e}{2a^2}$$
(9)
We emphasize that, in the units we have chosen, the flux $`\varphi _b`$ of a magnetic field $`\stackrel{}{B}`$ through a surface $`\mathrm{\Omega }`$ is obtained via the following formula:
$$\varphi _b=\frac{1}{2\pi }_\mathrm{\Omega }\stackrel{}{B}.\stackrel{}{d}S=\frac{1}{2\pi }_\mathrm{\Omega }\stackrel{}{A}.\stackrel{}{d}l$$
(10)
An extra factor $`1/2\pi `$ appears here because $`B`$ is given in units of $`\frac{\varphi _0}{4\pi \lambda ^2}`$, the surface in units of $`2\lambda ^2`$ and the flux in units of $`\varphi _0`$.
Since we are studying a superconductor in an applied external field, the relevant thermodynamic potential is the Gibbs free energy $`G`$ obtained from $`F`$ via a Legendre transformation:
$$G=FH_e_\mathrm{\Omega }B=FH_e2\pi \varphi _b=F4\pi a^2\varphi _e\varphi _b$$
(11)
In a normal sample, $`\psi =0`$ and $`B=H_e`$. Therefore, the Gibbs free energy $`G_N`$ of a normal sample is given by:
$$G_N=F_NH_e_\mathrm{\Omega }B=F_N2\pi a^2\varphi _e^2$$
(12)
At thermodynamic equilibrium, the superconductor selects the state of minimal Gibbs free energy. The quantity that we are interested in, and which is measured in experiments, is the magnetization $`M`$ of the superconductor due to the applied field given by $`4\pi M=BH_e`$. It is obtained, at thermodynamic equilibrium and up to a constant equal to the superconducting condensation energy, from the difference of the (dimensionless) Gibbs energies
$$𝒢=G_SG_N=+2\pi a^2\varphi _e^24\pi a^2\varphi _e\varphi _b$$
(13)
using the thermodynamic relation :
$$M=\frac{1}{2\pi }\frac{𝒢}{\varphi _e}$$
(14)
## III Free energy of a superconductor at the dual point
We now study the particular case of the dual point, defined by $`\kappa =\frac{1}{\sqrt{2}}.`$ For this value of the Ginzburg-Landau parameter, the free energy (2) of a two dimensional domain $`\mathrm{\Omega }`$ can be written as :
$$=_\mathrm{\Omega }\left(\frac{1}{2}\left(B1+|\psi |^2\right)^2+|𝒟\psi |^2\right)+_\mathrm{\Omega }(\stackrel{}{ȷ}+\stackrel{}{A}).\stackrel{}{d}l$$
(15)
where the operator $`𝒟`$ is defined as $`𝒟=_x+i_yi(A_x+iA_y)`$ and the second integral is over the boundary of the domain $`\mathrm{\Omega }.`$
### A The case of an infinite system
If we suppose that the domain $`\mathrm{\Omega }`$ is infinite and superconducting at large distances , i.e. $`|\psi |1`$ at infinity, then the boundary integral in (15) is identical to the fluxoid. Using the quantization property (7), we obtain
$$=2\pi n+_\mathrm{\Omega }\left(\frac{1}{2}(B1+|\psi |^2)^2+|𝒟\psi |^2\right)$$
(16)
The free energy is thus minimum when Bogomol’nyi equations are satisfied, that is when,
$`𝒟\psi `$ $`=`$ $`0`$ (17)
$`B`$ $`=`$ $`1|\psi |^2`$ (18)
Thus, the total free energy results only from the boundary term in (15) and is a purely topological number:
$$=2\pi n.$$
(19)
The free energy is proportional to the number of vortices: at the dual point, vortices do not interact with each other .
### B Finite size systems
In a finite system with boundaries, vortices do not interact with each other at the dual point but they are repelled by the edge currents. Therefore, at thermodynamic equilibrium, all vortices collapse into a giant vortex state. Since the superconductor under discussion has a circular cross-section, this giant vortex (or multi-vortex) is located at the center and the system is invariant under cylindrical symmetry. In a finite size mesoscopic superconductor at the dual point, the boundary integral, in (15), can not be identified with the fluxoid because $`|\psi |`$ is in general different from 1 on the boundary of the system. This quantity is no more a topological integer but a continuously varying real number. The two terms of (15) can not, therefore, be minimized separately to obtain the optimal free energy. In , we found a method to circumvent this difficulty: if the system is invariant under cylindrical symmetry, i.e. all the vortices are at the center of the disk, then the current density has only an azimuthal component $`ȷ_\theta `$. The current $`ȷ_\theta `$ has opposite signs near the center (where the vortex is located) and at the edge of the disk (where Meissner currents oppose the penetration of the external field). Hence, there exists a circle $`\mathrm{\Gamma }`$ on which $`ȷ_\theta `$ vanishes . Along $`\mathrm{\Gamma }`$, we have
$$\stackrel{}{ȷ}+\stackrel{}{A}=\frac{\stackrel{}{ȷ}}{|\psi |^2}+\stackrel{}{A}=\stackrel{}{}\chi \text{ and therefore }_\mathrm{\Gamma }(\stackrel{}{ȷ}+\stackrel{}{A}).\stackrel{}{d}l=2\pi n$$
(20)
The domain $`\mathrm{\Omega }`$ can thus be divided into two sub-regions, $`\mathrm{\Omega }=\mathrm{\Omega }_1\mathrm{\Omega }_2,`$ such that the boundary between $`\mathrm{\Omega }_1`$ and $`\mathrm{\Omega }_2`$ is the circle $`\mathrm{\Gamma }`$. By convention, we call $`\mathrm{\Omega }_1`$ the bulk and the annular ring $`\mathrm{\Omega }_2`$ the boundary region (see figure 1).
Numerical solutions of the Ginzburg-Landau equations in a two dimensional superconductor, with cylindrical symmetry, clearly show the separation of the sample cross-section into two distinct subdomains. In figure 2, we have plotted the order parameter and the magnetic field in a cylinder of radius $`R=10\lambda \sqrt{2}`$ with one vortex at the center. These two quantities vary only near the center and near the edge: there is a whole intermediate region in which $`|\psi |`$ and $`B`$ remain almost constant. When the system is large enough, these constant values are, to an excellent precision, identical to the asymptotic values of $`|\psi |`$ and $`B`$ in an infinite system. The current vanishes for a value of $`r`$ for which $`\frac{dB}{dr}=0`$, and this determines the radius of the circle $`\mathrm{\Gamma }.`$ In figure 2, this corresponds approximately to $`r5.5`$, though practically, the circle $`\mathrm{\Gamma }`$ can be placed anywhere in the saturation region where the current is infinitesimally small.
Although the Bogomol’nyi equations (17, 18) do not minimize the Ginzburg-Landau free energy , we notice that the behaviour of $`|\psi |`$ and $`B`$ in the bulk subdomain $`\mathrm{\Omega }_1`$ is still given by the relation $`B=1|\psi |^2,`$(18) which, as shown in figure 3, represents indeed an excellent approximation. Thus, using (15) and (20), we conclude that at the dual point the free energy of $`\mathrm{\Omega }_1`$ can be calculated as that of an infinite domain, namely
$$(\mathrm{\Omega }_1)=2\pi n$$
(21)
We also emphasize that the flux in $`\mathrm{\Omega }_1`$ is quantized and one has:
$$\frac{1}{2\pi }_{\mathrm{\Omega }_1}B=n$$
(22)
To calculate $`(\mathrm{\Omega }_2)`$, the identity (15) valid at the dual point is of no use anymore, since $`|\psi |`$ is in general different from 1 at the boundary, and the boundary integral in (15) can not be identified to the fluxoid. We therefore have to go back to the definition (2) of the Ginzburg-Landau free energy which becomes at the dual point:
$$(\mathrm{\Omega }_2)=_{\mathrm{\Omega }_2}\frac{B^2}{2}+(|\psi |)^2+|\psi |^2|\stackrel{}{}\chi \stackrel{}{A}|^2+\frac{(1|\psi |^2)^2}{2}$$
(23)
The assumption of cylindrical symmetry implies that $`\chi =n\theta `$ where $`\theta `$ is the polar angle and $`n`$ the number of vortices present at the center of the disk. Examining again figure 2, we observe that in $`\mathrm{\Omega }_2`$, the order parameter and the magnetic field vary from their values on the edge to their saturation values over a region of width $`\delta `$, which is of order 1 in units of $`\lambda \sqrt{2}`$<sup>*</sup><sup>*</sup>*Indeed, one has for a thick system $`\delta \lambda `$ at the dual point. For a thin film of thickness $`d,`$ $`\delta \lambda ^2/d`$ in the London limit . Since we are considering a mesoscopic regime in which $`d\lambda `$, both expressions indicate that $`\delta `$ is of order 1. The length $`\delta `$ therefore represents the typical distance over which the integrand in (23) has a non negligible value.
With the help of this observation, we shall estimate $`(\mathrm{\Omega }_2)`$ using a variational Ansatz: we shall consider that the modulus of the order parameter has a constant value $`\psi _0`$ over a ring of width $`\delta `$, included in $`\mathrm{\Omega }_2`$ and that $`\stackrel{}{A}`$ and $`\stackrel{}{B}`$ decay exponentially with a characteristic length $`\delta `$ from their boundary value to their bulk value. Clearly, our approximation will be valid only if the width of $`\mathrm{\Omega }_2`$ is large enough compared to 1. We first remark that our Ansatz is compatible with the boundary condition (6), which reduces here to $`\frac{d\psi }{dR}=0`$ and that it allows us to neglect the curvature term $`(|\psi |)^2`$ in (23). To evaluate the term proportional to the superfluid velocity $`v_s(r)`$ (5), we first notice that, due to the Meissner effect, it decreases from the boundary at $`r=R`$ with a behaviour well described by
$$v_s(r)=v_s(R)e^{(Rr)/\delta }$$
(24)
with $`v_s(R)=a(n\varphi _b)`$. To obtain the last equality we used that the boundary value of the vector potential is $`\stackrel{}{A}(R)=a\varphi _b\widehat{u}_\theta ,`$ where $`\varphi _b`$ is the total flux through the system. Hence, for a constant amplitude $`\psi _0`$ of the order parameter, we have
$$\frac{1}{2\pi }(\mathrm{\Omega }_2)=\frac{\delta }{2a}\left(\psi _0^2v_s^2(R)+(1\psi _0^2)^2\right)+\frac{1}{2\pi }_{\mathrm{\Omega }_2}\frac{B^2}{2}$$
(25)
The magnetic contribution in (23) is obtained from the typical magnitude $`\overline{B}`$ of the magnetic field in $`\mathrm{\Omega }_2`$ determined using (10) and (22) as
$$\varphi _b=\frac{1}{2\pi }_\mathrm{\Omega }B=\frac{1}{2\pi }_{\mathrm{\Omega }_1}B+\frac{1}{2\pi }_{\mathrm{\Omega }_2}B=n+\frac{\delta }{a}\overline{B}$$
(26)
Thus, using the fact that $`B^2`$ decrease exponentially with a characteristic length $`\delta /2`$, we estimate the contribution of the magnetic energy to $`(\mathrm{\Omega }_2)`$ as being:
$$\frac{1}{2\pi }_{\mathrm{\Omega }_2}\frac{B^2}{2}=\frac{\delta }{2a}\frac{\overline{B}^2}{2}=\frac{a}{4\delta }(n\varphi _b)^2$$
(27)
After substituting this expression into (25) we minimize $`(\mathrm{\Omega }_2)`$ with respect to $`\psi _0`$. The optimal variational value of $`\psi _0`$ is given by:
$$\psi _0^2=\{\begin{array}{cc}1\frac{1}{2}v_s^2(R)=1\frac{a^2}{2}(n\varphi _b)^2\hfill & \text{ if }|a(n\varphi _b)|\sqrt{2}\\ 0\hfill & \text{ if }|a(n\varphi _b)|>\sqrt{2}\end{array}$$
(28)
Inserting these expressions in (25), we obtain the variational free energy $`(\mathrm{\Omega }_2)`$:
$$\frac{1}{2\pi }(\mathrm{\Omega }_2)=\{\begin{array}{cc}Av_s^2(R)Bv_s^4(R)\hfill & \text{ if }|a(n\varphi _b)|\sqrt{2}\\ \frac{\delta }{2a}+\frac{1}{4a\delta }v_s^2(R)\hfill & \text{ if }|a(n\varphi _b)|>\sqrt{2}\end{array}$$
(29)
with $`A`$ and $`B`$ defined by
$`A`$ $`=`$ $`{\displaystyle \frac{\delta }{2a}}\left(1+{\displaystyle \frac{1}{2\delta ^2}}\right)`$ (30)
$`B`$ $`=`$ $`{\displaystyle \frac{\delta }{8a}}`$ (31)
The total free energy of the mesoscopic superconductor containing $`n`$ vortices, at the dual point, is thus:
$$\frac{1}{2\pi }(n,\varphi _b)=n+\{\begin{array}{cc}Av_s^2(R)Bv_s^4(R)\hfill & \text{ if }|a(n\varphi _b)|\sqrt{2}\\ \delta /2a+v_s^2(R)/4a\delta \hfill & \text{ if }|a(n\varphi _b)|>\sqrt{2}\end{array}$$
(32)
This energy is the sum of two contributions:
(i) a bulk term proportional to $`n`$ which is a topological quantity at the dual point.
(ii) A boundary term, reminiscent of the well-known ‘Little and Parks’ free energy (this boundary term can be given a geometric interpretation in terms of a geodesic curvature ).
## IV Free energy and magnetization of a cylinder at the dual point
We now apply the relations (32) to the simple case of an infinitely long superconducting cylinder of radius $`R>\lambda `$, lying in an external field $`H_e`$ directed along its axis. There are two contributions to the total flux $`\varphi _b`$: the flux of $`n`$ vortices present at the center of the sample and a fraction of the applied flux $`\varphi _e`$ localized near the boundary and proportional to $`\lambda /R`$ (due to the Meissner effect). Hence,
$$\varphi _b=n+2a\varphi _e\text{ with }\varphi _e=\frac{H_e}{2a^2}$$
(33)
The exact numerical coefficient in front of the term $`a\varphi _e`$ does not affect the result of our calculation; we take it equal to 2, the value obtained in the London limit . The total free energy, using (32) and the fact that $`v_s(R)=2a^2\varphi _e`$, is given by
$$\frac{1}{2\pi }(n,\varphi _b)=n+\{\begin{array}{cc}4a^4\left(A\varphi _e^24a^4B\varphi _e^4\right)\hfill & \text{ if }a^2|\varphi _e|1/\sqrt{2}\\ \frac{\delta }{2a}+\frac{a^3}{\delta }\varphi _e^2\hfill & \text{ if }|a^2|\varphi _b|>1/\sqrt{2}\end{array}$$
(34)
Using (11) and (34), the Gibbs free energy, $`𝒢(n,\varphi _e),`$ of a cylinder containing $`n`$ vortices at the dual point is given by:
$$\frac{1}{2\pi }𝒢(n,\varphi _e)=n\left(12a^2\varphi _e\right)+P(\varphi _e)$$
(35)
where $`P(\varphi _e)`$ is a polynomial in $`\varphi _e`$ that does not depend on $`n`$. Hence, all the curves $`𝒢(n,\varphi _e)`$ meet at
$$\varphi _c=\frac{1}{2a^2}$$
(36)
For values of $`\varphi _e`$ less than this critical value, the free energy is minimized if there are no vortices. At $`\varphi _e=\frac{1}{2a^2}`$ all vortices are nucleated simultaneously and the sample becomes normal. This value corresponds to a critical applied field $`H_e`$ which is equal to 1 in our units, or restoring the units back, and recalling that $`\kappa =1/\sqrt{2}`$
$$H_e=\frac{\varphi _0}{4\pi \lambda ^2}=\frac{\varphi _0}{2\sqrt{2}\pi \lambda \xi }$$
(37)
This is precisely the formula for the thermodynamic critical field of a superconductor (which, for a cylindrical superconductor with $`\kappa \frac{1}{\sqrt{2}},`$ is the same as the upper critical field). The magnetization $`M`$ of the cylinder satisfies the linear Meissner effect:
$$M=\frac{1}{2\pi }\frac{𝒢(n,\varphi _e)}{\varphi _e}=H_e(1ca)\text{ with }c=4\delta \left(1+\frac{1}{2\delta ^2}\right).$$
(38)
The macroscopic result is $`M=H_e`$; the finite-size correction to the susceptibility is proportional to $`R^1`$.
Thus, the well-known results for an infinite superconducting cylinder can easily be retrieved from the dual point approach. We now proceed to the study of the magnetic response of a thin disk.
## V A mesoscopic disk at the dual point
To modelize the experimental sample of , we consider a mesoscopic disk of thickness $`d`$ smaller than $`\xi `$ and $`\lambda `$. Because the disk is very thin, we take the order parameter and the magnetic field to be constant across the thickness $`d`$ of the sample . This enables us to study the disk as an effective two-dimensional system. However, unlike the case of a long cylindrical sample, strong demagnetization effects are present in a thin disk. The value of $`B`$ near the edge of the disk is larger than the applied field $`H_e`$ because geometric demagnetization effects induce a distortion of the flux lines . Hence the continuity condition $`B(R)=H_e`$ (33) valid for a long cylinder does not apply to describe a thin disk.
In order to find a more suitable choice for the boundary condition for a thin disk, we notice that the higher value of the magnetic field at the boundary, a feature which has been obtained from numerical computations , results from a demagnetization factor $`𝒩`$ close to one, such that $`H=\frac{H_e}{1𝒩}`$ in the Meissner phase. The flux lines are distorted by the sample and they pile up near the edge of the disk. To describe this, we shall thus take as boundary condition for a thin disk, the expression proposed in , which consists in taking the potential-vector at the edge of the disk equal to its applied value, i.e.
$$\stackrel{}{A}(R)=\varphi _ea\widehat{u}_\theta $$
(39)
or
$$\varphi _b=\varphi _e$$
(40)
Again, this relation does not mean that the field $`B`$ is uniform and equal to its external strength. A more refined value for the boundary condition could have been obtained by using the expression $`𝒩1\frac{\pi }{2}\frac{d}{R}`$ in the limit $`dR`$ of a flat disk. Then, $`H\frac{2R}{\pi d}H_e`$ or equivalently $`\varphi _b\frac{4\delta }{d}\varphi _e`$. But, since $`\delta d`$, we shall use for convenience the simpler boundary condition given above.
Substituting $`v_s(R)=a(n\varphi _e)`$ in (32), the free energy $`(n,\varphi _e)`$ of a thin disk containing $`n`$ vortices is found to be
$$\frac{1}{2\pi }(n,\varphi _b)=n+\{\begin{array}{cc}Aa^2(n\varphi _e)^2Ba^4(n\varphi _e)^4\hfill & \text{ if }a|(n\varphi _e)|\sqrt{2}\\ \frac{\delta }{2a}+\frac{a}{4\delta }(n\varphi _e)^2\hfill & \text{ if }a|(n\varphi _e)|>\sqrt{2}\end{array}$$
(41)
and the corresponding Gibbs free energy is obtained using (13). In our previous work , we obtained an expression which can be retrieved from (41) by taking $`\delta =1`$ and by neglecting the magnetic energy as well as the $`a^3`$ term. Despite these crude approximations, our analytical results agreed satisfactorily with experimental data, though they could not describe neither the behaviour of a disk with a radius smaller than $`\lambda `$ and $`\xi `$, nor its behaviour when $`R`$ is increased. We apply our present approach to a thin disk with a radius $`R`$ much smaller than $`\xi `$, and then we consider the case $`R>\xi .`$
### A Fractional fluxoid disk and Non-Linear Meissner Effect
We now consider a disk small enough so that no vortices can nucleate i.e. its radius $`R`$ is less than $`\xi `$ (such a system is sometimes called a fractional fluxoid disk ). If there are no vortices, the domain $`\mathrm{\Omega }_1`$ is empty and $`\mathrm{\Omega }=\mathrm{\Omega }_2`$. Since the radius of $`\mathrm{\Omega }`$ is small with respect to both $`\lambda `$ and $`\xi `$, we can no more use the expression (41) for the free energy, but we can assume that the amplitude $`|\psi |`$ of the order parameter has a uniform value $`\psi _0`$ all over the disk and that the magnetic field equals the external applied field $`B=H_e`$. Moreover, in the absence of vortices $`\stackrel{}{}\chi =0`$ and we can choose the Landau gauge $`A(r)=rB/2`$. Starting from (23), and after minimizing the free energy with respect to $`\psi _0`$, we find the difference between the free energies of the superconducting and the normal states to be:
$`{\displaystyle \frac{𝒢}{2\pi }}`$ $`=`$ $`{\displaystyle \frac{\varphi _e^2}{4}}\left(1{\displaystyle \frac{a^2}{4}}\varphi _e^2\right)\text{ if }a\varphi _e\sqrt{2}`$ (42)
$`{\displaystyle \frac{𝒢}{2\pi }}`$ $`=`$ $`0\text{ otherwise }`$ (43)
From (14) we deduce the magnetization $`M`$ of the sample:
$`M`$ $`=`$ $`{\displaystyle \frac{1}{2\pi }}{\displaystyle \frac{𝒢}{\varphi _e}}={\displaystyle \frac{1}{2}}\left(\varphi _e{\displaystyle \frac{a^2}{2}}\varphi _e^3\right)\text{ if }a\varphi _e\sqrt{2}`$ (44)
$`M`$ $`=`$ $`0\text{ otherwise }`$ (45)
The curve representing this magnetization is a cubic. The upper critical field is $`\varphi _e=1/a`$, i.e. $`H_eR^1`$; this scaling agrees with the linear analysis of in the limit $`R\xi `$. The transition between the superconducting phase and the normal phase is of second order. In figure 4, we plot the relation (45) for $`M`$ as a function of the external flux $`\varphi _e`$. The dots represent the experimental points obtained from . The analytical curve has been scaled so that the maximum value of the magnetization and the critical flux coincide with the corresponding experimental data.
### B Mesoscopic disk with vortices
We now consider a disk with $`R\xi `$. The Gibbs free energy difference $`𝒢(n,\varphi _e)`$ of the disk with $`n`$ vortices is given by (13). The entrance field $`H_n`$ of the $`n`$-th vortex is obtained by solving the equation $`𝒢(n,\varphi _e)=𝒢(n1,\varphi _e)`$ which, using (41), reduces to
$$\frac{2}{\delta }=a(1+\frac{1}{2\delta ^2})\left((n1\varphi _e)^2(n\varphi _e)^2\right)\frac{a^3}{4}\left((n1\varphi _e)^4(n\varphi _e)^4\right)$$
(46)
Using the following change of variable
$$\varphi _e=n\frac{1}{2}+\frac{y}{2a}$$
(47)
we obtain an equation for $`y`$
$$\frac{2}{\delta }=(1+\frac{1}{2\delta ^2})y\frac{y^3}{8}$$
(48)
(a term $`a^2/8`$ has been neglected in comparison to 1). The solution of (48) that satisfies $`y0`$ (because $`\varphi _e0`$) depends on the value of the parameter $`\delta `$. One can show that the polynomial $`P(y)=(1+\frac{1}{2\delta ^2})y\frac{y^3}{4}\frac{2}{\delta }`$ always has a positive root. We retain only the smaller positive root $`y_0`$ of (48) because in thermodynamic equilibrium, the system always chooses the state with minimal Gibbs free energy. Restoring the usual units, and using (47), the nucleation fields are found to be:
$`H_1`$ $`=`$ $`y_0{\displaystyle \frac{\varphi _0}{2\pi \sqrt{2}R\lambda }}+{\displaystyle \frac{\varphi _0}{2\pi R^2}}`$ (49)
$`H_{n+1}`$ $`=`$ $`H_1+n{\displaystyle \frac{\varphi _0}{\pi R^2}}`$ (50)
When the applied field $`H_e`$ lies between $`H_n`$ and $`H_{n+1}`$, the disk contains exactly $`n`$ vortices and its magnetization is calculated using (14). In figure 5, we have plotted the magnetization of a mesoscopic disk with $`R=10\lambda \sqrt{2}`$ both from exact numerical solutions of the Ginzburg-Landau equations and from the expression (14). The agreement is very satisfactory. For larger values of the number $`n`$ of vortices, a discrepancy between the theoretical and the numerical expressions appears which results from the interaction between the vortices and the edge currents that we have neglected until now.
The expression (32) is also in good agreement with previous experimental and numerical results . A non-linear Meissner behaviour still exists before the nucleation of the first vortex as well as between successive jumps. The field $`H_1`$ of nucleation of the first vortex scales as $`R^1`$. The transition between a state with $`n`$ vortices to a state with $`(n+1)`$ vortices is of first order since the entrance of a new vortex induces a jump in the magnetization. These jumps are of constant height and have a period $`\frac{\varphi _0}{\pi R^2}`$. If we use the experimental values of for $`R`$ and $`\lambda `$ we obtain a value for the period of the jumps which is in very good agreement with the experimental value.
If $`R`$ is smaller than a threshold value, the system is a fractional fluxoid disk with a second order phase transition. If $`R1`$, a vortex can nucleate in the disk and a first order transition occurs. When $`R`$ increases, the number of jumps increases (as $`R^2`$). These qualitative changes of behaviour with increasing $`R`$, which are the important features obtained from the present model, have been indeed observed in experiments carried out on disks of different sizes. In an earlier study , we obtained satisfactory values for the nucleation fields but the fractional fluxoid disk, and the different regimes obtained by increasing $`R`$ could not be explained because we neglected subdominant terms that are retained here.
It has been observed experimentally that the period and the height of the jumps cease to be constant when the number of vortices increases. These effects are related both to interactions between the vortices and between vortices and edge currents. The purpose of the next section is to take into account these interactions and to obtain a better estimate for the free energy and the magnetization of a mesoscopic disk.
## VI Weakly interacting vortices in the vicinity of the dual point
So far we have obtained analytical expressions for the free energy and the magnetization of a thin superconducting disk at the dual point. When the Ginzburg-Landau parameter has the special value, $`\kappa =\frac{1}{\sqrt{2}},`$ vortices do not interact. This fact, discussed in , implies that the bulk free energy does not depend on the location of the vortices. However, when $`\kappa `$ is away from the dual point, the vortices start interacting among themselves; therefore the bulk free energy ceases to be a purely topological integer $`n`$ and the vortex interaction energy must be taken into account. Because of this interaction the vortices are no longer necessarily placed at the center of the disk: in an equilibrium configuration, the cylindrical symmetry can be broken and the optimal free energy may correspond to geometrical patterns such as regular polygons, polygons with a vortex at the center, or even rings of polygons . It is the competition between the interaction amongst vortices and the interaction between vortices and edge currents that determines the shape of the equilibrium configuration.
Analytical studies were mostly carried out in the limit $`\kappa \mathrm{}`$ and were based on the London equation for which vortices are point-like and have a hard-core repulsion . We shall study a regime where $`\kappa `$ is slightly different than $`\frac{1}{\sqrt{2}}`$, i.e. a regime where vortices interact weakly. We shall determine, to the leading order in $`(\kappa \frac{1}{\sqrt{2}}),`$ the interaction energy of the vortices.
### A The interaction energy
In order to obtain an estimate for the free energy of a system of interacting vortices, we have solved numerically the Ginzburg-Landau equations for a cylindrically symmetric infinite system with $`n`$ vortices located at the center (these equations are explicitly written in Appendix A). The free energy per vortex is plotted in figure 6 as a function of $`\kappa `$, for $`n=1,2,3,5`$ and 10. At the dual point, the free energy per vortex is equal to 1 and is independent of $`n`$: all the curves pass through this point. When $`\kappa `$ is different from $`\frac{1}{\sqrt{2}}`$ the interaction between the vortices changes the value of the free energy. One can deduce from figure 6 that vortices attract each other for $`\kappa `$ less than $`\frac{1}{\sqrt{2}}`$ while they repel each other when $`\kappa \frac{1}{\sqrt{2}}.`$
From our numerical results we observed that in the vicinity of the dual point, the free energy $`(\kappa ,n)`$ satisfies the following scaling behaviour:
$$\frac{1}{2\pi }(\kappa ,n)=n(\kappa \sqrt{2})^{\alpha (n)}$$
(51)
We note that the relation (51) is exact at the dual point. For $`n=1`$, $`(\kappa ,1)`$ is nothing but the self energy $`𝒰_S`$ of a vortex. In the vicinity of the dual point we can write:
$$\frac{1}{2\pi }(\kappa ,1)=1+\alpha (1)(\kappa \sqrt{2}1)$$
(52)
The values of the function $`\alpha (n)`$, as determined from numerical computations, for $`n`$ ranging from 1 to 30 are given in the table I.
We can now derive an approximation for the free energy of a $`n`$ vortices configuration located at the center of the disk and for $`\kappa `$ close to $`\frac{1}{\sqrt{2}}.`$ Since this configuration is cylindrically symmetric, one can again use the circle $`\mathrm{\Gamma }`$ to separate the system into two subdomains $`\mathrm{\Omega }_1`$ and $`\mathrm{\Omega }_2`$ and then estimate separately the two contributions to the total free energy. From our numerical scaling result, we deduce a formula for the bulk free energy of a finite system which is valid for $`\kappa `$ close to the dual point. Expanding (51) in the vicinity of the dual point, we obtain:
$$\frac{1}{2\pi }(\mathrm{\Omega }_1)=n+(\kappa \sqrt{2}1)n\alpha (n)$$
(53)
and the boundary contribution, obtained via a variational Ansatz is now given by:
$$\frac{1}{2\pi }(\mathrm{\Omega }_2)=\{\begin{array}{cc}Av_s^2(R)B(\kappa )v_s^4(R)\hfill & \text{ if }|a(n\varphi _e)|2\kappa \\ \kappa ^2\delta /a+v_s^2/4a\delta \hfill & \text{ if }|a(n\varphi _e)|>2\kappa \end{array}$$
(54)
where $`A`$ is still given by the relation (31) while $`B(\kappa )`$ is now given by $`B(\kappa )=\delta /16a\kappa ^2`$.
The magnetization curve of figure 7 shows both the numerical results and a plot of the magnetization deduced from (54) using (14). We notice that the magnetization of a mesoscopic disk is modified when the interactions between vortices are taken into account. The period and the amplitude of the jumps are not constant anymore; besides, the non-linearity of the curve between two successive jumps is enhanced. These important features of the $`MH_e`$ curve were observed in previous experimental and numerical results . Here we have shown that these features are a consequence of vortex interactions.
### B Two-body interaction energy
The exponent $`\alpha (n)`$ in the relations (51) or (53) allows to describe the interacting potential between vortices. It is interesting to compare the result (51) with the energy of $`n`$ vortices obtained by assuming a two-body interaction. In this case the energy of the whole system of $`n`$ vortices can be written as a sum of two terms
$$\frac{1}{2\pi }=n𝒰_S+\frac{n(n1)}{2}𝒰_I(0)$$
(55)
where $`𝒰_S`$ represents, as noted before, the self-energy of a vortex and $`𝒰_I`$ the two body interaction potential. Using the data of , we can estimate these two energies to the leading order in $`(\kappa \sqrt{2}1)`$. We obtained:
$`𝒰_S`$ $`=`$ $`1+\beta _1(\kappa \sqrt{2}1)\text{ with }\beta _10.4`$ (56)
$`𝒰_I(r)`$ $`=`$ $`\beta _2(\kappa \sqrt{2}1)\mathrm{min}(1,\mathrm{exp}(C(r{\displaystyle \frac{1}{\kappa }})))`$ (57)
with $`\beta _2\frac{1}{4}`$ and $`C\frac{1}{2}`$. From this analysis, and assuming only two-body interaction, we derive an approximate value for the free energy of a configuration with $`n`$ vortices placed at the same point:
$$\frac{1}{2\pi }(\kappa ,n)=n𝒰_S+\frac{n(n1)}{2}𝒰_I(0)n+(\kappa \sqrt{2}1)n\left(\beta _1+\beta _2\frac{n1}{2}\right)$$
(58)
If we compare this relation to the previous expression (53) we find that instead of the sublinear function $`\alpha (n)`$ we have a linear behaviour $`\beta _1+\frac{n1}{2}\beta _2`$. Hence, the function $`\alpha (n)`$ takes into account not only two-body interactions among vortices but also multiple interactions which are present for values of $`\kappa `$ around the dual point unlike the large $`\kappa `$ limit where only the two-body contribution remains.
## VII vortex/edge interactions in system without cylindrical symmetry
In this section, we calculate the energy at the dual point of a system with only one vortex that is not located at the center of the disk. Such a configuration is not in thermodynamic equilibrium and its free energy can be related to a surface energy barrier (analogous to the classical Bean-Livingston barrier in the London limit). We first show that even when the cylindrical symmetry is broken, the system can still be separated into bulk and edge domains.
### A Bulk and edge domains. The curve $`\mathrm{\Gamma }`$
We have seen in section III B that when one or more vortices are located at the center of the disk, there exists a circle $`\mathrm{\Gamma }`$ on which the current vanishes identically. This circle allowed us to define a bulk and an edge domain and to identify the bulk energy with the fluxoid.
If all the vortices are not placed at the center of the disk (i.e. the configuration is not cylindrically symmetric) there is in general no curve of zero current. However the curve $`\mathrm{\Gamma }`$ has now the following property: at each point $`M`$ of $`\mathrm{\Gamma }`$ the current $`\stackrel{}{ȷ}`$ is normal to $`\mathrm{\Gamma }`$. The existence of such a curve is shown by the following argument. Consider a disk with only one vortex $`V`$ situated at a point different from the center of the disk. Take a line segment joining the vortex $`V`$ to the closest point $`S`$ on the boundary of the disk (see figure 8). The component of the current density normal to the $`VS`$ segment changes its sign when one goes from $`V`$ to $`S`$. Hence, there exists a point $`M`$ along this segment where the current either vanishes or is parallel to $`VS.`$ To draw the curve $`\mathrm{\Gamma }`$ we start from $`M`$ in a direction orthogonal to the $`VS`$ segment, and then $`\mathrm{\Gamma }`$ is constructed via infinitesimal steps by imposing that at a point $`M^{}=M+dM`$, very close to $`M`$, the direction of $`\mathrm{\Gamma }`$ is orthogonal to the direction of the current at $`M^{}.`$
Although we lack a general proof, we believe on topological grounds that for vortices at arbitrary positions, there always exists a $`\mathrm{\Gamma }`$ curve which is everywhere orthogonal to the current (one should note that $`\mathrm{\Gamma }`$ does not necessarily have only one connected component). In , we shall present a numerical construction of $`\mathrm{\Gamma }.`$ In the sequel of this work we assume that $`\mathrm{\Gamma }`$ exists, that it encircles all the vortices, and consists of one or many simple closed curves. We shall call the curve $`\mathrm{\Gamma }`$ the separatrix.
Using $`\mathrm{\Gamma },`$ the domain $`\mathrm{\Omega }`$ can be decomposed in two regions $`\mathrm{\Omega }_1`$ and $`\mathrm{\Omega }_2`$ such that:
(i) $`\mathrm{\Omega }_1\mathrm{\Omega }_2=\mathrm{\Omega }`$ ;
(ii) $`\mathrm{\Omega }_1`$ contains all the vortices ( $`\mathrm{\Omega }_1`$ may have multiply connected components);
(iiii) $`\mathrm{\Omega }_2`$ contains the edge of the disk;
(iv) the separatrix $`\mathrm{\Gamma }`$ is the boundary between $`\mathrm{\Omega }_1`$ and $`\mathrm{\Omega }_2`$ and is everywhere normal to the current density.
The remarkable property of the separatrix implies that along $`\mathrm{\Gamma }`$ one can write:
$$_\mathrm{\Gamma }(\stackrel{}{ȷ}+\stackrel{}{A}).\stackrel{}{d}l=_\mathrm{\Gamma }\left(\frac{\stackrel{}{ȷ}}{|\psi |^2}+\stackrel{}{A}\right).\stackrel{}{d}l=_\mathrm{\Gamma }\stackrel{}{}\chi .\stackrel{}{d}l$$
(59)
since along $`\mathrm{\Gamma }`$, $`\stackrel{}{ȷ}.\stackrel{}{d}l=0`$. Since the separatrix is the boundary of $`\mathrm{\Omega }_1`$, the property (59) ensures that the total magnetic flux through $`\mathrm{\Omega }_1`$ is quantized. Hence, at the dual point, we can again use the method of Bogomoln’yi and find the free energy of $`\mathrm{\Omega }_1`$ to be a purely topological number, just as for an infinite domain, even if the cylindrical symmetry is broken.
### B Free energy of one vortex: the surface energy barrier.
As before, we estimate the contribution $`(\mathrm{\Omega }_2)`$ to the total free energy via a variational Ansatz, taking the modulus of the order parameter to be constant. To obtain a qualitative result for the surface energy barrier we neglect the magnetic energy so that, at the dual point, we have:
$`{\displaystyle \frac{1}{2\pi }}(\mathrm{\Omega }_2)`$ $``$ $`{\displaystyle _{\mathrm{\Omega }_2}}|\psi |^2|\stackrel{}{}\chi \stackrel{}{A}|^2+{\displaystyle \frac{(1|\psi |^2)^2}{2}}`$ (60)
$``$ $`{\displaystyle \frac{\delta }{2a}}\left(\psi _0^2v_s^2+(1\psi _0^2)^2\right)`$ (61)
where
$$v_s^2=\frac{d\theta }{2\pi }|\stackrel{}{}\chi \stackrel{}{A}(R)|^2$$
(62)
is the superfluid velocity square averaged over the boundary of the disk. As before, we have replaced the integral over $`\mathrm{\Omega }_2`$ by a line integral along the boundary of the sample (i.e. the disk of radius $`R`$) multiplied by an effective length $`\delta .`$ The function $`\chi `$ appearing in (61) is the phase of the order parameter, and the vector potential is taken, as before, to its value on the boundary of the sample. Optimizing (61) with respect to $`\psi _0`$ we find that:
$`\psi _0^2`$ $`=`$ $`1{\displaystyle \frac{v_s^2}{2}}`$ (63)
$`{\displaystyle \frac{1}{2\pi }}(\mathrm{\Omega }_2)`$ $`=`$ $`{\displaystyle \frac{\delta }{2a}}\left(v_s^2{\displaystyle \frac{v_s^2^2}{4}}\right)`$ (64)
for $`v_s^2\sqrt{2}`$. The phase function $`\chi `$ and the vector potential near the edge of the disk are calculated in Appendix B. Using these results, we obtain (for n=1):
$$\frac{1}{2\pi }(\mathrm{\Omega }_2)=\frac{\delta }{2a}\left(a(1\varphi _e)^2\frac{a^3}{4}(1\varphi _e)^4\right)+f(x,a,\varphi _e1)\delta $$
(65)
The function $`f(x,a,\varphi _e1)`$ determines the dependance of the free energy on the position $`x`$ of the vortex; hence, it measures the interaction energy between the edge currents and the vortex as a function of its position. It is given by
$$f(x,a,\varphi _e1)=\frac{2ax^2}{1x^2}(\varphi _e1)^2\left(1a^2\frac{(\varphi _e1)^2}{1x^2}\right)$$
(66)
From this expression, we observe that the edge currents tend to confine the vortex inside the system. In figure 9 the surface energy as a function of the position $`x`$ of the vortex is plotted. According to (63), only the increasing part of the curve is physical. We nevertheless plot the curve defined by (66) in the whole range $`0x1`$ in order to emphasize the similarity between our result and the well-known Bean-Livingston surface barrier effect that was first derived using the London theory .
## VIII Conclusion
In this work, we have obtained analytical results for the free energy and the magnetization of a mesoscopic superconductor. We have used a known exact solution for the two dimensional Ginzburg-Landau equations in an infinite plane, valid at the dual point, to study a finite system with boundaries. With the help of numerical simulations, we have carried out a perturbative calculation in the vicinity of the dual point. This approach enabled us to study thermodynamically stable states but also metastable states (to obtain a surface energy barrier). This model gives theoretical insights into the physical mechanisms involved in the experimental results of and our analytical results agree quantitatively with experimental measurements. In fact, other related thermodynamic quantities such as the surface tension measuring the thermodynamic stability of vortex states can also be computed along this way and could generalize to two-dimensional systems previous results obtained in one dimension .
More generally, we believe that a theoretical study in the vicinity of the dual point provides a lot of information about the Ginzburg-Landau equations. Although one usually relies on exact results derived from London’s equation, one should be aware of the fact that these results agree with numerical simulations of Ginzburg-Landau equations only when $`\kappa `$ is large (typically $`\kappa 50).`$ We verified that the behaviour we found in the vicinity of the dual point, such as the scaling of the free energy, remains valid when $`\kappa `$ ranges from 0.1 to 10 and this interval of values is indeed relevant for many conventional superconductors.
Our study can be extended in many directions. The scaling results in the vicinity of $`\kappa =1/\sqrt{2}`$ were derived from numerical simulations: a systematic perturbative expansion around the dual point would put them on a more rigorous basis. Secondly, a linear stability analysis of the cylindrically symmetrical solution should allow to understand the fragmentation transition between a giant vortex and unit vortices. Since the separatrix $`\mathrm{\Gamma }`$ exists even for vortex configurations breaking cylindrical symmetry, our approach can be used to analyze hysteretic behaviour of metastable states, and to study polygonal vortex configurations found numerically in mesoscopic superconductors .
## IX Acknowledgments
It is pleasure to thank G. Dunne for numerous discussions. K.M. would like to express his gratitude to S. Mallick for his constant help during the preparation of this work and to thank A. Lemaitre for many interesting discussions. During the completion of this work, we received a preprint by G.S. Lozano et al. which contains results similar to ours for the case of non interacting vortices. It is a good opportunity to thank Gustavo Lozano for correspondence about his results.
This research was supported in part by the U.S.–Israel Binational Science Foundation (BSF), by the Minerva Center for Non-linear Physics of Complex Systems, by the Israel Science Foundation, by the Niedersachsen Ministry of Science (Germany) and by the Fund for Promotion of Research at the Technion.
## X Appendix A: The Ginzburg-Landau equations in a cylindrically symmetric system
For a cylindrically symmetric system, we can use $`\psi =f(r)e^{in\theta }`$ and $`\stackrel{}{A}=A(r)\widehat{e}_\theta `$ where $`n`$ is a non-negative integer which represents the number of vortices at the center of the system. We also define the superfluid velocity $`\stackrel{}{v}_s=v_s(r)\widehat{e}_\theta `$, where
$$v_s(r)=\left(\frac{n}{r}A(r)\right)$$
(67)
In this case the Ginzburg-Landau equations are:
$`{\displaystyle \frac{d^2f}{dr^2}}+{\displaystyle \frac{1}{r}}{\displaystyle \frac{df}{dr}}v_s^2f`$ $`=`$ $`2\kappa ^2f(1f^2)`$ (68)
$`{\displaystyle \frac{d}{dr}}\left({\displaystyle \frac{1}{r}}{\displaystyle \frac{d}{dr}}(rv_s)\right)`$ $`=`$ $`2v_sf^2`$ (69)
It is convenient to define the quantity $`p(r)=rv_s(r)`$. The magnetic field $`\stackrel{}{B}=B(r)\widehat{e}_z`$ is given in terms of $`p(r)`$ by
$$B(r)=\frac{1}{r}\frac{dp}{dr}$$
(70)
We obtain finally two coupled ordinary differential equations
$`f^{\prime \prime }`$ $`=`$ $`2\kappa ^2f(1f^2)+p^2f^2/r^2f^{}/r`$ (71)
$`p^{\prime \prime }`$ $`=`$ $`2pf^2+p^{}/r`$ (72)
with the following boundary conditions at $`r=a^1`$ for $`n0`$:
$$\begin{array}{cc}f(0)=0\hfill & f^{}(a^1)=0\hfill \\ p(0)=n\hfill & p(a^1)=n\varphi _e\hfill \end{array}$$
(73)
for a disk and
$$\begin{array}{cc}f(0)=0\hfill & f^{}(a^1)=0\hfill \\ p(0)=n\hfill & p^{}(a^1)=2a\varphi _e\hfill \end{array}$$
(74)
for a cylinder. These are the equations we have solved numerically using the relaxation method . From the analysis of the equations (72) we deduce the following behaviour in the vicinity of the center of the disk:
$$fr^n\text{ and }pr^2\text{ when }r0$$
The free energy (2) is then given in terms of the solution of (72) by
$$\frac{}{2\pi }=_0^{1/a}r𝑑r\left(\frac{B^2}{2}+\kappa ^2\left(1f^4\right)\right)$$
(75)
## XI Appendix B: phase and vector potential of an off centered configuration with one vortex
In this appendix we measure the distances in units of $`R`$, so the disk has unit radius. Suppose that the vortex is located at a distance $`x`$ from the center of the disk $`(0x<1)`$. The phase $`\chi (\rho ,\theta )`$ of the order parameter satisfies $`\mathrm{\Delta }\chi =0`$ everywhere on the disk except on the vortex with boundary condition $`\widehat{𝐧}.\stackrel{}{}\chi =0`$
Using the image method, the phase $`\chi (\rho ,\theta )`$ at a point located at a distance $`\rho `$ from the center of the disk (with $`0\rho 1`$) is given by :
$$\chi (\rho ,\theta )=\text{Im}\mathrm{ln}\left(\frac{\rho \mathrm{exp}(i\theta )x}{\rho \mathrm{exp}(i\theta )x^1}\right)$$
(76)
where Im denotes the imaginary part part of a complex-valued function. Or equivalently:
$$\mathrm{tan}\chi (\rho ,\theta )=\frac{1x^2}{1+x^2}\frac{\mathrm{sin}\theta }{\mathrm{cos}\theta \frac{\rho +\rho ^1}{x+x^1}}$$
(77)
On the boundary of the disk, $`\rho =1`$, and one finds that
$$\frac{\chi }{\theta }=\frac{1x^2}{1+x^2}\frac{1}{1\frac{2x}{1+x^2}\mathrm{cos}\theta },\frac{\chi }{\rho }=0$$
(78)
therefore
$$\frac{d\theta }{2\pi }|\stackrel{}{}\chi (1,\theta )|^2=a^2\frac{1+x^2}{1x^2}$$
(79)
The vector-potential $`\stackrel{}{A}(R)`$ at the boundary of the sample is a function of the polar angle $`\theta `$ since the vortex is not at the center of the disk. We determine $`\stackrel{}{A}(R)`$ from the following conditions:
$$\stackrel{}{}.\stackrel{}{A}=0,_\mathrm{\Omega }\stackrel{}{A}(R).d\stackrel{}{l}=\varphi _e$$
and on the boundary $`\stackrel{}{A}(R).\widehat{𝐧}=0`$. The following choice,
$$\stackrel{}{A}(R)=\varphi _e\stackrel{}{}\chi $$
(80)
valid near boundary of the system, satisfies these requirements.
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# PRE, submitted Analytic Calculation of the Anomalous Exponents in Turbulence: Using the Fusion Rules to Flush Out a Small Parameter
## I Introduction
The aim of this paper is to build on previous work to achieve a controlled evaluation of the anomalous exponents that characterize various correlation and structure function in Navier-Stokes turbulence, and in particular the exponents $`\zeta _n`$ that characterize $`n`$th order structure functions. The main result of this paper is that given a single experimental input (for example the value of the anomalous exponent of the second order structure function), the $`n`$ dependence of all the other exponents that were reliably measured in experiments and simulations can be calculated analytically.
Decades of experimental and theoretical attention (see for example ) have been devoted to two types of simultaneous correlation functions; the first type includes the structure functions of velocity differences,
$$S_n(𝑹)=|𝒖(𝒓+𝑹)𝒖(𝒓)|^n,$$
(1)
where $`\mathrm{}`$ stands for a suitably defined ensemble average. A second type of correlations include gradients of the velocity field. An important example is the rate $`ϵ(𝒓,t)`$ at which energy is dissipated into heat due to viscous damping. This rate is roughly $`\nu |𝒖(𝒓,t)|^2`$. An often-studied simultaneous correlation function of $`\widehat{ϵ}(𝒓,t)=ϵ(𝒓,t)\overline{ϵ}`$ is
$$K_{ϵϵ}(𝑹)=\widehat{ϵ}(𝒓+𝑹)\widehat{ϵ}(𝒓).$$
(2)
It has been hypothesized by Kolmogorov in 1941 (K41) and 1962 (K62) that statistical objects of this type exhibit power law dependence on $`R`$ :
$$S_n(R)R^{\zeta _n},K_{ϵϵ}(R)R^\mu .$$
(3)
In addition, the K41 theory predicted the values of $`\zeta _n`$ to be $`n/3`$. Experimental measurements and computer simulations show that in some aspects K41 was remarkably close to the truth. The major aspect of its predictions, that the statistical quantities depend on the length scale $`R`$ as power laws, is corroborated by experiments. On the other hand, the predicted exponents seem not to be exactly realized. The numerical values of $`\zeta _n`$ deviate progressively from $`n/3`$ when $`n`$ increases . K62 tried to improve on this prediction by taking into account the fluctuations in the rate of energy dissipation. On the basis of a phenomenological model, assuming the distributions function of energy dissipation to be lognormal, K62 reached the predictions
$$\zeta _n=\frac{n}{3}\frac{\mu n(n3)}{18}.$$
(4)
Besides the fact that these predictions did not follow from fluid mechanical considerations, it was pointed out that they are violating basic inequalities that do not allow the exponents $`\zeta _n`$ to decrease, something that always happens with (4) with $`n`$ large enough. The quest for computing the scaling exponents from the equations of fluid mechanics was long, arduous, and on the whole pretty unsuccessful.
In this paper we present an approach that is based on our own previous findings which culminates in the analytic calculation of exponents like $`\zeta _n`$ and $`\mu `$. At present the calculation is not completely from first principle. We need the input of one number from experiment, say $`\delta _2\zeta _22/3`$. Given this number we can calculate all the other exponents systematically with $`\delta _2`$ being a small parameter that organizes our calculations. To first order in $`\delta _2`$ we recapture (4). We will show that the result to $`O(\delta _2)`$ is universal, independent of the details of the calculations performed below. To second order we find the results
$`\zeta _n`$ $`=`$ $`{\displaystyle \frac{n}{3}}{\displaystyle \frac{n(n3)}{2}}\delta _2[1+2\delta _2(n2)b_2]+O(\delta _2^3),`$ (5)
$`\mu `$ $`=`$ $`9\delta _2(1+8b_2\delta _2)+O(\delta _2^3).`$ (6)
The curves of $`\zeta _n`$ vs. $`n`$ are shown in Fig. 1. We show the K41 prediction, the result of our calculation to 1-loop order, and the 2-loop result that is presented in Eq. (5). While the form of these results is universal, the numerical value of the dimensionless parameter $`b_2`$ depends on the details of the calculations; we find that $`b_2`$ is always negative and of the order of unity. Note that the 2-loop results correct the unwanted down curving of the 1-loop calculation (which is the same malaise as in K62).
In thinking about the strategy for this work we were led by some insights that developed in the context of understanding how to compute the scaling exponents of the Kraichnan model of passive scalar advection . In that model a scalar field $`T(𝒓,t)`$ is advected by a Gaussian velocity field $`𝒖(𝒓,t)`$ which is delta-correlated in time but which has a scaling exponent $`\zeta _2=ϵ`$. For $`ϵ=0`$ the advected scalar has trivial statistics, and for $`ϵ`$ small the model has a natural small parameter. It turned out that the calculation of the exponents can proceed along two ways. The first, which is non-perturbative, was pioneered in . It considers the differential equations that the $`n`$th order correlation functions satisfy, and identifies the anomalous scaling solutions as the zero modes of these differential equations. The calculation of the exponents then depends on the calculation of the zero modes themselves, a task that is not at all easy, and therefore such calculations were never done for any order but $`O(ϵ)`$. In this method the renormalization scale is the outer scale $`L`$, not the inner scale $`\eta `$, and the dimensionless ratio of scales that carries the anomalous part of the exponents of the structure functions is $`L/R`$ where $`R`$ is defined in Eq. (1). A second method that was discussed in detail in considers instead of the correlation functions the averages of higher moments of $`T`$, i.e. $`|T|^n`$. These quantities diverge as powers of $`L/\eta `$, and the exponents of this power are the same as the anomalous part of the exponent of the $`n`$th order structure function. The great advantage of the second method is that one can write a perturbative theory in $`ϵ`$ for the scaling exponents themselves, without any need to compute the zero modes or any other functions of many variables. Thus the second method allows easily computations to $`O(ϵ^2)`$ and with some more efforts to higher orders. The insight gained is that one needs to focus on a quantity that offers the most feasible calculation, exposing the anomalous part of the exponents as it appears with a dimensionless ratio of length scales.
In Navier Stokes turbulence the situation is similar. On the one hand we have a non-perturbative theory which in this case is the infinite hierarchy formed by the equations of motion of the correlation function . We can use this hierarchy to demonstrate that anomalous solutions exist, but the computation of the scaling exponents requires a calculation of the correlation functions themselves. This is a very difficult task that up to now was not accomplished in a satisfactory manner. The other approach will be described in this paper. It will be a perturbative theory for the scaling exponents themselves, not requiring the computation of the correlation function along the way. Similarly to the second method in the Kraichnan problem it will be based on considering limits of correlation functions when $`p`$ coordinates fuse. In that limit we create, even when all the distances are in the inertial range, a ratio of large and small lengthscales that appears raised to the power of the anomalous exponent $`\zeta _n`$.
The two previous findings that influence crucially the present formulation are the mechanism for anomalous scaling that was announced in Ref., and the fusion rules that were discussed in Ref.. In short, Ref. exposed the ladder diagrams which appear in the theory of turbulence. These diagrams contain logarithmic divergences that are summable to power laws with anomalous exponents. These ladder diagrams contain “rungs” of the ladder, that are actually vertices with four, six and more “legs”, representing 4-point, 6-point and higher order interaction amplitudes. In Ref. these objects were represented in terms of infinite series of diagrams that could not be resummed analytically. This is where the fusion rules are now called to save the day. The fusion rules determine the asymptotic properties of $`n`$-point correlation functions when subgroups of $`p`$ coordinates coalesce together. As such the fusion rules are non-perturbative, and are believed to be exact. We use the fusion rules to determine the asymptotic properties of the rungs. This is done such that a calculation of $`\zeta _2`$ from the theory will agree with the experimental value of $`\zeta _2`$. We then show that the knowledge of the asymptotics suffices for constructing a calculation of all the other scaling exponents, and in particular of $`\zeta _n`$ for $`n>2`$ and of $`\mu `$. The crux is that in the process of determining the analytic form of the rungs we discover that their amplitude contains powers of a dimensionless small parameter $`\delta _2=\zeta _22/30.03`$. Using this small parameter in the renormalized 4-point interaction allows us to develop a systematic expansion in orders of $`\delta _2`$. At the end of this paper we sketch a way to understand the remaining task regarding the origin of the small parameter $`\zeta _22/3`$.
In Sect. II we summarize past results that are necessary for the present developments. In Sect. III we show how the fusion rules can be used to determine the properties of the rungs in the ladder diagrams appearing in the $`n`$th order correlation functions. The first important new result is demonstrated in Sect. V \- the numerical coefficient contained in the 4-point rung is shown to be small, of the order of the anomaly of $`\zeta _2`$. This result is crucial since it allows (to our knowledge for the first time) to develop a perturbative calculation of the anomalous parts of all the other exponents. The physical reason for this result is that ( in Sect. IV) we are developing the theory around the K41 solution instead of the dissipative solution as was always attempted. In Sect. V we pave the way for the calculation of all the other exponents. In Sect. VI we calculate the scaling exponents by resumming the logarithmically divergent ladder diagrams up to $`O(\delta _2)`$ (which is known in the field-theoretic jargon as the “1-loop order” in the renormalized rungs). We find (admittedly to our surprise) that to this order the scaling exponents are identical to K62. Like the latter they suffer from the violation of the known requirement that $`\zeta _n`$ cannot decrease with $`n`$ . In Sect. VII we show that the 2-loop order cures the malaise of K62, and we present the result (5) for $`\zeta _n`$ that in our theoretical estimate is valid for $`n12`$. The exponent $`\mu `$ is also computed in this Section. If one wanted results for $`\zeta _n`$ with higher values of $`n`$ one would need to go to 3-loop order (and see Sect. VIII where the form of $`\zeta _n`$ to this order is presented), but the present experimental situation does not warrant a theoretical prediction of $`\zeta _n`$ for very high values of $`n`$. In Sect. VIII we summarize the paper, paying special attention to the range of validity of the theory and to demonstrating that no uncontrolled approximations were made.
## II Summary of pertinent previous results
In this Section we present a brief summary of some past work which is most pertinent. We refer to as Paper I, to as Paper II and to as Paper III.
### A The basic perturbation theory
The starting point of the analysis are the Navier-Stokes equations for the velocity field of an incompressible fluid with kinematic viscosity $`\nu `$ which is forced by an external force $`𝒇(𝒓,t)`$:
$$\left(\frac{}{t}\nu \mathrm{\Delta }^2\right)𝒖+\stackrel{}{𝑷}(𝒖\mathbf{})𝒖=\stackrel{}{𝑷}𝒇,$$
(7)
where $`\stackrel{}{𝑷}`$ is the transverse projection operator $`\stackrel{}{𝑷}\mathrm{\Delta }^2\mathbf{}\times \mathbf{}\times `$. It is well known, (and see for example Paper I) that developing a perturbative approach for the correlation functions and response functions in terms of the Eulerian velocity $`𝒖(𝒓,t)`$ results in a theory that is plagued with infra-red divergences. On the other hand one can transform to new variables, and after the transformation (which amounts to infinite partial resummations in the perturbation theory) one finds a renormalized perturbation theory that is finite, without any divergences in any order of the expansion (cf. and Paper I) . One can achieve such a theory using Lagrangian variables ; we find it technically simpler to employ the Belinicher-L’vov transformation,
$$𝒗[𝒓_0|𝒓,t]𝒖[𝒓+𝝆(𝒓_0,t),t)],$$
(8)
where $`𝝆(𝒓_0,t)`$ is the Lagrangian trajectory of a fluid point which has started at point $`𝒓=𝒓_0`$ at time $`t=t_0`$
$$𝝆(𝒓_0,t)=_0^t𝒖[𝒓+𝝆(𝒓_0,\tau ),\tau ]𝑑\tau .$$
(9)
The natural variables for a divergence free theory are the velocity differences
$$𝒘(𝒓_0|𝒓,t)𝒗[𝒓_0|𝒓,t]𝒗[𝒓_0|𝒓_0,t].$$
(10)
Since the averages of quantities that depend on one time only can be computed at $`t=0`$, it follows that the average moments of these BL-variables are the structure functions of the Eulerian field $`S_n(𝑹)`$ defined by Eq. (1). It was shown that these variables satisfy the Navier Stokes equations, and that one can develop (cf. Paper I) a perturbation theory of the diagrammatic type in which the natural quantities are the Green’s function $`G_{\alpha \beta }(𝒓_0|𝒓,𝒓^{},t,t^{})`$ and the correlation function $`F_{\alpha \beta }(𝒓_0|𝒓,𝒓^{},t,t^{})`$:
$`G_{\alpha \beta }(𝒓_0|𝒓,𝒓^{},t,t^{})`$ $`=`$ $`{\displaystyle \frac{\delta w_\alpha (𝒓_0|𝒓,t)}{\delta f_\beta (𝒓^{},t^{})}}|_{f0},`$ (11)
$`F_{\alpha \beta }(𝒓_0|𝒓,𝒓^{},t,t^{})`$ $`=`$ $`w_\alpha (𝒓_0|𝒓,t)w_\beta (𝒓_0|𝒓^{},t^{}).`$ (12)
Physically the Green’s function is the mean response of the velocity difference to the action of a vanishingly small forcing. In stationary turbulence these quantities depend on $`t^{}t`$ only, and we can denote this time difference as $`t`$. The quantities satisfy the well known and exact Dyson and Wyld coupled equations. The Dyson equation reads
$`[{\displaystyle \frac{}{t}}`$ $``$ $`\nu \mathrm{\Delta }]G_{\alpha \beta }(𝒓_0|𝒓,𝒓^{},t)=G^0_{\alpha \beta }(𝒓_0|𝒓,𝒓^{},0^+)\delta (t)`$ (13)
$`+`$ $`{\displaystyle 𝑑𝒓_2G_{\alpha \delta }^0(𝒓_0|𝒓,𝒓_2,0^+)𝑑𝒓_1_0^t𝑑t_1}`$ (14)
$`\times `$ $`\mathrm{\Sigma }_{\delta \gamma }(𝒓_0|𝒓_2,𝒓_1,t_1)G_{\gamma \beta }(𝒓_0|𝒓_1,𝒓^{},tt_1),`$ (15)
where $`G_{\alpha \beta }^0(𝒓_0|𝒓,𝒓^{},0^+)`$ is the bare Green’s function determined by Eq. (3.20) in Paper I. The Wyld equation has the form
$`F_{\alpha \beta }(𝒓_0|𝒓,𝒓^{},t)={\displaystyle 𝑑𝒓_1𝑑𝒓_2\underset{0}{\overset{\mathrm{}}{}}𝑑t_1𝑑t_2G_{\alpha \gamma }(𝒓_0|𝒓,𝒓_1,t_1)}`$ (16)
$`\times `$ $`\left[D_{\gamma \delta }(𝒓_1𝒓_2,tt_1+t_2)+\mathrm{\Phi }_{\gamma \delta }(𝒓_0|𝒓_1,𝒓_2,tt_1+t_2)\right]`$ (17)
$`\times `$ $`G_{\delta \beta }(𝒓_0|𝒓^{},𝒓_2,t_2).`$ (18)
In Eq. (14) the “mass operator” $`\mathrm{\Sigma }`$ is related to the “eddy viscosity” whereas in Eq. (18) the “mass operator” $`\mathrm{\Phi }`$ is the renormalized “nonlinear” noise which arises due to turbulent excitations. Both these quantities are dependent on the Green’s function and the correlator, and thus the equations are coupled.
The main result of Paper I is a demonstration of the property of “locality” in the Dyson and Wyld equations. This property means that given a value of $`|𝒓𝒓_0|`$ in Eq. (14), the important contribution to the integral on the RHS comes from that region where $`|𝒓_1𝒓_0|`$ and $`|𝒓_2𝒓_0|`$ are of the order of $`|𝒓𝒓_0|`$. In other words, all the integrals converge both in the upper and the lower limits. The same is true for the Wyld equation, meaning that in the limit of large $`L`$ and small $`\eta `$ these length scales disappear from the theory, and there is no natural cutoff in the integrals in the perturbative theory. In this case one cannot form a dimensionless parameter like $`L/r`$ or $`r/\eta `$ to carry dimensionless corrections to the K41 scaling exponents. For $`\eta |𝒓𝒓_0|L`$ scale invariance prevails, and one finds precisely the K41 scaling exponents:
$`G_{\alpha \beta }(\lambda 𝒓_0|\lambda 𝒓,\lambda 𝒓^{},\lambda ^zt)`$ $`=`$ $`\lambda ^{\beta _2}G_{\alpha \beta }(𝒓_0|𝒓,𝒓^{},t),`$ (19)
$`F_{\alpha \beta }(\lambda 𝒓_0|\lambda 𝒓,\lambda 𝒓^{},\lambda {}_{}{}^{z}t)`$ $`=`$ $`\lambda ^{\zeta _2}F_{\alpha \beta }(𝒓_0|𝒓,𝒓^{},t).`$ (20)
One can derive two scaling relations which hold order by order, i.e.
$$2z+\zeta _2=2,z+2\zeta _2=2.$$
(21)
The solution is $`z=\zeta {}_{2}{}^{}=2/3`$. It was also shown that the scaling exponent of the Green’s function (19) is $`\beta {}_{2}{}^{}=3`$. Extending such considerations to the higher order structure functions leads to the order-by order K41 prediction that $`\zeta _n=n/3`$.
Of course, the order by order result (19) which leads to (21) is not necessarily the correct one. If one could resum all the diagrammatic expansion one could find nonperturbative answers that may be different. The whole sum of diagrams may diverge when the outer scale goes to infinity or the inner scale to zero, allowing a renormalization scale to creep in even though the order-by-order theory is convergent. The difficulty is that no one knows how to resum the infinite expansion which exhibits no obvious small parameter.
In this paper we will propose a way out of this difficulty. The new thinking is based on the fusion rules. Instead of considering fully unfused correlation functions only, we will allow some coordinates to be much closer together, say within a distance $`𝒓`$, whereas the rest will be separated by a much larger distance, say of the order of $`𝑹`$ where $`rR`$. We will show that we can form a dimensionless ratio with $`R/r`$, and that such ratios carry anomalous exponents that are going to survive the process of fusion of coordinates in correlation functions when we make structure functions. We will thus be able to recognize the anomalous exponents even though at first sight there is no obvious renormalization scale.
To clarify how this mechanism works we need to remind ourselves of the appearance of ladder diagrams in the theory of correlations functions. Such diagrams appear in the most transparent way in nonlinear Green’s functions, and we review briefly our past results on these objects.
### B The Nonlinear Green’s Functions
The nonlinear Green’s function $`G_{2,2}(𝒓_0|x_1,x_2,x_3,x_4)`$ describes the response of a product of two velocity fields taken at different space-time coordinates to the action of two forces $`𝒇`$:
$`G_{2,2}^{\alpha \beta \gamma \delta }(𝒓{}_{0}{}^{}|x{}_{1}{}^{},x{}_{2}{}^{},x{}_{3}{}^{},x{}_{4}{}^{})`$ $`=`$ $`{\displaystyle \frac{\delta ^2w{}_{\alpha }{}^{}(𝒓{}_{0}{}^{}|x{}_{1}{}^{})w{}_{\gamma }{}^{}(𝒓_0|x_2)}{\delta f{}_{\beta }{}^{}(𝒓{}_{0}{}^{}|x{}_{3}{}^{})\delta f{}_{\delta }{}^{}(𝒓{}_{0}{}^{}|x{}_{4}{}^{})}},`$ (22)
where for brevity we use the notation $`x_j\{𝒓{}_{j}{}^{},t{}_{j}{}^{}\}`$. Similarly, one defines the nonlinear Green’s functions $`𝑮_{p,p}`$ as the response of the product of $`p`$ velocity differences between $`p`$ distinct points and $`𝒓_0`$ to the action of $`p`$ forces in different points.
In a Gaussian theory (which ours is not) $`𝑮_{2,2}`$ would be the products of the linear Green’s functions like $`G{}_{}{}^{\alpha \beta }(𝒓{}_{0}{}^{}|x{}_{1}{}^{},x{}_{3}{}^{})G{}_{}{}^{\gamma \delta }(𝒓{}_{0}{}^{}|x{}_{2}{}^{},x{}_{4}{}^{})`$. In a non-Gaussian theory it is natural to assume that this quantity is a homogeneous function of its arguments when they are in the scaling regime. This means that
$`G_{2,2}^{\alpha \beta \gamma \delta }(𝒓{}_{0}{}^{}|\lambda 𝒓{}_{1}{}^{},\lambda {}_{}{}^{z}t{}_{1}{}^{},\lambda 𝒓{}_{2}{}^{},\lambda {}_{}{}^{z}t{}_{2}{}^{},\lambda 𝒓{}_{3}{}^{},\lambda {}_{}{}^{z}t{}_{3}{}^{},\lambda 𝒓{}_{4}{}^{},\lambda {}_{}{}^{z}t{}_{4}{}^{})`$ (24)
$`=`$ $`\lambda ^{\beta _4}G_2^{\alpha \beta \gamma \delta }(𝒓{}_{0}{}^{}|x{}_{1}{}^{},x{}_{2}{}^{},x{}_{3}{}^{},x{}_{4}{}^{})`$ (25)
From the Gaussian decomposition of this quantity we would guess that $`\beta {}_{4}{}^{}=2\beta {}_{2}{}^{}=6`$. The proof of locality in Paper I means that there is no perturbative mechanism to change this scaling index. On the other hand, this quantity, which is a function of four space-time coordinates $`x_i`$ has scaling properties that are not exhausted by the overall scaling exponent $`\beta _4`$. We have shown in Paper II that when we consider its dependence on ratios of space-time coordinates in their asymptotic regimes w e pick up a set of anomalous scaling exponents. The main result of Paper II was that in the regime $`r_1r_2r{}_{3}{}^{}r_4`$ the diagrammatic expansion of this object produces logarithms like ln$`(r{}_{3}{}^{}/r{}_{1}{}^{})`$ to some power. It was explained that the sum of such logarithmically large contributions is given by $`(r{}_{3}{}^{}/r{}_{1}{}^{})^\mathrm{\Delta }`$ with some anomalous exponent $`\mathrm{\Delta }`$. To make the appearance of anomalous exponents evident we review the simplest object that resums to logarithms, i.e. the series of “ladder diagrams”.
The diagrammatic representation of the nonlinear Green’s function (LABEL:c1) is shown in Fig. 2a, where the notation of the diagrams is explained shortly in the figure legend and at length in Appendix A. The 4-point “rungs” appearing in the ladders were represented in Paper II as an infinite series expansion whose beginning is shown in Fig. 2b. The rest of the expansion involves exceedingly complicated diagrams that we failed to resum analytically. Nevertheless, just from general properties one could show that the ladder with $`n`$ rungs contains a contribution of order $`\left[\mathrm{\Delta }\mathrm{ln}(r{}_{3}{}^{}/r{}_{1}{}^{})\right]^n/n!`$. The summation of all these contributions gives a term proportional to $`(r_3/r{}_{1}{}^{})^\mathrm{\Delta }`$, and this is the observation that we want to build on in this paper.
## III Ladder diagrams in $`n`$’th order correlation functions
In this Section we demonstrate how the anomalous exponents of higher order correlation functions can be related to resummed ladder diagrams. The idea is to consider a typical $`n`$th order correlation function and to almost fuse $`p`$ coordinates, $`p<n`$, chosen from the available $`n`$ coordinates. The point to observe is that the diagrammatic theory allows us to write, upon observation, all the topologically possible diagrams appearing in the expansion of a given object. Thus for example consider Fig. 3 where we represent a general $`n`$th order correlation function, in which two coordinates are a distance $`r`$ from each other, and all the rest are a distance $`R`$ from them and from each other. While coalescing the two coordinates we pull out a contribution that reads $`G_{2,2}`$ connected to the rest of the diagram with $`n2`$ coordinates. We know that this contribution exists since it is allowed topologically. But this fragment is represented in its turn by the sum of the ladder diagrams that we present in Fig. 2. We thus state that whenever we are about to fuse two coordinates in any $`n`$th order correlation function we can expose a series of diagrams that are the same as those that appear in the expansion of the nonlinear Green’s function $`G_{2,2}(𝒓_0|x_1,x_2,x_3,x_4)`$, together with the logarithmic divergences that are associated with them. In Appendix A we explain that in doing so we really take into account all the necessary contributions, leaving nothing uncontrolled.
Now we employ the fusion rules. These tell us that as a function of the two fusing coordinates the $`n`$th order correlation function is a homogeneous function whose homogeneity exponent is $`\zeta _2`$. Accordingly, if we succeed to resum the ladder diagrams in the limit $`rR`$ we should find an anomalous part of $`\zeta _2`$ from the dimensionless power $`\left(R/r\right)^\mathrm{\Delta }`$.
Similarly, we can almost fuse 3, 4 or $`p`$ coordinates, and accordingly pull out of the diagram for the $`n`$th order correlation functions more complex ladders with 3, 4 or $`p`$ struts. For example in Fig. 3b we show how the fusion of 3 coordinates singles out a fragment that is $`G_{3,3}`$ whose ladder diagrams have three struts cf. Fig. 4. The fusion rules guarantee that the $`n`$th order correlation is a homogeneous function of the $`p`$ fusing coordinates with $`\zeta _p`$ being the homogeneity exponent. We will show that these more complex ladders also resum into power laws in $`R/r`$, being responsible for the anomalous parts of $`\zeta _p`$.
At this point all this is a bit formal, since we do not have an explicit form of the rungs in the ladder diagrams, and we can compute nothing without this knowledge. In the next two Sections we will address this issue and demonstrate that a judicious use of the fusion rules dictates enough knowledge of the rungs to take us through a useful calculation.
## IV Building the theory on the background of K41
In this Section we reorganize the theory such that Kolmogorov’s 41 theory serves as its “free” limit. In other words, we aim at achieving a theory in which resummations of divergent contributions would directly give the anomalous parts of the scaling exponents; the K41 parts should be obvious order-by-order. This is done in two steps, that are correspondingly presented in Subsects. IV A and IV B.
### A Resummation into K41 propagators
It was explained in Subsect. II A that our theory is developed in the BL-representation, to eliminate spurious IR divergences that stem from the sweeping interactions. The main result of Paper I was that after line-resummation each diagram in the BL-diagrammatic expansion of the propagators (Green’s function and double correlation function) converged in the infrared and the ultraviolet regimes. Accordingly, K41 scaling is a solution of the order-by-order theory. Nevertheless, the propagators in the BL-representation lose translational invariance, and are therefore not diagonal in Fourier space. For the purpose of actual calculations it is extremely advantageous to rearrange the theory such that the BL-propagators become again diagonal in Fourier space.
The actual resummation that is necessary is presented in Appendix B. It results in a diagrammatic theory that is topologically exactly the same as the standard Wyld diagrammatic expansion before line resummation. There are two differences as explained in the Appendix B. For the purposes of our considerations below the main issue is the simple form of the the propagators that appear as lines in the diagrams: they exhibit K41 scaling exponents:
$`G_{\alpha \beta }(𝒌,\omega )`$ $`=`$ $`P_{\alpha \beta }(𝒌)g(k,\omega ),g(k,\omega )={\displaystyle \frac{1}{\omega +i\gamma (k)}},`$ (26)
$`F_{\alpha \beta }(𝒌,\omega )`$ $`=`$ $`P_{\alpha \beta }(𝒌)f(k,\omega ),f(k,\omega )={\displaystyle \frac{\varphi (k)}{\omega ^2+\gamma ^2(k)}}.`$ (27)
In these formulae the scaling exponents are carried by
$$\gamma (k)=c_\gamma \overline{ϵ}^{1/3}k^{2/3},\varphi (k)=c_\varphi \overline{ϵ}k^3,$$
(28)
where $`c_\gamma `$ and $`c_\varphi `$ are dimensionless constants.
### B Renormalization to K41 4-point rung
In this Subsection we determine the form of the 4-point rungs of the ladder diagrams in two steps. These two steps are based on the following observation: the diagrammatic expansion of the rung includes many diagrams, some of which contain in them additional subsets of ladder diagrams. In the first step we will consider the rungs as if all the diagrams appearing in their infinite series were resummed, except for their own internal subsets of ladder diagrams. In the second step we will consider also the ladder diagrams appearing in the series for the rung. We aim to a situation in which all the ladders that appear in the theory, like in Fig. 2, contain already renormalized rungs. However, instead of evaluating the rungs from actual resummations we are going to determine their form using the fusion rules. Thus in the first step we find the form of the rung that results, upon fusion, in K41 scaling exponents. In the second step we recognize that the rungs themselves have ladders, leading to an anomalous correction in the scaling properties of the rungs themselves. This being accomplished, we will have our final form of the rung. Then we turn to the ladder diagrams appearing in the fused correlation functions, using the rung as a basic building block of the theory. All anomalies of all the measurable statistical objects will result from resummations of the remaining ladder diagrams.
Consider Fig. 2a, in which the rung appears as an object. It is given in terms of an infinite series of diagrams in Fig. 2b. It is in fact a 4-point vertex depending on four $`𝒌`$-vectors and four frequencies. As a first step we consider the value of the rung when all the frequencies are zero, denoting it in this limit as $`𝑹(𝒌_a,𝒌_b,𝜿_c,𝜿_d)`$. At a later point we will explain that this is sufficient for our purposes. The bare value of this object can be read directly from diagram (1) in Fig. 2b, with two bare BL-vertices $`𝚪`$ and one double correlation function. The answer is
$`R_0^{\alpha \beta \gamma \delta }(𝒌_a,𝒌_b,𝜿_c,𝜿_d)`$ (29)
$`=\delta _0\overline{ϵ}^{1/3}{\displaystyle \frac{\mathrm{\Gamma }^{\alpha \gamma \sigma }(𝒌_a,𝜿_c,𝒌_e)\mathrm{\Gamma }^{\beta \delta \sigma }(𝒌_b,𝜿_d,𝒌_e)}{k_e^{13/3}}},`$ (30)
where $`𝒌_e𝒌_a𝜿_c=𝜿_d𝒌_b`$, and $`\delta _0`$ is a dimensionless constant.
We demonstrate now that if we use this bare form of the rung the fusion rules would predict dissipative exponents, $`\zeta _n=n`$. We first demonstrate this in the context of $`\zeta _2`$. Consider a general $`n`$ order correlation function as in Fig. 3a, and fuse two coordinates, pulling out the fragment of $`𝑮_{2,2}`$ as shown on the RHS of Fig. 3a. We will now compute the scaling exponent by finding the $`r`$ dependence of this fragment when the two coordinates approach each other to a small distance $`rR`$ where $`R`$ is the typical distance between all the other coordinates. To find the $`r`$ dependence we must integrate according to the explanation in Appendix C, and to this aim we introduce the object $`T_2(r,𝜿)`$:
$`T_2(r,𝜿)={\displaystyle \frac{d𝒌_a}{(2\pi )^3}4\mathrm{sin}^2(\frac{1}{2}𝒌_a𝒓)}`$ (31)
$`\times {\displaystyle }{\displaystyle \frac{d\omega _a}{2\pi }}G_{2,2}(𝒌_a,\omega _a,𝒌_a,\omega _a,𝜿,0,𝜿,0).`$ (32)
We are interested in the $`r`$ dependence of this object in the limit $`r\kappa 1`$. To calculate $`T_2(r,𝜿)`$ in this limit we return to Fig. 2a. Obviously the Gaussian contribution diagram (1) is irrelevant in this limit. The skeleton diagram (2) contributes the following integral:
$`T_2^\mathrm{s}(r,𝜿){\displaystyle \frac{d\omega _a}{2\pi }\frac{d𝒌_a}{(2\pi )^3}4\mathrm{sin}^2(\frac{1}{2}𝒌_a𝒓)}`$ (33)
$`\times g(k_a,\omega _a)g(k_a,\omega _a)R_0(𝒌_a,𝒌_a,𝜿,𝜿),`$ (34)
where the tensor indices of the rung were contracted for the longitudinal contribution. The superscript “s” is used here and below to denote skeleton contributions. We note that in the limit $`k\kappa `$ the BL vertices are proportional to the smallest wavevector $`\kappa `$. Thus the rung is proportional to $`\kappa ^2/k_a^{13/3}`$. Integrating over the frequencies of the two Green’s functions $`g(k_a,\omega _a)`$ in this rung \[cf. Eqs. (26, 28)\] results in the evaluation $`1/[\gamma (k_a)k_a^{13/3}]k_a^5`$. Thus the $`r`$ dependence of $`T_2^\mathrm{s}`$ is given by
$`T_2^\mathrm{s}(r,𝜿){\displaystyle \frac{d𝒌_a}{(2\pi )^3}4\mathrm{sin}^2(\frac{1}{2}𝒌_a𝒓)\frac{1}{k_a^5}}.`$ (35)
Up to logarithmic corrections this integral is proportional to $`r^2`$ which is the dissipative solution. Similarly, if we use the bare rung in the diagram in Fig. 2b to determine $`\zeta _3`$ we will find $`\zeta _3=3`$. In general we will find $`\zeta _n=n`$ instead of the K41 value of $`n/3`$. Now one could think that the correct values of the inertial range exponents may be obtained from resumming all the ladder diagrams with the bare rungs. This was the point of view proposed in Paper I. In such a case the sought after correction to the scaling exponents is of the order of unity, and it is unclear how to develop a controlled resummation. In this paper we point out a new way, based on the existence of a renormalized rung which gives, upon fusion, K41 exponents before ladder resummations. The characteristics of the renormalized rung in such a scheme are dictated by the fusion rules.
Next we want to determine the renormalized form of the rung. To this aim we repeat the exercise of integrating over the two Green’s functions and the rung, with the vertices determined as before in the limit $`k_ak_b\kappa _c\kappa _d`$. But now we leave the exponent of $`k_a`$ in the asymptotic evaluation of the rung free, and demand that the result of the integration will be $`k_a^{5/3}`$. We find that this requires that $`R_0k_a^3`$.
We are now in a position to propose a renormalized form of the rung which in a proper calculation could be obtained by summing up all the non-ladder diagrams that contribute to this rung. This conforms with our basic hypothesis that all non-ladder diagrams contribute to K41, whereas the ladders are responsible for the anomalous scaling. Since K41 does not allow $`L`$ renormalization we propose the form
$`R^{\alpha \beta \gamma \delta }(𝒌_a,𝒌_b,𝜿_c,𝜿_d)`$ (36)
$`=\delta \overline{ϵ}^{1/3}{\displaystyle \frac{\mathrm{\Gamma }^{\alpha \gamma \sigma }(𝒌_a,𝜿_c,𝒌_e)\mathrm{\Gamma }^{\beta \delta \sigma }(𝒌_b,𝜿_d,𝒌_e)}{k_e^3\kappa _c^{2/3}\kappa _d^{2/3}}}.`$ (37)
This form gives the K41 overall scaling exponent (which is the same as in the bare rung $`𝑹_0`$, and in addition agrees with the fusion rules for the second order correlation function with a K41 scaling exponent. In addition it is symmetric, as it should be, for exchanging the indices $`a`$ and $`b`$ together with $`c`$ and $`d`$. Note that $`\delta `$ is now a renormalized unknown dimensionless parameter which will be determined later.
To proceed, we note that our actual calculation (see below) depends really only on the asymptotic properties of the rung, which are rigidly determined by the fusion rules. We can thus attempt to simplify the form of the rung as much as possible, preserving the asymptotic and parity properties unchanged. In particular we note that the BL vertices $`𝚪(𝒌_a,𝜿_c,𝒌_e)`$ have complicated structure which makes calculations involving them rather difficult. Therefore we propose to use instead Eulerian vertices $`𝑽(𝒌_a,𝜿_c,𝒌_e)`$ corrected by a factor $`2(𝒌_b𝜿_c)/[k_a^2+k_b^2+\kappa _c^2]`$. The correction is aimed at reproducing the asymptotic behavior of the BL vertex $`\mathrm{\Gamma }^{\alpha \gamma \sigma }(𝒌_a,𝜿_c,𝒌_e)\mathrm{min}\{k_a,k_b,\kappa _c\}`$. Thus instead of (37) one has
$`R^{\alpha \beta \gamma \delta }(𝒌_a,𝒌_b,𝜿_c,𝜿_d)={\displaystyle \frac{4\delta \overline{ϵ}^{1/3}(𝜿_c𝒌_e)(𝜿_d𝒌_e)}{[k_a^2+\kappa _c^2+k_e^2][k_b^2+\kappa _d^2+k_e^2]}}`$ (38)
$`\times `$ $`{\displaystyle \frac{V^{\alpha \gamma \sigma }(𝒌_a,𝜿_c,𝒌_e)V^{\beta \delta \sigma }(𝒌_b,𝜿_d,𝒌_e)}{k_e^3(\kappa _c\kappa _d)^{2/3}}}.`$ (39)
As a further simplification of the actual calculations we will use a 1-dimensional reduction of the problem (preserving the asymptotic scaling properties and parity) in which instead of 3-dimensional integrations $`d^3k/(2\pi )^3`$ we will use 1-dimensional one $`_{\mathrm{}}^{\mathrm{}}𝑑k/2\pi `$. Then we can disregard the vector indices and replace $`V^{\alpha \gamma \sigma }(𝒌_a,𝜿_c,𝒌_e)k_a`$, $`(𝒌𝒌^{})kk^{}`$ (keeping the signs) and $`k_e^3`$ in the denominator by $`|k_e|`$ (because we replaced 3-d by 1-d integration). The 1-dimensional version of the rung (39) turns into
$$R(k_a,k_b,\kappa _c,\kappa _d)=\frac{4\delta \overline{ϵ}^{1/3}k_ak_b\kappa _c\kappa _d|k_e|}{[k_a^2+\kappa _c^2+k_e^2][k_b^2+\kappa _d^2+k_e^2]|\kappa _c\kappa _d|^{2/3}}.$$
(40)
Note that here $`k_a,k_b,\kappa _c,\kappa _d`$ are in the interval $`\pm \mathrm{}`$ and that they carry signs in order to preserve the parity of the rungs. $`k_a,k_b`$ are incoming wave vectors and $`\kappa _c,\kappa _d`$ are outgoing, and they conserve momentum,
$$k_a+k_b=\kappa _c+\kappa _d.$$
(41)
Substituting into (40) $`k_e=k_a\kappa _c=k_b\kappa _d`$ one gets finally:
$`R(k_a,k_b,\kappa _c,\kappa _d)`$ (42)
$`=`$ $`\delta \overline{ϵ}^{1/3}{\displaystyle \frac{k_ak_b\kappa _c\kappa _d|k_a\kappa _c|}{[k_a^2k_a\kappa _c+\kappa _c^2][k_b^2k_b\kappa _d+\kappa _d^2]|\kappa _c\kappa _d|^{2/3}}}.`$ (43)
In particular,
$`R(k,k,\kappa ,\kappa )={\displaystyle \frac{\delta \overline{ϵ}^{1/3}k^2\kappa ^{2/3}|k\kappa |}{[k^2k\kappa +\kappa ^2]^2}},`$ (44)
$`R(k,k,\kappa ,\kappa ^{})={\displaystyle \frac{\delta \overline{ϵ}^{1/3}\mathrm{sign}(\kappa \kappa ^{})|\kappa \kappa ^{}|^{1/3}}{|k|}}`$ (45)
$`\text{for}\kappa ,\kappa ^{}k.`$ (46)
To check that we get the right K41 scaling exponents with the new renormalized rung we need to recalculate the 1-dimensional version of Eq. (34) with (42) for the rung:
$`T_2^\mathrm{s}(r,\kappa )`$ $`=`$ $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle \frac{dk}{2\pi }}4\mathrm{sin}^2\left({\displaystyle \frac{kr}{2}}\right)R(k,k,\kappa ,\kappa )`$ (47)
$`\times `$ $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle \frac{d\omega }{2\pi }}g(k,\omega )g(k,\omega ).`$ (48)
Using Eq. (26) the frequency integral yields $`1/2\gamma (k)`$. Thus
$$T_2^\mathrm{s}(r,\kappa )=\frac{\delta \kappa ^{2/3}}{c_\gamma \pi }_{\mathrm{}}^{\mathrm{}}𝑑k\frac{|k|^{4/3}|k\kappa |\mathrm{sin}^2\left(\frac{kr}{2}\right)}{[k^2k\kappa +\kappa ^2]^2}.$$
(49)
In the limit $`\kappa 0`$ the integral simplifies to
$$T_2^\mathrm{s}(r,\kappa )=2\stackrel{~}{\delta }\kappa ^{2/3}_0^{\mathrm{}}\frac{dk}{k^{5/3}}\mathrm{sin}^2\left(\frac{kr}{2}\right),$$
(50)
where
$$\stackrel{~}{\delta }\frac{\delta }{\pi c_\gamma }.$$
(51)
This integral is elementary, reading
$$T_2^\mathrm{s}(r,\kappa )=\frac{1}{2}\stackrel{~}{\delta }(\kappa r)^{2/3}\mathrm{\Gamma }\left(\frac{2}{3}\right)2\stackrel{~}{\delta }(\kappa r)^{2/3},$$
(52)
where $`\mathrm{\Gamma }(x)`$ is the gamma-function.
The point to notice is that in the asymptotic limit $`\kappa 0`$ the only properties of the rung that guaranteed the appearance of the scaling exponent 2/3 are the asymptotic properties that we preserved in the series of simplifications leading to (42). In general, we will show below that the series of simplifications of the model form of the rung are of absolutely no import also for the calculation of the anomalous scaling exponents up to 1-loop order. We will show below that we get precisely the same exponents in this order with any arbitrary analytic form of the rung, with tensor indices or without, in 3-d form or 1-d form or whatever, as long as the asymptotics are preserved, as they are. In 2-loop order this is no longer true. The actual numbers obtained in the 2-loop order are model dependent. We will show however that the sensitivity of the predicted exponents $`\zeta _n`$ to the model for the rung is small as long as $`n<8`$. We need the 2-loop order mainly to make sure that it corrects for some unacceptable properties of the 1-loop results for higher order correlation functions.
Before we proceed we need to check the self consistency of our approach. We need to make sure that all higher order non-linear Green’s function $`𝑮_{p,p}`$ (the response of $`p`$ velocities to $`p`$ forcing) yield, upon fusion, the correct K41 exponent for $`p`$th order correlation functions $`\zeta _p^{^{\mathrm{K41}}}=p/3`$. For this aim we have to consider the so called skeleton diagrams which are the lowest order connected diagrams without loops. For $`𝑮_{2,2}`$ this is diagram (2) in Fig. 2a. for $`𝑮_{3,3}`$ the skeleton contribution are shown as diagrams (3) on Fig. 4a. We must make sure that the skeleton diagrams, upon fusion, yield K41 scaling $`\zeta _p^{^{\mathrm{K41}}}=p/3`$ for the appropriate $`p^{\mathrm{th}}`$-order correlation functions, since our grand hypothesis is that the anomalous scaling comes only from ladder resummations.
This test of self consistency is presented in Appendix C. The important conclusion of this Appendix is that the skeleton diagrams for $`G_{p,p}(\{k_j,\kappa _j\})`$ (with asymptotics of the rung defined by the two-point fusion rules with $`\zeta _2^{^{\mathrm{K41}}}=2/3`$) automatically reproduces the K41 scaling exponent $`\zeta _p^{^{\mathrm{K41}}}=p/3`$ when $`p`$ points are fused.
## V Sanding the floor in the one-loop order
In this Section we demonstrate the most important new property of the resummed theory, i.e. that the rungs in the ladder diagrams appear with a small parameter. This will allow us to develop a controlled ladder resummation, something that to our knowledge has never been available before. In fact, we will lay out in this Section all that is needed to calculate the scaling exponents in the 1-loop order. In Subsect. V A we demonstrate that the prefactor $`\delta `$ of the rung is the order of $`\delta _2`$ which is the anomalous part of $`\zeta _2`$ and thus small. In Subsect. V B we consider the anomalous exponent of the rung itself, denoted as $`\delta _a`$, and stemming from ladder resummations within the rung infinite series representation. In Subsect. V C we reconsider the contribution of the skeleton diagrams to the scaling exponents upon fusion, taking into account the anomaly of the rung. In Subsect. V D we throw in the inputs: the fact that $`\zeta _3=1`$ and the experimental value of $`\zeta _2`$. The result is Eq. (81) which states that in the 1-loop order all the unknown parameters are numerically identical. From this point the calculation of all the other scaling exponents in the 1-loop order is straightforward.
### A The 4-point rung is small!
Here we show that the coefficient $`\delta `$ in front of the renormalized 4-point rung (37) is of the order of the correction to K41 of the scaling exponent $`\zeta _2`$:
$$\delta _2=\zeta _2\zeta _2^{^{\mathrm{K41}}}0.03.$$
(53)
To this aim consider the 1-dimensional version of the quantity $`T_2(r,𝜿)`$ of (32):
$`T_2(r,\kappa )={\displaystyle \underset{\mathrm{}}{\overset{\mathrm{}}{}}}{\displaystyle \frac{dk_a}{2\pi }}4\mathrm{sin}^2({\displaystyle \frac{1}{2}}k_ar)`$ (54)
$`\times {\displaystyle \underset{\mathrm{}}{\overset{\mathrm{}}{}}}{\displaystyle \frac{d\omega _a}{2\pi }}G_{2,2}(k_a,\omega _a,k_a,\omega _a,\kappa ,0,\kappa ,0).`$ (55)
We will examine the ratio of the contributions of the 1-loop diagram (3), denoted below as $`T_2^{(1)}(r,\kappa )`$, to the contribution of the skeleton diagram Eq. (48) \[diagram (2) in Fig. 2a\]. After performing the frequency integrals the 1-dimensional form of $`T_2^{(1)}(r,\kappa )`$ \[cf. diagram (3) in Fig. 2a\] is:
$`T_2^{(1)}(r,\kappa )={\displaystyle \underset{\mathrm{}}{\overset{\mathrm{}}{}}}{\displaystyle \frac{dk_a}{\pi \gamma (k_a)}}\mathrm{sin}^2({\displaystyle \frac{1}{2}}k_ar)`$ (56)
$`\times {\displaystyle \underset{\mathrm{}}{\overset{\mathrm{}}{}}}{\displaystyle \frac{dq}{2\pi \gamma (k_a)}}R(k_a,k_a,q,q)R(q,q,\kappa ,\kappa ).`$ (57)
In the asymptotic limit defined by $`\kappa r0`$ the main contribution to the integral comes from the 2 symmetric regions of the q-integration in which $`|\kappa ||q||k_a|`$. In these regions we use the form (46). We calculate
$$T_2^{(1)}(r,\kappa )2\stackrel{~}{\delta }^2\kappa ^{2/3}\underset{0}{\overset{\mathrm{}}{}}\frac{dk_a}{k_a^{5/3}}\mathrm{sin}^2(\frac{1}{2}k_ar)\underset{\kappa }{\overset{k_a}{}}\frac{dq}{q}.$$
(58)
As expected the loop integral over $`q`$ produces a logarithmic contribution. At this point we use the asymptotical identity
$$\underset{\kappa r0}{lim}_0^{\mathrm{}}d(kr)f(kr)\mathrm{ln}\left(\frac{k}{\kappa }\right)=\mathrm{ln}\left(\frac{1}{r\kappa }\right)_0^{\mathrm{}}d(kr)f(kr),$$
(59)
which produces, upon comparison with Eq. (50) the final result
$$T_2^{(1)}(r,\kappa )=\stackrel{~}{\delta }\mathrm{ln}\left[\frac{1}{\kappa r}\right]T_2^{(\mathrm{s})}(r,\kappa ).$$
(60)
The factor $`\stackrel{~}{\delta }`$ that was introduced in Eq. (51) reappears here in front of the logarithm as the effective parameter of expansion.
Analogously one computes the leading contribution of the two-loop diagram (4) in Fig. 2a. This is done explicitly in Subsect. VII A and Appendix D:
$$T_2^{(2)}(r,\kappa )=\frac{1}{2}\left[\stackrel{~}{\delta }\mathrm{ln}\left(\frac{1}{r\kappa }\right)\right]^2T_2^{(\mathrm{s})}(r,\kappa ),$$
(61)
and, in general the leading contribution of the $`n`$-loop diagram:
$$T_2^{(n)}(r,\kappa )=\frac{1}{n!}\left[\stackrel{~}{\delta }\mathrm{ln}\left(\frac{1}{r\kappa }\right)\right]^nT_2^{(\mathrm{s})}(r,\kappa ).$$
(62)
The sum of all these contributions is as follows:
$`T_2(r,\kappa )`$ $`=`$ $`T_2^{(\mathrm{s})}(r,\kappa )+{\displaystyle \underset{m=1}{\overset{\mathrm{}}{}}}T_2^{(m)}(r,\kappa )`$ (63)
$`=`$ $`T_2^{(\mathrm{s})}(r,\kappa ){\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{m!}}\left[\stackrel{~}{\delta }\mathrm{ln}\left({\displaystyle \frac{1}{r\kappa }}\right)\right]^m={\displaystyle \frac{T_2^{(\mathrm{s})}(r,\kappa )}{(r\kappa )^{\stackrel{~}{\delta }}}}.`$ (64)
We see that, as usual, resummation of the leading contributions from the logarithmic ladder diagrams results in the power function with the exponent $`\stackrel{~}{\delta }`$ which is the prefactor of the logarithm in the one-loop diagram, see Eq. (60). Because the expected correction $`\delta _2`$ to the K41 exponent $`\zeta _2^{^{\mathrm{K41}}}`$ is small ( $`\delta _20.03`$ ) we conclude that the prefactor of the rung $`\stackrel{~}{\delta }`$ is small as well! This allows us to begin with the one-loop approximation in computing the higher order scaling exponents $`\zeta _p`$ with $`p>2`$.
### B Anomalous correction of the rung asymptotics
So far we have disregarded the explicit appearance of ladder diagrams in the infinite series that defines the rung itself. As pointed in Paper II the same kind of ladder resummation that is responsible for the anomaly of the exponents of the nonlinear Green’s functions will also contribute an anomalous part to the scaling properties of the rung. Nevertheless the outer and inner scale do not appear in the rung either, and therefore the anomaly is explicit only in the asymptotic regime where we have a ratio of large and small scales. In this Subsection we flush out this anomaly.
Instead of (46) we expect
$`R_\mathrm{a}(k,k,\kappa ,\kappa ^{})={\displaystyle \frac{\delta \overline{ϵ}^{1/3}\mathrm{sign}(\kappa \kappa ^{})|\kappa \kappa ^{}|^{1/3+\delta _\mathrm{a}}}{|k|^{1+2\delta _\mathrm{a}}}}`$ (65)
$`\text{for}k\kappa ,\kappa ^{},`$ (66)
with some anomalous exponent $`\delta _\mathrm{a}`$ which is expected (and later demonstrated) to be of the order of $`\stackrel{~}{\delta }`$ as it stems from the same origin. This correction to the asymptotics may be achieved, for example by the following model form of the rung (42):
$$R_\mathrm{a}(k_a,k_b,\kappa _c,\kappa _d)=R(k_a,k_b,\kappa _c\kappa _d)\left(\frac{|\kappa _c\kappa _d|}{|k_a\kappa _c|^2}\right)^{\delta _\mathrm{a}}.$$
(67)
As before, we will argue that the exact analytic form of the rung is not important for our calculations, and only the asymptotic scaling form is essential. This statement will be shown to be exact in the 1-loop order. We thus need at this point only to preserve the essential properties, i.e. that the outer scale $`L`$ cannot appear due to locality, that the rung has to be symmetric with respect $`a,bc,d`$, etc.
### C Contributions of the skeleton diagrams with the anomalous 4-point rung
In this Subsection we reconsider the skeleton diagrams appearing in the nonlinear Green’s functions $`G_{p,p}`$ but taking into account the anomaly of the rung. In other words we are going to compute the scaling exponents accounting only for the ladder resummation inside the rung, but not the ladder resummation with the anomalous renormalized rungs. This final step will be done in Sects. VI and VII.
We are interested in the scaling exponents of structure functions in $`r`$ representation, and these are obtained from the correlation functions in $`k`$ representation as detailed in Appendix C. Upon fusion we obtain automatically contributions behaving like nonliner Green’s functions. Accordingly the objects of interest in the analysis below are the nonlinear Green’s functions in which the $`k`$ dependence of the fusing coordinates is transformed to $`r`$ representation. The outgoing wavevectors $`\kappa `$ are left as are, and the outgoing frequencis can be put to zero with impunity. This results in objects defined in a mixed $`r,\kappa `$ representation, which we denote as $`T_p(r,\{\kappa _j^{}\})`$:
$`T_p(r,\{\kappa _j^{}\})`$ $`=`$ $`{\displaystyle \underset{\mathrm{}}{\overset{\mathrm{}}{}}}{\displaystyle \underset{i=1}{\overset{p}{}}}{\displaystyle \frac{d\omega _idk_i}{(2\pi )^2}}\delta (\omega _1+\mathrm{}\omega _p)\delta (k_1+\mathrm{}k_p)`$ (68)
$`\times `$ $`f_p(r,\{k_j\})G_{p,p}(\{k_j,\omega _j,\kappa _j,0\}).`$ (69)
Here $`f_p(r,\{k_j\})`$ are one-dimensional versions of the functions $`f_p(𝒓,\{𝒌_j\})`$ defined by Eq. (C9). The set $`\{\kappa _j\}`$ denotes all the outgoing wavevectors.
We consider the skeleton contributions to the nonlinear Green’s function $`G_{p,p}`$, denoted by $`G_{p,p}^\mathrm{s}`$. Similarly to the definition (48) we introduce $`T_p^\mathrm{s}(r,\{\kappa _j\})`$ as
$`T_p^\mathrm{s}(r,\{k_j^{}\})`$ $`=`$ $`{\displaystyle \underset{\mathrm{}}{\overset{\mathrm{}}{}}}{\displaystyle \underset{i=1}{\overset{p}{}}}{\displaystyle \frac{d\omega _idk_i}{(2\pi )^2}}\delta (\omega _1+\mathrm{}\omega _p)\delta (k_1+\mathrm{}k_p)`$ (70)
$`\times `$ $`f_p(r,\{k_j\})G_{p,p}^\mathrm{s}(\{k_j,\omega _j,\kappa _j,0\}).`$ (71)
Repeating the calculation of Appendix C in the asymptotic regime $`\kappa _jr1`$ but with the redefined rung (65) one gets
$$T_p(r,\{\kappa _j\})=C_p(\overline{ϵ}r)^{p/3}r^{p\delta _\mathrm{a}}\underset{j=1}{\overset{p}{}}|\kappa _j|^{1/3+\delta _\mathrm{a}}.$$
(72)
Here $`C_p`$ are dimensionless constants that absorb all the numerical factors. In fact, the result (72) could be guessed directly by recognizing that every rung which is connected with the “outgoing” Green’s functions $`G(\kappa _j,0)`$ contributes to $`T_p(r,\{\kappa _j\})`$ a factor $`|\kappa _j|^{1/3+\delta _\mathrm{a}}`$. All together they give $`_{j=1}^p|\kappa _j|^{1/3+\delta _\mathrm{a}}`$. Convergence of the integrals over $`\kappa _j`$ and $`\omega _j`$ implies that neither inner nor outer scales may appear, and therefore dimensional consideration require a factor $`r^{p(\frac{1}{3}+\delta _\mathrm{a})}`$.
### D The 2nd and 3rd order correlation functions: relations between $`\delta _\mathrm{a}`$, $`\delta _2`$ and $`\stackrel{~}{\delta }`$
Consider first the scaling exponent $`\zeta _2`$. In Subsect. V A we showed that the resummation of the ladder diagrams leads to an anomalous correction to the exponent $`\zeta _2^{^{\mathrm{K41}}}=\frac{2}{3}`$ which is $`\stackrel{~}{\delta }`$ (cf. Eq. (63). But according to Eq. (72) the ladder in the skeleton contribution brings in an additional correction $`2\delta _\mathrm{a}`$. Altogether we have in the 1-loop order
$$\zeta _2=\frac{2}{3}\stackrel{~}{\delta }+2\delta _\mathrm{a}.$$
(73)
Therefore the exponent $`\delta _2`$ defined by (53) may be expressed as follows:
$$\delta _2=2\delta _\mathrm{a}\stackrel{~}{\delta }.$$
(74)
Another relation between the exponents will follow from the analysis of the fusion of three points. To 1-loop order the nonlinear Green’s function $`G_{3,3}`$ has the skeleton contribution diagram 3 in Fig. 4a, and the 1-loop diagrams in Fig. 4b. The skeleton contribution can be read directly from Eq. (72):
$$T_3^\mathrm{s}(r,\{\kappa _j\})=C_3\overline{ϵ}rr^{3\delta _\mathrm{a}}|\kappa _1\kappa _2\kappa _3|^{1/3+\delta _\mathrm{a}}.$$
(75)
To discuss the other contributions we refer to Fig. 4 in which all the diagram of $`G_{3,3}`$ with zero, one and two rungs are represented. We have one diagram with no rung, three with one, nine with 2. The multiplicity of 3 in the diagrams of type (2) represent the three possible connections of two struts by two rungs. The multiplicity of 6 in the diagrams of type (3) represent the different pair-permutations of three struts. In general there are $`3^n`$ diagrams with $`n`$ rungs, out of which three will have one disconnected strut. Diagrams with disconnected struts will not contribute in the asymptotic regime that interests us here. Thus out of the diagrams in Fig. 4a only the skeleton diagram (3) remains in the asymptotic regime.
In general, with $`n`$ rungs we have $`3^n3`$ fully linked diagrams. This number is $`6[(3^{n1}1)/2]`$, and the number $`[(3^{n1}1)/2]`$ counts the topologically distinct fully linked diagrams with $`n`$ rungs. Thus for example we represent in Fig. 4b the four topologically distinct contributions with three rungs. These are all the 1-loop ladder diagrams contributing to the 3’rd order correlation function. We show now that of these four terms diagram (1a) does not contribute a logarithmic divergence, whereas the other three contribute the same logarithmic term. In fact, this is the beginning of a systematic rule: the only diagrams that contribute logarithmic terms in the 1-loop order are those in which the last rung appears to the right of the skeleton diagrams. Similar rules will be established below for higher loop contributions.
Consider then the 1-loop diagram (1a) in Fig. 4b in which this rule is not obeyed. We focus on the loop made by the two rungs and the Green’s functions $`q`$ and $`q^{}`$, considering the asymptotic regime $`k_a,k_bq\kappa _a,\kappa _b`$. This is the only regime in which a logarithmic divergence is possible. In this regime $`q^{}q^{\prime \prime }k_a+k_b`$. Thus the rung $`R_a(k_a,k_b,q,q^{})`$ contributes to the loop $`q^{1/3+\delta _a}`$. The rung $`R_{q,q^{},\kappa _a,q^{\prime \prime }}`$ and the Green’s function $`G(q^{})`$ do not contribute any $`q`$ dependence to the integrand. The $`\omega `$ integration over the product of the Green’s functions $`G(q)`$ and $`G(q^{})`$ gives approximately $`1/\gamma (q^{})`$ and again contributes no $`q`$ dependence. Finally we have the evaluation
$$T_3^{(1a)}_\kappa ^{k_a}𝑑qq^{1/3+\delta _a}.$$
(76)
Clearly, this diagram does not exhibit a logarithmic divergence and as such it does not contribute to the renormalization of the scaling exponent.
The other three diagrams in Fig. 4b (namely 1b,1c and 1d) are different, they all have a logarithmic divergence. The reason for the difference is that in these three diagrams there are four Green’s functions, instead of five in diagram (1a), which carry large wavevectors. This is the same situation as in the skeleton diagram (3) in Fig. 2a. In the loop we have now two Green’s functions, instead of one in diagram (1a), that carry small wavevectors $`q`$. This difference leads to a different $`q`$ dependence in the loop, and to a logarithmic divergence. We demonstrate this explicitly in the next paragraph, but we already draw the conclusion which is general: 1-loop ladder diagrams with logarithmic divergences are those in which the additional rung (compared to skeleton diagram) has been positioned to the right of the skeleton structure.
Explicitly, consider diagram (1b) in Fig. 4b. The rung $`R_a(k_a,k_b,q,q^{})`$ contributes $`q^{1/3+\delta _a}`$ as before. But now also the rung $`R_a(q^{},k_c,q,\kappa _c)`$ contributes the same $`q`$-dependence. On the other hand the rung $`R_a(q,q,\kappa _a,\kappa _b)`$ contributes $`|q|^{12\delta _a}`$. The $`\omega `$ integration with the product of the two Green’s functions $`G(q,\omega )G(q,\omega )`$ is the same as (48) leading to $`1/2\gamma (q)`$. In total we have a logarithmic integral.
The diagram (1d) is very similar to (1b); it has the same rung structure at the left, and the rightmost rung is $`R_a(q,q,\kappa _a,\kappa _c)`$. This makes no difference to the $`q`$ dependence and thus to the logarithmic divergence or to the factor in front of the logarithm. Diagram (1c) is slightly different, having the third rung on the same ladder as the second rung. Nevertheless the rung $`R_a(q^{},k_c,q,q))`$ contributes exactly the same $`q`$ dependence as the two rungs in diagrams (1b) or (1d). Thus it yields at the end the same factor with the same logarithm. Finally, comparing 1c to diagram 3 in Fig. 2a we see that the loop structures are identical in both, and thus if diagram 3 had a prefactor $`\stackrel{~}{\delta }`$, we can immediately conclude that the three diagrams (1b, 1c and 1d) will result in a total prefactor of $`3\stackrel{~}{\delta }`$:
$$T_3^{(1)}(r,\{\kappa _j\})=3\stackrel{~}{\delta }\mathrm{ln}\left[\frac{1}{r\kappa }\right]T_3^{(s)}(r,\{\kappa _j\}),$$
(77)
where $`\kappa [\kappa _1\kappa _2\kappa _3]^{1/3}`$. The leading contribution from the higher loop diagrams can be seen to contribute higher order terms in the series of a power law, similarly to the mechanism displayed in Eqs. (6063):
$$T_3(r,\{\kappa _j\})=\frac{T_3^{(s)}(r,\{\kappa _j\})}{[r\kappa ]^{3\stackrel{~}{\delta }}}.$$
(78)
Substituting (75) we find finally
$$T_3(r,\{\kappa _j\})=C_3\overline{ϵ}[r\kappa ]^{1+3\delta _a3\stackrel{~}{\delta }}.$$
(79)
Accordingly to 1-loop order we write
$$\zeta _3=1+3\delta _a3\stackrel{~}{\delta }.$$
(80)
At this point we use the exact, nonperturbative result that $`\zeta _3=1`$ to find the relationship between $`\delta _a`$ and $`\stackrel{~}{\delta }`$: $`\delta _a=\stackrel{~}{\delta }`$. Together with (74) we get the important conclusion that all our $`\delta `$’s are the same:
$$\delta _2=\delta _a=\stackrel{~}{\delta }=\zeta _2\zeta _2^{^{\mathrm{K41}}}0.03.$$
(81)
We should stress that this important result is obtained using only the asymptotic scaling properties of the rung. Changing the explicit form of the rung without ruining the asymptotics will affect only the subleading terms in the analysis. The leading logarithmic terms are insensitive to the details of the analytic form of the rung.
## VI anomalous scaling exponents in the 1-loop approximation: Surprise, surprise
We are poised to compute now the anomalous corrections to all the scaling exponents of the $`p`$-order correlation functions in the 1-loop approximation. Start with the 4th order nonlinear Green’s function, and consider the skeleton diagrams in Fig. 5. In the one loop order, to obtain a logarithmic divergence in the asymptotic regime we must add the additional rung on the right of the skeleton structure. The combinatorics are elementary: Each skeleton diagram can host a new rung on the right in six different ways. Once a rung has been put in place the leading (logarithmic) contribution to the loop integral is the same as the loop integrals considered in the last Section. It gives the same logarithm with the same prefactor. The only difference is in the combinatorics. We can thus write by inspection
$$T_4^{(1)}(r,\{\kappa _j\})=6\stackrel{~}{\delta }\mathrm{ln}\left[\frac{1}{r\kappa }\right]T_4^{(s)}(r,\{\kappa _j\}),$$
(82)
where $`\kappa `$ is the geometric mean of all the $`\kappa _j`$. Resumming the leading contributions of the higher order loop diagrams results in the power law
$$T_4(r,\{\kappa _j\})=\frac{T_4^{(s)}(r,\{\kappa _j\})}{[r\kappa ]^{6\stackrel{~}{\delta }}}.$$
(83)
Using Eq. (72) we reach the final result
$$\zeta _4=4/3+4\delta _a6\stackrel{~}{\delta }=4/32\delta _2,\text{1-loop order.}$$
(84)
The analysis of the 1-loop order contribution to the anomalous exponents of the $`p`$-order correlation functions is as straightforward. There are $`p(p1)/2`$ possibilities to append an additional rung to the right of the skeleton structure of the $`p`$-order nonlinear Green’s function. All these diagrams contribute identical leading order logarithmic terms, with the same prefactor, summing up to an anomalous correction to the scaling exponent of the skeleton diagrams which is $`\stackrel{~}{\delta }p(p1)/2`$. According to Eq. (72) the scaling exponent of the skeleton contribution itself is corrected with respect to K41 by $`p\delta _a`$. Thus altogether
$$\zeta _p=\frac{p}{3}+p\delta _a\stackrel{~}{\delta }\frac{p(p1)}{2}.$$
(85)
Using Eq. (81)
$$\zeta _p=\frac{p}{3}\delta _2\frac{p(p3)}{2},\text{1-loop order}.$$
(86)
We note that this formula, which is valid in our case to 1-loop order only, is identical in prediction to Kolmogorov’s log-normal phenomenological model (known as K62). This is interesting, as it stems from the nontrivial topology of the ladder diagrams, in which only the most leading were considered. The present authors find the connection between lognormality and ladder diagrams unexpected.
Nevertheless we should recognize that in the present approach this result has a limited region of validity. The analysis of the 2-loop order which is provided below will show that Eq. (4) is only valid when $`p\delta _21`$. The 2-loop order will contribute positive terms of the order of $`\delta _2^2p^2(p3)`$, reducing the negative tendency of the correction to K41. Accordingly the present theory will not suffer from the well know deficiencies of the K62 log-normal model which for us is only a first order result.
## VII Anomalous Scaling Exponents in the 2-loop approximation: K62 is cured
In this Section we calculate the 2-loop contributions to the scaling exponents $`\zeta _p`$. Even though these contributions are very small when $`p\delta _2`$ is small (for, say, $`p<6`$), they become important for larger values of $`p`$ where K62 begins to turn down the $`n`$ dependence of $`\zeta _n`$. In addition this calculation allows to present clear ranges of validity for the 1-loop and 2-loop calculations.
### A 2-loop contributions to $`\zeta _2`$
We consider the 2-loop diagram (4) in Fig. 2a. Substituting it instead of $`G_{2,2}`$ in Eq. (55) we obtain the quantity $`T_2^{(2)}(r,\kappa )`$. We want to compute the correction that this diagram gives to the skeleton diagram (2), and to this aim we divide it by $`T_2^\mathrm{s}(r,\kappa )`$. In the asymptotic regime $`\kappa r1`$ the loop integrals over $`q_1`$ and $`q_2`$ contribute mostly in the range $`kq_1,q_2\kappa `$. In this regime the integrals over $`k_a,\omega _a`$ cancel from the ratio of $`T_2^{(2)}(r,\kappa )/T_2^\mathrm{s}(r,\kappa )`$. In addition the Green’s functions $`G(\kappa _c)`$ and $`G(\kappa _d)`$ cancel. Thus this ratio can be read from the ratio of the corresponding diagrams for $`G_{2,2}`$, or taking the diagram (4) and amputating the incoming and outgoing legs. We still need to divide by the rung in diagram (2) where $`k_a`$ is replaced by $`1/r`$:
$`{\displaystyle \frac{T_2^{(2)}(r,\kappa )}{T_2^\mathrm{s}(r,\kappa )}}={\displaystyle \frac{1}{R(1/r,1/r,\kappa ,\kappa )}}{\displaystyle \underset{\mathrm{}}{\overset{\mathrm{}}{}}}{\displaystyle \frac{dq_1dq_2J(q_1)J(q_2)}{(2\pi )q_1q_2}}`$ (87)
$`\times R(r^1,r^1,q_1,q_1)R(q_1,q_1,q_2,q_2)R(q_2,q_2,\kappa ,\kappa ),`$ (88)
$`J(q)={\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle \frac{d\omega }{2\pi }}G(q,\omega )G(q,\omega )={\displaystyle \frac{1}{2\gamma (q)}}.`$ (89)
In Appendix D we analyze this integral in the asymptotic limit $`\kappa r1`$ with the final result
$$T_2^{(2)}(r,\kappa )=\stackrel{~}{\delta }^2\left[\frac{1}{2}\mathrm{ln}^2\left(\frac{1}{\kappa r}\right)+b_1\mathrm{ln}\left(\frac{1}{\kappa r}\right)\right]T_2^\mathrm{s}(r,\kappa ),$$
(90)
where $`b_1`$ is a dimensionless constant
$$b_10.434.$$
(91)
In Eq. (90) the $`\mathrm{ln}^2`$ term accounts for the exponentiation of the 1-loop contribution, whereas the $`\mathrm{ln}`$ term provides the 2-loop correction to the scaling exponent $`\zeta _2`$. Instead of Eq. (74) we now read
$$\delta _2=2\delta _a\stackrel{~}{\delta }b_1\stackrel{~}{\delta }^2.$$
(92)
A second relation between these exponents will be derived in the next subsection.
### B 2-loop contributions to $`\zeta _3`$
The calculation to $`O(\delta _2^2)`$ of the contributions to $`\zeta _3`$ and of higher order $`\zeta _p`$ due to ladder resummations introduces for the first time 6-point irreducible interactions amplitudes. These appear in the ladder diagrams as rungs with six legs, arising from diagrams that due to their topology cannot be resummed into reducible contributions consisting of two 4-point rungs and one Green’s function. The 6-point rung is discussed in Appendices A and G. In particular in Appendix G we explain why the functional dependence of $`\zeta _p`$ on $`p`$ can be understood completely on the basis of the analysis of ladders with 4-point rungs. This stems from the fact that the reducible and irreducible contributions to the 6-point rung are of the same order, and their combinatorical factors are identical.
There are many possible two loop diagrams involving 4-point rungs that appear in the expansion of $`G_{3,3}`$. However, we are only interested in those contributing a logarithmic divergence in the asymptotic regime. As before, to get the relevant diagrams we need to append the last rung to the right of the existing 1-loop structure. Thus, we begin with the 3 logarithmic diagrams in Fig. 4b (i.e. 1b, 1c and 1d) and consider all the diagrams that are obtained by adding an additional rung on the right which connects two struts. The nine resulting diagrams are shown in Appendix F. These diagrams are subdivided into two groups: three diagrams in which the last rung connects the same struts as the previous rung, and six diagrams in which the last rung connects different struts. In Appendix F we explain that the first group of diagrams gives exactly the same asymptotic integral as the 2-loop contribution to $`\zeta _2`$. This statement should be reiterated because of its importance to the structure of the theory: the integrals are different, but once the limits $`kq_1,q_2\kappa `$ are taken, the resulting integrals coalesce with those computed in the previous Subsection. Thus the contribution to $`\zeta _3`$ from these three diagrams will be $`3b_1\stackrel{~}{\delta }^2`$ (which is 3 times larger than the corresponding contribution to $`\zeta _2`$).
The six diagrams of the second group look topologically different, but again in the asymptotic regime coalesce into an identical integral Eq. (E1) with $`\stackrel{~}{\mathrm{\Psi }}\stackrel{~}{\mathrm{\Psi }}_2`$, where
$`\stackrel{~}{\mathrm{\Psi }}_2(q_1,q_2)`$ (93)
$`=`$ $`{\displaystyle \frac{q_1|q_1|^{1/3}|q_1+q_2|^{4/3}\mathrm{sign}(q_2)}{(|q_1|^{2/3}+|q_2|^{2/3}+|q_1+q_2|^{2/3})(q_1^2+q_2q_1+q_2^2)}}.`$ (94)
Following the procedure outlined in Appendix E we find the coefficients of expansion
$$a_2=1,b_20.55.$$
(95)
Finally we get the 2-loop form of $`\zeta _3`$:
$$\zeta _3=1+3\delta _a3\stackrel{~}{\delta }\stackrel{~}{\delta }^2(3b_1+6b_2).$$
(96)
Demanding again $`\zeta _3=1`$ we find from Eqs. (92, 96):
$`\stackrel{~}{\delta }`$ $`=`$ $`\delta _2(b_1+4b_2)\delta _2^2+O(\delta _2^3),`$ (97)
$`\delta _a`$ $`=`$ $`\stackrel{~}{\delta }[1+(b_1+2b_2)\stackrel{~}{\delta }].`$ (98)
These results are used in the next Subsection to calculate $`\zeta _n`$ for $`n3`$ to 2-loop order.
### C 2-loop contributions to $`\zeta _p`$, $`p4`$
The calculation of the contribution of 4-point rungs to $`\zeta _p`$ for higher values of $`p`$ does not necessitate the evaluation of new integrals. In Appendix F we explain that all the 2-loop integrals appearing in the ladders of $`G_{4,4}`$ and higher order nonlinear Green’s functions are identical in the asymptotic regime to one of the two integrals appearing in the 3-order quantity. The only differences are in the combinatorial factors that account for how many ways we can choose the rungs to connect between $`p`$ struts.
If the second rung is connecting the same struts as the rung before it we have the same combinatorial factor as in the 1-loop order, namely $`p(p1)/2`$. This provides a contribution to $`\zeta _p`$ which is $`p(p1)\stackrel{~}{\delta }^2b_1/2`$. If the second rung is not connecting the same struts as the rung before it, we have $`p(p1)(p2)`$ contributions. This is due to the existence of $`p(p1)/2`$ ways connect two struts with the first rung, and then $`2(p2)`$ ways to connect one of these two struts with the remaining $`(p2)`$ struts. This leads to a contribution $`p(p1)(p2)b_2\stackrel{~}{\delta }^2`$. We should stress that the loops must have a joint strut to give a $`\mathrm{ln}`$ contribution. Two disconnected loops lead only to $`\mathrm{ln}^2`$ contributions, which do not enter the 2-loop corrections to the scaling exponents. In total we find
$`\zeta _p`$ $`=`$ $`{\displaystyle \frac{n}{3}}+p\delta _a{\displaystyle \frac{p(p1)}{2}}[\stackrel{~}{\delta }+b_1\stackrel{~}{\delta }^2]`$ (99)
$``$ $`p(p1)(p2)b_2\stackrel{~}{\delta }^2+O(\stackrel{~}{\delta }^3).`$ (100)
Substituting Eqs. (97, 98) we obtain finally
$$\zeta _p=\frac{p}{3}\frac{p(p3)}{2}\delta _2[1+2\delta _2b_2(p2)]+O(\delta _2^3)$$
(101)
We should stress that the functional form presented in this equation is solid. It is shown in Appendix G that the contribution coming from 6-point irreducible rungs is only renormalizing the value of $`b_2`$ which anyway depends on the precise analytic form of the 4-point rung which is not available at the present time. We estimate the range of validity of this order of the calculation by the Hölder inequalities, which disallow a nonlinear increase in the $`\zeta _p`$ as a function of $`p`$. The inflection point where this requirement is violated may serve as a good estimate for the range of validity. This inflection point occurs at $`p1.41/(6b_2\delta _2)12`$. In Fig. 1 we show, within this range, the K41 prediction, the 1-loop approximation (equivalent to K62) and our 2-loop final result. It is obvious that the 2-loop loop prediction goes considerably beyond the range of validity of the K62 formula which has an unphysical maximum at $`p11.5`$. We believe that all the reliably measured values of $`\zeta _p`$ agree very well with this prediction.
Using the bridge relation $`\mu =2\zeta _6`$ we predict
$$\mu =9\delta _2(1+8b_2\delta _2).$$
(102)
Plugging in the numbers we get $`\mu =0.235+O(\delta _2^3)`$. This is to be contrasted with the K62 prediction $`\mu 0.27`$. We conclude that the 2-loop contribution is very significant for experimentally measured exponents. If one wishes to obtain theoretical results for $`\zeta _p`$ with higher values of $`p`$ one needs to consider the 3-loop contributions, which pose no further conceptual difficulties. Nevertheless the experimental situation does not warrant at the present time the effort needed to accomplish such a calculation.
One should stress before closing this Section that the form of Eq. (101) is universal, stemming from the structure of the ladder diagrams and from combinatorics only. However the numerical value of $`b_2`$ is model dependent. We have checked that changing the form of the rung keeping the asymptotics unchanged results in $`b_2`$ remaining negative while its value not changing by more than a factor of 2 or so. At this moment in time one can determine $`b_2`$ using the value of $`\zeta _4`$ from experiments, allowing us then to predict accurate values of $`\zeta _n`$ for $`n`$ up to 12. It is our plan however to develop in the near future a theoretical equations for the 4-point and 6-point rungs, leading to an ab-inito determination of their analytic forms, and with them of the parameters $`\stackrel{~}{\delta }`$ and $`b_2`$.
## VIII Summary and discussion
The main steps of this and previous papers leading to the present results have been as follows:
* The theory is developed using BL-velocities to eliminate the spurious infrared divergences that are due to sweeping effects when Eulerian velocities are employed.
* The Dyson-Wyld perturbation theory was line resummed in order to achieve order by order convergent perturbation theory with K41 propagators as the lines in the theory. At this point the objects of the theory are two 2-point propagators (Green’s function and correlator) and one 3-point vertex. The 3-point vertex is in no way “small”, and renormalizing it does not change this fact .
* Multipoint correlation functions are considered when $`p`$ coordinates coalesce together. In the fusion limit $`\kappa r0`$ it is advantageous to reorganize the theory in terms of one propagator (K41 Green’s function), and 4-point, 6-point vertices etc. (the rungs). The series of diagrams contributing to the fusion limit are then simple ladder diagrams.
* The crucial step of the theory is achieved by two requirements: (i) the 4-point rung should be consistent at the level of the skeleton diagrams with the fusion rules with K41 scaling exponents. (ii) The resummation of the ladder diagrams that appear when 2 coordinates fuse together should lead to the correct value of $`\zeta _2`$. These double requirements accomplish two things in one go: (i) the theory is now developed around the K41 limit, leading to the appearance of the small parameter $`\delta _2`$ in front of the 4-point rung, and (ii) all the anomalies are coming from the ladder resummations. The 6-point rung is shown explicitly (Appendix G) to be of second order in the small parameter, 8-point rungs are of third order, etc.
* We computed the anomalous exponents in 1-loop order, inputting the value of $`\zeta _2`$ and requiring that $`\zeta _3=1`$. The result is that the scaling exponents are predicted to this order to agree with the log normal model K62. We showed that to this order the result is universal, independent of the simplifications and of the model form of the rung.
* We computed the anomalous exponents in 2-loop order. The malaise of K62 is cured, the 2-loop contribution has a sign that lifts up the exponents from the down curve of the K62 parabola. While the form of the 2-loop result is universal, the numerical value of the parameter $`b_2`$ appearing in the final result is model dependent, with contributions for the 4-point and 6-point rungs.
To improve upon the present theory one needs to develop a theory for the 4-point and 6-point interaction amplitudes. Here we determined only the asymptotic properties of the 4-point rung, and this allowed us to predict the form of the scaling exponents, but an input of the value of the anomalous part of $`\zeta _2`$ was needed to nail the 1-loop order. In fact we could use the value of $`\zeta _4`$ to fix the value of $`b_2`$ and gain a solid prediction of all the exponents to 2-loop order. Such a prediction for $`\zeta _n`$ would be valid up to $`n12`$. It is very easy to generalize the result that we have to 3-loop order, with the introduction of yet one more parameter associated with the 3-loop integrals, say $`b_3`$, which included also contributions from the irreducible 8-point rung. The result would read
$`\zeta _n`$ $`=`$ $`{\displaystyle \frac{n}{3}}{\displaystyle \frac{n(n3)}{2}}\delta _2[1+2\delta _2(n2)b_2`$ (103)
$`+`$ $`6\delta _2^2b_3(n1)(n2)]+O(\delta _2^4).`$ (104)
We stress that this form stems from the structure of the ladder diagrams, and we consider it very solid. From one point of view we can now use the value of $`\zeta _5`$ to fix $`b_3`$ to provide a prediction that is valid for any $`n`$ within experimental reach for quite some time. But this is not the main point. The main point is that we have identified the coefficients appearing in this formula with particular objects, i.e the 4-point and higher order vertices which appear in the theory as the rungs of the ladders. Obviously, a calculation of the renormalized rungs from first principle would remove the need to input experimental information altogether, affording us a complete theory of the scaling exponents of isotropic turbulence. At this point this is still not in the cards.
###### Acknowledgements.
It is a pleasure to thank Anna Pomyalov for her patient help with the diagrams in this paper. We thank her, Yoram Cohen, Ayse Erzan and Massimo Vergassola for useful comments on the manuscript. This work has been supported in part by the Israel Science Foundation, the German-Isreali Foundation, the European Commission under contract HPRN-CT-2000-00162 (“Nonideal Turbulence”), and the Naftali and Anna Backenroth-Bronicki Fund for Research in Chaos and Complexity.
## A Explanation of the diagrammatic expansion in Figs. 2, 3
It is important to stress that in the present theory we take into account all the necessary contributions. To understand this we need to say a few more words about the representation of $`G_{2,2}`$ in terms of simple ladders only, and of the $`n`$th order correlation function with $`p`$ fused coordinates as shown in Fig. 3.
The natural objects in the straightforward perturbation theory are the 2-point correlation function and Green’s function (11), and the 3-point vertex resulting from the Navier-Stokes nonlinearity $`𝒖\mathbf{}𝒖`$. After line resummation the theory contains “dressed” correlator and Green’s function. The 3-point vertex is protected by Galilean invariance and is not effected much by dressing . Thus, when we write the expansion of $`G_{2,2}`$ many diagrams involving these objects appear. The strategy that leads to the simple ladder expansion of Fig. 2 is as follows: Every diagram that contributes to the series is inspected for its cross section, or in other words what are the kind of objects that intersect a line cutting across the diagram. The line is put at the left of the diagram, and is moved to the right. Every time that the line intersects two Green’s functions that are oriented as shown in Fig. 2 we mark that position, and move the line further to the right, until we intersect again two Green’s functions, etc. For every pair of such intersections we now sum up all the topologically allowed diagrams that can be inflated ad infinitum from the fragments appearing between the two pairs of Green’s functions. This infinite resummation is the representation of the rung, which is actually a 4-point vertex. This procedure is flawless, taking into account all the possible diagram in the series of $`G_{2,2}`$ except for one subseries. This is the subseries of diagrams in which the cross section contains exactly two 2-point correlation functions. It was shown in Paper II that this can happen only once per diagram, and therefore we cannot resum such contributions to the rung, since this will lead to one rung differing from all the others. Thus the series shown in Fig. 2a contains in the $`+\mathrm{}`$ also ladders in which the struts contain two correlators above each other. With these we account for all the possible diagrams in $`G_{2,2}`$.
The presence of these diagrams also complicates the discussion of the fusion rules and Fig. 3. When we pool out the fragment containing the two coalescing coordinates we will generate in the series expansion of this fragment also the diagrams containing two correlators in the cross section. To understand their roles we will always consider the fragment of the diagrams to the right of this special cross section as belonging to the main body, see Fig. 6. But now the fragment pulled out becomes a 4th order correlation function instead of $`G_{2,2}`$. One can analyze the role of this diagram in the fusion limit, and the conclusion is that when $`rR`$ it contributes the same power $`(R/r)^{\zeta _2}`$ as the diagrams considered in the body of the paper. We thus conclude that the procedure followed in the body of the paper is amply sufficient for the calculation of $`\zeta _2`$.
The same type of considerations apply when we fuse $`p`$ coordinates. In the series expansion for $`G_{p,p}`$ we will have however ladder diagrams whose topology requires resummation into 2$`p`$-point irreducible interaction amplitudes which serve as new types of rungs. The first one appears at the level of $`G_{3,3}`$, and is the 6-point rung that is discussed explicitly in Appendix G. It will be argued that the 6-point rung is of second order in the smallness $`\delta _2`$ that characterizes the 4-point rung. Respectively the 2$`p`$-point rungs are of $`O(\delta _2^{p1})`$ and thus rungs with $`p>3`$ will not affect the analysis of this paper.
On top of the diagrams considered in the body of the paper one needs to consider those having $`p`$ correlators in the cross section. The conclusion is however the same - the diagrams considered give all the necessary information, no new (or more divergent) information is available in the diagrams that we do not consider explicitly.
At the end of the day the procedure of calculating the scaling exponents to $`O(\delta _2^2)`$ leaves us with only three dressed objects in the theory, the dressed 2-point Greens’ function and new 4-point and 6-point dressed vertices which we call the rungs. We will show that the rungs are small and therefore these three objects suffice for a consistent and controlled theory of the scaling exponents.
A reader who is an expert in diagrammatic theories may feel worried that the procedure advocated here mixes into the 4-point vertex some 2-particle reducible diagrams. This is indeed so. In other words (and see Paper III for more details) there are contributions in which the cross section includes two propagators, like one Green’s function and one correlator, or two Green’s functions oriented oppositely to the ones plotted in Figs. 2 and 3, which we resum into our 4-point vertex. Usually one would prefer not to include 2-particle reducible diagrams in the 4-point vertex, but rather to distinguish different 4-point vertices. In our case one would need three different 4-point rungs, all presented in detail in Paper III. Indeed, if we planned to compute our 4-point vertices from summing up diagrams it would be very advisable to distinguish different types of rungs. In this paper we determine however the properties of the 4-point vertex from the fusion rules, and it makes no difference at all how we classify the diagrams. For the sake of clear presentation it is much better to have just one type of rung. This rung is 2-particle reducible and hides three types of 2-particle irreducible rungs as a single object. The conclusions of the analysis are independent of this simplification.
## B Resummation into diagonal K41 propagators
The starting point of this rearrangement are the mass operators in $`k,\omega `$ representation $`\mathrm{\Sigma }_{\alpha \beta }(𝒓_0|𝒌_1,𝒌_2,\omega )`$ and $`\mathrm{\Phi }_{\alpha \beta }(𝒓_0|𝒌_1,𝒌_2,\omega )`$. Define the “diagonal” part of the mass operators as
$`\sigma _{\alpha \beta }({\displaystyle \frac{𝒌_1+𝒌_2}{2}})`$ $``$ $`{\displaystyle \frac{d(𝒌_1𝒌_2)}{(2\pi )^3}\mathrm{\Sigma }_{\alpha \beta }(𝒌_1,𝒌_2,0)},`$ (B1)
$`\varphi _{\alpha \beta }({\displaystyle \frac{𝒌_1+𝒌_2}{2}})`$ $``$ $`{\displaystyle \frac{d(𝒌_1𝒌_2)}{(2\pi )^3}\mathrm{\Phi }_{\alpha \beta }(𝒌_1,𝒌_2,0)}.`$ (B2)
In these definitions $`𝒓_0`$ disappears. The reason is that for objects which are time independent the Eulerian and BL-representations are equivalent and the designation $`𝒓_0`$ is unneeded. Here we have objects with $`\omega =0`$, or time-integrated quantities. It was shown in ref. that time integrated quantities are related to simultaneous correlations, and as such they lose the $`𝒓_0`$ designation.
Denote the rest of the mass operators as
$`\stackrel{~}{\mathrm{\Sigma }}_{\alpha \beta }(𝒓_0|𝒌_1,𝒌_2,\omega )`$ $``$ $`\mathrm{\Sigma }_{\alpha \beta }(𝒓_0|𝒌_1,𝒌_2,\omega )\sigma _{\alpha \beta }({\displaystyle \frac{𝒌_1+𝒌_2}{2}}),`$ (B3)
$`\stackrel{~}{\mathrm{\Phi }}_{\alpha \beta }(𝒓_0|𝒌_1,𝒌_2,\omega )`$ $``$ $`\mathrm{\Phi }_{\alpha \beta }(𝒓_0|𝒌_1,𝒌_2,\omega )\varphi _{\alpha \beta }({\displaystyle \frac{𝒌_1+𝒌_2}{2}}).`$ (B4)
For translationally invariant tensors in homogeneous and incompressible turbulence one can write:
$`\sigma _{\alpha \beta }(𝒌)=P_{\alpha \beta }(𝒌)\sigma (k),`$ (B5)
$`\varphi _{\alpha \beta }(𝒌)=P_{\alpha \beta }(𝒌)\varphi (k),`$ (B6)
where $`P_{\alpha \beta }(𝒌)`$ is the transverse projector,
$$P_{\alpha \beta }(𝒌)=\delta _{\alpha \beta }\frac{k_\alpha k_\beta }{k^2}.$$
(B7)
It is known that $`\sigma (k)`$ (which is the mass operator taken at $`\omega =0`$) is purely imaginary
$$\sigma (k)=i\gamma (k),$$
(B8)
with $`\gamma (k)`$ real positive. On the other hand $`\varphi (k)`$ is purely real. The diagrammatic series expansion of both $`\gamma (k)`$ and $`\varphi (k)`$ converge order by order, and using scaling relations as shown in (21) one can find their scaling behavior. The order-by-order theory dictates a K41 evaluation of these objects which is
$`\gamma (k)`$ $`=`$ $`c_\gamma [\overline{ϵ}k]^{2/3},`$ (B9)
$`\varphi (k)`$ $`=`$ $`c_\varphi \overline{ϵ}k^3,`$ (B10)
where $`c_\gamma `$ and $`c_\varphi `$ are dimensionless constants.
The Dyson-Wyld equations can be written shortly as
$`(\omega +i\nu k^2)𝑮=𝑷+𝚺𝑮,`$ (B11)
$`𝑭=𝑮(𝚽+𝑫)𝑮,`$ (B12)
where $`\nu `$ is the molecular viscosity, $`𝑷`$ is the transverse projector, and $`𝑫`$ is the correlation function of the external force which is localized in the energy containing interval. The symbol $``$ stands for summation over tensor indices and integration over intermediate $`𝒌`$. Substituting $`𝚺`$ from Eq. (B3) into the Dyson equation we rewrite:
$$[\omega +i\nu k^2+i\gamma (k)]𝑮=𝑷+\stackrel{~}{𝚺}𝑮.$$
(B13)
In the bulk of the inertial interval we can neglect $`\nu k^2`$ with impunity. The zero order solution of this equation is obtained by neglecting $`\stackrel{~}{𝚺}`$:
$`G_{\alpha \beta }g_{\alpha \beta }(𝒌,\omega )`$ $`=`$ $`P_{\alpha \beta }(𝒌)g(k,\omega ),`$ (B14)
$`g(k,\omega )`$ $`=`$ $`{\displaystyle \frac{1}{\omega +i\gamma (k)}}.`$ (B15)
The zero order solution of $`𝑭`$ is obtained in three steps: first replace $`𝚽`$ by $`\mathit{\varphi }`$, secondly neglect $`𝑫`$ in the inertial interval in comparison with $`\mathit{\varphi }`$, and lastly substitute $`𝒈`$ instead of $`𝑮`$ in Eq. (B12). The result is
$`F_{\alpha \beta }f_{\alpha \beta }(𝒌,\omega )`$ $`=`$ $`P_{\alpha \beta }(𝒌)f(k,\omega ),`$ (B16)
$`f(k,\omega )`$ $`=`$ $`{\displaystyle \frac{\varphi (k)}{\omega ^2+\gamma ^2(k)}}.`$ (B17)
Iterating Eqs. (B11, B12) without the bare forcing and viscosity results in a new diagrammatic series, which topologically is exactly the same as the old Wyld diagrammatic expansion before line re-summation. The difference is twofold. First, instead of bare propagators we have K41 propagators $`𝒈`$ and $`𝒇`$, and every 1-particle reducible fragment of any diagram will have a counter term which subtracts its “diagonal” part. This counter term is of no consequence for our procedure here since the diagrams involving it are resummed in the 4-point vertices (the rungs) together with all the other contributions as explained in Appendix A. The resulting topological structure of the ladder diagrams is thus unchanged in the new formulation.
## C Self consistency at the level of K41
Before establishing this self consistency we need to pass from correlation functions in $`𝒌,\omega `$ representation to structure functions. The theory is done naturally in $`𝒌,\omega `$ representation but the experimental scaling exponents are measured in simultaneous structure functions. We first transform from $`\omega `$-representation of $`p^{\mathrm{th}}`$-order correlation function $`𝓕_p(\{𝒌_j,\omega _j\})`$ to simultaneous correlation function $`𝑭_p(\{𝒌_j\})`$ by the integration:
$$𝑭_p(\{𝒌_j\})=\underset{\mathrm{}}{\overset{\mathrm{}}{}}\underset{i=1}{\overset{p}{}}\frac{d\omega _i}{2\pi }\delta (\omega _1+\mathrm{}\omega _p)𝓕_p(\{𝒌_j,\omega _j\}).$$
(C1)
Here $`\{𝒌_j,\omega _j\}`$ and $`\{𝒌_j\}`$ are sets of corresponding variables with $`j=1,\mathrm{}p`$. The transformation from $`𝒌`$ representation of $`𝑭_p(\{𝒌_j\})`$ to the $`p^{\mathrm{th}}`$-order structure function is done as follows: define the longitudinal component of the velocity as
$$S_p(r)=\left\{\left[𝒖(\frac{𝒓}{2})𝒖(\frac{𝒓}{2})\right]\frac{𝒓}{r}\right\}^p.$$
(C2)
Each of the factors is Fourier transformed according to
$`[𝒖({\displaystyle \frac{𝒓}{2}})𝒖({\displaystyle \frac{𝒓}{2}})]={\displaystyle \frac{d𝒌_j}{(2\pi )^3}\widehat{𝒖}(𝒌_j)}`$ (C3)
$`\times `$ $`\left[\mathrm{exp}\left(i{\displaystyle \frac{𝒌_j𝒓}{2}}\right)\mathrm{exp}\left(i{\displaystyle \frac{𝒌_j𝒓}{2}}\right)\right],`$ (C4)
Accordingly,
$`S_p(r)`$ $`=`$ $`(2\pi )^3{\displaystyle \underset{\mathrm{}}{\overset{\mathrm{}}{}}}{\displaystyle \underset{i=1}{\overset{p}{}}}{\displaystyle \frac{d𝒌_i}{(2\pi )^3}}\delta (𝒌_1+\mathrm{}𝒌_p)`$ (C6)
$`\times `$ $`f_p(𝒓,\{𝒌_j\})F_p(\{𝒌_j\}).`$ (C7)
Here
$$(2\pi )^3F_p(\{𝒌_j\})\delta (𝒌_1+\mathrm{}𝒌_p)=\underset{j=1}{\overset{p}{}}\widehat{𝒖}(𝒌_j)\frac{𝒓}{r}.$$
(C8)
The functions $`f_p(𝒓,\{𝒌_j\})`$ are seen from Eq. (LABEL:fdiff) to be:
$$f_p(𝒓,\{𝒌_j\})=\underset{j=1}{\overset{p}{}}[2i\mathrm{sin}(\frac{1}{2}𝒌_j𝒓)].$$
(C9)
In the limit $`𝒓0`$
$$f_p(𝒓,\{𝒌_j\})\underset{j=1}{\overset{p}{}}(𝒌_j𝒓).$$
(C10)
The K41 scaling exponents $`y_p`$ associated with $`p^{\mathrm{th}}`$-order correlation function $`𝓕_p(\{𝒌_j,\omega _j\})k^{y_p}`$ in $`(𝒌,\omega )`$-representation is
$$y_p=4p11/3.$$
(C11)
This corresponds to $`𝑺_p(𝒓)r^{p/3}`$ under the condition of convergence of integrals (C1, C6).
Next consider the 3rd order Green’s function, $`G_{3,3}(\{k_j,\kappa _j\})`$ in which we denoted by $`k_j`$ the set of incoming wave vectors and by $`\kappa _j`$ the set of outgoing wave vectors. The skeleton diagram of $`G_{3,3}^\mathrm{s}(\{k_j,\kappa _j\})`$ which involves 4-point rungs is shown as diagram (3) in Fig. 4a. (The contribution of 6-point rungs to the skeleton is considered in Appendix G and shown not to change the present considerations). This skeleton has two rungs, and we consider it in the limit that the incoming $`k_j`$ vectors are much larger than the outgoing $`\kappa _j`$. In this limit we have four Green’s functions with large $`k`$, contributing $`\gamma _k^4`$, and one vertex with all $`k`$ large, contributing $`k`$. The two rungs have large $`k`$ vector in them \[$`k_5`$ in Eq. (30)\], giving $`k^6`$. Finally, one of the rungs has large $`k`$ coming and going, and Eq. (30) requires for it a $`k^{2/3}`$. Altogether this gives $`G_{3,3}(\{k_j,\kappa _j\})k^{x_3}`$ with $`x_3=25/3`$ which is equal to $`y_3`$ given by Eq. (C11). This means that the skeleton diagrams for $`G_{3,3}(\{k_j,\kappa _j\})`$ (with asymptotics of the rung defined by the two-point fusion rules) automatically reproduces the K41 scaling exponent $`\zeta _3=1`$ in the three-point fusion. This is true subject to the condition that the integrals (C1, C6) for $`p=3`$ converge. That this is so may be shown by a direct calculation. For future purposes it is extremely important to note that the principal contribution to the $`𝒌`$ integral (C6) comes from the region where $`k_1k_2k_31/r`$.
Now let us compare diagram (3) in Fig. 4a and Fig. 5 with the skeleton diagrams for $`G_{3,3}`$ and $`G_{4,4}`$. One recognizes that in general for $`G_{p,p}`$, we will have $`(p1)`$ rungs with large incoming $`k`$, contributing $`k^{3(p1)}`$ \[originating from $`k_e`$ in Eq. (30)\]. We will have also $`2p2`$ Green’s functions with large $`k`$ contributing $`k^{(2p2)2/3}`$. Next we will have $`p2`$ outgoing legs with large $`k`$ contributing $`k^{(p2)2/3}`$ from Eq. (9). Finally we will have $`2p2`$ vertices having incoming and outgoing large $`k`$ vectors, contributing $`k^{p2}`$. All together we find that $`G_{p,p}(\{k_j,\kappa _j\})k^{x_p}`$ with $`x_p=4p11/3`$ which is equal to $`y_p`$ given by Eq. (C11). Convergence of the $`𝒌`$ integral (C6) for $`p=4`$ may be shown by direct calculations. A proof of convergence of the $`𝒌`$ integrals (C6) for $`p>4`$ is a tedious exercise which nevertheless may be done, for example, iteratively. It is readily demonstrated that the integral converges when all $`k_j`$-vectors are of the same order of magnitude (say, $`k`$). Then $`G_{p,p}k^{11/34p}`$. After $`(p1)`$ $`\omega `$-integrations (each of them giving a factor $`k^{2/3}`$) one has $`k^{3(p1)p/3}`$ which is enough for convergence of $`(p1)`$ $`d^3k`$ integrals in the UV region $`k_jk1/r`$. In the IR region $`k_jk1/r`$ the functions $`f_p`$ provide the integral with additional $`k^p`$ factor \[according to (C10)\] which guarantees the convergence.
The considerations of the 6-point and higher order rungs leave these conclusions invariant.
## D Analysis of two loop integrals contributing to $`\zeta _2`$
The integrand in the integral (89) is a function of $`q_1`$ and $`q_2`$ and it depends on $`k`$ and $`\kappa `$ as parameters. The integration range is the $`q_1q_2`$ infinite plane, but in the limit $`k\kappa `$ the main contribution comes form the four finite quadrants $`\kappa <|q_1|,|q_2|<k`$. Well inside the quadrants we are allowed to use the asymptotic form in which $`\kappa |q_1|,|q_2|k`$. In this regime the integrand is $`k,\kappa `$-independent, and the dependence of the integrals on $`k,\kappa `$ appears only via the limits of integration. By changing the dummy variables $`q_1`$ and $`q_2`$ we can now project all four quadrants into one of them, say $`q_1`$ and $`q_2`$ positive. In this asymptotic regime we can use for the rungs in the integrand of (89) that include either $`k`$ or $`\kappa `$ their asymptotic form (46). This results in
$`K(k,\kappa )`$ $`=`$ $`\stackrel{~}{\delta }^2{\displaystyle _p^k}{\displaystyle \frac{dq_1}{q_1}}{\displaystyle _p^k}{\displaystyle \frac{dq_2}{q_2}}\mathrm{\Psi }(q_1,q_2),`$ (D1)
$`\mathrm{\Psi }(q_1,q_2)`$ $`=`$ $`\stackrel{~}{\mathrm{\Psi }}(q_1,q_2)\stackrel{~}{\mathrm{\Psi }}(q_1,q_2).`$ (D2)
In Appendix E we show how to analyze this kind of integral with the aim of extracting the coefficients of the leading and first subleading logarithmic terms, i.e.
$$K_1(k,\kappa )=\frac{a_1}{2}\mathrm{ln}^2(k/\kappa )+b_1\mathrm{ln}(k/\kappa )$$
(D3)
Using the results there (E5) with $`\stackrel{~}{\mathrm{\Psi }}(q_1,q_2)=\stackrel{~}{\mathrm{\Psi }}_1(q_1,q_2)`$,
$$\stackrel{~}{\mathrm{\Psi }}_1(q_1,q_2)=\frac{q_1^3|q_1q_2|\mathrm{sign}(q_2)}{2(q_1^2q_1q_2+q_2^2)^2},$$
(D4)
one find immediately $`aa_1=1`$ as required by the anticipated expansion employed in Eqs. (61)-(63). To compute $`b_1`$ we examine the integral $`b_1(A)`$ numerically, see Fig. 7. We see that the requested limit exists and that $`b_10.434`$.
## E Extraction of the subleading logarithmic term from the 2-loop integrals
The 2-loop integrals have the characteristic structure appearing in (D1)
$$I(A)=\stackrel{~}{\delta }^2_1^A\frac{dq_1}{q_1}_1^A\frac{dq_2}{q_2}\mathrm{\Psi }(q_1,q_2),$$
(E1)
where $`A1`$ and $`\mathrm{\Psi }(q_1,q_2)`$ is homogeneous function of degree zero: $`\mathrm{\Psi }(\lambda q_1,\lambda q_2)=\mathrm{\Psi }(q_1,q_2)`$. When $`\mathrm{\Psi }(q_1,q_2)=1`$ then $`I(A)=\mathrm{ln}^2A`$. In general only the leading term of $`I(A)`$ is proportional to $`\mathrm{ln}^2A`$ and we expect the following subleading terms:
$$I(A)=\stackrel{~}{\delta }^2[\frac{a}{2}\mathrm{ln}^2A+b\mathrm{ln}A+c+\frac{d}{A}+\mathrm{}]$$
(E2)
Our goal is to find the coefficient $`b`$ in the limit $`A\mathrm{}`$. Taking the first derivative of (E2) with respect to $`A`$ and multiplying by $`A`$ we find
$`a\mathrm{ln}A+b{\displaystyle \frac{d}{A}}\mathrm{}={\displaystyle _1^A}{\displaystyle \frac{dq_2}{q_2}}\mathrm{\Psi }(1,q_2)+{\displaystyle _1^A}{\displaystyle \frac{dq_1}{q_1}}\mathrm{\Psi }(q_1,1)`$ (E3)
$`={\displaystyle _{1/A}^1}{\displaystyle \frac{dx}{x}}\mathrm{\Psi }(x,1)+{\displaystyle _{1/A}^1}{\displaystyle \frac{dy}{y}}\mathrm{\Psi }(1,y),`$ (E4)
where we changed the dummy variables $`q_1=xA`$ and $`q_2=yA`$. Taking another derivative and multiplying by $`A`$ we find for large $`A`$
$$a=\mathrm{\Psi }(1,\frac{1}{A})+\mathrm{\Psi }(\frac{1}{A},1).$$
(E5)
Substituting this result in (E4), and representing $`\mathrm{ln}A`$ as $`_{1/A}^1𝑑x/x`$ we find
$`b`$ $`=`$ $`\underset{A\mathrm{}}{lim}b(A),`$ (E6)
$`b(A)`$ $`=`$ $`{\displaystyle _{1/A}^1}{\displaystyle \frac{dx}{x}}\left[\mathrm{\Psi }(x,1)+\mathrm{\Psi }(1,x)\mathrm{\Psi }({\displaystyle \frac{1}{A}},1)\mathrm{\Psi }(1,{\displaystyle \frac{1}{A}})\right].`$ (E7)
If the expansion assumed in Eq. (E2) is valid, this limit must exist.
## F The nine 2-loop diagrams of $`G_{3,3}`$
Consider diagram (1a) in Fig. 8. We are interested in the ratio of $`T_{3,1a}^{(2)}/T_3^\mathrm{s}`$, where $`T_{3,1a}^{(2)}`$ is obtained by substituting the diagram (1a) instead of $`G_{3,3}`$ in Eq. (69). In the asymptotic regime $`\kappa r1`$ the loop integrals over $`q_1`$ and $`q_2`$ contribute mostly in the regime $`kq_1,q_2\kappa `$. In this regime the integrals over $`k_a,\omega _a,k_b,\omega _b`$ cancel in the desired ratio. Similarly the Green’s functions $`G(\kappa _d)`$, $`G(\kappa _e)`$ and $`G(\kappa _f)`$ also cancel in the ratio. Accordingly
$`T_{3,1a}^{(2)}/T_3^\mathrm{s}`$ can be calculated from the amputated diagram (2) in Fig. 10, in which the explicit dependence on $`k_j`$ and $`\kappa _j`$ has disappeared. These wavevectors remain only in the limits of the integrals over $`q_1`$ and $`q_2`$, with $`k`$ replaced by $`1/r`$. In this diagram every black dot contributes a factor of $`q_j^{1/3+\delta _a}`$ where $`q_j`$ is the wavevector on the right of the black dot. This is a remnant of the corresponding rung before the amputation. The thin line connecting these dots is just a reminder that we have loop integrals to perform.
The point to understand now is that if we use diagrams (1b) and (1c) in Fig. 8 to form $`T_{3,1b}^{(2)}`$ and $`T_{3,1c}^{(2)}`$, the ratio of these to $`T_3^\mathrm{s}`$ can be again calculated from the amputation of their own diagrams. This will lead to the identical amputated diagram (2) of Fig. 10. In addition, and most importantly, the integral that needs to be computed is the same as Eq. (89). Thus one recaptures Eq. (90) but with the combinatorial factor 3 in front of the RHS:
$`T_{3,1a+1b+1c}^{(2)}(r,\kappa )`$ $`=`$ $`3\stackrel{~}{\delta }^2[{\displaystyle \frac{1}{2}}\mathrm{ln}^2\left({\displaystyle \frac{1}{\kappa r}}\right)`$ (F1)
$`+`$ $`b_1\mathrm{ln}\left({\displaystyle \frac{1}{\kappa r}}\right)]T_3^\mathrm{s}(r,\kappa ),`$ (F2)
with $`b_1`$ of Eq. (91).
The second group of six diagrams (2a – 3c) shown in Fig. 9 yields to a similar analysis, but the amputated diagram is shown as diagram (3) in Fig. 10. All six diagrams result in the very same amputation, up to permutations of the three struts. Analyzing the amputated diagram (3) one brings it to the canonical form (D1) with $`\stackrel{~}{\mathrm{\Psi }}(q_1,q2)`$ given by (93). Accordingly we write
$$T_{3,2a,\mathrm{}3c}^{(2)}(r,\kappa )=6\stackrel{~}{\delta }^2\left[\frac{1}{2}\mathrm{ln}^2\left(\frac{1}{\kappa r}\right)+b_2\mathrm{ln}\left(\frac{1}{\kappa r}\right)\right]T_3^\mathrm{s}(r,\kappa ),$$
(F3)
with $`b_2`$ of (95).
The analysis of the 2-loop diagrams that involve 4-point rungs in the context of $`G_{p,p}`$ follows exactly the same lines, with the amputated diagrams being those of Fig. 10. The only thing to mind is the combinatorics, which are presented explicitly in Fig. 10, leading to the numbers in Eq. (100).
## G Resummed equations for the 4-point and 6-point rungs
In this Appendix we sketch a theory for the 4-point and 6-point rungs. Our main aim here is to explain why the 6-point rung is quadratic in the smallness, but we use the opportunity to indicate how a future theory of these objects may be formulated.
Consider the beginning of the series expansion of the 4-point rung which is shown in Fig. 2b. Diagram (2) contains a cross of correlators each attached to two 3-point vertices. This is exactly diagram (1), and therefore the equation can lend itself to resummation resulting in the equation shown in Fig. 11a. We note that this is not the full equation for the 4-point rung even in 1-loop order
since we did not take into account the ladders with a correlator and Green’s function in a cross section. Taking into account all the needed contributions is not difficult, but is not the main point of this Appendix, and we proceed for simplicity without the additional terms.
In the asymptotic regime the bare contribution diagram (1) in Fig. 11a is negligible. With this contribution discarded, the remaining equation is homogeneous, calling for finding a zero mode of the equation. Since we have already demonstrated that the 4-point rung is small, of the order of $`\delta _2`$ we can conclude that the loop integral which we denote as $`\mathrm{}_1`$ must be large, or the order of $`1/\delta _2`$. (the homogeneous equation can be only solved if $`\delta _2\delta _2^2\mathrm{}_1`$). In fact, in the future it would be extremely worthwhile to solve the full equation in the 1-loop order and demonstrate that this is the case, and thus to lend further weight to the theory presented in this paper. Of course solving such an equation will also supply us with a functional form of the 4-point rung, and with it a substantial part of the value of the parameter $`b_2`$ which appears in the final result for the scaling exponents.
In Fig. 11b we present the resummed form of the equation of the 6-point rung, to the same level of qualitative discussion. Again we discard in the asymptotic limit the bare contribution of diagram (1), but we cannot neglect diagram (2) since it has the same asymptotic behavior as the resummed 6-point rung. Diagram (2) is of the order of $`\delta _2^2`$. Diagram (3) is of the order of $`ϵ\delta _2\mathrm{}_2`$ where $`\mathrm{}_2`$ is the loop integral. This integral is very similar to $`\mathrm{}_1`$, and we therefore estimate $`\mathrm{}_2\mathrm{}_11/\delta _2`$, and thus diagram (3) is of the order of the LHS. Diagram (4) is of the order of $`ϵ^2\mathrm{}_3\mathrm{}_4`$ where $`\mathrm{}_3`$ and $`\mathrm{}_4`$ each refers to one of the loop integrals. With the same level of approximation we estimate it thus to be of $`O(ϵ^2/\delta _2^2)`$. Denoting $`xϵ/\delta ^2`$ we thus represent the order of magnitude relations that result from panel b by the equation
$$x=1+ax+bx^2,$$
(G1)
where $`a`$ and $`b`$ are dimensionless constants of $`O(1)`$. It is obvious that only $`x1`$ is a consistent solution of this equation, and we thus conclude that the 6-point rung is quadratic in the smallness $`\delta _2`$.
We therefore understand that the 6-point rung appears in our considerations only at the level of the $`O(\delta ^2)`$ order. In this order it appears in addition to the 2-loops integrals which are formed by two 4-point rungs, as discussed in detail in the text of the paper. But since the 6-point rung connects three struts, exactly like the structure made of two 4-point rungs, the combinatorical factors appearing in the $`p`$th order scaling exponents are identically the same. Accordingly we understand that the effect of the 6-point rung is only in renormalizing the value of the parameter $`b_2`$ which anyway is model dependent.
Similar consideration apply to the 8-point rung which begins to affect the theory only in $`O(\delta _2^3)`$. It will renormalize the value of the parameter $`b_3`$ in Eq. (104). Higher order rungs are even less relevant for the calculation at hand.
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# Realizations of quantum hom-spaces, invariant theory and quantum determinantal ideals
## Introduction
Let $`V`$ be a vector space of finite dimension over a field of characteristic zero. There are two ways of interpreting the symmetric tensor algebra $`𝖲(V)`$ over $`V`$. The first one is to consider it as a factor algebra of the tensor algebra $`𝖳(V)`$ over $`V`$. Thus we have a relation between two tensors, like $`ab=ba`$. The second way is to consider $`𝖲(V)`$ as a subspace of $`𝖳(V)`$ consisting of symmetrized tensors, e.g., for $`a,bV`$, the element $`ab+ba`$ belongs to the symmetric tensor algebra on $`V`$. We are thus speaking of two realizations of the symmetric tensor algebras over $`V`$. The analogous procedure applies also for the exterior (anti-symmetric) tensor algebra.
A Hecke operator $`R`$ on a vector space $`V`$ of finite dimension is an invertible operator on $`VV`$ that satisfies the Yang-Baxter equation and the Hecke equation $`(x+1)(xq)=0`$. To a Hecke operator there is associated a quantum space, given in terms of a pair of quadratic algebras (Section 1). The latter algebras are quantum analogues (or deformations) of the symmetric and anti-symmetric tensor algebras over a vector space.
For a quantum space associated to a Hecke operator, the two realizations for its (quantum) symmetric and anti-symmetric tensor algebras was first obtained by Gurevich . While the first realization was taken as the definition, the second realizations followed from the general theory of Hecke algebras. In this paper we give the second realization for the (anti-) symmetric tensor algebras on the quantum semi-group of endomorphism associated to a Hecke operator and the quantum hom-space associated to a pair of Hecke operators.
The main idea in giving the second realization is to construct a projector on the tensor power of $`V`$. For the classical case, it is the (anti-) symmetrizer operator, e.g., $`ab(ab+ba)/2`$. For the quantum space, it is the quantum (anti-) symmetrizer, constructed in terms of the trivial and signature representations of the Hecke algebras. In Section 1, we recall the definition of the quantum exterior and the quantum symmetric algebras associated to a Hecke operator and their second realization, due to Gurevich (Equation (8)). Together, these algebras determine a quantum space. Then, we recall the definition of the matrix bialgebra associated to the Hecke operator, which is considered as the function algebra of the quantum semi-group of endomorphism of the quantum space. Unlike the classical case, where a matrix can also be considered as a vector, the matrix bialgebra generally cannot be defined as a quantum symmetric algebra associated to a Hecke operator. In fact, it can still be defined analogously in terms of a Yang-Baxter operator, but this operator has the minimal polynomial of degree 3. As a result, only an analogue of the quantum anti-symmetrizer was defined (the operator $`\mathrm{\Phi }^n`$ in Equation (12)), which is no more a projector. There is not strightforward analogue of the quantum symmetrizer.
There is a simple solution by the following remark. Since the Yang-Baxter operator defining the matrix bialgebra is a tensor product of the ordinary Hecke operator with the inverse of its dual, our operator $`\mathrm{\Phi }^n`$ is a homomorphic image of a Casimir element in $`_n_n`$, where $`_n`$ is the Hecke algebra of type $`A_n`$. Using various dual bases in $`_n`$, we found eigenvalues of $`\mathrm{\Phi }^n`$. So that we can modify it to obtain a projector $`\overline{\mathrm{\Phi }}^n`$ (Subsection 2.1). Now, the choice of a quantum symmetrizer becomes clear. The remark above also suggest us define an operator $`\mathrm{\Psi }^n`$, which play the role of the quantum symmetrizer. We find its eigenvalues and modify it to get a projector $`\overline{\mathrm{\Psi }}^n`$. It is the operators $`\overline{\mathrm{\Psi }}^n`$, $`n=1,2,\mathrm{}`$, which give the second realization for the matrix bialgebra.
The results can be generalized for the function algebra on the space of homomorphisms of two quantum spaces or quantum hom-space. This quadratic algebra is introduced in Section 3, it is defined in terms of a pair of Hecke operators. Thus, on it coact the matrix bialgebras associated to these Hecke operators. Our second realization for quantum spaces of homomorphisms implies a quantum analogue of Cauchy’s decomposition (21). Although in the classical case, the operator $`\overline{\mathrm{\Psi }}^n`$ reduces to the ordinary symmetrizer operator, its relationship with Cauchy’s decomposition is new. Moreover, the second realization also implies interesting results in invariant theory.
Before describing the invariant theory for quantum groups of type $`A`$, let me briefly recall the classical theory. Let $`M(m,n)`$ denote the space of $`m\times n`$-matrices. Let $`\mu `$ denote the matrix multiplication map $`\mu :M(m,t)\times M(t,n)M(m,n)`$, $`\mu (A,B)=AB`$, $`AM(m,t),BM(t,n)`$. On the variety $`M(m,t)\times M(t,n)`$ acts the general linear group $`GL(t)`$, $`g(A,B)=(Ag^1,gB)`$. It this easy to see that elements of an orbit of $`GL(t)`$ have the same image under $`\mu `$. The above action of $`GL(t)`$ induces an action on the polynomial ring on $`M(m,t)\times M(t,n)`$, $`𝒪(M(m,t)\times M(t,n))`$. The classical invariant theory for general linear groups studies the subring of invariant polynomials in $`𝒪(M(m,t)\times M(t,n))`$. Let $`m_j^i`$ the $`(i,j)`$ coordinate function on $`𝒪(M(m,n))`$ and $`a_j^i`$ be the composition of $`\mu `$ with $`m_j^i`$, which are then polynomial functions on $`𝒪(M(m,t)\times M(t,n))`$. The first fundamental theorem of invariant theory states that any invariant polynomial on $`𝒪(M(m,t)\times M(t,n))`$ can be represented as a polynomial functions on the functions $`a_j^i`$. The second fundamental theorem states that the relations between the functions $`a_j^i`$ are exactly the minors of ranks $`t+1`$ in the matrix $`(a_j^i)`$.
Using the associated homomorphism of algebras $`\mu ^{}:𝒪(M(m,n))𝒪(M(m,t)\times M(t,n))`$, we can reformulate the above theorems as follows:
1. A polynomial in $`𝒪(M(m,t)\times M(t,n))`$, invariant under the action of $`GL(t)`$, is contained in the image of $`\mu ^{}`$.
2. The kernel of $`\mu ^{}`$ is the ideal in $`𝒪(M(m,n))`$, generated by minors of degree $`t+1`$ in the matrix $`(e_j^i)`$.
The characteristic free proof of these theorems, due to DeConcini-Procesi , uses the notion of standard basis and has close relationship with combinatorics. A generalization of these results for standard quantum general linear groups was obtained by Goodearl et.al . Their proof closely follows DeConcini-Procesi’s proof.
In the quantum setting, the variety $`M(m,n)`$ is replaced by a quantum hom-space. Given two Hecke operators $`R,S`$, the quantum hom-space associated to $`R,S`$ is denoted by $`𝖬_{SR}`$. On the algebra $`𝖬_{SR}`$ coact the bialgebras $`𝖤_R`$ on the right and $`𝖤_S`$ on the left. In Section 4, we study two sided ideals in $`𝖬_{SR}`$ which are invariant with respect to these actions. In the classical case, these ideals were studied by DeConcini-Eisenbud-Procesi . We show that there is a one-one correspondence between invariant ideals and diagram ideals in the sense of (note that our notation here slightly differes from the notion in , a partition is replaced by its conjugate partition). In the classical case, those invariant ideals that correspond to diagram ideals of the form $`(1^k)=\{\lambda |\lambda _1k\}`$ are determinantal ideals, i.e., generated by minors of degree $`k`$. In our general case, the notions of quantum determinant and quantum minor is not defined. We show, however, that in the case of standard quantum general linear groups, where these notions are defined, the quantum determinantal ideals, introduced by Goodearl et.al., are precisely those corresponding to $`(1^k)`$, for some $`k`$.
The setting for invariant theory of quantum groups of type $`A`$ involves three Hecke operators $`R,S`$ and $`T`$, of which $`R`$ is also a Hecke symmetries. The morphism $`\mu ^{}`$ mentioned above becomes an algebra morphism
$$\mu ^{}:𝖬_{TS}𝖬_{TR}𝖬_{RS}.$$
The corresponding quantum group is the Hopf algebra associated to $`R`$, $`𝖧_R`$. This Hopf algebra coacts on the source and the target of $`\mu ^{}`$ and the formulation of the first and the second fundamental theorems can be made analogously as in the classical case. Since $`𝖧_R`$ is not commutative, there are more than one coactions of it on $`𝖬_{RS}`$, which yield different versions of the fundamental theorems.
The method of our proof is new. It relies mainly on the second realization of function algebras on quantum spaces homomorphisms. In the case of standard quantum general linear groups, our result is precisely those obtained in . On the other hand, our assumption about the quantum groups also covers the case of standard quantum general linear supergroups. Thus, we have particularly proved the fundamental theorems for quantum general linear supergroups, which have as a special case the fundamental theorems for general linear supergroups. Unlike the case of (quantum) general linear groups, in the super case, the kernel of $`\mu ^{}`$ is not generated by (quantum) minors. It is an interesting problem to study such ideals.
## 0. Preliminaries
Throughout this paper, we work over an algebraically closed field 𝕂 of characteristic zero.
### 0.1. Partitions.
A partitions $`\lambda `$ of $`n\text{}`$ is a sequence $`\lambda =(\lambda _1,\lambda _2,\mathrm{})`$ of non-increasing non-negative integers, whose sum is $`n`$, we write $`\lambda n`$ or $`|\lambda |=n`$. The maximal number $`r`$, for which $`\lambda _r0`$, is called the length of $`\lambda `$. The diagram $`[\lambda ]`$ associated to $`\lambda `$ is a matrix, whose first row contains $`\lambda _1`$ elements, called nodes, second row contains $`\lambda _2`$ elements, and so on. The conjugated to $`\lambda `$, denoted by $`\lambda ^{}`$, is the one, whose diagram $`[\lambda ^{}]`$ is obtained from $`[\lambda ]`$ by rotating it $`180^{}`$ along its diagonal.
A standard $`\lambda `$-tableau is the diagram $`[\lambda ]`$, filled by numbers $`1,2,\mathrm{},|\lambda |`$, in such a way that they increase along rows and columns. The number of standard tableaux is denoted by $`d_\lambda `$.
### 0.2. Symmetric Groups.
The symmetric group $`\text{S}_n`$ consists of permutations of the set $`1,2,\mathrm{},n`$. It can be regarded as the group generated by transposition $`v_i=(i,i+1)`$, $`1in1`$, subject to the relations: $`v_i^2=1,v_iv_{i+1}v_i=v_{i+1}v_iv_{i+1}`$ and $`v_iv_j=v_jv_i`$ if $`|ij|2`$. The length $`l(w)`$ of an element $`w`$ is the minimal length of the words in $`v_i`$ expressing $`w`$. It equals the number of pairs $`1i<jn`$ for which $`iw>jw`$.
### 0.3. Hecke Algebras .
The Hecke algebra $`_n=_{n,q}`$ is a $`q`$-analogue of the group algebra $`\text{𝕂}[\text{S}_n]`$. It is generated over 𝕂 by 1 and the elements $`T_i,1in1`$, subject to the relations $`T_i^2=(q1)T_i+q`$, $`T_iT_{i+1}T_i=T_{i+1}T_iT_{i+1}`$ and $`T_iT_j=T_jT_i`$ if $`|ij|2`$. When $`q=1`$, $`_{n,1}`$ reduces to $`\text{𝕂}[\text{S}_n]`$, $`T_iv_i`$. $`_n`$ has a basis consisting of $`T_w,w\text{S}_n`$, $`T_1:=1`$, $`T_{v_i}:=T_i`$, $`T_wT_u=T_{wu}`$ if $`l(w)+l(u)=l(wu)`$.
We shall always assume that $`q^n1,n>1`$. In this case $`_{n,q}`$ is semisimple.
The embedding $`_l_m_{l+m}`$, mapping $`_lT_i`$ to $`T_i_{l+m}`$ and $`_mT_j`$ to $`I_{m+j}_{l+m}`$, is called the standard embedding.
### 0.4. Representations of the Hecke Algebras .
Representations of $`_{n,q}`$, for $`q`$ not being root of unity, can be parameterized by partitions of $`n`$. Let $`S_\lambda `$ be the simple representation of $`_n`$, corresponding to $`\lambda `$, then the dimension of $`S_\lambda `$ is $`d_\lambda `$ ($`d_\lambda `$ is defined in 0.1). Let $`𝒜_\lambda `$ be the block, i.e. a minimal two sided ideal, in $`_n`$, that corresponds to $`\lambda n`$. Let $`E_\lambda ^{ij},1i,jd_\lambda `$ be a basis of $`𝒜_\lambda `$ such that $`E_\lambda ^{ij}E_\lambda ^{kl}=\delta _k^jE_\lambda ^{il}.`$ Thus, $`E_\lambda ^{ii}`$ are mutually orthogonal primitive idempontents of $`_n`$. Notice that for $`\lambda \mu `$, $`E_\lambda ^{ij}E_\mu ^{kl}=0`$.
We shall also consider the algebra $`_n^{\mathrm{op}}`$. Its simple comodules are canonically identified with $`S_{\lambda }^{}{}_{}{}^{}`$, the dual vector space to $`S_\lambda `$: if $`\varphi _\lambda `$ is the representation of $`_n`$ on $`S_\lambda `$, then the representation of $`_n^{\mathrm{op}}`$ on $`S_\lambda `$ is given by $`\overline{\varphi }_\lambda (w):=\varphi _\lambda (w)^{}.`$
### 0.5. A bilinear form on $`_n`$ .
There exists a non-degenerate, symmetric, associative bilinear form on $`_n`$, defined as follows: $`(T_u,T_w):=q^{l(u)}\delta _w^{u^1}.`$ Thus, we see that $`\{q^{l(w)}T_{w^1},w\text{S}_n\}`$ is the dual basis to $`\{T_w,w\text{S}_n\}`$ with respect to this bilinear form. Therefore, the Casimir element $`_{w\text{S}_n}q^{l(w)}T_wT_{w^1}`$ is central:
$$\underset{w\text{S}_n}{}q^{l(w)}T_iT_wT_{w^1}=\underset{w\text{S}_n}{}q^{l(w)}T_wT_{w^1}T_i.$$
Since $`𝒜_\lambda `$ is simple, the bilinear form restricted on $`𝒜_\lambda `$ should satisfy $`(E_\lambda ^{ij},E_\lambda ^{kl})=\delta _k^j\delta _l^ik_\lambda `$, for certain coefficient $`k_\lambda 0`$. Thus $`\{E_\lambda ^{ij},\lambda n,1i,jd_\lambda \}`$ and $`\{k_\lambda ^1E_\lambda ^{ji},\lambda n,1i,jd_\lambda \}`$ are dual bases. Therefore
(1) $`{\displaystyle \underset{w\text{S}_n}{}}q^{l(w)}T_wT_{w^1}={\displaystyle \underset{\genfrac{}{}{0pt}{}{\lambda n}{1i,jd_\lambda }}{}}k_\lambda ^1E_\lambda ^{ij}E_\lambda ^{ji}.`$
Denote $`\mathrm{\Pi }_\lambda :=_{1i,jd_\lambda }k_\lambda ^1E_\lambda ^{ij}E_\lambda ^{ji}.`$ Then
$$\mathrm{\Pi }_\lambda ^2=d_\lambda k_\lambda ^1\mathrm{\Pi }_\lambda .$$
The coefficients $`k_\lambda `$ can be explicitly given (cf. )
$$k_\lambda =q^{_i\lambda _i(i1)}\underset{k=1}{\overset{n}{}}\frac{1}{[r_\lambda (k)+r]_q}\frac{[\lambda _i\lambda _j+ji]_q}{[ji]_q}.$$
### 0.6. Hecke Operators.
Let $`V`$ be a finite dimensional vector space over a field 𝕂 and $`R`$ be a invertible operator on $`VV`$. $`R`$ is called Hecke operator if the following conditions are satisfied:
$`R_1R_2R_1=R_2R_1R_2\text{where }R_1:=R\text{id}_V,R_2:=\text{id}_VR,`$
$`(R+1)(Rq)=0.`$
A Hecke operator induces a right representation $`\rho `$ of $`_{n,q}`$ on $`V^n`$, for $`n>1`$: $`\rho (T_i)=R_i:=\text{id}_V^{i1}R\text{id}_V^{ni1}.`$ For any $`w\text{S}_n`$, we denote $`R_w:=\rho (T_w)`$.
### 0.7. Hopf algebras.
We assume that the reader is familiar with the notions of bialgebras and Hopf algebras and their (co)modules. The reader may consult or \[19, Chapter1\] for basic notions of bialgebras and Hopf algebras. For a coalgebra $`C`$, $`C^{\mathrm{cop}}`$ denotes $`C`$ with the opposite coproduct. Similary, for a bialgebra $`B`$, $`B^{\mathrm{cop}}`$ denotes the $`B`$ with the opposite product and coproduct, it is a bialgebra too. If $`M`$ is a right $`C`$ comodule then $`M^{}`$ is a left $`C`$-comodule in a canonical way, namely, if $`\delta `$ denotes the coaction of $`C`$ on $`M`$, then the left coaction $`\lambda `$ of $`C`$ on $`M^{}`$ is given by $`\lambda (\varphi )(m)=\varphi (\delta (m))`$ for all $`mM,\varphi M^{}`$. Hence, $`M^{}`$ is a right $`C^{\mathrm{cop}}`$-comodule.
### 0.8. The dual space to a tensor product.
To a vector space $`V`$, the dual vector space is defined to be $`V^{}:=\text{Hom }(V,\text{𝕂})`$. If $`V`$ is finite dimensional over 𝕂, then so is $`V^{}`$ and they have the same dimension. For finite dimension vector spaces, there is the following equivalent definition of dual vector spaces which is more suitable for further generalization (cf. ). The dual vector space to a vector space $`V`$ (of finite dimension) is a pair $`(V^{},\text{ev}_V:V^{}V\text{𝕂})`$, such that there exists a linear map $`\text{db}_V:\text{𝕂}VV^{}`$, satisfying the following conditions ($`\text{id}_V`$ is the identity map on $`V`$):
$$(\text{ev}\text{id}_M^{})(\text{id}_M^{}\text{db})=\text{id}_M^{},(\text{id}_M\text{ev})(\text{db}\text{id}_M)=\text{id}_M.$$
The dual vector space is determined uniquely up to an isomorphism by these conditions. For the generalization for monoidal categories see \[3, Chap. 1\] or . Notice that ev plays the role of the pairing between $`V`$ and $`V^{}`$.
Although the vector space $`(VW)^{}:=\text{Hom }(VW,\text{𝕂})`$ is uniquely determined, there are more than one way of specifying a basis for this space based on given bases on $`V`$ and $`W`$. This can be seen from the point of view of the second definition above as the existence of more than one choices of the 𝕂-linear mapping $`\text{ev}_{VW}`$. More precisely, one can set $`(VW)^{}=V^{}W^{}`$, which is the usual way, or $`(VW)^{}=W^{}V^{}`$ which is actually a more standard way. In this paper we shall use both ways of identifying $`(VW)^{}`$. Note that for longer tensor products, we shall identify their duals in the above two ways (although there are more), namely $`(VW\mathrm{}U)^{}=V^{}W^{}\mathrm{}U^{}`$ or $`(VW\mathrm{}U)^{}=U^{}\mathrm{}W^{}V^{}`$. For the case of $`V^n=VV\mathrm{}V`$, the above notation may be confusing, so we shall use the notations $`V^n`$ and $`V^n`$ to refer to the first and the second identification, respectively.
One of the reasons to specify the dual spaces is for describing the matrices of adjoint operators. Recall that to each operator $`R:UV`$ on finite dimensional vector spaces, there corresponds an adjoint operator $`R^{}:U^{}V^{}`$, which is uniquely determined. For operators on tensor products, the determination of an adjoint operator obviously depends on the choice of the specification of dual spaces. In this paper we shall use the following notation for the matrices of adjoint operators: $`{}_{}{}^{t}R`$ for the usual identification $`(VW\mathrm{}U)^{}=V^{}W^{}\mathrm{}U^{}`$ and $`R^{}`$ for the “standard” identification $`(VW\mathrm{}U)^{}=U^{}\mathrm{}W^{}V^{}`$.
## 1. Matrix Quantum Semigroups of Type $`A`$
Let 𝕂 be an algebraically closed field of characteristic zero, which will be fixed through out this paper. Let $`q\text{𝕂}^\times `$ which is not a root of unity of order greater that 1. $`q`$ will also be fixed through out the paper.
### 1.1. Quadratic Algebras
Let $`V`$ be a vector space over 𝕂 of finite dimension. Let $`R`$ be a subspace of $`VV`$. Let $`𝖠=𝖠(V,R)`$ be the quotient algebra of the tensor algebra on $`V`$ by the two-sided ideal generated by elements of $`R`$:
$$𝖠(V,R):=𝖳(V)/(R).$$
Such an algebra is called quadratic algebra \[19, Chap. 1\]. $`R`$ is called the space of relations of $`𝖠`$. The two-sided ideal, generated by $`R`$, is usually denoted by $`R(𝖠)`$. Since the relations on $`𝖠`$ are homogeneous, $`𝖠`$ inherits a grading from $`𝖳(V)`$
$`𝖠={\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}𝖠_n,𝖠_n=V^n/R^n(𝖠)`$
$`R^n(𝖠)={\displaystyle \underset{i=1}{\overset{n1}{}}}R(𝖠)_i^n,R(𝖠)_i^n:=V^{i1}RV^{ni1},1in1.`$
The Poincaré (or Hilbert) series of $`𝖠`$ is by definition the formal power series $`P_𝖠(t):=_{k=0}^{\mathrm{}}\text{dim }_\text{𝕂}𝖠_k.`$ The dual quadratic algebra to $`𝖠`$, $`𝖠^!`$ is defined to be $`𝖠^!:=𝖳(V^{})/(R(𝖠)^{}).`$ It is easy to see that $`(𝖠_n^!)^{}=_{i=1}^{n1}R(𝖠)_i^n.`$
Example. Assume that $`V`$ has dimension $`d`$. Then $`𝖳(V)`$ is canonically isomorphic to the free non-commutative algebra with $`d`$ generators: $`kx_1,x_2,\mathrm{},x_d`$. The polynomial ring in $`d`$ indeterminates is the quotient of this algebra by the ideal generated by elements of the form $`x_ix_jx_jx_i`$, thus a quadratic algebra. The Poincaré series of the polynomial ring is equal to $`(1t)^d`$ as a formal series. The dual quadratic algebra is the polynomial algebra on anti-commuting indeterminates $`(\xi ^1,\xi ^2,\mathrm{},\xi ^d)`$, it is the quotient of the free non-commutative algebra $`k\xi ^1,\xi ^2,\mathrm{},\xi ^d`$ by the ideal generated by elements of the form $`\xi ^i\xi ^j+\xi ^j\xi ^i`$.
### 1.2. Quantum Spaces and Quantum Endomorphism Rings
Let $`R=R_q`$ be a Hecke operator on $`V`$, where $`q\text{𝕂}^\times `$ will be assumed not to be a root of unity of order greater than 1. We define the following quadratic algebras
(2) $`𝖲=𝖲_R:=𝖳(V)/(\text{Im }(Rq)),`$
(3) $`=_R:=𝖳(V)/(\text{Im }(R+1)).`$
$`𝖲_R`$ and $`_R`$ are considered as the function algebra and the exterior algebra on a quantum space. The function algebra on the quantum semi-group of endomorphisms of this quantum space is defined to be a quadratic algebra on $`V^{}V`$, given by
(4) $`𝖤=𝖤_R:=𝖳(V^{}V)/(\text{Im }(\overline{R}1)),`$
where $`\overline{R}:=s_{(23)}(R^1R)s_{(23)}`$, acting on $`(V^{}V)^2`$, $`s_{(23)}`$ interchanges the second and the third components in a tensor product, (see \[8, Section 1.2\]). One can check that
$$R(𝖤)=s_{(23)}\left(R(𝖲)^{}R(𝖲)R()^{}R()\right).$$
Using $`s_{(23)}`$ we shall identify $`\overline{R}`$ with $`{}_{}{}^{t}R_{}^{1}R`$ acting on $`V^2V^2`$.
Let $`\text{ev}_V`$ denote the linear map $`\text{ev}_V:V^{}Vk`$, $`\varphi x\varphi (x)`$. Then there exists uniquely a morphism $`\text{db}_V:\text{𝕂}VV^{}`$, subject to the following conditions:
$$(\text{ev}\text{id}_V^{})(\text{id}_V^{}\text{db})=\text{id}_V^{},(\text{id}_V\text{ev})(\text{db}\text{id}_V)=\text{id}_V.$$
In fact, if $`x_1,x_2,\mathrm{},x_n`$ form a basis of $`V`$ and $`\xi ^1,\xi ^2,\mathrm{},\xi ^n`$ form a basis of $`V^{}`$, such that $`\text{ev}_V(\xi ^ix_j)=\delta _j^i`$, then $`\text{db}_V`$ is given by $`\text{db}_V(1)=_k\xi ^kx_k`$. Conversely, the dual vector space $`V^{}`$ can be given in terms of the linear map $`\text{ev}_V`$ and $`\text{db}_V`$. Then $`𝖳_nV^nV^n`$ is a coalgebra with the coproduct
$$\mathrm{\Delta }_n=\text{id}\text{db}_{V^n}\text{id}:V^nV^nV^nV^nV^nV^n.$$
The direct sum of $`\mathrm{\Delta }_n`$ defines a coalgebra structure on $`𝖳`$ making it a bialgebra. It turns out that the ideal generated by $`R(𝖤)`$ is a biideal (cf. \[19, Chap. 2\]), hence $`𝖤=𝖳/R(𝖤)`$ is a bialgebra too. Set $`e_j^i:=\xi ^ix_j`$. Then $`\{e_j^i:=\xi ^ix_j,1i,jd\}`$ is a basis of $`V^{}V`$. Then the coproduct on $`𝖤`$ is given by
$$\mathrm{\Delta }(e_j^i)=\underset{k}{}e_k^ie_j^k.$$
We shall also write $`\mathrm{\Delta }(E)=E\dot{}E,`$ for $`E=(e_j^i)`$.
The relation on $`𝖤`$ can be written in terms of the matrix $`E`$ as follows (cf ):
(5) $`RE_1E_2=E_1E_2R,\text{where }E_1:=E\text{id}(d),E_2:=\text{id}(d)E\text{.}`$
$`V`$ is a right comodule over $`𝖳`$, with the coaction $`\delta _V=\text{db}_V\text{id}_V:VVV^{}V`$, hence a right comodule over $`𝖤`$. Since $`𝖤`$ is a bialgebra, $`V^n`$ is also a right comodule over $`𝖤`$. Its dual, $`(V^n)^{}`$ is a left comodule over $`𝖤`$ in a canonical way, hence a right comodule over $`𝖤^{\mathrm{cop}}`$. The relation in (5) implies that $`R`$ is a morphism of $`𝖤`$-comodules. Hence, $`_n`$, $`𝖲_n`$ are right $`𝖤`$-comodules, they are factor comodules of $`V^n`$. Actually, $`𝖤_n`$ is a subcoalgebra of $`𝖤`$ and the cocation of $`𝖤`$ on $`V^n,_n,𝖲_n`$ factorizes through $`𝖤_n`$. Since $`𝖤_n`$ is finite dimensional, $`𝖤_{n}^{}{}_{}{}^{}`$ is an algebra and its left modules are in one to one correspondence with right $`𝖤_n`$-comodules. Therefore $`𝖤_{n}^{}{}_{}{}^{}`$ acts on $`V^n,_n,𝖲_n`$. On the other hand, the operator $`R`$ induces a right action $`\rho =\rho _R`$ on $`V^n`$ of the Hecke algebra $`_n`$. We have a quantum analogue of Schur’s double centralizers theorem \[8, Theorem 2.1\]:
(6) $`𝖤_{n}^{}{}_{}{}^{}`$ $``$ $`\text{End }__n(V^n),`$
(7) $`\rho (_n)`$ $``$ $`\text{End }_{𝖤_{n}^{}{}_{}{}^{}}(V^n)=\text{End }^{𝖤_n}(V^n).`$
As a consequence, simple $`𝖤_n`$ comodules are parameterized by primitive idempotents of $`_n`$, which, in their order, are parameterized by partitions of $`n`$. For each primitive idempotents $`E_\lambda `$ of $`_n`$, $`\rho (E_\lambda )`$ is either zero of a simple $`𝖤_n`$-comodule, conjugated idempotents define isomorphic comodules. In particular, as $`𝖤_n`$-comodules
$$_n\text{Im }\rho (Y_n),𝖲_n\text{Im }\rho (X_n),$$
where $`X_n=([n]_q!)^1_{w\text{S}_n}T_w`$ , $`Y_n=([n]_{q^1}!)^1_{w\text{S}_n}q^{l(w)}T_w`$, the $`q`$-symmetrizer and $`q`$-anti-symmetrizer operators, where $`[n]_q:=(q^n1)/(q1)`$. Moreover we have isomorphisms of algebras
(8) $`𝖲{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}\rho (X_n),{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}\rho (X_n),`$
where the product on the spaces on the right-hand sides of these isomorphisms are given by
(9) $`x\rho (X_n),y\rho (X_m)xy:=\rho _{m+n}(X_{m+n})(xy)`$
(10) $`x\rho (Y_n),y\rho (Y_m)xy:=\rho _{m+n}(Y_{m+n})(xy).`$
This is the second realization of $`𝖲`$ and $``$, first considered by Gurevich . In other words, let $`X=_{n=0}^{\mathrm{}}\rho (X_n)`$. Then $`X`$ is a projection on $`𝖳(V)`$, which carries the algebra structure of $`𝖳(V)`$ to its image $`\text{Im }X`$ making this space an algebra. The second realization states that $`\text{Im }X`$ is isomorphic to $`𝖲_R`$. Further, we have
(11) $`\text{Im }\rho (X_n)={\displaystyle \underset{i=1}{\overset{n1}{}}}R()_i^n,\text{Im }\rho (Y_n)={\displaystyle \underset{i=1}{\overset{n1}{}}}R(S)_i^n.`$
This means $`𝖲`$ is isomorphic to $`^!`$, $``$ is isomorphic to $`𝖲^!`$ as graded space, where means graded dual.
Example. Let $`R_d`$ be Drinfel’d-Jimbo’s solution of the Yang-Baxter equation of type $`A_{d1}`$. Explicitly, $`R_d`$ is given as follows, with respect to a basis $`x_1,x_2,\mathrm{},x_d`$,
$`R_{d}^{}{}_{ij}{}^{kl}={\displaystyle \frac{q^2q^{2\epsilon _{ij}}}{1+q^{2\epsilon _{ij}}}}\delta _{ij}^{kl}+{\displaystyle \frac{q^{\epsilon _{ij}}(q^2+1)}{1+q^{2\epsilon _{ij}}}}\delta _{ij}^{lk},1i,j,k,ld,\epsilon _{ij}:=\text{sign }(ji).`$
The Hecke equation of $`R`$ is $`(R+1)(Rq^2)=0`$. $`\text{Im }(Rq^2)`$ has a basis consisting of elements of the form $`x_ix_jq^1x_jx_i,i<j`$. $`\text{Im }(R+1)`$ has a basis consisting of elements of the form $`x_ix_j+qx_jx_i,ij`$. Thus the relations on the algebra $`𝖲_{R_d}`$ are $`x_ix_jq^1x_jx_i=0,i<j`$ and the relations on the algebra $`_{R_d}`$ are $`x_ix_j+qx_jx_i=0,ij`$. The second realization of $`𝖲`$ and $``$ says that their elements can be realized as $`q`$-symmetrized tensors in $`kx_1,x_2,\mathrm{},x_d`$:
$`\text{for }𝖲\text{ }:x_{i_1}x_{i_2}\mathrm{}x_{i_k}{\displaystyle \underset{s\text{S}_k}{}}q^{l(w)}x_{i_1s}x_{i_2s}\mathrm{}x_{i_ks}`$
$`\text{for }\text{ }:x_{i_1}x_{i_2}\mathrm{}x_{i_k}{\displaystyle \underset{s\text{S}_k}{}}(q)^{l(w)}x_{i_1s}x_{i_2s}\mathrm{}x_{i_ks}.`$
The situation for the algebra $`𝖤`$ is more complicated. In studying the Poincaré series of $`𝖤`$, Sudbery introduced an operator $`\mathrm{\Phi }`$, acting on $`𝖤_1^3`$, which is an analogue of the symmetrizer $`X_n`$. The study was further developed in , where a complete description of the Poincaré series of $`𝖤`$ was given in terms of the Poincaré series of $``$. Define for each $`n0`$ an operator
(12) $`\mathrm{\Phi }^n:={\displaystyle \underset{w\text{S}_n}{}}(1)^{l(w)}\overline{R}_w.`$
As it has been mentioned, we identify $`\overline{R}_w`$ with $`{}_{}{}^{t}R_{}^{1}{}_{w}{}^{}R_w`$, acting on $`V^nV^n`$. For $`n=3`$, $`\mathrm{\Phi }^3`$ is the operator introduced by Sudbery. It was shown that
(13) $`\text{Im }\mathrm{\Phi }^n={\displaystyle \underset{i=1}{\overset{n1}{}}}R(𝖤)_i^n,`$
for transcendent $`q`$ \[8, Thm. 2.6\], from which follows
(14) $`P_𝖤(t)=P_𝖲P_{},`$
$``$” is the product in the $`\lambda `$-ring $`\text{}_0[[t]]`$ of power series with constant coefficient equal to 1 (see, e.g. ).
## 2. Another Realization for The Algebra $`𝖤`$
Unlike $`X_n`$ and $`Y_n`$, the operator $`\mathrm{\Phi }^n`$ introduced in (13) is not a projection. This is the main difficulty in giving the second realization for the algebra $`𝖤`$. If we know the eigenvalues of this operator, we can modify it to get a projection.
### 2.1. The Operator $`\mathrm{\Phi }^n`$
Let us denote $`R^{}:=qR^1`$. Then $`R^{}`$ also satisfies the equation $`(x+1)(xq)=0`$; hence induces a representation of $`_n`$ on $`V^n`$. Let $`\sigma =\sigma _R`$ be the representation of $`_n^{\mathrm{op}}`$ on $`(V^n)^{}`$: $`\sigma (T_w)=(R_w^{})^{}`$. If we identify $`(V^n)^{}`$ with $`V^n`$ then $`(R_w^{})^{}={}_{}{}^{t}R_{w^1}^{}.`$ Applying $`\sigma \rho `$ on the identity (1), we get
(15) $`{\displaystyle \underset{w\text{S}_n}{}}q^{l(w)}{}_{}{}^{t}R_{w}^{}R_w={\displaystyle \underset{w\text{S}_n}{}}\sigma (E_\lambda ^{ij})\rho (E_\lambda ^{ji}).`$
Recall that $`\overline{R}=q^1{}_{}{}^{t}R_{}^{}R`$. Therefore,
(16) $`{\displaystyle \underset{w\text{S}_n}{}}(1)^{l(w)}\overline{R}_w`$ $`=`$ $`{\displaystyle \underset{w\text{S}_n}{}}q^{l(w)}{}_{}{}^{t}R_{w^1}^{}R_w`$
$`=`$ $`{\displaystyle \underset{w\text{S}_n}{}}k_\lambda ^1\sigma (E_\lambda ^{ij})\rho (E_\lambda ^{ji}).`$
Denote $`\mathrm{\Phi }_\lambda :=(\sigma \rho )\mathrm{\Pi }_\lambda ={\displaystyle \underset{1i,jd_\lambda }{}}k_\lambda ^1\sigma (E_\lambda ^{ij})\rho (E_\lambda ^{ji}).`$ Then we have
(17) $`\mathrm{\Phi }_\lambda ^2=d_\lambda k_\lambda ^1\mathrm{\Phi }_\lambda .`$
Thus, the spectrum of $`\mathrm{\Phi }^n`$ is contained in $`\{d_\lambda k_\lambda ^1|\lambda n\}0`$. Note that if $`q=1`$, $`d_\lambda k_\lambda ^1=n!.`$
The isomorphisms in (6) and (7) imply the following decomposition of $`V^n`$ as $`𝖤_{n}^{}{}_{}{}^{}_n`$-bimodule:
$$V^n\underset{\lambda n}{}\text{Im }\rho (E_\lambda )S_\lambda ,$$
where $`S_\lambda `$ is the simple right $`_n`$-module, isomorphic to the right ideal in $`_n`$, spanned by $`\{E_\lambda ^{jm}|1md_\lambda \}`$, for each fixed $`j`$. Analogously, as $`𝖤_{n}^{}{}_{}{}^{\mathrm{op}}_n^{\mathrm{op}}`$-bimodules,
$$V^n\underset{\lambda }{}\text{Im }\sigma (E_\lambda )S_{\lambda }^{}{}_{}{}^{},$$
where $`S_{\lambda }^{}{}_{}{}^{}`$ is the simple right $`_n^{\mathrm{op}}`$-module, isomorphic to the right ideal in $`_n^{\mathrm{op}}`$, spanned by $`\{E_\lambda ^{mj}|1md_\lambda \}`$, for each fixed $`j`$. Therefore,
(18) $`V^nV^n{\displaystyle \underset{\lambda ,\mu n}{}}\text{Im }\sigma (E_\lambda )\text{Im }\rho (E_\mu )S_{\lambda }^{}{}_{}{}^{}S_\mu ,`$
as $`𝖤_{n}^{}{}_{}{}^{\mathrm{op}}𝖤_{n}^{}{}_{}{}^{}_n^{\mathrm{op}}_n`$-bimodules. In particular, for any elements $`x\text{Im }\sigma (E_\mu ),y\rho (E_\nu )`$, $`xyS_{\mu }^{}{}_{}{}^{}S_\nu `$ is an invariant space of $`\mathrm{\Phi }^n`$ and the action of $`\mathrm{\Phi }^n`$ on this space does not depend on $`x,y`$. Moreover, we have:
###### Lemma 2.1.
Let $`x\text{Im }\sigma (E_\mu )`$ and $`y\text{Im }\rho (E_\nu )`$. Then the operator $`\mathrm{\Phi }_\lambda `$ vanishes on $`xyS_{\mu }^{}{}_{}{}^{}S_\nu `$, if $`\mu \lambda `$ or $`\nu \lambda `$; and if $`\mu =\nu =\lambda `$, it has rank 1 on $`xyS_{\lambda }^{}{}_{}{}^{}S_\lambda `$.
Proof. Since $`S_\nu `$ (resp. $`S_{\mu }^{}{}_{}{}^{}`$) is simple over $`_n`$ (resp. $`_n^{\mathrm{op}}`$), the action of $`\mathrm{\Phi }^\lambda `$ on $`xyS_{\mu }^{}{}_{}{}^{}S_\nu `$ is equivalent to the action of $`\mathrm{\Pi }_\lambda `$. Thus, it is zero if $`\mu \lambda `$ or $`\nu \lambda `$.
Fix $`i,j`$ and let $`\{E_\lambda ^{mj}E_\lambda ^{in}|1m,nd_\lambda \}`$ be a basis of $`S_{\lambda }^{}{}_{}{}^{}S_\lambda `$. We have, as elements in $`_n^{\mathrm{op}}_n`$,
$`(E_\lambda ^{jm}E_\lambda ^{in})\mathrm{\Pi }_\lambda `$ $`=`$ $`(E_\lambda ^{jm}E_\lambda ^{in}){\displaystyle \underset{1i,jd_\lambda }{}}k_\lambda ^1E_\lambda ^{ij}E_\lambda ^{ji}`$
$`=`$ $`\delta _n^mk_\lambda ^1{\displaystyle \underset{1ld_\lambda }{}}E_\lambda ^{lj}E_\lambda ^{il}`$
That is, the action of $`\mathrm{\Pi }_\lambda `$ from the right on $`S_{\lambda }^{}{}_{}{}^{}S_\lambda `$ has rank 1.
###### Corollary 2.2.
As $`𝖤^{\mathrm{cop}}𝖤`$-comodules
$`\text{Im }(\mathrm{\Phi }_\lambda )\text{Im }\sigma (E_\lambda )\text{Im }\rho (E_\lambda ),`$
(19) $`\text{Im }(\mathrm{\Phi }^n){\displaystyle \underset{\lambda n}{}}\text{Im }\sigma (E_\lambda )\text{Im }\rho (E_\lambda ).`$
Not all $`\lambda n`$ contribute in the second decomposition of (19). For some $`\lambda `$, $`\sigma (E_\lambda )`$ or $`\rho (E_\lambda )`$ my vanish. See Theorem 2.6 for the vanishing condition of $`\sigma (E_\lambda )`$ and $`\rho (E_\lambda )`$.
### 2.2. The Operator $`\mathrm{\Psi }^n`$.
The operator $`\overline{\mathrm{\Phi }}^n`$ is not a right projection that we need for describing $`𝖤`$. The discussion in the previous subsection suggests us a new operator. According to (4), $`𝖤_2`$ is isomorphic to $`\text{Im }(\stackrel{~}{R}+q)`$, where $`\stackrel{~}{R}:=s_{(23)}\left({}_{}{}^{t}RR\right)s_{(23)}`$. We shall also identify $`\stackrel{~}{R}`$ with $`{}_{}{}^{t}RR`$ acting on $`V^nV^n`$. Let us define
$$\mathrm{\Psi }^n:=\underset{w\text{S}_n}{}q^{l(w)}\stackrel{~}{R}_w.$$
Let $`\tau `$ be the representation of $`_n^{\mathrm{op}}`$ on $`(V^n)^{}`$, induced by $`R`$. Thus, $`\tau (T_w)={}_{}{}^{t}R_{w^1}^{}`$. In analogy to (16), we have
$`{\displaystyle \underset{w\text{S}_n}{}}q^{l(w)}\stackrel{~}{R}_w={\displaystyle \underset{w\text{S}_n}{}}k_\lambda ^1\tau (E_\lambda ^{ij})\rho (E_\lambda ^{ji}).`$
Notice that $`\rho (E_\lambda )^{}`$, considered as $`_n^{\mathrm{op}}`$-module, is isomorphic to $`\tau (E_\lambda )`$. This is because $`E_\lambda `$ is also a primitive idempotent in $`_n^{\mathrm{op}}`$. Therefore, we have a decomposition of $`𝖤_{n}^{}{}_{}{}^{\mathrm{op}}𝖤_{n}^{}{}_{}{}^{}_n^{\mathrm{op}}_n`$-bimodules:
(20) $`V^nV^n{\displaystyle \underset{\lambda ,\mu n}{}}\text{Im }\rho (E_\lambda )^{}\text{Im }\rho (E_\lambda )S_{\lambda }^{}{}_{}{}^{}S_\lambda .`$
Set $`\mathrm{\Psi }_\lambda =(\rho ^{}\rho )\mathrm{\Pi }_\lambda .`$ An analogue of Lemma 2.1 holds for $`\mathrm{\Psi }_\lambda `$. Consequently, we have
###### Corollary 2.3.
There exist an isomorphism of $`𝖤^{\mathrm{cop}}𝖤`$-comodules
(21) $`\text{Im }\mathrm{\Psi }^n{\displaystyle \underset{\lambda n}{}}\text{Im }\tau (E_\lambda )\text{Im }\rho (E_\lambda ){\displaystyle \underset{\lambda n}{}}\text{Im }\rho (E_\lambda )^{}\text{Im }\rho (E_\lambda ).`$
On the other hand, according to (7), there is an algebra isomorphism
$$𝖤_{n}^{}{}_{}{}^{}\underset{\lambda n}{}\text{End }_{𝒜_\lambda }(V^n)\underset{\lambda n}{}\text{Im }\rho (E_\lambda )\text{Im }\rho (E_\lambda )^{}.$$
Consequently, we have an isomorphism of coalgebras
(22) $`𝖤_n{\displaystyle \underset{\lambda n}{}}\text{Im }\rho (E_\lambda )^{}\text{Im }\rho (E_\lambda ).`$
Therefore, as $`𝖤_n^{\mathrm{cop}}𝖤_n`$-comodules,
(23) $`𝖤_n\text{Im }\mathrm{\Psi }^n.`$
The operator $`\mathrm{\Psi }^n`$ is not a projector. However we can slightly modify it to have a projection
$$\overline{\mathrm{\Psi }}^n:=\underset{\lambda }{}k_\lambda d_\lambda ^1\tau (E_\lambda ^{ij})^{}\rho (E_\lambda ^{ji}).$$
Set $`\mathrm{\Psi }=_{n=0}^{\mathrm{}}\overline{\mathrm{\Psi }}^n.`$ Then $`\mathrm{\Psi }`$ is a projector on $`𝖳(V^{}V)`$, which then induces an algebra structure on its image $`\text{Im }\mathrm{\Psi }`$: for $`aV^nV^n`$ and $`bV^mV^m`$, $`ab:=\overline{\mathrm{\Psi }}^{n+m}(ab)`$.
###### Theorem 2.4.
Assume that the parameter $`q`$ is not a root of unity of order greater that 1. Then the projection $`\mathrm{\Psi }`$ induces an algebra isomorphism from $`𝖤=`$ to $`\text{Im }\mathrm{\Psi }`$.
Proof. We have $`(\stackrel{~}{R}+q)(\overline{R}1)=(\overline{R}1)(\stackrel{~}{R}+q)=0`$. Therefore, if $`x\text{Im }(\overline{R}_i1)`$ then $`(\stackrel{~}{R}+q)x=0`$, consequently, $`\mathrm{\Psi }^n`$ vanishes on $`_{i=1}^{n1}R(𝖤)_i^n`$; hence $`\overline{\mathrm{\Psi }}^n`$ vanishes on $`_{i=1}^{n1}R(𝖤)_i^n`$, too. Taking (23) into account, we conclude that $`_{i=1}^{n1}R(𝖤)_i^n`$ is precisely the Kernel of $`\overline{\mathrm{\Psi }}^n`$. Thus, the linear map $`\mathrm{\Psi }:𝖳(V^{}V)_{n=0}^{\mathrm{}}\text{Im }\overline{\mathrm{\Psi }}^n`$, $`a_{n=0}^{\mathrm{}}\overline{\mathrm{\Psi }}^n(a)`$ has the following properties
(26) $`\begin{array}{c}aR(𝖤)\mathrm{\Psi }(a)=0\hfill \\ a=\overline{\mathrm{\Psi }}(a)modR(𝖤).\hfill \end{array}`$
That is $`R(𝖤)=\text{Ker }\mathrm{\Psi }`$. Therefore, $`\mathrm{\Psi }`$ factorizes through $`R(𝖤)`$ to a linear isomorphism $`\psi :𝖤_{n=0}^{\mathrm{}}\text{Im }\overline{\mathrm{\Psi }}^n`$.
Let $`a(V^{}V)^m`$, $`b(V^{}V)^n`$. According to (26), we have
$$\overline{\mathrm{\Psi }}^m(a)=amodR(𝖤)^m\overline{\mathrm{\Psi }}^n(b)=amodR(𝖤)^n,$$
hence
$$\overline{\mathrm{\Psi }}^m(a)\overline{\mathrm{\Psi }}^n(b)=abmodR(𝖤)^{m+n}.$$
Consequently,
(27) $`\overline{\mathrm{\Psi }}^{m+n}(\overline{\mathrm{\Psi }}^m(a)\overline{\mathrm{\Psi }}^n(b))=\overline{\mathrm{\Psi }}^{m+n}(ab),`$
meaning that $`\psi `$ is an algebra homomorphism and therefore isomorphism $`𝖤\text{Im }\mathrm{\Psi }`$.
We now proceed to show that, for $`q`$ not a root of unity of order greater than 1 (cf. Eq. (13)),
(28) $`\text{Im }\mathrm{\Phi }^n={\displaystyle \underset{i=1}{\overset{n=1}{}}}R(𝖤)_i^n.`$
The inclusion “$``$” is obviously, for we have $`\mathrm{\Phi }^n=(\overline{R}_i1)P_i`$ for certain operator $`P_i`$, $`i=1,2,\mathrm{},n1`$. To show the equality, we compare the dimensions of $`𝖤_{}^{!}{}_{n}{}^{}`$ and $`\text{Im }\mathrm{\Phi }^n`$. Let $`l_\lambda =\text{dim }_k\text{Im }\rho (E_\lambda )`$. From the definition of the operator $`S`$, we see that $`\sigma (E_\lambda )\rho (E_\lambda ^{})`$ as vector spaces. Therefore $`\text{dim }\text{Im }\mathrm{\Phi }^n=_{\lambda n}l_\lambda ^{}l_\lambda `$.
On the other hand, since $`𝖤`$ is a Koszul algebra (with the above assumption on $`q`$), (cf. \[8, Thm 2.5\])
(29) $`P_𝖤(t)P_{𝖤^!}(t)=1.`$
Hence, according to (14), we have $`P_{𝖤^!}(t)=P_{}(t)P_{}(t)`$, that is dim$`𝖤_{}^{!}{}_{n}{}^{}=_{\lambda n}l_\lambda ^2,`$ . Therefore
(30) $`\text{dim }{\displaystyle \underset{i=1}{\overset{n1}{}}}R(𝖤)_i^n=\text{dim }(𝖤_{}^{!}{}_{n}{}^{}{}_{}{}^{})=\text{dim }E_n^!={\displaystyle \underset{\lambda n}{}}l_\lambda ^{}l_\lambda .`$
Thus, $`\text{dim }\text{Im }\mathrm{\Phi }^n=\text{dim }𝖤_{}^{!}{}_{n}{}^{}`$. Whence (28) follows.
Finally, let us denote by $`𝖥`$ the quadratic algebra on $`V^{}V`$ with relation $`R(𝖥):=\text{Im }(\stackrel{~}{R}+q)`$. Thus, $`𝖥`$ can be considered as the quantum exterior algebra over the matrix quantum semi-group. We have
(31) $`R(𝖥)=s_{(23)}\left(R(𝖲)^{}R(){\displaystyle R()^{}R(𝖲)}\right).`$
In terms of the matrix $`E`$, the relations can be given as follows:
$$RE_1E_2=qE_1E_2R.$$
The vector space $`V^{}V`$ is self dual with respect to the pairing $`(\varphi x,\psi y):=(\varphi ,y)(\psi ,x)`$. With respect to this pairing, $`𝖥`$ is canonically isomorphic to $`𝖤^!`$. Therefore $`𝖥`$ is Koszul algebra and
(32) $`P_𝖥(t)=P_𝖤(t)^1=P_{}(t)P_𝖲(t),`$
or equivalently
(33) $`\text{dim }𝖥_n={\displaystyle \underset{\lambda n}{}}l_\lambda l_\lambda ^{}.`$
Decompose $`(V^{}V)^n`$ (which is identified with $`V^nV^n`$) into simple $`_n^{\mathrm{op}}_n`$-modules (with the action given by $`\tau \rho `$). Notice that, as in the case of $`𝖤`$, $`𝖥_n`$ and $`R(𝖥)^n`$ are $`_n^{\mathrm{op}}_n`$-modules. Recall that the operator $`\mathrm{\Phi }^n`$ also has rank $`_{\lambda n}l_\lambda l_\lambda ^{}`$ on $`V^nV^n`$ and vanishes on $`R(𝖥)^n`$. Comparing the dimension we see that $`R(𝖥)^n`$ is precisely the kernel of $`\mathrm{\Phi }^n`$. Analogously, the image of $`\mathrm{\Psi }^n`$ is $`_{i=1}^{n1}R(𝖥)_i^n`$. Thus we proved an analogue of Theorem 2.4:
###### Theorem 2.5.
The algebra $`𝖥`$ is a Koszul algebra and can be realized as $`\text{Im }\mathrm{\Phi }`$, $`\mathrm{\Phi }:=_{n=0}^{\mathrm{}}\overline{\mathrm{\Phi }}^n`$.
The question when $`\rho (E_\lambda )`$ is zero can be answered by knowing the Poincaré series of $`𝖲_R`$. More precisely, is is proved that the Poincaré series of $`𝖲_R`$ is a rational function having only negative roots and positive pole (as a complex function). Let $`r`$ denote the number of poles and $`s`$ denote the number of roots of $`P_𝖲(t)`$. We call the pair $`(r,s)`$ the birank of $`R`$. For example, the birank of the operator $`R_d`$ in Section 0. Preliminaries is $`(d,0)`$.
Let $`\mathrm{\Gamma }_{r,s}`$ denote the set of partitions $`\lambda `$ such that $`\lambda _{r+1}s`$.
###### Theorem 2.6.
\[12, Theorem 5.1\] Assume that the operator $`R`$ has the birank $`(r,s)`$. Then the comodule $`V_\lambda :=\text{Im }\rho (E_\lambda )`$ is non-zero if and only if $`\lambda \mathrm{\Gamma }_{r,s}`$.
## 3. The Quantum Spaces of Homomorphisms
The notion of $`𝖤_R`$ as an “endomorphism ring” of a quantum space can be generalized to the notion of “space of homomorphisms” of two quantum spaces. Let $`R`$ and $`S`$ be Hecke operators on $`V`$ and $`W`$, respectively. We define the quadratic algebra $`𝖬=𝖬_{SR}`$ on $`W^{}V`$, whose relation is
(34) $`R(𝖬):=s_{(23)}((R(𝖲_S)^{}R(𝖲_R)R(_S)^{}R(_R)),`$
where, as usual, $`R(𝖲_S):=\text{Im }(Sq)`$, $`R(𝖲_R):=\text{Im }(Rq)`$, and so on. We have
(35) $`R(𝖬)=\text{Im }(s_{(23)}({}_{}{}^{t}S_{}^{1}R)s_{(23)}\text{id}).`$
As in the previous section, we shall identify the two vector space $`(W^{}V)^n`$ and $`W^nV^n`$.
The algebra $`𝖬`$ can be interpreted as the function algebra on the quantum space of homomorphisms (or quantum hom-space) from the quantum space associated to $`R`$ to the one associated to $`S`$. Let $`x_1,\mathrm{},x_m`$ be a basis of $`V`$ and $`\eta ^1,\mathrm{},\eta ^n`$ be a basis of $`W^{}`$. Then $`\{m_j^i:=\eta ^ix_j\}`$ form a basis of $`W^{}V`$. $`𝖬`$ is then isomorphic to
(36) $`\text{𝕂}m_j^i|1in,1jm/(SM_1M_2M_1M_2R)`$
where $`M=(m_j^i)`$, $`M_1:=M\text{id}(n)`$, $`M_2:=\text{id}(m)M`$. $`𝖬`$ has the following properties: Let $`A:=(a_1,a_2,\mathrm{},a_n)`$ be a point of $`𝖲_S`$, i.e., $`(AA)(Rq)=0`$, and $`T=(t_j^i)`$ be a point of $`𝖬_{SR}`$, such that $`a_i`$ and $`t_l^k`$ commute. Then $`A\dot{}M`$ is a point of $`𝖲_R`$, where $`(A\dot{}M)_i:=_ka_km_i^k`$. Analogously, if $`B`$ is a point of $`_S`$ commuting with $`T`$ then $`B\dot{}T`$ is a point of $`_R`$.
There is also an interpretation of $`𝖬_{SR}`$ from the categorical view-point. The bialgebra $`𝖤`$ can be constructed as the Coend of the functor F from the braided monoidal category $`𝒱`$, generated by one object $`v`$ and one morphism $`\tau :v^2v^2`$, into the category of vector space, such that $`\text{F}(v)=V`$ and $`\text{F}(\tau )=R`$ (cf. ). That is, for any vector space $`X`$,
(37) $`\text{Nat }(\text{F},\text{F}X)\text{Hom }_\text{𝕂}(𝖤,X)`$
where $`\text{Nat }(\text{F},\text{G})`$ denotes the set of natural transformation between functors F and G, $`\text{F}X`$ is the functor that sends $`w`$ to $`\text{F}(w)X`$ and sends $`f`$ to $`\text{F}(f)\text{id}_X`$, $`v𝒱,f\text{Mor }(𝒱)`$. Let us now consider another functor G, with $`\text{G}(v)=W`$ and $`\text{G}(\tau )=T`$. Then we have
(38) $`\text{Nat }(\text{F},\text{G}X)\text{Hom }(𝖬_{SR},X).`$
The exterior algebra on the quantum hom-space is defined to be
(39) $`𝖭=𝖭_{SR}:=𝖳(W^{}V)/\text{Im }(s_{(23)}({}_{}{}^{t}SR)s_{(23)}+q\text{id}).`$
The bialgebra $`𝖤_R`$ coacts on $`𝖬`$ and $`𝖭`$ from the right. The coaction is induced from the one on $`W^{}V`$: $`\delta (m_j^i)=_km_k^ie_j^k.`$ Analogously, $`𝖤_S`$ coacts on $`𝖬`$ and $`𝖭`$ from the left, with the coaction induced from $`\delta (m_j^i)=_le_l^im_j^l.`$ Thus, $`𝖬`$ and $`𝖭`$ are right $`𝖤_S^{\mathrm{cop}}𝖤_R`$-comodule algebras.
We show in this section that $`𝖬`$ and $`𝖭`$ are Koszul algebras, compute their Poincaré series and give a second realization. Since $`N_{SR}=𝖬_{S^{}R}`$, $`S^{}=qS^1`$, it is sufficient to study $`𝖬_{SR}`$.
As in (18), we have a decomposition of $`(W^{}V)^n`$ as an $`𝖤_{S}^{}{}_{}{}^{\mathrm{cop}}𝖤_R_n^{\mathrm{op}}_n`$-bimodule
(40) $`W^nV^n{\displaystyle \underset{\lambda ,\mu n}{}}\text{Im }\sigma (E_\lambda )\text{Im }\rho (E_\mu )S_{\lambda }^{}{}_{}{}^{}S_\mu .`$
The subspace $`R(𝖬)^n`$ of $`(W^{}V)^n`$ is an $`𝖤_{S}^{}{}_{}{}^{\mathrm{cop}}𝖤_R`$-comodule, hence
$$R(𝖬)^n=\underset{\lambda ,\mu n}{}\left(R(𝖬)^n\text{Im }\sigma _S(E_\lambda )\text{Im }\rho _R(E_\mu )S_{\lambda }^{}{}_{}{}^{}S_\mu \right).$$
Remark. if the action of $`{}_{}{}^{t}S_{i}^{}R_i`$ on $`S_{\lambda }^{}{}_{}{}^{}S_\mu `$ is not zero then this action does not depend on $`S`$ and $`R`$. In fact, the action of $`{}_{}{}^{t}S_{i}^{}R_i`$ on $`S_{\lambda }^{}{}_{}{}^{}S_\mu `$ in this case is the action of $`T_iT_i`$.
Define the operators $`\mathrm{\Phi }_{SR}`$ and $`\mathrm{\Psi }_{SR}`$ as in (12) with $`\overline{R}=q{}_{}{}^{t}S_{}^{1}R`$ and $`\overline{R}={}_{}{}^{t}SR`$, respectively. The corresponding projectors $`\overline{\mathrm{\Phi }}_{SR}`$ and $`\overline{\mathrm{\Psi }}_{SR}`$ are defined similarly. Notice that $`\mathrm{\Phi }_{SR}=\mathrm{\Psi }_{S^{}R}`$. From the proof of Theorem 2.4 and using the above remark we have
$$R(𝖬)^n=\text{Ker }\overline{\mathrm{\Psi }}_{SR}^n.$$
Consider the action of $`\overline{\mathrm{\Psi }}_{SR}^n`$ on a module $`xyS_{\lambda }^{}{}_{}{}^{}S_\mu `$, $`x\sigma (E_\lambda ),y\rho (E_\mu )`$. This has rank 1 if $`\lambda =\mu `$ and 0 otherwise. Therefore, for $`k_\lambda :=\text{dim }\text{Im }\rho _S(E_\lambda ),l_\lambda :=\text{dim }\text{Im }\rho _R(E_\lambda )`$, we have
$$\text{dim }_k𝖬_n=\text{dim }_k(𝖬_1^n/R(𝖬)^n)=\underset{\lambda n}{}l_\lambda k_\lambda .$$
Since $`\overline{\mathrm{\Psi }}_{SR}^n`$ is a projector, so is the map $`\mathrm{\Psi }_{SR}:_{n=0}^{\mathrm{}}\overline{\mathrm{\Psi }}_{SR}^n`$, which induces an isomorphism of algebras $`𝖬_{n=0}^{\mathrm{}}\text{Im }\mathrm{\Psi }_{SR}`$. This map is also a homomorphism of $`𝖤_S^{\mathrm{cop}}𝖤_R`$-comodules because each $`\overline{\mathrm{\Psi }}_{SR}`$ is.
The remark above implies that the lattice induced by $`R(𝖬)_i^n(S_{\lambda }^{}{}_{}{}^{}S_\mu ),1in1`$ is distributive in $`(S_{\lambda }^{}{}_{}{}^{}S_\mu )`$ (see for more details). Consequently, the lattice generated by $`R(𝖬)_i^n,1in1`$ is distributive in $`W^nV^n`$, that is $`𝖬_{SR}`$ is a Koszul algebra.
###### Theorem 3.1.
Let $`R=R_q`$ and $`S=S_q`$ be Hecke operators, where $`q\text{𝕂}^\times `$ is not a root of unity of order greater than 1. Then the algebra $`𝖬_{SR}`$ is a Koszul algebra, its Poincaré series are given by
(41) $`P_𝖬(t)=P_{𝖲_R}P_{𝖲_S}(t).`$
There is a realization of $`𝖬`$ a subspace of $`𝖳(W^{}V)`$: the following is an isomorphism of $`𝖤_S𝖤_R`$-bicomodule algebras
(42) $`𝖬{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}\text{Im }\overline{\mathrm{\Phi }}_{SR}^n`$
where, for each $`n`$, as $`𝖤_S𝖤_R`$-bicomodules
$$𝖬_n\text{Im }\overline{\mathrm{\Psi }}_{SR}^n=\underset{\lambda n}{}\text{Im }\rho _S(E_\lambda )^{}\text{Im }\rho _R(E_\lambda ).$$
Further, we have
(43) $`\text{Im }(\overline{\mathrm{\Psi }}_{SR}^n)={\displaystyle \underset{i=1}{\overset{n1}{}}}R(𝖬)_i^n.`$
## 4. Quantum Determinantal Ideals
In commutative algebra, the ideal $`I_k`$, generated by the $`k\times k`$-minors in the coordinate ring of the varieties $`M_k(m,n)`$ is called determinantal ideal ($`0k\mathrm{min}(m,n)`$). This ideal is invariant with respect to a natural action of the group $`G=GL_k(m)\times GL_k(n)`$ on $`𝒪(M_k(m,n))`$. It is proved to be prime, see, e.g. . The variety determined by $`I_k`$ is called determinantal variety.
In the quantum setting, we have a coaction of the bialgebra $`𝖦=𝖤_S^{\mathrm{cop}}𝖤_R`$ on the algebra $`𝖬_{SR}`$. A subspace (resp. two-sided ideal) in $`𝖬_{SR}`$, which is invariant with respect to the coaction of $`𝖦`$ will be called invariant subspace (resp. invariant ideal). Let $`M_\lambda `$ denote $`\text{Im }\tau (E_\lambda )\text{Im }\rho (E_\lambda )`$. Then any invariant subspace of $`𝖬`$ is a direct sum of some $`M_\lambda `$, $`\lambda 𝒫`$, ($`𝒫`$ is the set of all partitions).
Let us denote by $`𝒫_{SR}`$ the set of partitions $`\lambda `$, such that $`k_\lambda l_\lambda 0`$ (see the previous section). This set can be fully described using Theorem 2.6. In this section we show that there is a one-one correspondence between $`𝖦`$-invariants ideals in $`𝖬`$ and D-ideals in $`𝒫_{SR}`$. The latter is defined as follows: a subset $`J`$ of $`𝒫_{SR}`$ is called a D-ideal if for any $`\sigma J`$ and any $`\tau 𝒫_{SR}`$, such that $`\tau \sigma `$, one has $`\tau J`$ (cf. ). For any subset $`J`$ of $`𝒫_{SR}`$, let $`I(J)`$ denote the subspace $`_{\sigma J}M_\sigma `$.
We need the following key lemma.
###### Lemma 4.1.
Let $`\lambda `$ and $`\mu `$ be partitions, such that all the Littlewood-Richardson coefficients $`c_{\lambda \mu }^\gamma `$ is at most 1. Let $`C_{\lambda ,\mu }`$ denote the set of partitions $`\gamma `$ such that $`c_{\lambda ,\mu }^\gamma =1`$. Then the image of the product of $`𝖬`$, restricted on $`M_\lambda M_\mu `$, is $`I(C_{\lambda ,\mu })`$.
Reference for the Littlewood-Richardson coefficients is . The proof of this lemma will be given at the end of this section.
Notice that for any partition $`\lambda `$, $`\lambda `$ and $`(1^k)`$ satisfies the condition of the lemma above. In fact, using the Littlewood-Richardson rule $`c_{\lambda (1^k)}^\gamma `$ is equal to 1 if and only if $`\gamma \lambda `$ and $`\gamma _j\lambda _i1`$, $`|\gamma ||\lambda |=k`$; otherwise $`c_{\lambda (1^k)}^\gamma =0.`$ Thus we have
###### Corollary 4.2.
For any partition $`\lambda `$ and any integer $`k`$, $`M_\lambda M_{(1^k)}=I(C_{\lambda ,(1^k)})`$.
Following , we denote by $`I_\sigma `$ the ideal in $`𝖬`$, generated by $`M_\sigma `$. We have
###### Theorem 4.3.
$`I_\sigma =_{\tau \sigma }M_\tau .`$
Proof. Since $`M_\sigma M_{(1)}=I(C_{\sigma (1)})`$ and $`\gamma C_{\sigma (1)}`$ implies $`\gamma \sigma `$, we have $`I_\sigma _{\tau \sigma }𝖬_\tau .`$
Assume $`\tau \sigma `$ and $`|\tau |=|\sigma |+1`$. Then $`\tau C_{\lambda (1)}`$. Hence, $`M_\tau `$ appears in $`M_\sigma M_{(1)}`$, thus $`M_\tau I_\sigma `$. By induction, $`_{\tau \sigma }M_\tau I_\sigma .`$
For each set $`J𝒫`$, let $`J`$ denote the smallest D-ideal in $`𝒫`$, containing $`J`$. The following facts follow immediately from Theorem 4.3.
###### Corollary 4.4.
1. The ideal generated by $`I(S)`$ for some subset $`S𝒫`$ is $`I(S)`$.
2. $`I_\sigma I_\tau `$ if and only if $`\tau \sigma `$.
3. Let $`J𝒫`$. Then $`I(J)`$ is an ideal in $`𝖬`$ if and only if $`J`$ is a D-ideal in $`𝒫`$.
We now describe the product of determinantal ideals. We need the following order on the set of partitions (cf. ). Firstly, we identify a partition with the diagram it determines. Define functions $`\beta _k`$, $`k=1,2,\mathrm{}`$, on the set of diagram as follows. For a partition $`\sigma =(\sigma _1,\sigma _2,\mathrm{})`$, $`\beta _k(\sigma ):=\sigma _1^{}+\sigma _2^{}+\mathrm{}+\sigma _k^{}`$, where $`\sigma ^{}`$ is the conjugate partition of $`\sigma `$. In other words, $`\beta _k(\sigma )`$ is the number of boxes in the first $`k`$ columns of the diagram determined by $`\sigma `$. We say that $`\tau \sigma `$ if for all $`k`$, $`\beta _k(\tau )\beta _k(\sigma )`$.
###### Proposition 4.5.
Let $`\sigma =(\sigma _1,\sigma _2,\mathrm{})`$ be a partition. Then the product of ideals $`I_{\sigma _i}`$ is $`I(D_\sigma ^{})`$ where $`D_\sigma ^{}`$ is the set of partitions $`\tau `$ such that $`\tau \sigma ^{}`$ in the above defined order.
Proof. We use induction. If $`\sigma `$ has only one non-zero component, the claim is obvious. Let $`\sigma =(\sigma _1,\sigma _2,\mathrm{},\sigma _r)`$. Denote $`\sigma _{}`$ the partition $`\sigma =(\sigma _1,\sigma _2,\mathrm{},\sigma _{r1})`$. Using the induction assumption, one reduces it to showing that
$$I(D_\sigma )=I(D_\sigma _{})I_{\sigma _r}.$$
Let $`M_\lambda I(D_\sigma _{})I_{\sigma _r}.`$ Then $`\lambda C_{\gamma (1^{\sigma _r})}`$, for some $`\gamma D_\sigma _{}`$, by Corollary 4.2. In this case $`\lambda `$ contains $`\gamma `$ and satisfies $`\lambda _i\gamma _i1`$, $`|\lambda ||\gamma |=\sigma _r.`$ Hence, for any $`k=1,2,\mathrm{},r1`$, $`\beta _k(\lambda )\beta _k(\gamma )\beta _k(\sigma _{})=\beta _k(\sigma )`$, and for $`k=r`$, $`\beta _r(\lambda )=\beta _{r1}(\gamma )+\sigma _r\beta _{r1}(\sigma _{})+\sigma _r=\beta _r(\sigma )`$. Thus, $`\lambda 𝒟_\sigma `$, that is $`I(D_\sigma _{})I_{\sigma _r}I(D_\sigma ).`$
To show the converse inclusion we need the following lemma.
###### Lemma 4.6.
With the above notations, if $`\lambda \sigma ^{}`$ then there exists $`\gamma \lambda `$, such that $`\lambda \sigma _{}^{}{}_{}{}^{}`$ and $`\lambda C_{\gamma (1^{\sigma _r})}`$.
Proof. Instead of considering partitions $`\lambda `$ and $`\sigma `$, we consider their diagrams using the same notation. By assumption, the number of boxes in the first $`k`$ columns of $`\lambda `$ is greater than the corresponding number in $`\sigma `$, for $`k=1,2,\mathrm{}`$ Let $`b(\lambda )`$ denote the set of boxes lying in the rightest place in each row. Thus the number of boxes in $`b(\lambda )`$ is $`\beta _1(\lambda )`$. Since $`\lambda \sigma ^{}`$, $`|b(\lambda )|b_1(\sigma ^{})=\sigma _1\sigma _r`$. To obtain $`\gamma `$, we remove from $`\lambda `$ $`\sigma _r`$ boxes in $`b(\lambda )`$. Starting from the lowest one in the rightest column we remove all the boxes bottom up (all these boxes belong to $`b(\lambda )`$). Then we remove those boxes of $`b(\lambda )`$ lying in the second rightest column in the same way and keep doing this further to the left. We remove as many boxes as the number of boxes in the rightest column of $`\sigma ^{}`$, i.e. $`\sigma _r`$ boxes. The diagram obtained is $`\gamma `$.
To see that $`\gamma \sigma _{}^{}{}_{}{}^{}`$, we proceed as follows. Each time when we remove a box from $`\lambda `$, we also remove a box from the rightest column of $`\sigma ^{}`$ in the order bottom up. Thus we have two sequences of diagram denoted by $`\lambda =\lambda (0),\lambda (1),\mathrm{},\lambda (\sigma _r)=\gamma `$ and $`\sigma ^{}=\sigma ^{}(0),\sigma ^{}(1),\mathrm{},\sigma ^{}(\sigma _r)=\sigma _{}^{}{}_{}{}^{}`$. It is easy to see that if $`\lambda (t)\sigma ^{}(t)`$ then $`\lambda (t+1)\sigma ^{}(t+1)`$. Since $`\lambda \sigma ^{}`$, we conclude that $`\gamma \sigma _{}^{}{}_{}{}^{}`$.
On the other hand the set $`C_{\gamma (1^{\sigma _r})}`$ contains such partitions $`\lambda `$ that $`0\tau _k\gamma _k1`$ and $`|\tau ||\gamma |=\sigma _r`$. Thus, the construction of $`\gamma `$ above also implies that $`\lambda C_{\gamma (1^{\sigma _r})}`$. The Lemma is therefore proved.
Assume now that $`\lambda D_\sigma `$, by Lemma, there exist $`\gamma \lambda `$ such that $`\gamma \sigma _{}`$ and $`\lambda C_{\gamma (1^{\sigma _r})}`$. That is $`\gamma D_\sigma _{}`$ and, by Corollary 4.2, $`M_\lambda I(D_\sigma _{}^{}{}_{}{}^{})I_{\sigma _r}`$. That is, $`I(D_\sigma ^{})I(D_\sigma _{})I_{\sigma _r}.`$ The proposition is proved.
The next interesting question is whether the primeness, radicals, preserve through the correspondence $`JI(J)`$. This problem is completely open. The following property is the only one in this direction that I am able to prove.
###### Proposition 4.7.
Let $`J`$ be an invariant ideal in $`𝖬`$. Then $`\sqrt{J}`$ is invariant, too.
We first recall the definition of the radical of a (two-sided) ideal in a non-commutative ring. An $`m`$-system in a ring $`R`$ is a set of elements such that for each two of its elements $`a,b`$, there exists an element $`r`$ from the ring, such that $`arb`$ belong to this set. Let $`I`$ be an ideal. The radical $`\sqrt{I}`$ of $`I`$ is defined to be the set of elements $`s`$, such that any $`m`$-system containing $`s`$ intersects non-trivially with $`I`$.
Proof. We have to show that for any element $`a\sqrt{I}`$, the $`𝖦`$-invariant space generated by $`a`$ also belong to $`\sqrt{I}`$. To simplify the discussion, we introduce the algebra $`𝖦^{}`$. Recall that $`𝖦=𝖤_S^{\mathrm{op}}𝖤_R`$ is a cosemisimple bialgebra. Hence its dual is the completion of a direct sum of endomorphism rings on certain vector spaces
$$𝖦^{}=\overline{\underset{\lambda \mu }{}\text{End }(V_{\lambda \mu })}\underset{\lambda \mu }{}\text{End }(V_{\lambda \mu }),$$
where the product on the right-hand side is defined componentwise. Let $`\pi _{\lambda \mu }`$ be the projection on the $`\lambda \mu `$-component. It is obvious that an element $`f`$ of $`𝖦^{}`$ is invertible if and only if $`\pi _{\lambda \mu }f`$ is invertible in $`\text{End }(V_{\lambda \mu })`$, for all $`\lambda ,\mu `$. Let $`U`$ be the of invertible elements in $`𝖦^{}`$. Since the set of invertible elements in $`\text{End }(V_{\lambda \mu })`$ spans this vector space, the set $`Um`$ spans the submodule generated by $`m`$. Let now $`m\sqrt{I}`$. It is thus sufficient to check that $`Um\sqrt{I}`$. This is an easy consequence of the fact that $`U`$ consists of invertible elements and the product on $`𝖬`$ is $`𝖦^{}`$-equivariant. The proposition is proved.
Proof of Lemma 4.1. We use the characterization of the product of $`𝖬`$ given in Theorem 2.4. Thus for $`aM_\lambda `$, $`bM_\mu `$, $`ab=\overline{\mathrm{\Psi }}^{l+m}(ab)`$, where $`l=|\lambda |,m=|\mu |`$. It is to show that for any given $`\gamma C_{\lambda \mu }`$ there exist such $`a`$ and $`b`$, that $`\overline{\mathrm{\Psi }}^\gamma (ab)0`$ . According to the fact that the product is $`𝖦`$-equivariant and that $`M_\sigma `$ is a simple $`𝖦`$-comodule the assertion of Lemma will then follow.
Denote $`V_\lambda :=\text{Im }\rho (E_\lambda )`$. We have an isomorphism of $`𝖤_{n}^{}{}_{}{}^{\mathrm{op}}𝖤_{n}^{}{}_{}{}^{}_n^{\mathrm{op}}_n`$ modules (cf. Eq. (20))
$$V^nV^nV_\lambda ^{}V_\lambda S_{\lambda }^{}{}_{}{}^{}S_\mu .$$
For any fixed $`i`$ and $`j`$, $`1i,jd_\lambda `$, we can assume that the set $`\{E_\lambda ^{mi}|1md_\lambda \}`$ and $`\{E_\lambda ^{jm}|1md_\lambda \}`$ are bases of $`S_{\lambda }^{}{}_{}{}^{}`$ and $`S_\lambda `$, respectively. In this setting, an element of $`M_\lambda `$ can be represented as a linear combination of elements of the form $`v(_mE_\lambda ^{mi}E_\lambda ^{jm})`$ for some $`vV_\lambda ^{}V_\lambda `$. That is, as a subspace of $`V_\lambda ^{}V_\lambda S_{\lambda }^{}{}_{}{}^{}S_\lambda `$,
$$M_\lambda =V_\lambda ^{}V_\lambda (\underset{m}{}E_\lambda ^{mi}E_\lambda ^{jm}).$$
We have, with the assumption on $`\lambda `$ and $`\mu `$
(44) $`V_\lambda V_\mu {\displaystyle \underset{\sigma C_{\lambda ,\mu }}{}}V_\sigma ,V_\lambda ^{}V_\mu ^{}{\displaystyle \underset{\sigma C_{\lambda ,\mu }}{}}V_\sigma ^{},`$
as $`𝖤_{l+m}`$-comodules, where the coaction of $`𝖤_{l+m}`$ on the left-hand side is induced by the product on $`𝖤`$, $`𝖤_l𝖤_m𝖤_{m+l}`$. Therefore,
$$M_\lambda M_\mu =\underset{\sigma ,\eta C_{\lambda ,\mu }}{}V_\sigma ^{}V_\eta \left(\underset{m,n}{}E_\lambda ^{mi}E_\lambda ^{jm}E_\mu ^{nl}E_\mu ^{kn}\right),$$
so that, for each $`\gamma C_{\lambda ,\mu }`$, we can choose $`vV_\lambda ^{}V_\lambda `$ and $`wV_\mu ^{}V_\mu `$, such that the projection $`\pi _\gamma (vw)`$ of $`vw`$ on $`V_\gamma ^{}V_\gamma `$ through the above isomorphism is not zero. The crucial point here is that in the decomposition (44) is multiplicity-free.
On the other hand, if we embed $`_l_m`$ into $`_{l+m}`$ in the standard way, then $`S_\gamma `$, considered as $`_l_m`$-module, contains $`S_\lambda S_\mu `$ as a simple subcomodule. Therefore, the space
$$V_\gamma ^{}V_\gamma \left(\underset{m,n}{}E_\lambda ^{mi}E_\lambda ^{jm}E_\mu ^{nl}E_\mu ^{kn}\right)$$
is a subspace of $`V_\gamma ^{}V_\gamma S_{\gamma }^{}{}_{}{}^{}S_\gamma `$. The element $`(ab)\overline{\mathrm{\Psi }}^\gamma `$ is then
$$\pi _\gamma (vw)\left(\underset{m,n}{}E_\lambda ^{mi}E_\lambda ^{jm}E_\mu ^{nl}E_\mu ^{kn}\right)\overline{\mathrm{\Psi }}_\gamma .$$
Here, the element $`(_mE_\lambda ^{mi}E_\lambda ^{jm})`$ is considered as an element of $`_{l+m}`$ by the embedding $`_l_m_{l+m}`$. Thus, it remains to show that (in $`_n^{\text{op}}_n`$)
(45) $`\left({\displaystyle \underset{m,n}{}}E_\lambda ^{mi}E_\lambda ^{jm}E_\mu ^{nl}E_\mu ^{kn}\right)\mathrm{\Pi }_\gamma 0.`$
Since the sets $`\{E_\lambda ^{jm},1md_\lambda \}`$, for $`i=i_1`$ and $`i=i_2`$ can be obtained from each other by multiplying with $`E_\lambda ^{i_1i_2}`$ on the left, (45) does not depend on $`i,j`$ and $`k,l`$. That means, if (45) holds (or fails) for some $`i,j,k,l`$, it should holds (or fails) for all $`i,j,k,l`$. So that, if on the left-hand side of (45) we set $`i=j`$, $`k=l`$ and summing it up after these indices and show that this element is not zero, we will be done. Thus, we have to show that
$`\left({\displaystyle \underset{mi}{}}E_\lambda ^{m,i}E_\lambda ^{im}E_\mu ^{nk}E_\mu ^{kn}\right)\mathrm{\Pi }_\gamma 0,`$
or $`\mathrm{\Pi }_\lambda \mathrm{\Pi }_\mu \mathrm{\Pi }_\gamma 0`$, for any $`\gamma C_{\lambda ,\mu }`$, here, $`\mathrm{\Pi }_\lambda `$ and $`\mathrm{\Pi }_\mu `$ are considered as elements of $`_{l+m}^{\mathrm{op}}_{l+m}`$. Let $`F_\lambda `$ be the minimal central primitive element in $`_l`$, corresponding to $`\lambda `$. Then $`F_\lambda `$ is also a central primitive element in $`_l^{\mathrm{op}}`$. Hence from the definition of $`\mathrm{\Pi }_\lambda `$ and (1), we have
(46) $`\mathrm{\Pi }_\lambda =k_\lambda ^1(F_\lambda F_\lambda )\left({\displaystyle \underset{w\text{S}_l}{}}q^{l(w)}R_wR_{w^1}\right).`$
Analogous equalities hold for $`\mathrm{\Pi }_\mu `$ and $`\mathrm{\Pi }_\gamma `$. Let $`𝒟_\lambda `$ be the set of left coset representatives of $`\text{S}_l`$ in $`\text{S}_{l+m}`$, such that $`l(w)l(t)=l(wt)`$ for $`w\text{S}_l,t𝒟_\lambda `$. Then we have, as elements of $`_{l+m}^{\mathrm{op}}_{l+m}`$,
$$\underset{v\text{S}_{l+m}}{}q^{l(v)}R_vR_{v^1}=\left(\underset{w\text{S}_l}{}q^{l(w)}R_wR_{w^1}\right)\left(\underset{t𝒟_\lambda }{}q^{l(t)}R_tR_{t^1}\right).$$
According to (46), we have
$$\mathrm{\Pi }_\lambda \underset{w\text{S}_l}{}q^{l(w)}R_wR_{w^1}=d_\lambda k_\lambda ^1(F_\lambda F_\lambda )\underset{w\text{S}_l}{}q^{l(w)}R_wR_{w^1}$$
Therefore,
$`\mathrm{\Pi }_\lambda \mathrm{\Pi }_\gamma =\mathrm{\Pi }_\lambda \left({\displaystyle \underset{w\text{S}_l}{}}q^{l(w)}R_wR_{w^1}\right)\left({\displaystyle \underset{t𝒟_\lambda }{}}q^{l(t)}R_tR_{t^1}\right)(F_\gamma F_\gamma )`$
$`=`$ $`d_\lambda k_\lambda ^1(F_\lambda F_\lambda )\left({\displaystyle \underset{w\text{S}_l}{}}q^{l(w)}R_wR_{w^1}\right)\left({\displaystyle \underset{t𝒟_\lambda }{}}q^{l(t)}R_tR_{t^1}\right)(F_\gamma F_\gamma )`$
$`=`$ $`d_\lambda k_\lambda ^1(F_\lambda F_\lambda )(F_\gamma F_\gamma ){\displaystyle \underset{w\text{S}_{l+m}}{}}q^{l(w)}R_wR_{w^1}.`$
Thus, it is led to showing that $`F_\lambda F_\mu F_\gamma 0`$. This is obvious by the assumption, that $`\gamma C_{\lambda ,\mu }`$. The Lemma is therefore proved.
## 5. Invariant theory
Let $`m,n,t`$ be positive integers. The group $`GL(t)`$ acts on the variety $`M(m,t)\times M(t,n)`$ in the following way:
(47) $`g(A,B)=(Ag^1,gB),gGL(t),AM(m,t),BM(t,n).`$
This action induces an action of $`GL(t)`$ on the coordinate ring on $`M(m,t)\times M(t,n)`$, which is a polynomial ring in $`mt+tn`$ variables. The classcial invariant theory studies the ideal of polynomials, which are invariant under this action. Let $`\mu `$ be the natural morphism of affine varieties $`M(m,t)\times M(t,n)M(m,n),`$ $`(A,B)AB`$, inducing a morphism of algebras
$$\mu ^{}:𝒪(M(m,n)𝒪(M(m,t)\times M(t,n))𝒪(M(m,t))𝒪(M(t,n)).$$
Let $`m_k^i`$ (resp. $`n_j^i,p_k^j`$) be the standard generators of $`𝒪(M(m,n))`$, (resp. $`𝒪(M(m,t)),𝒪(M(t,n))`$), such that $`\mu ^{}`$ is given by
$$\mu ^{}(m_k^i)=\underset{k}{}n_j^ip_k^j.$$
The first and the second fundamental theorems for general linear groups state that
1. Any invariant polynomial on $`M(m,t)\times M(t,n)`$ can be obtained by composing a polynimial on $`M(m,n)`$ with $`\mu `$, or, equivalently, the ideal of invariant polynomials is precisely $`\text{Im }\mu ^{}`$. Thus, it is the quotient of $`𝒪(M(m,n))`$ by $`\text{Ker }\mu ^{}`$.
2. The ideal $`\text{Ker }\mu ^{}`$ in $`𝒪(M(m,n))`$ is generated by the minors of rank $`(t+1)\times (t+1)`$, i.e., it is the ideal $`I_{t+1}`$.
In this section we formulate and prove a quantum analogue of the above theorems for arbitrary Hecke operators $`S,R`$ and $`T`$, acting on $`U,V`$ and $`W`$ respectively. Thus $`M(m,t)`$ (resp. $`M(t,n),M(m,n),GL(t)`$) will be replaced by $`𝖬_{RS}`$ (resp. $`𝖬_{TR},𝖬_{TS},𝖧_R`$). Here, $`𝖧_R`$ is a Hopf algebra associated to $`R`$. The left action in (47) is replaced by a right coaction of $`𝖧_R`$ on $`𝖬_{TR}𝖬_{RS}`$. The set of polynomials on $`M(m,t)\times M(t,n)`$, invariant with the action of $`GL(t)`$ now corresponds to the set of coinvariants of the coaction $`\delta _{RST}`$, i.e. the set of $`x𝖬_{TR}𝖬_{RS}`$ such that $`\delta _{RST}(x)=x1`$.
In the quantum case, the are (at least) two ways to define the action of the Hopf algebra $`𝖧_R`$ on the algebra $`𝖬_{RS}`$, which coincide when the Hecke operators reduce to the ordinary flip operators. Therefore there are (at least) two versions of the fundamental theorems, depending on the way we define the coaction and on the way we define the algebra structure on $`𝖬_{TR}𝖬_{RS}`$. In 5.1 we define an algebra structure on $`𝖬_{TR}𝖬_{RS}`$ in a usual way and choose an appropriate coaction of $`𝖧_R`$ on it. In this setting $`𝖬_{TR}`$ is an $`𝖧_R`$-comodule algebra, $`𝖬_{RS}`$ is an $`𝖧_R`$-comodule but not comodule algebra. Another setting is considered Subsection 5.2, in which we modify $`𝖬_{RS}`$ so that it becomes $`𝖧_R`$-comodule algebra (in fact, $`R`$ is replaced by $`\widehat{R}`$). The algebra structure on $`𝖬_{TR}𝖬_{\widehat{R}S}`$ is also modified making the new algebra an $`𝖧_R`$-comodule algebra.
### 5.1. Fundamental theorems for quatum groups of type $`A`$, the first version
In formulating the fundamental theorems for quantum groups of type $`A`$, we need to intoduce the multiplication map $`\mu ^{}`$ and the coaction of the quantum group $`𝖧_R`$.
We first define the morphism $`\mu ^{}`$. The linear map
$$\theta _1=\text{id}\text{db}_V\text{id}:W^{}UW^{}VV^{}U𝖳(W^{}V)𝖳(V^{}U)$$
induces an algebra morphism
$$\theta :𝖳(W^{}U)𝖳(W^{}V)𝖳(V^{}U).$$
Note that $`𝖳(W^{}V)𝖳(V^{}U)`$ is the tensor product of the algebras$`𝖳(W^{}V)`$ and $`𝖳(V^{}U)`$, that is, the elements from the latter algebras commute in $`𝖳(W^{}V)𝖳(V^{}U)`$. In other words, we identify the two vector spaces
(48) $`(W^{}VV^{}U)^n(W^{}V)^n(V^{}U)^n,`$
by means of the usual flip operator (that changes orders of tensor components). The restriction of $`\theta `$ on the $`n^{\mathrm{th}}`$ component is then the $`n^{\mathrm{th}}`$ tensor power of $`\theta _1`$, taking in account the above identification.
Fix bases of $`U,V,W`$ and then define their dual bases on $`U^{},V^{},W^{}`$. These bases define bases for $`W^{}U,W^{}V,V^{}U`$, which will be denoted by $`M=(m_k^i),N=(n_j^i),L=(l_k^j),`$ respectively. For convennience, we shall omit all tensor signs when describe an element of the algebra $`𝖳(W^{}V)𝖳(V^{}U)`$. Then we have $`\theta (M)=NL`$. Since in the algebra $`𝖳(W^{}V)𝖳(V^{}U)`$, the entries of $`N`$ and $`L`$ commute, i.e. $`L_1N_2=N_2L_1,`$ where $`L_1=L\text{id},L_2=\text{id}L`$, we have
$$\theta (M_1M_2)=(NL)_1(NL)_2=N_1N_2L_1L_2.$$
Notice that $`(MN)_1=M_1N_1,(MN)_2=M_2N_2`$.
Combining $`\theta `$ with the quotient map $`𝖳(W^{}V)𝖳(V^{}U)𝖬_{TR}𝖬_{RS}`$, we obtain an algebra morphism
$$\overline{\theta }:𝖳(W^{}U)𝖬_{TR}𝖬_{RS}.$$
On $`𝖬_{TR}𝖬_{RS}`$, we have
$$T(NL)_1(NL)_2=TN_1N_2L_1L_2=N_1N_2RL_1L_2=N_1N_2L_1L_2S=(NL)_1(NL)_2S.$$
Thus, $`\overline{\theta }(TM_1M_2M_1M_2S)=0`$. Hence, it factorizes to a morphism
$$\mu ^{}:𝖬_{TS}𝖬_{TR}𝖬_{RS}.$$
Next, we define the coaction of the Hopf algebra $`𝖧_R`$. The Hopf algebra $`𝖧_R`$ is by definition the Hopf envelope of the bialgebra $`𝖤_R`$ , that is, there exists uniquely a bialgebra morphism $`i:𝖤_R𝖧_R`$ such that any bialgebra morphism $`f:𝖤_RH`$ to a Hopf algebra $`H`$ factorizes uniquely through as a composition of $`i`$ and a Hopf algebra morphism $`j:𝖧_RH`$, $`f=ji`$. The Hopf envelope of any bialgebras exists, hence we can define $`𝖧_R`$ for any Hecke operators $`R`$. If $`R`$ satisfies the Yang-Baxter equation (so for example, when $`R`$ is a Hecke operator) it is known that $`𝖤_R`$ is a coquasitriangular bialgebra . However, in general, the coquasitriangular structure on $`𝖤_R`$ cannot be extended on $`𝖧_R`$.
We shall assume that $`R`$ is a Hecke symmetry which means that the operator $`R^\mathrm{\#}:=(\text{ev}_V\text{id}_{VV^{}})(\text{id}_V^{}R\text{id}_V^{})(\text{id}_{V^{}V}\text{db}_V):V^{}VVV^{}`$ is invertible, this condition provides the coquasitriangular structure on $`𝖤_R`$ can be extended on $`𝖧_R`$, the antipode on $`𝖧_R`$ is bijective , and the map $`i`$ is injective \[9, Thm. 2.3.5\].
By means of the injective bialgebra morphism $`i`$, we identify $`𝖤_R`$ with a subbialgebra of $`𝖧_R`$. Each (simple) $`𝖤_R`$-comodule becomes then (simple) $`𝖧_R`$-comodule. Further, since the antipode of $`𝖧_R`$ is an anti-homomorphism of (co)algebras, a left (resp. right) $`𝖤_R`$-comodule becomes a right (left) $`𝖧_R`$ by composing the coaction with the antipode. In particular, since $`(V^n)^{}`$ is a left $`𝖤_R`$-comodule, it is a right $`𝖧_R`$-comodule. The coaction is explicitly given as follows. Let $`\{x_i\}_{i=1}^d`$ be a basis of $`V`$ and $`\{\xi ^i\}_{i=1}^d`$ be the dual basis for $`V^{}`$. Then $`\{e_i^j:=\xi ^jx_i\}_{i,j=1}^d`$ form a multiplicative matrix in $`𝖤_R`$ and the left coaction of $`𝖤_R`$ on $`V^{}`$ is $`\lambda _V^{}(\xi ^i)=_ke_k^i\xi ^k.`$ The right coaction of $`𝖧_R`$ on $`(V^n)^{}`$ is then
(49) $`\delta (\xi ^{i_1}\mathrm{}\xi ^{i_n})={\displaystyle \underset{k_1,\mathrm{},k_n}{}}\xi ^{k_1}\mathrm{}\xi ^{k_n}\mathrm{SS}(e_{k_1}^{j_1}\mathrm{}e_{k_n}^{j_n}),\text{ }\mathrm{SS}\text{ denotes the antipode}.`$
Since $`\mathrm{SS}`$ is injective, we also have
$$\text{End }^{𝖧_R}((V^n)^{})={}_{}{}^{𝖤_R}\text{End }((V^n)^{}).$$
In particular, $`({}_{}{}^{t}R)_w`$ are endomorphisms of right $`𝖧_R`$-comodules, for all $`w\text{S}_n`$. Therefore, $`𝖬_{RS}`$ is a right $`𝖧_R`$-subcomodule of $`𝖳(V^{}U)`$. Notice that $`𝖳(V^{}U)`$ and hence $`𝖬_{RS}`$, is not an $`𝖧_R`$-comodule algebra. The reason is that the usual isomorphism $`V^nV^mV^{m+n}`$ is not an $`𝖧_R`$-comodule morphism.
We are now ready to formulate a quantum analogue of the first and the second fundamental theorems for general linear groups.
###### Theorem 5.1.
Let $`S,T`$ be Hecke operators $`R`$ be a Hecke symmetry with the parameter $`q`$ not a root of unity of order greater that 1. Let $`\delta _{RST}`$ be the coaction of $`𝖧_R`$ on $`𝖬_{TR}𝖬_{RS}`$, which is the tensor product of the coactions of $`𝖧_R`$ on $`𝖬_{TR}`$ and $`𝖬_{RS}`$ given above. Then:
1. The set of coinvariants in $`𝖬_{TR}𝖬_{RS}`$ with respect to the coaction $`\delta _{RST}`$ is precisely $`\text{Im }\mu ^{}`$.
2. The kernel of $`\mu ^{}`$ in $`𝖬_{TS}`$ is the ideal $`I(((r+1)^{s+1}))`$ where $`(r,s)`$ is the birank of $`S`$.
Proof. From Section 2, we know that $`𝖬_{TS}`$ decomposes into a direct sum of $`𝖤_T^{\mathrm{cop}}𝖤_S`$-comodules and $`𝖬_{TR}𝖬_{RS}`$ decomposes into a direct sum of $`𝖤_T^{\mathrm{cop}}𝖤_R𝖤_R^{\mathrm{cop}}𝖤_S`$-comodules:
(52) $`\begin{array}{ccc}\hfill 𝖬_{TS}& & _\lambda W_\lambda ^{}U_\lambda ,\hfill \\ \hfill 𝖬_{TR}𝖬_{RS}& & _{\lambda \mu }W_\lambda ^{}V_\lambda V_\mu ^{}U_\mu .\hfill \end{array}`$
On the other hand, $`V_\lambda V_\mu ^{}`$ is a comodule over $`𝖧_R`$ and thus $`𝖬_{TR}𝖬_{RS}`$ can be considered as an $`𝖤_T^{\mathrm{cop}}𝖧_R𝖤_S`$-comodule. Consider the trivial coaction of $`𝖧_R`$ on $`𝖬_{TS}`$, i.e., consider $`𝖬_{TS}`$ as a direct sum of copies of 𝕂 which is $`𝖧_R`$-comodule by mean of the unit element $`\delta (1_\text{𝕂})=1_\text{𝕂}1_{𝖧_R}`$.
###### Lemma 5.2.
With the coaction of $`𝖧_R`$ described above, $`\mu ^{}`$ is a morphism of $`𝖤_T^{\mathrm{cop}}𝖧_R𝖤_S`$-comodules, the restriction of $`\mu ^{}`$ on $`W_\lambda ^{}U_\lambda `$ can be given by the map $`\text{db}_{V_\lambda }`$:
$$\mu ^{}|_{W_\lambda ^{}U_\lambda }=\text{id}\text{db}_{V_\lambda }\text{id}:W_\lambda ^{}U_\lambda W_\lambda ^{}V_\lambda V_\lambda ^{}U_\lambda .$$
Assume that Lemma is true. Since the map $`\text{db}_{V_\lambda }:\text{𝕂}V_\lambda V_\lambda ^{}`$ is injective unless $`V_\lambda =0`$, the restriction of $`\mu `$ on $`W_\lambda ^{}U_\lambda `$ in injective unless $`V_\lambda =0`$. This implies that $`\text{Ker }\mu ^{}`$ is the set $`{\displaystyle \underset{\lambda ,V_\lambda =0}{}}W_\lambda ^{}U_\lambda `$, which is precisely the set $`{\displaystyle \underset{\lambda ((r+1)^{s+1})}{}}W_\lambda ^{}U_\lambda =I(((r+1)^{s+1})),`$ by Theorem 2.6. The first claim of Theorem is proved.
Since $`V_\lambda `$ is simple over $`𝖧_R`$,
$$\text{Hom }_{𝖧_R}(\text{𝕂},V_\lambda V_\lambda ^{})=\text{Hom }_{𝖧_R}(V_\lambda ,V_\lambda )=\text{𝕂}.$$
Therefore, the image of $`\text{db}_{V_\lambda }`$ is the subspace of $`𝖧_R`$-invariants in $`V_\lambda V_\lambda ^{}`$. Consequenlty, the set of invariants in $`W_\lambda ^{}V_\lambda V_\lambda ^{}U_\lambda `$ is the image of $`W_\lambda ^{}U_\lambda `$. On the other hand, if $`\lambda \mu `$ then
$$\text{Hom }_{𝖧_R}(\text{𝕂},V_\lambda V_\mu ^{})=\text{Hom }_{𝖧_R}(V_\lambda ,V_\mu )=0.$$
Hence, for $`\lambda \mu `$, the comodule $`W_\lambda ^{}V_\lambda V_\mu ^{}U_\mu `$ does not contains non-zero invariants. Taking the direct sum for all $`\lambda `$ we prove the second claim of Theorem.
Proof of Lemma 5.2. Recall that the map $`\theta `$ is given in terms of the inclusion $`W^{}U𝖳(W^{}V)𝖳(V^{}U)`$, which in its order is given by the map $`\text{db}_V`$. The restriction of $`\theta `$ on $`(W^{}U)^n`$ is
$$\theta _n:W^nU^n(W^{}U)^n\stackrel{\theta _1^n}{}(W^{}VV^{}U)^nW^nV^nV^nU^n.$$
Hence
$$\theta _n=\text{id}_{V^n}\text{db}_{V^n}\text{id}_{U^n}:W^nU^nW^nV^nV^nU^n.$$
Thus, is easy to see that $`\theta `$ is a morphism of $`𝖤_T^{\mathrm{cop}}𝖤_S`$-comodules. On the other hand, with respect to the coaction of $`𝖧_R`$ on $`(V^n)^{}`$ given in (49), $`\text{db}_{V^n}`$ is an $`𝖧_R`$-comodules morphism. Therefore $`\theta _n`$ is a morphism of $`𝖤_T^{\mathrm{cop}}𝖧_R𝖤_S`$-comodules.
The restriction of $`\theta _n`$ on $`W_\lambda ^{}U_\mu `$ is then the map
$$W_\lambda ^{}U_\mu W_\lambda ^{}V^nV^nU_\mu \underset{\gamma }{}W_\lambda ^{}V_\gamma V_\gamma ^{}U_\mu .$$
The map $`\mu ^{}`$ is obtained by passing to quotients. According to Theorem 2.4, we can identify $`𝖬_{TS}`$ with a subspace of $`𝖳(W^{}U)`$ such that the quotient map $`𝖳(W^{}U)𝖬_{TS}`$ is given by the sum $`\mathrm{\Psi }`$ of the projectors $`\mathrm{\Psi }_n`$. Thus $`\mu ^{}`$ can be considered as the composition $`(\mathrm{\Psi }_{TR}\mathrm{\Psi }_{RS})\theta \mathrm{\Psi }_{TS}`$. Therefore, the restriction of $`\mu ^{}`$ on $`W_\lambda ^{}U_\lambda `$ is the map
$$\text{id}\text{db}_{V_\lambda }\text{id}:W_\lambda ^{}U_\lambda W_\lambda ^{}V_\lambda V_\lambda ^{}U_\lambda .$$
Since the map $`\text{db}_{V_\lambda }`$ is a morphism of $`𝖧_R`$-comodules, the above map is a morphism of $`𝖤_T^{\mathrm{cop}}𝖧_R𝖤_S`$-comodules. Lemma 5.2 is therefore proved.
### 5.2. Fundamental theorems for quantum groups of type $`A`$, the second version
Since $`R`$ is a Hecke symmetry, $`𝖧_R`$ is a coquasitriangular Hopf algebra, i.e., the category of $`𝖧_R`$-comodules is braided (see, e.g., \[14, Chapter XI\]). We can modify the struture above so that the morphism $`\mu ^{}`$ is an $`𝖧_R`$-comodule algebra morphism. To do it, we identify the two vector spaces in (48) by means of an $`𝖧_R`$-comodule isomorphism $`\tau _{V^{}V}:V^{}VVV^{}`$, $`\xi ^ix_jx_k\xi ^lR_{}^{1}{}_{jl}{}^{ik}`$. More precisely, in defining an isomorphism from $`(W^{}VV^{}U)^n`$ to $`(W^{}V)^n(V^{}U)^n`$, whenever we have to interchange $`V^{}`$ and $`V`$ we shall use the comodule isomorphism $`\omega _{V^{}V}`$ above. We therefore modify the algebra structure on $`𝖳(W^{}V)𝖳(V^{}U)`$ replacing the commuting relation of $`N`$ and $`P`$ by the following relation
(53) $`L_1N_2=N_2R^1PL_1,`$
where $`P`$ is the matrix of the usual flip operator $`xyyx`$, with respect to the basis $`x_1,x_2,\mathrm{},x_d`$, $`P_{kl}^{ij}=\delta _l^i\delta _k^j`$. Let $`\omega _n`$ denote the isomorphism from $`(VV^{})^nV^nV^n`$ obtained by using $`\omega _{V^{}V}`$, so, for example, $`\omega _1=\text{id}`$, $`\omega _2=\text{id}\omega _{V^{}V}\text{id}`$.
Since the algebra structure on $`𝖳(W^{}V)𝖳(V^{}U)`$ is modified, the map $`\theta `$ should also be modified. The restriction of $`\theta `$ on $`(W^{}U)^n`$ is now
$`\theta _n:(W^{}U)^n\stackrel{\theta _1^n}{}(W^{}VV^{}U)^n\stackrel{\omega _n}{}W^nV^nV^nU^n,`$
If we identify $`(W^{}U)^n`$ with $`W^nU^n`$, then we can consider $`\theta _n`$ as a morphism
$`\theta _n=\text{id}_{V^n}\omega _n\text{db}_{V}^{}{}_{}{}^{n}\text{id}_{U^n}:W^nU^nW^nV^nV^nU^n.`$
Thus, if we define the coaction of $`𝖧_R`$ on $`V^n`$ in such a way that this comodule is isomorphic to $`(V^{})^n`$ as $`𝖧_R`$-comodules then $`\omega _n\text{db}_{V}^{}{}_{}{}^{n}`$ is a comodule morphism and hence so is $`\theta _n`$. Explicitly, the coaction is given by
$$\delta (\xi ^{i_1}\xi ^{i_2}\mathrm{}\xi ^{i_n})=\underset{k_1,k_2,\mathrm{},k_n}{}\xi ^{k_1}\xi ^{k_2}\mathrm{}\xi ^{k_n}\mathrm{SS}(e_{k_1}^{i_1})\mathrm{SS}(e_{k_2}^{i_2})\mathrm{}\mathrm{SS}(e_{k_n}^{i_n}).$$
Note that with respect to these coactions, the operators $`\widehat{R}_w,w\text{S}_n`$, where $`\widehat{R}=PRP`$, are comodules endomorphisms of $`V^n`$. Hence the quotient $`𝖬_{\widehat{R}S}`$ of $`𝖳(V^{}U)`$ is a comodule over $`𝖧_R`$. Further, since the usual identification $`(V^{})^m(V^{})^n(V^{})^{m+n}`$ is a morphism of $`𝖧_R`$-comodules, $`𝖬_{\widehat{R}S}`$ is an $`𝖧_R`$-comodule algebra.
On the other hand, we can check that $`\theta `$ factorizes to an algebra morphism
$$\mu _m^{}:𝖬_{ST}𝖬_{SR}_𝗆𝖬_{\widehat{R}T}$$
where in $`𝖬_{SR}_𝗆𝖬_{\widehat{R}T}`$, $`N`$ and $`L`$ commute by the rule in (53). Indeed, on $`𝖬_{SR}_𝗆𝖬_{\widehat{R}T}`$ we have
$`T(NL)_1(NL)_2`$ $`=`$ $`TN_1N_2R^1PL_1L_2=N_1N_2RR^1PL_1L_2=N_1N_2R^1P\widehat{R}L_1L_2`$
$`=`$ $`N_1N_2R^1PL_1L_2S=(NL)_1(NL)_2S.`$
Thus, we obtain a coaction of $`𝖧_R`$ on $`𝖬_{SR}_𝗆𝖬_{\widehat{R}T}`$, for which the morphism $`\mu _𝗆^{}`$ is an $`𝖤_T^{\mathrm{cop}}𝖧_R𝖤_S`$-comodule morphism.
Although we have modified the algebra structure on $`𝖬_{SR}_𝗆𝖬_{\widehat{R}T}`$, this does not, affect the decomposition (52). That is, we still have an isomorphism of $`𝖤_T^{\mathrm{cop}}𝖧_R𝖤_S`$-comodules
$$𝖬_{TR}_𝗆𝖬_{RS}\underset{\lambda \mu }{}W_\lambda ^{}V_\lambda V_\mu ^{}U_\mu .$$
Therefore an anlogue of Lemma 5.2 can be easily obtained, whence one gets an analogue of Theorem 5.1
###### Theorem 5.3.
Let $`S,T`$ be Hecke operators $`R`$ be a Hecke symmetry. Then $`\mu _𝗆^{}`$ is an $`𝖧_R`$-comodule algebra morphism:
1. The set of coinvariants in $`𝖬_{TR}_𝗆𝖬_{RS}`$ with respect to the coaction $`\delta _{RTS}`$ is precisely $`\text{Im }\mu _𝗆^{}`$.
2. The kernel of $`\mu _𝗆^{}`$ in $`𝖬_{TS}`$ is the ideal $`I(((r+1)^{s+1}))`$ where $`(r,s)`$ is the birank of $`T`$.
Proof. It remains to check that $`𝖬_{SR}_𝗆𝖬_{\widehat{R}T}`$ is an $`𝖧_R`$-comodule algebra. In fact, what we have done above is to define an algebra struture on the tensor product of two $`𝖧_R`$-comodule algebras $`𝖬_{TR}`$ and $`𝖬_{\widehat{R}S}`$ (see, e.g., ). Namely, the isomorphism $`\omega _{V^{},V}`$ gives rise to an $`𝖧_R`$-comodule isomorphism $`𝖬_{\widehat{R}S}𝖬_{TR}𝖬_{TR}𝖬_{\widehat{R}S}`$. It follows from the standard argument that $`𝖬_{TR}_𝗆𝖬_{\widehat{R}S}`$ is an $`𝖧_R`$-comodule algebra.
## 6. Example: Standard quantum general linear groups
The notion of quantum determinantal ideals presented here does not seem to have relationship with any quantum determinants. In fact, for an arbitrary Hecke operator, we aren’t able to define any quantum determinant. However, in the case of standard $`R`$-matrix (see Subsection 1.2), the quantum minors are definable, and our notion of quantum determinantal ideals can be given in terms of these quantum minors. The quantum determinantal ideals associated to standard quantum matrix of type $`A`$ were studied by Goodearl, Lenagan and Rigal in , where the primeness was particularly proved. Here we show that in the case of standard $`R`$-matrix, our notion of quantum determinant ideal coincides with the notion give by Goodearl and Lenagan.
Recall that the standard $`R`$-matrix of type $`A_{n1}`$ has, with respect to a certain basis $`x_1,x_2,\mathrm{},x_n`$, the following form:
$`R_{n}^{}{}_{ij}{}^{kl}={\displaystyle \frac{q^2q^{2\epsilon _{ij}}}{1+q^{2\epsilon _{ij}}}}\delta _{ij}^{kl}+{\displaystyle \frac{q^{\epsilon _{ij}}(q^2+1)}{1+q^{2\epsilon _{ij}}}}\delta _{ij}^{lk},1i,j,k,ln,\epsilon _{ij}:=\text{sign }(ji)`$
The Hecke equation for $`R_n`$ is $`(R_nq^2)(R_n+1)=0`$. The quantum exterior algebra $`_{R_n}`$ is the factor algebra of the non-commutative algebra $`\text{𝕂}x_1,x_2,\mathrm{},x_n`$ by the relations $`x_ix_j=qx_jx_i`$ for $`ij`$, it can be realized as subalgebra of $`\text{𝕂}x_1,x_2,\mathrm{},x_n`$ spanned by $`q`$-symmetrized tensors
$$x_{i_1}x_{i_2}\mathrm{}x_{i_k}:=\underset{\sigma \text{S}_k}{}(q)^{l(\sigma )}x_{i_1\sigma }x_{i_2\sigma }\mathrm{}x_{i_k\sigma },$$
for any sequence $`(i_1<i_2<\mathrm{}<i_k)`$ of elements from $`\{1,2,\mathrm{},n\}`$, for $`k=1,2,\mathrm{},n.`$ In particular, $`_{R_n}^{}{}_{k}{}^{}`$ is spanned by $`x_{i_1}x_{i_2}\mathrm{}x_{i_k}`$.
Analogously, assume that $`_{R_m}^{}{}_{k}{}^{}`$ has a basis consisting of $`q`$-antisymmetrized tensors $`y_{j_1}y_{j_2}\mathrm{}y_{j_k}`$, for any sequence $`(j_1<j_2<\mathrm{}<j_k)`$ of elements from $`\{1,2,\mathrm{},m\}`$, for $`k=1,2,\mathrm{},m`$. Then the space $`(_{R_m}^{}{}_{k}{}^{})^{}=\text{Im }\tau _k(E_{(1^k)})`$ is canonically spanned by the set
$$\xi ^{i_1}\xi ^{i_2}\mathrm{}\xi ^{i_k}:=\underset{\sigma \text{S}_k}{}(q)^{l(\sigma )}\xi ^{i_1\sigma }\xi ^{i_2\sigma }\mathrm{}\xi ^{i_k\sigma },$$
where $`\xi ^1,\xi ^2,\mathrm{},xi^m`$ is the dual basis to $`y_1,y_2,\mathrm{},y_m`$.
Denote $`e_i^j:=\xi ^jx_i`$. Then the algebra $`𝖬_{R_mR_n}`$ can be considered as a subspace of $`\text{𝕂}e_1^1,\mathrm{},e_n^m`$ spanned by the elements
$`e_{i_1i_2\mathrm{}i_k}^{j_1j_2\mathrm{}j_k}`$ $`:=`$ $`{\displaystyle \underset{\sigma ,\tau \text{S}_k}{}}{\displaystyle \frac{(q)^{l(\sigma )}}{(q)^{l(\sigma )}}}e_{i_1\sigma }^{j_1\tau }e_{i_2\sigma }^{j_2\tau }\mathrm{}e_{i_k\sigma }^{j_k\tau }`$
$`=`$ $`k!{\displaystyle \underset{\sigma \text{S}_k}{}}(q)^{l(\sigma )}e_{i_1\sigma }^{j_1}e_{i_2\sigma }^{j_2}\mathrm{}e_{i_k\sigma }^{j_k}`$
The element $`_{\sigma \text{S}_k}(q)^{l(\sigma )}e_{i_1\sigma }^{j_1}e_{i_2\sigma }^{j_2}\mathrm{}e_{i_k\sigma }^{j_k}`$ is precisely the quantum determinant of the submatrix of $`Z=(e_i^j)`$ formed on rows $`i_1,i_2,\mathrm{},i_k`$ and columns $`j_1,j_2,\mathrm{},j_k`$.
The fundamental theorems for standard quantum groups were proved by Goodearl, Lenagan and Rigal in \[5, Theorem 2.5\] and \[6, Theorem 4.5\]. We would like to mention that the setting of these theorems is slightly different from our setting here, namely, in defining the coaction of the quantum group and the algebra structure. In the language of our paper, the algebra structure of $`𝖬_{TR}𝖬_{RS}`$ considered in \[loc.cit\] is the ordinary algebra structure, i.e., as in the setting of Subsection 5.1. The coaction of $`𝖧_R`$ on $`𝖬_{RS}`$ considered in \[loc.cit\] corresponds however to the coaction given in Subsection 5.2. One can do that because for standard deformation we have $`\widehat{R}_n:=PR_nP=R_n`$.
## 7. Example: The standard quantum general linear supergroups
There are two main differences between the coalgebras $`𝒪(GL(n))`$ and $`𝒪(GL(m|n))`$. Firstly, as a coalgebra, $`𝒪(GL(n))`$ is cosemisimple while $`𝒪(GL(m|n))`$ is not. Secondly, the determinant is a polynomial function while the super determinant is not. Looking at the more general construction of quantum groups of type $`A`$, we see that the quantum determinant is determined in terms of the quantum exterior algebra, which should be finite dimensional. Thus, if a Hecke opertor $`R`$ produces a finite dimensional quantum exterior algebra, we call the corresponding bialgebra $`𝖤_R`$ a quantum semi-group of type $`A_n`$, where $`n`$ is the of the Poincaré series of the quantum exterior algebra. In this case, $`R`$ is called even Hecke operator. It may happen that the quantum symmetric algebra has finite dimension, in this case $`R`$ is called odd Hecke operator . A typical of non-even and non-odd Hecke operator is the flip operators in super geometry. Hence it is natural to suggest non-even non-odd Hecke operators define analogies of the general linear supergroups. It this then also natural to consider them as operators in the category of vector superspaces. It turns out, however, that the basis category does not play a great role. That is, it doesn’t matter whether or not we consider $`R`$ as an operator in the category of vector superspaces and define $`𝖤_R`$ as a superbialgebra, many properties of $`𝖤_R`$ remain unchanged. In other words, many properties of $`𝖤_R`$ depend only on the intrinsic properties of $`R`$.
For example, let $`R_d`$ be the standard matrix considered in the previous section. If we assume that some of the basis vectors $`x_1,x_2,\mathrm{},x_d`$ have odd parity and the other have even parity, thus, $`V`$ is a vector superspace, then $`R`$ is an operator in the category of vector superspaces. In this category, the associated bialgebra $`𝖤_R`$ is defined differently but it remains cosemisimple. On the other hand, if we take a flip operator on a strict vector superspace and consider it as an operator in the category of (non-super) vector spaces, the associated bialgebra $`𝖤_R`$ remains non-cosemisimple.
All results in this paper hold in the category of vector superspaces. In fact, all what we have to do in the category of vector superspaces is to replace the ordinary flip operator by it super counterpart, i.e. to insert signs at some places.
Let now $`R_{r|s}`$ denote the super analogue of the standard $`R`$-matrices of type $`A`$. Explicitly, with respect to some homogeneous basis $`x_1,x_2,\mathrm{},x_d`$, $`d=r+s`$, where $`\widehat{x}_i:=\overline{0}`$ if $`ir`$ and $`\overline{1}`$ if $`ir+1`$, the operator $`R_{r|s}`$ has the following form:
$`R_{r|s}^{}{}_{ij}{}^{kl}={\displaystyle \frac{q^2q^{2\epsilon _{ij}}}{1+q^{2\epsilon _{ij}}}}\delta _{ij}^{kl}+\widehat{i}\widehat{j}{\displaystyle \frac{q^{\epsilon _{ij}}(q^2+1)}{1+q^{2\epsilon _{ij}}}}\delta _{ij}^{lk},\mathrm{\hspace{0.33em}1}i,j,k,ln=r+s,\epsilon _{ij}:=\text{sign }(ji),`$
where $`\widehat{i}`$ denotes the parity of $`x_i`$. Then $`R_{r|s}`$ is a Hecke symmetry of birank $`(r,s)`$. The associated super bialgebra $`𝖤_R`$ and Hopf algebra $`𝖧_R`$ are called the function aglebras on the standard quantum super semi-group $`M_q(r|s)`$ and the standard quantum supergroup $`GL_q(r|s)`$, respectively, see, e.g. .
Theorem 5.1 applied to this case gives us the fundamental theorems for standard quantum supergroups.
###### Theorem 7.1.
Let $`M_q(m|n,r|s)`$ denote the super bialgebra $`𝖬_{R_{r|n}R_{r|s}}`$ and let $`GL_q(m|n)`$ denote the Hopf superaglebra $`𝖧_{R_{m|n}}`$. Let $`\mu ^{}`$ be the algebra morphism $`M_q(m|n,u|v)M_q(m|n,r|s)M_q(r|s,u|v)`$ induced from the map
$$e_j^i\underset{k=1}{\overset{r+s}{}}e_k^ie_j^k$$
where $`\{e_j^i\}_{i=1,j=1}^{m+nu+v}`$ is the standard generators of $`M_q(m|n,u|v)`$, similary, $`n_k^i`$ and $`p_j^k`$ are generators for $`M_q(m|n,r|s)`$ and $`M_q(r|s,u|v)`$. Assume that $`q`$ is not a root ou unity of order greater than 1. Then,
1. the set of coinvariants of $`M_q(m|n,r|s)M_q(r|s,u|v)`$ with respect ot the coaction of $`GL_q(r|s)`$ is precisely $`\text{Im }\mu ^{}`$,
2. the kernel of $`\mu ^{}`$ in $`M_q(m|n,u|v)`$ is the ideal $`I(((s+1)^{r+1}))=_{\sigma _{u+1}v+1}I_\sigma `$.
Setting $`q=1`$ in the above theorem, we obtain the fundamental theorems for general linear super groups. These theorems can be formulated in the classical way. Let $`\mu `$ denote the multiplication of supermatrices
$$\mu :M(m|n,r|s)\times M(r|s,u|v)M(m|n,r|s).$$
Then $`\mu `$ induces a morphism $`\mu ^{}`$
$$𝒪(M(m|n,r|s))𝒪(M(m|n,r|s)\times M(r|s,u|v))𝒪(M(m|n,r|s))𝒪(M(r|s,u|v)).$$
Consider the natural action of $`GL(r|s)`$ on $`M(m|n,r|s)\times M(r|s,u|v)`$: $`(A,B)(Ag^1,gB)`$, which induces a natural coaction of $`GL(r|s)`$ on the function algebra on $`M(m|n,r|s)\times M(r|s,u|v)`$.
###### Theorem 7.2.
We have:
1. A polynomial in $`𝒪(M(m|n,r|s)\times M(r|s,u|v))`$, invariant with the action of $`GL(r|s)`$, can be obtained by composing a polynomial in $`𝒪(M(m|n,u|v))`$ with $`\mu `$.
2. The kernel of $`\mu ^{}`$ is the ideal $`I(((s+1)^{r+1}))`$.
Remark. Except for the case $`n=s=v=0`$, the ideal $`I(((s+1)^{r+1}))`$ is not a determinantal ideal. It is an interesting problem to describe this ideal more explicitly.
Acknowledgment
The author should like to thank Professor C. Procesi for explaining him invariant theory. He also should like to thank Professor K. Goodearl for useful discussions.
The first part of the work was done at the Max-Planck Intitut für Mathematik, Bonn and appeared as Preprint MPI-99/12. The work was completed during the author’s stay at the Mathematical Sciences Research Institute, Berkeley. The author would like to thank these Institutes for the excellent working condition and financial support.
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# Lie Algebras Associated to Fiber-Type Arrangements
## 1. Introduction
Two classical constructions of interest in group theory and topology are:
1. The Lie algebra arising from the filtration quotients associated to the descending central series of a discrete group $`G`$; and
2. The Lie algebra of primitive elements in the singular homology of the loop space of a space $`X`$, for certain topological spaces $`X`$.
The purpose of this article is to illustrate that these two a priori unrelated Lie algebras are, in fact, isomorphic in certain natural cases. This work is motivated by recent results relating the Lie algebras of (i) and (ii) arising in the context of classical configuration spaces, and resolves a conjecture of the second two authors concerning the generalization of these results to spaces arising from certain hyperplane arrangements.
The main result here asserts that the Lie algebra associated to the fundamental group $`G`$ of the complement of a fiber-type hyperplane arrangement is, up to regrading, isomorphic to the Lie algebra of primitive elements in the homology of the loop space of the complement of a higher dimensional analogue of the arrangement. The main theorem is, in fact, stronger. The Samelson product for the loop space gives rise to a graded Lie algebra given by the homotopy groups modulo torsion. This Lie algebra is, again up to regrading, also isomorphic to the Lie algebra associated to the descending central series quotients of $`G`$. In addition, after looping further, there are natural related Poisson algebras arising from the homology of associated iterated loop spaces.
Given a discrete group $`G`$, let $`G_n`$ be the $`n`$-th stage of the descending central series, defined inductively by $`G_1=G`$, and $`G_{n+1}=[G_n,G]`$ for $`n1`$, and let $`E_0^n(G)=G_n/G_{n+1}`$ be the $`n`$-th associated quotient. Let $`E_0^{}(G)=_{n1}E_0^n(G)`$ be the Lie algebra obtained from the descending central series of $`G`$, with Lie algebra structure induced by the commutator map $`G\times GG`$, $`(x,y)xyx^1y^1`$. For each positive integer $`k`$, use the ungraded Lie algebra $`E_0^{}(G)`$ to define a related graded Lie algebra as follows.
###### Definition 1.1.
For a group $`G`$, let $`E_0^{}(G)_k`$ be the graded Lie algebra given by
$$E_0^q(G)_k=\{\begin{array}{cc}E_0^n(G)\hfill & \text{if }q=2nk\text{,}\hfill \\ 0\hfill & \text{otherwise,}\hfill \end{array}$$
with Lie bracket structure induced by that of the Lie algebra $`E_0^{}(G)`$ obtained from the descending central series of $`G`$ in the obvious manner.
A theorem relating the Lie algebras of (i) and (ii) above is described next. Let $`P_n`$ be the Artin pure braid group, the fundamental group of the configuration space $`F(,n)`$. The results on configuration spaces alluded to above, due to Fadell and Husseini and Cohen and Gitler , may be summarized as follows.
###### Theorem 1.2.
For $`k1`$, the homology of the loop space of the configuration space $`F(^{k+1},n)`$ is isomorphic to the universal enveloping algebra of the graded Lie algebra $`E_0^{}(P_n)_k`$. Moreover,
1. The image of the Hurewicz homomorphism
$$\pi _{}(\mathrm{\Omega }F(^{k+1},n))H_{}(\mathrm{\Omega }F(^{k+1},n);)$$
is isomorphic to $`E_0^{}(P_n)_k`$; and
2. The Hurewicz homomorphism induces isomorphisms of graded Lie algebras
$$\pi _{}(\mathrm{\Omega }F(^{k+1},n))/\mathrm{Torsion}\mathrm{Prim}H_{}(\mathrm{\Omega }F(^{k+1},n);)E_0^{}(P_n)_k,$$
where $`\mathrm{Prim}`$ denotes the module of primitive elements, and the Lie algebra structure of the source is induced by the classical Samelson product.
The Lie algebra arising in the above theorem is the “universal Yang-Baxter Lie algebra” $`(n)`$, the quotient of the free Lie algebra on a free abelian group of rank $`\left(\genfrac{}{}{0pt}{}{n}{2}\right)`$ by relations recorded in (4.1) below, and known variously as the “infinitesimal pure braid relations” or the “horizontal four-term relations and framing independence.” Furthermore, the homology of an iterated loop space of configuration space, $`\mathrm{\Omega }^qF(^{k+1},n)`$ for $`q>1`$, admits the structure of a graded Poisson algebra, see . The associated relations are called the “universal infinitesimal Poisson braid relations.” This Poisson algebra structure on $`H_{}(\mathrm{\Omega }^qF(^{k+1},n))`$ has recently been used in the context of algebraic groups by Lehrer and Segal .
The universal Yang-Baxter Lie algebra, and infinitesimal pure braid relations, arise in a number of contexts. These include the classification of pure braids by Vassiliev invariants, see Kohno , and the Knizhnik-Zamolodchikov differential equations from conformal field theory, where the relations appear as integrability conditions on the associated Gauss-Manin connection, see Varchenko . Moreover, any finite dimensional representation of the Lie algebra $`(n)`$ induces a representation of the pure braid group $`P_n`$ on the same vector space, see Kapovich and Millson .
An important ingredient in the proof of Theorem 1.2 is a classical result of Fadell and Neuwirth which shows that configuration spaces admit iterated bundle structure. Similar results are known to hold for certain orbit configuration spaces , which admit analogous bundle structure, and are described in more detail below. All of these spaces fit in the following general framework.
For each natural number $`\mathrm{}`$, let $`X_{\mathrm{}}`$ be a functor from Euclidean spaces, with morphisms restricted to endomorphisms, to topological spaces. For a Euclidean space $`𝔼`$, let $`𝒬_{\mathrm{}}(𝔼)`$ be a discrete subset of $`𝔼`$ of fixed (possible infinite) cardinality depending on $`\mathrm{}`$. Assume that there are natural transformations $`X_{\mathrm{}}(𝔼)X_\mathrm{}1(𝔼)`$ which satisfy the following conditions.
1. The space $`X_1(𝔼)=𝔼𝒬_1(𝔼)`$ is the complement of a discrete subset of $`𝔼`$.
2. The map $`X_{\mathrm{}}(𝔼)X_\mathrm{}1(𝔼)`$ is a fiber bundle projection, with fiber $`𝔼𝒬_{\mathrm{}}(𝔼)`$.
3. Each bundle $`X_{\mathrm{}}(𝔼)X_\mathrm{}1(𝔼)`$ admits a cross-section.
4. If $`𝔼`$, the fundamental group of $`X_\mathrm{}1(𝔼)`$ acts trivially on the homology of the fiber $`𝔼𝒬_{\mathrm{}}(𝔼)`$.
The prototypical examples are given by the configuration spaces $`X_{\mathrm{}}(𝔼)=F(𝔼,\mathrm{})`$, where $`𝔼=^k`$. Further examples are given below.
It seems likely that for many choices of $`X_{\mathrm{}}`$, the Lie algebras associated to $`X_{\mathrm{}}(𝔼)`$ as $`𝔼`$ varies are related in a manner analogous to those arising in Theorem 1.2. If $`𝔼`$, conditions (1) and (2) imply that $`X_{\mathrm{}}(𝔼)`$ is a $`K(G,1)`$ space, where $`G=\pi _1(X_{\mathrm{}}(𝔼))`$ is the fundamental group of $`X_{\mathrm{}}(𝔼)`$, as is readily seen from the homotopy sequence of a bundle. In this case, condition (3) further implies that the group $`G`$ admits the structure of an iterated semidirect product of free groups, and condition (4) restricts the type of free group automorphisms arising in this structure. These conditions determine the additive structure of the Lie algebra $`E_0^{}(G)`$, see and Section 4. For higher dimensional $`𝔼`$, conditions (1)–(3) imply that the homology of the loop space of $`X_{\mathrm{}}(𝔼)`$ is isomorphic to the universal enveloping algebra of the Lie algebra $`\pi _{}(\mathrm{\Omega }X_{\mathrm{}}(𝔼))/\mathrm{Torsion}`$, and determine the additive structure of $`\mathrm{Prim}H_{}(\mathrm{\Omega }X_{\mathrm{}}(𝔼);)`$, see and Section 5. For higher dimensional $`𝔼`$, these conditions have analogous implications for the homology of an iterated loop space $`\mathrm{\Omega }^qX_{\mathrm{}}(𝔼)`$ with $`q>1`$, and the Poisson algebra structure admitted by this homology, see and Section 6.
A brief indication how one may analyze and compare the Lie algebras arising for various choices of $`𝔼`$ is given next. First, there is a variant of the classical Freudenthal suspension, relating reduced suspensions and loop spaces as indicated below, where the maps are induced by (homology) suspensions.
$$\begin{array}{ccc}H_{2k2}(\mathrm{\Omega }X_{\mathrm{}}(^k))& & H_{2k1}(X_{\mathrm{}}(^k))\\ & & & & \\ & & H_{2k}(\mathrm{\Sigma }X_{\mathrm{}}(^k))& & H_{2k}(\mathrm{\Omega }X_{\mathrm{}}(^{k+1}))\end{array}$$
If $`k2`$, conditions (1)–(3) above imply that these maps are all (additive) isomorphisms. In the case $`k=1`$, these maps yield an additive isomorphism $`E_0^1(G)=H_1(X_{\mathrm{}}())H_2(\mathrm{\Omega }X_{\mathrm{}}(^2))`$ where $`G=\pi _1(X_{\mathrm{}}())`$. While this comparison does not in general preserve the structures of these Lie algebras, it does provide a geometric way to compare indecomposable elements in these Lie algebras.
To determine the Lie algebra structure, let $`S`$ be a sphere of appropriate dimension and $`A:SX_{\mathrm{}}(𝔼)`$ a map representing a (reduced) homology generator in minimal degree. Consider the pullback $`\xi (𝔼)`$ of the bundle $`X_{\mathrm{}+1}(𝔼)X_{\mathrm{}}(𝔼)`$ along the map $`A`$:
$$\begin{array}{ccccc}𝔼𝒬_{\mathrm{}}(𝔼)& & \xi (𝔼)& & S\\ & & & & A& & \\ 𝔼𝒬_{\mathrm{}}(𝔼)& & X_{\mathrm{}+1}(𝔼)& & X_{\mathrm{}}(𝔼)\end{array}$$
These bundles admit compatible cross-sections by condition (3). There is consequently a morphism of extensions of Lie algebras
$$\begin{array}{ccccccccc}0& & (𝔼𝒬_{\mathrm{}}(𝔼))& & (\xi (𝔼))& & (S)& & 0\\ & & \mathrm{id}& & & & A_{}& & \\ 0& & (𝔼𝒬_{\mathrm{}}(𝔼))& & (X_{\mathrm{}+1}(𝔼))& & (X_{\mathrm{}}(𝔼))& & 0\end{array}$$
where $`()`$ denotes the Lie algebra obtained from the descending central series of the fundamental group if $`𝔼`$, and the graded Lie algebra of primitive elements in the homology of the loop space for higher dimensional $`𝔼`$. Knowledge of the extension $`0(𝔼𝒬_{\mathrm{}}(𝔼))(\xi (𝔼))(S)0`$ and the map $`A_{}:(S)(X_{\mathrm{}}(𝔼))`$ for all homology generators completely determines the structure of the Lie algebra $`(X_{\mathrm{}+1}(𝔼))`$. In favorable situations, one can show that the extensions of Lie algebras which arise as $`𝔼`$ varies are, apart from grading, isomorphic by carefully combining these considerations with the aforementioned comparison of indecomposables.
Several natural families of examples which fit in the framework described above are given next. These examples either may be or have been studied using (variants of) the techniques sketched above. Let $`M`$ be a manifold, and $`\mathrm{\Gamma }`$ a group which acts properly discontinuously on $`M`$. The orbit configuration space $`F_\mathrm{\Gamma }(M,\mathrm{})`$ consists of all $`\mathrm{}`$-tuples of points in $`M`$, no two of which lie in the same $`\mathrm{\Gamma }`$-orbit.
First, consider orbit configuration spaces of the form $`F_\mathrm{\Gamma }(𝔼\times ^n,\mathrm{})`$, where $`\mathrm{\Gamma }`$ operates diagonally of $`𝔼\times ^n`$, and trivially on $`^n`$. Relevant examples include the following.
1. A parameterized lattice $`\mathrm{\Gamma }`$ acting on $`𝔼=`$, so that the orbit space is an elliptic curve. The orbit configuration spaces associated to the action of the standard integral lattice were the subject of , where it is shown that the analogue of Theorem 1.2 holds for these spaces.
2. A discrete group $`\mathrm{\Gamma }`$ acting properly discontinuously on the upper half-plane $`𝔼=`$, so that the orbit space is a complex curve.
3. A torsion free subgroup of $`\mathrm{\Gamma }<\mathrm{Sp}(2g,)`$ acting properly discontinuously on Siegel upper half-space $`𝔼=^g`$.
4. A torsion free subgroup $`\mathrm{\Gamma }`$ of the mapping class group for genus $`g`$ surfaces, acting on Teichmuller space $`𝔼`$.
Second, let $`M=^k\{0\}`$ and let $`\mathrm{\Gamma }=/p`$ act freely on $`M`$ by rotations. The orbit configuration spaces $`F_\mathrm{\Gamma }(M,\mathrm{})`$ were the subject of and , the results of which combine to show that the analogue of Theorem 1.2 also holds for these spaces.
In the instances where conditions (1)–(4) hold, one obtains generalizations of the universal Yang-Baxter Lie algebra, parameterized by the group $`\mathrm{\Gamma }`$. This is the case for the family $`F_\mathrm{\Gamma }(M,\mathrm{})`$ of orbit configuration spaces above, where $`M=^k\{0\}`$ and $`\mathrm{\Gamma }=/p`$. As noted by D. Matei, the resulting generalized Yang-Baxter Lie algebra with cyclic symmetry is of use in constructing Vassiliev invariants of links in the lens space $`L(p,1)`$. The Lie algebras arising from other families of orbit configuration spaces may be of similar use for other three-manifolds, among other potential applications. The orbit configuration spaces $`F_{/p}(^k\{0\},\mathrm{})`$, and the classical configuration spaces $`F(^k,\mathrm{})`$, may be realized as complements of finite hyperplane or subspace arrangements. This led to speculation in that similar results may hold for fiber-type arrangements whose complements, like configuration spaces, admit iterated bundle structure.
Let $`𝒜`$ be a hyperplane arrangement in $`^{\mathrm{}}`$, a finite collection of codimension one affine subspaces, with complement $`M(𝒜)=^{\mathrm{}}_{H𝒜}H`$. See Orlik and Terao as a general reference on arrangements. Given a hyperplane $`H^{\mathrm{}}`$, let $`H^k`$ be the codimension $`k`$ affine subspace of $`^k\mathrm{}=(^{\mathrm{}})^k`$ consisting of all $`k`$-tuples of points in $`^{\mathrm{}}`$, each of which lies in $`H`$. For each positive integer $`k`$, the elements of the hyperplane arrangement $`𝒜`$ may be used in this way to obtain an arrangement $`𝒜^k`$ of complex codimension $`k`$ subspaces in $`^k\mathrm{}`$, with complement $`M(𝒜^k)=^k\mathrm{}_{H𝒜}H^k`$.
When $`𝒜`$ is a fiber-type hyperplane arrangement, the behavior of the family of spaces $`\{X_{\mathrm{}}(^k)=M(𝒜^k),k1\}`$ is reminiscent of that of the family $`\{F(^k,n),k1\}`$ of configuration spaces. Let $`G=\pi _1(M(𝒜))`$ be the fundamental group of the complement of the fiber-type arrangement $`𝒜`$ in $`^{\mathrm{}}`$, and let $`E_0^{}(G)`$ be the Lie algebra obtained from the descending central series of $`G`$. The main result of this article is as follows.
###### Theorem 1.3.
For $`k1`$, the homology of $`\mathrm{\Omega }M(𝒜^{k+1})`$, the loop space of the complement of the subspace arrangement $`𝒜^{k+1}`$ in $`^{(k+1)\mathrm{}}`$, is isomorphic to the universal enveloping algebra of the graded Lie algebra $`E_0^{}(G)_k`$. Moreover,
1. The image of the Hurewicz homomorphism
$$\pi _{}(\mathrm{\Omega }M(𝒜^{k+1}))H_{}(\mathrm{\Omega }M(𝒜^{k+1});)$$
is isomorphic to $`E_0^{}(G)_k`$; and
2. The Hurewicz homomorphism induces isomorphisms of graded Lie algebras
$$\pi _{}(\mathrm{\Omega }M(𝒜^{k+1}))/\mathrm{Torsion}\mathrm{Prim}H_{}(\mathrm{\Omega }M(𝒜^{k+1});)E_0^{}(G)_k,$$
where the Lie algebra structure of the source is induced by the Samelson product.
The paper is organized as follows.
1. Given a hyperplane arrangement $`𝒜^{\mathrm{}}`$, there is an associated arrangement of codimension $`k`$ subspaces $`𝒜^k^k\mathrm{}`$. The combinatorics and topology of the subspace arrangement $`𝒜^k`$ are studied in this section.
2. The topology of the subspace arrangement $`𝒜^k`$, in the instance where the underlying hyperplane arrangement $`𝒜`$ is fiber-type, is further studied in this section.
3. The (known) structure of the Lie algebra $`E_0^{}(G)`$ associated to the descending central series of the fundamental group $`G=\pi _1(M(𝒜))`$ of the complement of a fiber-type hyperplane arrangement $`𝒜`$ is analyzed in this section.
4. The structure of the Lie algebra of primitive elements in the homology of the loop space of the complement of the subspace arrangement $`𝒜^k`$ is analyzed in this section, and the isomorphisms of graded Lie algebras asserted in Theorem 1.3 are established.
5. The Poisson algebra structure on the homology of an iterated loop space of the complement of the subspace arrangement $`𝒜^k`$ is briefly analyzed in this section.
## 2. Redundant Arrangements
Let $`H`$ be an affine hyperplane in $`^{\mathrm{}}`$, an affine subspace of codimension one. For each positive integer $`k`$, there is an affine subspace $`H^k`$ of codimension $`k`$ in $`^k\mathrm{}`$ obtained from $`H`$ in the following manner. Choose coordinates $`𝐱=(x_1,\mathrm{},x_{\mathrm{}})`$ on $`^{\mathrm{}}`$, and $`(𝐱_1,\mathrm{},𝐱_k)`$ on $`^k\mathrm{}=^{\mathrm{}}\times \mathrm{}\times ^{\mathrm{}}`$, where for each $`i`$, $`𝐱_i=(x_{i,1},\mathrm{},x_{i,\mathrm{}})^{\mathrm{}}`$. Then, if the hyperplane $`H`$ in $`^{\mathrm{}}`$ is given by $`H=\{𝐱^{\mathrm{}}_{j=1}^{\mathrm{}}a_jx_j=b\}`$, define a codimension $`k`$ affine subspace $`H^k`$ in $`^k\mathrm{}`$ by $`H^k=\{(𝐱_1,\mathrm{},𝐱_k)^k\mathrm{}_{j=1}^{\mathrm{}}a_jx_{i,j}=b,1ik\}`$.
Now let $`𝒜`$ be a hyperplane arrangement in $`^{\mathrm{}}`$, a finite collection of (affine) hyperplanes. Via the above process, there is an arrangement $`𝒜^k=\{H^kH𝒜\}`$ of codimension $`k`$ affine subspaces in $`^k\mathrm{}`$ obtained from $`𝒜`$. For evident reasons, call the subspace arrangement $`𝒜^k`$ redundant. A brief description of the relationship between the combinatorics and topology of the hyperplane arrangement $`𝒜=𝒜^1`$ and the redundant subspace arrangement $`𝒜^k`$ is given in this section.
###### Example 2.1.
Let $`𝒜_n`$ be the braid arrangement in $`^n`$, consisting of the hyperplanes $`H_{i,j}=\{𝐱^nx_i=x_j\}`$. As is well known, the complement $`M(𝒜_n)=F(,n)`$ is the configuration space of $`n`$ points in $``$.
For each positive integer $`k`$, the associated redundant arrangement $`𝒜_n^k`$ consists of subspaces $`H_{i,j}^k=\{(𝐱_1,\mathrm{},𝐱_k)(^n)^kx_{p,i}=x_{p,j},1pk\}`$. These subspaces may be realized as $`H_{i,j}^k=\{(𝐲_1,\mathrm{},𝐲_n)(^k)^n𝐲_i=𝐲_j\}`$. Thus the complement $`M(𝒜_n^k)=F(^k,n)`$ is the configuration space of $`n`$ points in $`^k`$.
For an arbitrary hyperplane arrangement $`𝒜`$, and for each $`k`$, let $`𝖫(𝒜^k)`$ be the intersection poset of the arrangement $`𝒜^k`$, the partially ordered set of non-empty multi-intersections of elements of $`𝒜^k`$. Order $`𝖫(𝒜^k)`$ by reverse inclusion, and include the ambient space $`^k\mathrm{}`$ in $`𝖫(𝒜^k)`$ as the minimal element, corresponding to the intersection of no elements of $`𝒜^k`$. For the hyperplane arrangement $`𝒜=𝒜^1`$, it is known that $`𝖫(𝒜)`$ is a geometric poset, see \[20, Section 2.3\]. This need not be the case for an arbitrary subspace arrangement. However, for redundant arrangements, the following holds.
###### Proposition 2.2.
If $`𝒜`$ is a hyperplane arrangement, then $`𝖫(𝒜)𝖫(𝒜^k)`$ for all $`k`$.
###### Proof.
It will be shown that the bijection between $`𝒜`$ and $`𝒜^k`$ given by $`HH^k`$ induces an isomorphism of posets $`𝖫(𝒜)𝖫(𝒜^k)`$. To establish this, it suffices to show that to each codimension $`r`$ flat $`X𝖫(𝒜)`$ there corresponds a codimension $`kr`$ flat $`X^k𝖫(𝒜^k)`$.
Write $`𝒜=\{H_1,\mathrm{},H_n\}`$, where $`H_i=\{𝐱^{\mathrm{}}_{j=1}^{\mathrm{}}a_{i,j}x_j=b_i\}`$, and let $`X=H_1\mathrm{}H_m`$. The flat $`X`$ may be realized as the set of solutions of the system of equations $`A𝐱=𝐛`$, where $`A=(a_{i,j})`$ is $`m\times \mathrm{}`$ and $`𝐛=(b_1,\mathrm{},b_m)`$. Then, $`X`$ has codimension $`r`$ in $`^{\mathrm{}}`$ if and only if $`\mathrm{rank}[A𝐛]=\mathrm{rank}A=\mathrm{}r`$ if and only if
$$\mathrm{rank}\left[\begin{array}{cccccc}A& & & & |& 𝐛\\ & A& & & |& 𝐛\\ & & & & |& \\ & & & A& |& 𝐛\end{array}\right]=\mathrm{rank}\left[\begin{array}{cccc}A& & & \\ & A& & \\ & & & \\ & & & A\end{array}\right]=k(\mathrm{}r)$$
if and only if $`X^k=H_1^k\mathrm{}H_m^k`$ has codimension $`kr`$ in $`^k\mathrm{}`$. ∎
For each $`k`$, let $`M(𝒜^k)=^k\mathrm{}_{H^k𝒜^k}H^k`$ denote the complement of the (subspace) arrangement $`𝒜^k`$. In the case $`k=1`$, the cohomology of the hyperplane complement $`M(𝒜)=M(𝒜^1)`$ is well known. It is isomorphic to the Orlik-Solomon algebra $`𝖠(𝒜)`$, see \[20, Sections 3.1, 3.2\]. A family of algebras which includes $`𝖠(𝒜)`$ is defined next.
For each positive integer $`k`$, let $`𝖤_{2k1}[k]=_{H𝒜}e_H^k`$ be a free $``$-module generated by degree $`2k1`$ elements $`e_H^k`$ in one-to-one correspondence with the hyperplanes of $`𝒜`$. Let $`𝖤[k]=𝖤_{2k1}[k]`$ be the exterior algebra of $`𝖤_{2k1}[k]`$, and denote by $`𝖨[k]`$ the ideal of $`𝖤[k]`$ generated by the homogeneous elements
$`{\displaystyle \underset{p=1}{\overset{q}{}}}(1)^{p1}e_{H_1}^k\mathrm{}\widehat{e_{H_p}^k}\mathrm{}e_{H_q}^k`$ $`\text{if}0\mathrm{codim}H_1\mathrm{}H_q<q,`$
$`e_{H_1}^k\mathrm{}e_{H_q}^k`$ $`\text{if}H_1\mathrm{}H_q=\mathrm{}.`$
Let $`𝖠[k]=𝖤[k]/𝖨[k]`$. The Orlik-Solomon algebra is then given by $`𝖠(𝒜)=𝖠[1]`$.
Proposition 2.2 may be used to determine the cohomology of $`M(𝒜^k)`$ for $`k>1`$ in terms of that of $`M(𝒜)`$. Let $`P(𝒜^k,t)=_{q0}b_q(M(𝒜^k))t^q`$ be the Poincaré polynomial of $`M(𝒜^k)`$, where $`b_q(X)`$ is the $`q`$-th Betti number of $`X`$. Results of Goresky and MacPherson , and Yuzvinsky , see also Feichtner and Ziegler , together with Proposition 2.2, yield the following.
###### Corollary 2.3.
Let $`𝒜`$ be a hyperplane arrangement in $`^{\mathrm{}}`$.
1. For each $`k`$, the integral (co)homology of $`M(𝒜^k)`$ is torsion free, and we have $`P(𝒜^k,t)=P(𝒜,t^{2k1})`$.
2. For each $`k`$, the cohomology algebra of $`M(𝒜^k)`$ isomorphic to the algebra $`𝖠[k]`$, $`H^{}(M(𝒜^k);)𝖠[k]`$.
An explicit basis for the first non-zero (reduced) homology group, $`H_{2k1}(M(𝒜^k);)`$, of the complement of the subspace arrangement $`𝒜^k`$ is recorded next. Let $`L^{\mathrm{}}`$ be a complex line that is transverse to the hyperplane arrangement $`𝒜`$. Write $`L=\{t𝐮+𝐯\}`$ where $`𝐮,𝐯^{\mathrm{}}`$ are fixed and $`t`$ varies. For each hyperplane $`H`$ of $`𝒜`$, the intersection $`LH`$ is a point, say $`𝐪_H=\tau _H𝐮+𝐯`$ for some $`\tau _H`$. The following is immediate.
###### Lemma 2.4.
The subspace $`L^k=\{(t_1𝐮+𝐯,\mathrm{},t_k𝐮+𝐯)t_1,\mathrm{},t_k\}`$ of $`^k\mathrm{}`$ is transverse to the subspace arrangement $`𝒜^k^k\mathrm{}`$. For each subspace $`H^k`$ of $`𝒜^k`$, the intersection $`L^kH^k`$ is the point $`(𝐪_H,\mathrm{},𝐪_H)=(\tau _H𝐮+𝐯,\mathrm{},\tau _H𝐮+𝐯)`$.
Let $`S^{2k1}`$ be the unit sphere in $`^k`$. For $`ϵ>0`$ sufficiently small, the point
$$((\tau _H+ϵ^{}z_1)𝐮+𝐯,\mathrm{},(\tau _H+ϵ^{}z_k)𝐮+𝐯)L^k$$
lies in the intersection $`L^kM(𝒜^k)`$ for all $`ϵ^{}`$, $`0<ϵ^{}ϵ`$. Fix such an $`ϵ`$, and define a map $`c_H^k:S^{2k1}L^kM(𝒜^k)`$ using the above formula:
(2.1)
$$c_H^k(𝐳)=c_H^k(z_1,\mathrm{},z_k)=((\tau _H+ϵz_1)𝐮+𝐯,\mathrm{},(\tau _H+ϵz_k)𝐮+𝐯).$$
Let $`\iota _{2k1}`$ be the fundamental class of $`H_{2k1}(S^{2k1};)`$, and denote the image of $`(c_H^k)_{}(\iota _{2k1})H_{2k1}(L^kM(𝒜^k);)`$ under the map induced by the natural inclusion $`L^kM(𝒜^k)M(𝒜^k)`$ by $`C_H^kH_{2k1}(M(𝒜^k);)`$.
###### Proposition 2.5.
The classes $`\{C_H^kH𝒜\}`$ form a basis for $`H_{2k1}(M(𝒜^k);)`$.
###### Proof.
For $`H𝒜`$, define $`p_H^k:L^kM(𝒜^k)S^{2k1}`$ by
$$p_H^k(t_1𝐮+𝐯,\mathrm{},t_k𝐮+𝐯)=\frac{𝐭\tau _H𝐞}{𝐭\tau _H𝐞},$$
where $`𝐭=(t_1,\mathrm{},t_k)`$ and $`𝐞=(1,\mathrm{},1)`$ are in $`^k`$. It is then readily checked that $`p_H^kc_H^k=\mathrm{id}:S^{2k1}S^{2k1}`$ is the identity map. Furthermore, if $`H^{}H`$ is another hyperplane of $`𝒜`$, the composition $`p_H^kc_H^{}^k`$ is given by
$$p_H^kc_H^{}^k(𝐳)=\frac{𝐳+\frac{1}{ϵ}(\tau _H^{}\tau _H)𝐞}{𝐳+\frac{1}{ϵ}(\tau _H^{}\tau _H)𝐞},$$
so is null-homotopic. Consequently, the classes $`(c_H^k)_{}(\iota _{2k1})H_{2k1}(L^kM(𝒜^k);)`$ form a basis. Finally, using stratified Morse theory, one can show that the relative homology group $`H_i(M(𝒜^k),L^kM(𝒜^k);)`$ vanishes for $`i<4k2`$, see \[12, Parts II, III\]. It follows that the natural inclusion $`L^kM(𝒜^k)M(𝒜^k)`$ induces an isomorphism $`H_{2k1}(L^kM(𝒜^k);)\stackrel{}{}H_{2k1}(M(𝒜^k);)`$. So the classes $`C_H^k`$ form a basis for $`H_{2k1}(M(𝒜^k);)`$ as asserted. ∎
###### Remark 2.6.
The cohomology classes $`(C_H^k)^{}H^{2k1}(M(𝒜^k);)`$ dual to the classes $`C_H^kH_{2k1}(M(𝒜^k);)`$ generate the cohomology algebra $`H^{}(M(𝒜^k);)`$. Let $`a_H^k𝖠[k]`$ denote the image of $`e_H^k𝖤[k]`$ under the natural projection. Then the map $`H^{2k1}(M(𝒜^k);)𝖠_{2k1}[k]`$, $`(C_H^k)^{}a_H^k`$, induces an isomorphism of algebras $`H^{}(M(𝒜^k);)\stackrel{}{}𝖠[k]`$, see Corollary 2.3.
## 3. Linearly Fibered Arrangements
In this section, the topology of those redundant arrangements arising from strictly linearly fibered and fiber-type hyperplane arrangements is studied further. Recall the definition of arrangements of the former type from .
###### Definition 3.1.
A hyperplane arrangement $`𝒜`$ in $`^{\mathrm{}+1}`$ is strictly linearly fibered if there is a choice of coordinates $`(𝐱,z)=(x_1,\mathrm{},x_{\mathrm{}},z)`$ on $`^{\mathrm{}+1}`$ so that the restriction, $`p`$, of the projection $`^{\mathrm{}+1}^{\mathrm{}}`$, $`(𝐱,z)𝐱`$, to the complement $`M(𝒜)`$ is a fiber bundle projection, with base $`p(M(𝒜))=M()`$, the complement of an arrangement $``$ in $`^{\mathrm{}}`$, and fiber the complement of finitely many points in $``$. Refer to the hyperplane arrangement $`𝒜`$ as strictly linearly fibered over $``$.
The complements of hyperplane arrangements of this type are closely related to configuration spaces, as we now illustrate. For each hyperplane $`H`$ of $`𝒜`$, let $`f_H`$ be a linear polynomial with $`H=\mathrm{ker}f_H`$. Then $`Q(𝒜)=_{H𝒜}f_H`$ is a defining polynomial for $`𝒜`$. From the definition, if $`𝒜`$ is strictly linearly fibered over $``$ and $`|𝒜|=||+n`$, there is a choice of coordinates for which a defining polynomial for $`𝒜`$ factors as
(3.1)
$$Q(𝒜)=Q()\varphi (𝐱,z),$$
where $`Q()=Q()(𝐱)`$ is a defining polynomial for $``$, and $`\varphi (𝐱,z)`$ is a product
$$\varphi (𝐱,z)=(zg_1(𝐱))(zg_2(𝐱))\mathrm{}(zg_n(𝐱)),$$
with $`g_j(𝐱)`$ linear. Define $`g:^{\mathrm{}}^n`$ by
(3.2)
$$g(𝐱)=(g_1(𝐱),g_2(𝐱),\mathrm{},g_n(𝐱)),$$
Since $`\varphi (𝐱,z)`$ necessarily has distinct roots for any $`𝐱M()`$, the restriction of $`g`$ to $`M()`$ takes values in the configuration space $`F(,n)`$. The following result was proven by the first author, see \[2, Theorem 1.1.5, Corollary 1.1.6\].
###### Theorem 3.2.
Let $``$ be an arrangement of $`m`$ hyperplanes, and let $`𝒜`$ be an arrangement of $`m+n`$ hyperplanes which is strictly linearly fibered over $``$. Then the bundle $`p:M(𝒜)M()`$ is equivalent to the pullback of the bundle of configuration spaces $`p_{n+1}:F(,n+1)F(,n)`$ along the map $`g`$. Consequently, the bundle $`p:M(𝒜)M()`$ admits a cross-section and has trivial local coefficients in homology.
Since it is relevant to the subsequent discussion, a proof is included.
###### Proof.
Denote points in $`F(,n+1)`$ by $`(𝐲,z)`$, where $`𝐲=(y_1,\mathrm{},y_n)F(,n)`$ and $`z`$ satisfies $`zy_j`$ for each $`j`$. Similarly, denote points in $`M(𝒜)`$ by $`(𝐱,z)`$, where $`𝐱M()`$ and $`\varphi (𝐱,z)0`$. In this notation, we have $`p_{n+1}(𝐲,z)=𝐲`$ and $`p(𝐱,z)=𝐱`$. Let $`E=\{(𝐱,(𝐲,z))M()\times F(,n+1)g(𝐱)=𝐲\}`$ be the total space of the pullback of $`p_{n+1}:F(,n+1)F(,n)`$ along the map $`g`$. It is then readily checked that the map $`h:M(𝒜)E`$ defined by $`h(𝐱,z)=(𝐱,(g(𝐱),z))`$ is an equivalence of bundles.
Since the bundle $`p_{n+1}:F(,n+1)F(,n)`$ admits a cross-section, so does the pullback $`p:M(𝒜)M()`$. Furthermore, the structure group of the latter is the pure braid group $`P_n`$. Consequently, the action of the fundamental group of the base $`M()`$ on that of the fiber $`\{n\text{points}\}`$ is by pure braid automorphisms. As such, this action is by conjugation (see for instance or ), hence is trivial in homology. ∎
It is now shown that redundant strictly linearly fibered arrangements admit (linear) fibrations, just as their codimension one progenitors do.
###### Theorem 3.3.
Let $`𝒜`$ be a hyperplane arrangement in $`^{\mathrm{}+1}`$ which is strictly linearly fibered over $``$, with projection $`p:M(𝒜)M()`$ induced by the map $`^{\mathrm{}+1}^{\mathrm{}}`$ given by $`(x_1,\mathrm{},x_{\mathrm{}},z)(x_1,\mathrm{},x_{\mathrm{}})`$. Then for each $`k`$, the map $`^{k(\mathrm{}+1)}^k\mathrm{}`$ given by $`(𝐱_1,\mathrm{},𝐱_{\mathrm{}},𝐳)(𝐱_1,\mathrm{},𝐱_{\mathrm{}})`$ induces a fiber bundle projection $`p^k:M(𝒜^k)M(^k)`$. Furthermore, the bundle $`p^k:M(𝒜^k)M(^k)`$ admits a cross-section.
###### Proof.
By the previous result, the bundle $`p:M(𝒜)M()`$ is equivalent to the pullback of $`p_{n+1}:F(,n+1)F(,n)`$ along the map $`g`$ of (3.2). An analogous result for the complements of the subspace arrangements $`𝒜^k`$ and $`^k`$ is established next. For $`k2`$, view $`^k\mathrm{}`$ as $`(^{\mathrm{}})^k`$ and $`^{kn}`$ as $`(^k)^n`$. Denote points in the configuration space $`F(^k,n+1)`$ by $`(𝐲_1,\mathrm{},𝐲_n,𝐳)`$, where $`(𝐲_1,\mathrm{},𝐲_n)F(^k,n)`$ and $`𝐳𝐲_j`$ for each $`j`$. Define $`g^k:^k\mathrm{}^{kn}`$ by
(3.3)
$$g^k(𝐱_1,\mathrm{},𝐱_k)=((g_1(𝐱_1),\mathrm{},g_1(𝐱_k)),\mathrm{}\mathrm{},(g_n(𝐱_1),\mathrm{},g_n(𝐱_k))).$$
where $`(g_i(𝐱_1),\mathrm{},g_i(𝐱_k))^k`$ for each $`i`$. It is readily checked that the restriction of $`g^k`$ to $`M(^k)`$ takes values in the configuration space $`F(^k,n)`$. Let $`\pi ^k:E^kM(^k)`$ be the pullback of the bundle $`p_{n+1}^k:F(^k,n+1)F(^k,n)`$ along this restriction, with total space $`E^k`$ consisting of all points
$$((𝐱_1,\mathrm{},𝐱_k),(𝐲_1,\mathrm{},𝐲_n,𝐳))M(^k)\times F(^k,n+1)$$
for which $`g^k(𝐱_1,\mathrm{},𝐱_k)=p_{n+1}^k(𝐲_1,\mathrm{},𝐲_n,𝐳)=(𝐲_1,\mathrm{},𝐲_n)`$.
Since the hyperplane arrangement $`𝒜`$ is strictly linearly fibered over $``$, the complement of the subspace arrangement $`𝒜^k`$ may be realized as
$$M(𝒜^k)=\{(𝐱_1,\mathrm{},𝐱_k,𝐳)M(^k)\times ^k𝐳(g_i(𝐱_1),\mathrm{},g_i(𝐱_k))\text{for}1in\}.$$
Define $`h^k:M(𝒜^k)E^k`$ by $`h^k(𝐱_1,\mathrm{},𝐱_k,𝐳)=((𝐱_1,\mathrm{},𝐱_k),(g^k(𝐱_1,\mathrm{},𝐱_k),𝐳))`$.
The map $`h^k`$ is a homeomorphism. Moreover, the following diagram commutes.
(3.4)
$$\begin{array}{ccc}M(𝒜^k)& \stackrel{h^k}{}& E^k\\ p^k& & \pi ^k& & \\ M(^k)& \stackrel{\mathrm{id}}{}& M(^k)\end{array}$$
It follows that $`p^k:M(𝒜^k)M(^k)`$ is a bundle which is equivalent to the pullback of the bundle of configuration spaces $`p_{n+1}^k:F(^k,n+1)F(^k,n)`$ along the map $`g^k:M(^k)F(^k,n)`$, and therefore has a cross-section. ∎
An analysis of map in homology induced by the map $`g^k:M(^k)F(^k,n)`$ defined in (3.3) is given next. For $`1i<jn`$, define $`p_{i,j}:F(^k,n)S^{2k1}`$ by $`p_{i,j}(𝐲_1,\mathrm{},𝐲_n)=(𝐲_j𝐲_i)/𝐲_j𝐲_i`$. Recall that $`\iota _{2k1}H_{2k1}(S^{2k1};)`$ denotes the fundamental class. The classes $`p_{i,j}^{}(\iota _{2k1})`$ form a basis for $`H^{2k1}(F(^k,n))`$, and generate the cohomology algebra $`H^{}(F(^k,n))`$, see . Denote the dual classes in $`H_{2k1}(F(^k,n)`$ by $`A_{i,j}`$, $`1i<jn`$. Note that the classes $`A_{i,j}`$ may be represented by spheres linking the subspaces $`H_{i,j}^k=\{𝐲_i=𝐲_j\}`$ in $`^{kn}`$, as in (2.1).
As in Section 2, let $`L=\{t𝐮+𝐯\}^{\mathrm{}}`$ be a line transverse to the hyperplane arrangement $``$, and $`L^k`$ the corresponding codimension $`k`$ subspace of $`^k\mathrm{}`$, transverse to the subspace arrangement $`^k`$. Recall the maps $`c_H^k:S^{2k1}L^kM(^k)`$ from (2.1), and the resulting basis $`\{C_H^kH\}`$ for $`H_{2k1}(M(^k))`$ exhibited in Proposition 2.5.
###### Proposition 3.4.
Let $`^{\mathrm{}}`$ be an arrangement of complex hyperplanes, and let $`g:^{\mathrm{}}^n`$ be an affine transformation whose restriction, $`g:M()F(,n)`$, to the complement of $``$ takes values in the configuration space $`F(,n)`$. Then for every $`k1`$, the induced map $`(g^k)_{}:H_{2k1}(M(^k);)H_{2k1}(F(^k,n);)`$ is given by $`(g^k)_{}(C_H^k)=A_{i,j}`$ for each hyperplane $`H`$ of $``$, where the sum is over all distinct $`i`$ and $`j`$ for which $`g(H)`$ is contained in the hyperplane $`H_{i,j}=\{y_i=y_j\}`$ in $`^n`$.
###### Proof.
For each hyperplane $`H`$ of $``$, let $`\stackrel{~}{c}_H^k:S^{2k1}M(^k)`$ denote the composition of $`c_H^k:S^{2k1}L^kM(^k)`$ and the natural inclusion $`L^kM(^k)M(^k)`$. It will be shown that the composition $`p_{i,j}g^k\stackrel{~}{c}_H^k:S^{2k1}S^{2k1}`$ induces the identity in homology if $`g(H)H_{i,j}`$, and induces the trivial homomorphism if $`g(H)H_{i,j}`$, thereby establishing the result.
For $`𝐱^{\mathrm{}}`$, write $`g(𝐱)=(g_1(𝐱),\mathrm{},g_n(𝐱))`$ as in (3.2). Then $`g^k:^k\mathrm{}^{kn}`$ is given by $`g^k(𝐱_1,\mathrm{},𝐱_k)=(𝐲_1,\mathrm{},𝐲_n)`$, where $`𝐲_i=(g_i(𝐱_1),\mathrm{},g_i(𝐱_k))`$, see (3.3). Since the restriction of $`g`$ to $`M()`$ takes values in $`F(,n)`$, the restriction of $`g^k`$ to $`M(^k)`$ takes values in $`F(^k,n)`$.
From (2.1), the map $`\stackrel{~}{c}_H^k:S^{2k1}M(^k)`$ is given by $`\stackrel{~}{c}_H^k(𝐳)=(𝐰_1,\mathrm{},𝐰_k)`$, where $`𝐰_j=(\tau _H+ϵz_j)𝐮+𝐯`$, and $`LH`$ is the point $`𝐪_H=\tau _H𝐮+𝐯`$. Let $`\alpha _i=g_i(𝐪_H)`$, and define $`\beta _i`$ by the equation
$$g_i(𝐰_j)=g_i((\tau _H+ϵz_j)𝐮+𝐯)=g_i(𝐪_H+ϵz_j𝐮)=\alpha _i+ϵ\beta _iz_j.$$
Then, a calculation yields $`g^k\stackrel{~}{c}_H^k(𝐳)=(\alpha _1𝐞+ϵ\beta _1𝐳,\mathrm{},\alpha _n𝐞+ϵ\beta _n𝐳)`$ and
$$p_{i,j}g^k\stackrel{~}{c}_H^k(𝐳)=\frac{ϵ(\beta _j\beta _i)𝐳+(\alpha _j\alpha _i)𝐞}{ϵ(\beta _j\beta _i)𝐳+(\alpha _j\alpha _i)𝐞},$$
where, as before, $`𝐞=(1,\mathrm{},1)`$. Recall that $`ϵ>0`$ was chosen sufficiently small so as to insure that the point $`(𝐰_1^{},\mathrm{},𝐰_k^{})`$, where $`𝐰_j^{}=(\tau _H+ϵ^{}z_j)𝐮+𝐯`$, lies in $`L^kM(^k)`$ for all $`ϵ^{}`$, $`0<ϵ^{}ϵ`$. Since $`g^k:M(^k)F(^k,n)`$, it follows that $`g^k\stackrel{~}{c}_H^k(𝐳)F(^k,n)`$ for all $`𝐳S^{2k1}`$. In other words, $`ϵ(\beta _j\beta _i)𝐳+(\alpha _j\alpha _i)𝐞\mathrm{𝟎}`$ for all distinct $`i`$ and $`j`$.
If $`g(H)H_{i,j}`$, then $`g(𝐪_H)H_{i,j}`$ since $`𝐪_H=LH`$ is generic in $`H`$. Thus, $`\alpha _i=g_i(𝐪_H)g_j(𝐪_H)=\alpha _j`$, and the point $`(\alpha _i𝐞+ϵ^{}\beta _i𝐳,\alpha _j𝐞+ϵ^{}\beta _j𝐳)`$ lies in the configuration space $`F(^k,2)`$ for all $`ϵ^{}ϵ`$, including $`ϵ^{}=0`$. It follows that $`p_{i,j}g^kc_H^k`$ is trivial in homology in this instance.
If, on the other hand, $`g(H)H_{i,j}`$, then $`\alpha _i=\alpha _j`$ and thus $`\beta _j\beta _i`$ is necessarily non-zero. In this instance, $`p_{i,j}g^k\stackrel{~}{c}_H^k(𝐳)=\lambda 𝐳`$, where $`\lambda S^1^{}`$ is given by $`\lambda =(\beta _j\beta _i)/|\beta _j\beta _i|`$, which clearly induces the identity in homology. ∎
These results extend immediately to fiber-type arrangements, defined next.
###### Definition 3.5.
An arrangement $`𝒜=𝒜_1`$ of finitely many points in $`^1`$ is fiber-type. An arrangement $`𝒜=𝒜_{\mathrm{}}`$ of hyperplanes in $`^{\mathrm{}}`$ is fiber-type if $`𝒜`$ is strictly linearly fibered over a fiber-type hyperplane arrangement $`𝒜_\mathrm{}1`$ in $`^\mathrm{}1`$.
Let $`𝒜`$ be a fiber-type hyperplane arrangement in $`^{\mathrm{}}`$. Then there is a choice of coordinates $`(x_1,\mathrm{},x_{\mathrm{}})`$ on $`^{\mathrm{}}`$ so that a defining polynomial for $`𝒜`$ factors as $`Q(𝒜)=_{j=1}^{\mathrm{}}Q_j(x_1,\mathrm{},x_j)`$, see (3.1). Write $`Q_j=_{i=1}^{d_j}\left(x_jg_{i,j}(x_1,\mathrm{},x_{j1})\right)`$, where $`d_j`$ is the degree of $`Q_j`$ and each $`g_{i,j}`$ is linear. The non-negative integers $`\{d_1,\mathrm{},d_{\mathrm{}}\}`$ are called the exponents of $`𝒜`$. For each $`j\mathrm{}`$, the polynomial $`_{i=1}^jQ_i`$ defines a fiber-type arrangement $`𝒜_j`$ in $`^j`$ with exponents $`\{d_1,\mathrm{},d_j\}`$, and $`𝒜_j`$ is strictly linearly fibered over $`𝒜_{j1}`$. Furthermore, the map $`g_j=(g_{1,j},\mathrm{},g_{d_j,j}):^{j1}^{d_j}`$ gives rise to maps $`g_j^k:M(𝒜_{j1}^k)F(^k,d_j)`$ for each $`k`$. Theorems 3.2 and 3.3 yield
###### Theorem 3.6.
Let $`𝒜`$ be a fiber-type hyperplane arrangement in $`^{\mathrm{}}`$ with defining polynomial $`Q(𝒜)=_{j=1}^{\mathrm{}}Q_j`$. Then, for each $`j`$, $`2j\mathrm{}`$, and each $`k1`$, the projection $`^j^{j1}`$, $`(x_1,\mathrm{},x_{j1},x_j)(x_1,\mathrm{},x_{j1})`$, gives rise to a bundle map $`p_j^k:M(𝒜_j^k)M(𝒜_{j1}^k)`$. This bundle is equivalent to the pullback of the bundle of configuration spaces $`F(^k,d_j+1)F(^k,d_j)`$ along the map $`g_j^k:M(𝒜_{j1}^k)F(^k,d_j)`$. Consequently, the bundle $`p_j^k:M(𝒜_j^k)M(𝒜_{j1}^k)`$ admits a cross-section, has trivial local coefficients in homology, and the fiber is the complement of $`d_j`$ points in $`^k`$.
Proposition 3.4 also extends to fiber-type arrangements. The specific statement is omitted.
## 4. The Descending Central Series
In this section, the structure of the Lie algebra $`E_0^{}(G)`$ associated to the descending central series of the fundamental group $`G`$ of the complement of a fiber-type arrangement is analyzed. Additively, this structure is given by well known results of Falk and Randell stated below. Moreover, as shown by Jambu and Papadima , this Lie algebra is isomorphic to the (integral) holonomy Lie algebra of the arrangement $`𝒜`$ defined by Kohno . An alternate description of $`E_0^{}(G)`$, which will facilitate comparison with the Lie algebra of primitives in the homology of the loop space of the complement of the subspace arrangement $`𝒜^k`$ in Section 5, is given here.
###### Example 4.1.
Let $`P_n`$ be the Artin pure braid group, the fundamental group of the configuration space $`F(,n)`$. The structure of the Lie algebra $`E_0^{}(P_n)`$ was first determined rationally by Kohno . As observed by many authors, the following description holds over the integers as well. For each $`n2`$, let $`L[n]`$ be the free Lie algebra generated by elements $`A_{1,n+1},\mathrm{},A_{n,n+1}`$. Then the Lie algebra $`E_0^{}(P_n)`$ is additively isomorphic to the direct sum $`_{j=1}^{n1}L[j]`$, and the Lie bracket relations in $`E_0^{}(P_n)`$ are the infinitesimal pure braid relations, given by
(4.1) $`[A_{i,j}+A_{i,k}+A_{j,k},A_{m,k}]`$ $`=0\text{for }m=i\text{ or }m=j\text{, and}`$
$`[A_{i,j},A_{k,l}]`$ $`=0\text{for }\{i,j\}\{k,l\}=\mathrm{}\text{.}`$
Note that this description realizes the Lie algebra $`E_0^{}(P_{n+1})`$ as the semidirect product of $`E_0^{}(P_n)`$ by $`L[n]`$ determined by the Lie homomorphism $`\theta _n:E_0^{}(P_n)\mathrm{Der}(L[n])`$ given by $`\theta _n(A_{i,j})=\mathrm{ad}(A_{i,j})`$. From the infinitesimal pure braid relations, one has
(4.2)
$$\mathrm{ad}(A_{i,j})(A_{m,n+1})=\{\begin{array}{cc}[A_{m,n+1},A_{i,n+1}+A_{j,n+1}]\hfill & \text{if }m=i\text{ or }m=j\text{,}\hfill \\ 0\hfill & \text{otherwise.}\hfill \end{array}$$
This extension of Lie algebras arises topologically from the bundle of configuration spaces $`F(,n+1)F(,n)`$.
The additive structure noted above may be obtained from the following result of Falk and Randell .
###### Theorem 4.2.
Let $`1HGK1`$ be a split extension of groups such that the conjugation action of $`K`$ is trivial on $`H_1/H_2`$. Then there is a short exact sequence of Lie algebras $`0E_0^{}(H)E_0^{}(G)E_0^{}(K)0`$ which is split as a sequence of abelian groups. Furthermore, if the descending central series quotients of $`H`$ and $`K`$ are free abelian, then so are those of $`G`$.
Now let $`𝒜=𝒜_{\mathrm{}}`$ be a fiber-type hyperplane arrangement in $`^{\mathrm{}}`$. The complement of $`𝒜_{\mathrm{}}`$ sits atop a tower of fiber bundles
$$M(𝒜_{\mathrm{}})\stackrel{p_{\mathrm{}}}{}M(𝒜_\mathrm{}1)\stackrel{p_\mathrm{}1}{}\mathrm{}\stackrel{p_2}{}M(𝒜_1)=\{d_1\text{ points}\},$$
where the fiber of $`p_j`$ is homeomorphic to the complement of $`d_j`$ points in $``$. For simplicity, write $`=𝒜_\mathrm{}1`$ and $`n=d_{\mathrm{}}`$. Then $`𝒜`$ is strictly linearly fibered over $``$, and by Theorem 3.2, the bundle $`p:M(𝒜)M()`$ is equivalent to the pullback of the configuration space bundle $`p_{n+1}:F(,n+1)F(,n)`$ along the map $`g`$ of (3.2).
Application of the homotopy exact sequence of a bundle (and induction) shows that $`M(𝒜)`$ is a $`K(G,1)`$ space, where $`G=G(𝒜)=\pi _1(M(𝒜))`$. In light of Theorem 3.2, there is also a commutative diagram
(4.3)
$$\begin{array}{ccccccccc}1& & 𝔽_n& & G(𝒜)& & G()& & 1\\ & & \mathrm{id}& & & & g_\mathrm{\#}& & \\ 1& & 𝔽_n& & P_{n+1}& & P_n& & 1\end{array}$$
where $`g_\mathrm{\#}:G()P_n`$ is induced by $`g:M()F(,n)`$, and the fundamental group of the fiber $`\{n\text{points}\}`$ is identified with the free group $`𝔽_n`$ on $`n`$ generators. Since the underlying bundles admit cross-sections, the rows in the diagram above are split exact.
###### Theorem 4.3.
Let $`𝒜`$ be a fiber-type hyperplane arrangement. If the exponents of $`𝒜`$ are $`\{d_1,\mathrm{},d_{\mathrm{}}\}`$, then $`E_0^{}(G(𝒜))L[d_1]\mathrm{}L[d_{\mathrm{}}]`$ as abelian groups.
###### Proof.
The proof is by induction on $`\mathrm{}`$.
In the case $`\mathrm{}=1`$, $`𝒜`$ is an arrangement of $`d=d_1`$ points in $``$, the fundamental group of the complement is $`𝔽_d`$, the free group on $`d`$ generators, and it is well known that $`E_0^{}(𝔽_d)`$ is isomorphic to the free Lie algebra $`L[d]`$, see for instance \[21, Chapter IV\].
In general, assume that the fiber-type arrangement $`𝒜`$ is strictly linearly fibered over $``$ and that $`d_{\mathrm{}}=n`$ as above. Then there is a split, short exact sequence of groups $`1𝔽_nG(𝒜)G()1`$, and by Theorem 3.2, the action of $`G()`$ on $`𝔽_n`$ is by pure braid automorphisms. As such, this action is by conjugation, hence is trivial on $`H_{}(𝔽_n;)`$. By Theorem 4.2, the descending central series quotients of $`G(𝒜)`$ are free abelian, and there is a short exact sequence of Lie algebras
(4.4)
$$0E_0^{}(𝔽_n)E_0^{}(G(𝒜))E_0^{}(G())0,$$
which splits as a sequence of abelian groups. The result follows by induction. ∎
The additive decomposition provided by this result does not, in general, preserve the underlying Lie algebra structure. An inductive description of the Lie algebra structure of $`E_0^{}(G(𝒜))`$ is given next.
###### Theorem 4.4.
Let $`𝒜`$ and $``$ be fiber-type hyperplane arrangements with $`|𝒜|=||+n`$, and suppose that $`𝒜`$ is strictly linearly fibered over $``$. Then the Lie algebra $`E_0^{}(G(𝒜))`$ is isomorphic to the semidirect product of $`E_0^{}(G())`$ by $`L[n]`$ determined by the Lie homomorphism $`\mathrm{\Theta }=\theta _ng_{}:E_0^{}(G())\mathrm{Der}(L[n])`$, where $`g_{}:E_0^{}(G())E_0^{}(P_n)`$ is induced by the map $`g:M()F(,n)`$, and $`\theta _n:E_0^{}(P_n)\mathrm{Der}(L[n])`$ is given by $`\theta _n(A_{i,j})=\mathrm{ad}(A_{i,j})`$.
###### Proof.
From the exact sequence of Lie algebras (4.4) noted above, it follows that $`E_0^{}(G(𝒜))`$ is isomorphic to the semidirect product of $`E_0^{}(G())`$ by $`L[n]`$ determined by the Lie homomorphism $`\mathrm{\Theta }:E_0^{}(G())\mathrm{Der}(L[n])`$ given by $`\mathrm{\Theta }(b)=\mathrm{ad}_{L[n]}(b)`$ for $`bE_0^{}(G())`$. It suffices to show that the homomorphism $`\mathrm{\Theta }`$ factors as asserted.
From the diagram (4.3), and the results of Falk and Randell stated in Theorem 4.2, there is a commutative diagram of Lie algebras with split exact rows
$$\begin{array}{ccccccccc}0& & L[n]& & E_0^{}(G(𝒜))& & E_0^{}(G())& & 0\\ & & \mathrm{id}& & & & g_{}& & \\ 0& & L[n]& & E_0^{}(P_{n+1})& & E_0^{}(P_n)& & 0\end{array}$$
Via the splittings, view $`E_0^{}(G())`$ and $`E_0^{}(P_n)`$ as Lie subalgebras of $`E_0^{}(G(𝒜))`$ and $`E_0^{}(P_{n+1})`$ respectively. Then for $`aL[n]`$ and $`bE_0^{}(G())`$, we have $`[b,a]=[g_{}(b),a]`$ in $`L[n]`$. Thus $`\mathrm{ad}_{L[n]}(b)=\mathrm{ad}_{L[n]}(g_{}(b))`$ in $`\mathrm{Der}(L[n])`$ and $`\mathrm{\Theta }=\theta _ng_{}`$. ∎
This result, together with Proposition 3.4, provides an inductive description of the Lie bracket structure of $`E_0^{}(G(𝒜))`$. Recall the basis $`\{C_H^1H\}`$ for $`H_1(M();)=E_0^1(G())`$ exhibited in Proposition 2.5, and recall that the free Lie algebra $`L[n]`$ is generated by $`A_{1,n+1},\mathrm{},A_{n,n+1}`$.
###### Corollary 4.5.
For generators $`C_H^1`$ of $`E_0^1(G())`$ and $`A_{m,n+1}`$ of $`L[n]`$, one has
$$\mathrm{\Theta }(C_H^1)(A_{m,n+1})=\underset{g(H)H_{i,j}}{}[A_{i,j},A_{m,n+1}].$$
###### Proof.
By Proposition 3.4, one has $`g_{}(C_H^1)=A_{i,j}`$, where the sum is over all $`i`$ and $`j`$ for which $`g(H)H_{i,j}`$. The result follows. ∎
This corollary can be used to explicitly record the Lie bracket relations in $`E_0^{}(G(𝒜))`$, and to show that these relations are combinatorial, that is, completely determined by the intersection poset $`𝖫(𝒜)`$. The Lie algebra $`E_0^{}(G(𝒜))`$ is generated by $`\{C_H^1H𝒜\}`$. For a flat $`X𝖫(𝒜)`$, write $`C_X^1=_{XH}C_H^1`$. The following was proven by Jambu and Papadima , see also .
###### Theorem 4.6.
If $`𝒜`$ is a fiber-type hyperplane arrangement with exponents $`\{d_1,\mathrm{},d_{\mathrm{}}\}`$, then the Lie bracket relations in $`E_0^{}(G(𝒜))`$ are given by
$$[C_X^1,C_H^1]=0,$$
for codimension two flats $`X𝖫(𝒜)`$ and hyerplanes $`H𝒜`$ containing $`X`$.
###### Proof.
The proof is by induction on $`\mathrm{}`$.
In the case $`\mathrm{}=1`$, there is nothing to show since $`G(𝒜)`$ is a free group on $`d=d_1`$ generators, $`E_0^{}(G(𝒜))`$ is isomorphic to the free Lie algebra $`L[d]`$, and there are no codimension two flats in $`𝖫(𝒜)`$.
In general, assume that $`𝒜`$ is strictly linearly fibered over $``$ and that $`d_{\mathrm{}}=n`$ as before. Then $`𝒜`$ has a defining polynomial of the form $`Q(𝒜)=Q()_{j=1}^n(zg_j(𝐱))`$, see (3.1). View $``$ as a subarrangement of $`𝒜=\{HH\}\{H_j1jn\}`$, where $`H_j=\mathrm{ker}(zg_j(𝐱))`$. Then the set $`\{C_H^1H\}\{C_{H_j}^11jn\}`$ generates $`E_0^{}(G(𝒜))`$, where the generators $`C_{H_j}^1`$ correspond to the hyperplanes $`H_j`$ of $`𝒜`$, and to the generators $`A_{j,n+1}`$ of the free Lie algebra $`L[n]`$ under the additive isomorphism $`E_0^{}(G(𝒜))E_0^{}(G())L[n]`$.
By Theorem 4.4, $`E_0^{}(G(𝒜))`$ is isomorphic to an extension of $`E_0^{}(G())`$ by $`L[n]`$. Consequently, the Lie bracket relations in $`E_0^{}(G(𝒜))`$ consist of those of $`E_0^{}(G())`$, and those arising from the extension. By induction, the Lie bracket relations in $`E_0^{}(G())`$ are given by $`[C_X^1,C_H^1]=0`$ for codimension two flats $`X`$ contained only in hyperplanes $`H𝒜`$. So it remains to analyze those relations in $`E_0^{}(G(𝒜))`$ arising from the extension. These are given implicitly in Corollary 4.5.
Recall from (4.1) that $`[A_{i,j},A_{m,n+1}]=[A_{m,n+1},A_{i,n+1}+A_{j,n+1}]`$ if $`m\{i,j\}`$, and is zero otherwise. Thus the results of Corollary 4.5 may be recorded as
$$\mathrm{\Theta }(C_H^1)(A_{m,n+1})=[C_H^1,A_{m,n+1}]=\underset{g(H)H_{i,j}}{}[A_{m,n+1},(\delta _{i,m}+\delta _{j,m})(A_{i,n+1}+A_{j,n+1})],$$
where $`C_H^1E_0^{}(G())E_0^{}(G(𝒜))`$ and $`\delta _{i,m}`$ is the Kronecker delta. Note that the expression on the right lies in $`L[n]`$. Under the above identifications, these relations take the form
(4.5)
$$[C_H^1,C_{H_m}^1]=\underset{g(H)H_{i,j}}{}[C_{H_m}^1,(\delta _{i,m}+\delta _{j,m})(C_{H_i}^1+C_{H_j}^1)]$$
Now one can check that $`g(H)H_{i,j}`$ if and only if $`HH_iH_j`$ is a codimension two flat in $`𝖫(𝒜)`$ if and only if $`H_iH_jH`$. Using this observation, the relation (4.5) may be expressed as
$$[C_H^1,C_{H_m}^1]=[C_{H_m}^1,\underset{H_mH_jH}{}C_{H_j}^1].$$
A calculation then shows that this relation is equivalent to $`[C_X^1,C_{H_m}^1]=0`$, where $`X`$ is the codimension two flat in $`𝖫(𝒜)`$ contained in $`H`$ and $`H_m`$. Since this relation holds for all $`H_m𝒜`$ for which $`XHH_m`$, it follows that $`[C_X^1,C_H^1]=0`$ as well. ∎
## 5. Homology of the Loop Space
The structure of the Lie algebra of primitive elements in the homology of the loop space of the complement of a redundant subspace arrangement associated to a fiber-type hyperplane arrangement is analyzed in this section. In analogy with the previous section, begin by recalling this structure for the classical configuration spaces $`F(^{k+1},n)`$ for $`k1`$.
###### Example 5.1.
The integral homology of the loop space $`\mathrm{\Omega }F(^{k+1},n)`$ was calculated by Fadell and Husseini . The structure of the Lie algebra $`\mathrm{Prim}H_{}(\mathrm{\Omega }F(^{k+1},n);)`$ was subsequently determined by Cohen and Gitler . For brevity, denote this Lie algebra by $`(n)_k`$. The structure of $`(n)_k`$ may be described as follows.
For each $`n2`$, let $`L[n]_k`$ denote the free Lie algebra generated by elements $`B_{i,n+1}`$, $`1in`$, of degree $`2k`$. Then $`(n)_k`$ is additively isomorphic to the direct sum $`_{j=1}^{n1}L[j]_k`$, and the Lie bracket relations in $`(n)_k`$ are given by the infinitesimal pure braid relations on the $`B_{i,j}`$, see (4.1). Thus, there is an isomorphism of graded Lie algebras $`(n)_kE_0^{}(P_n)_k`$, see Definition 1.1, Theorem 1.2, and Example 4.1.
Furthermore, as is the case for the descending central series of the pure braid group, the Lie algebra $`(n+1)_k`$ is isomorphic to the semidirect product of $`(n)_k`$ by $`L[n]_k`$ determined by the Lie homomorphism $`\theta _n^k:(n)_k\mathrm{Der}(L[n]_k)`$ given by $`\theta _n^k(B_{i,j})=\mathrm{ad}(B_{i,j})`$. From the infinitesmial pure braid relations, there is a formula for $`\mathrm{ad}(B_{i,j})`$ analogous to that given in (4.2). As before, this extension of Lie algebras arises topologically from the bundle of configuration spaces $`F(^{k+1},n+1)F(^{k+1},n)`$.
Now let $`𝒜`$ be a fiber-type hyperplane arrangement in $`^{\mathrm{}}`$ with exponents $`\{d_1,\mathrm{},d_{\mathrm{}}\}`$. Then, for each $`k`$, there is a tower of fiber bundles
$$M(𝒜_{\mathrm{}}^k)\stackrel{p_{\mathrm{}}^k}{}M(𝒜_\mathrm{}1^k)\stackrel{p_\mathrm{}1^k}{}\mathrm{}\stackrel{p_2^k}{}M(𝒜_1^k)=^k\{d_1\text{ points}\},$$
where the fiber of $`p_j^k`$ is homeomorphic to the complement of $`d_j`$ points in $`^k`$, see Theorem 3.6. Furthermore, each of the fiber bundles $`p_j^k:M(𝒜_j^k)M(𝒜_{j1}^k)`$ involving the complements of the redundant subspace arrangements $`𝒜_j^k^{jk}`$ admits a cross-section, and, as indicated above, $`M(𝒜_1^k)`$ is the complement of $`d_1`$ points in $`^k`$. By work of the second two authors \[6, Theorem 1\], the following holds.
###### Theorem 5.2.
Let $`𝒜`$ be a fiber-type hyperplane arrangement in $`^{\mathrm{}}`$ with exponents $`\{d_1,\mathrm{},d_{\mathrm{}}\}`$. Then, for each $`k1`$,
1. There is a homotopy equivalence $`\mathrm{\Omega }M(𝒜^{k+1})_{j=1}^{\mathrm{}}\mathrm{\Omega }(^{k+1}\{d_j\text{ points}\})`$;
2. The integral homology of $`\mathrm{\Omega }M(𝒜^{k+1})`$ is torsion free, and is isomorphic to the tensor product $`_{j=1}^{\mathrm{}}H_{}(\mathrm{\Omega }(^{k+1}\{d_j\text{ points}\});)`$ as a coalgebra;
3. The module of primitives in the integral homology of $`\mathrm{\Omega }M(𝒜^{k+1})`$ is isomorphic to $`\pi _{}(\mathrm{\Omega }M(𝒜^{k+1}))`$ modulo torsion as a Lie algebra.
###### Remark 5.3.
The homotopy groups of a loop space admit a bilinear pairing, which satisfies the axioms for a graded Lie algebra in case there is no $`2`$ or $`3`$ torsion in the homotopy groups. The graded analogue of the symmetry law can fail in case $`2`$-torsion is present, while the graded analogue of the Jacobi identity can fail if $`3`$-torsion is present. Thus, forming the quotient of the homotopy groups by the torsion gives a graded module which satisfies the axioms for a graded Lie algebra. Analogous properties of iterated loop spaces yield a graded Poisson algebra, see Section 6 below.
###### Proof of Theorem 5.2.
Given a fibration $`F\stackrel{𝑖}{}E\stackrel{}{}B`$ with a section $`\sigma `$, there is a homotopy equivalence $`\mathrm{\Omega }E\mathrm{\Omega }B\times \mathrm{\Omega }F`$ given by the composite:
$$\mathrm{\Omega }B\times \mathrm{\Omega }F\stackrel{\mathrm{\Omega }\sigma \times \mathrm{\Omega }i}{}\mathrm{\Omega }E\times \mathrm{\Omega }E\stackrel{𝜇}{}\mathrm{\Omega }E,$$
where $`\mu `$ is the loop space multiplication and such that the inclusions of $`\mathrm{\Omega }B`$ and $`\mathrm{\Omega }F`$ into $`\mathrm{\Omega }E`$ are maps of $`H`$-spaces. Morever, if the spaces involved have torsion free homology then $`H_{}(\mathrm{\Omega }E)H_{}(\mathrm{\Omega }B)H_{}(\mathrm{\Omega }F)`$. By a theorem of Milnor and Moore, one obtains
(5.1)
$$\mathrm{Prim}H_{}(\mathrm{\Omega }E)\mathrm{Prim}H_{}(\mathrm{\Omega }B)\mathrm{Prim}H_{}(\mathrm{\Omega }F)$$
upon passing to the Lie algebra of primitives. This result is a topological analogue of Theorem 4.2 as the underlying Lie algebra structure is “twisted.”
Now apply these considerations to the fiber bundle $`p_j^{k+1}:M(𝒜_j^{k+1})M(𝒜_{j1}^{k+1})`$. The fiber in this case is $`F=^{k+1}\{d_j\text{ points}\}_{d_j}S^{2k+1}`$. Assertion (a) follows by induction, and then (b) by the Künneth theorem. By the Bott-Samelson Theorem, $`H_{}(\mathrm{\Omega }F)`$ is isomorphic to $`T[d_j]_k`$, a tensor algebra on $`d_j`$ generators of degree $`2k`$.
Thus the module of primitive elements is generated as a Lie algebra by the primitive elements in degree $`2k`$ which are in the image of the Hurewicz map. Since the Hurewicz map takes values in the the module of primitive elements, that module is generated as a Lie algebra by those spherical classes given by the homology classes of degree $`2k`$.
Next notice that the homology groups here are torsion free. Hence the Hurewicz map factors through $`\pi _{}\mathrm{\Omega }(M(𝒜^{k+1}))/\mathrm{Torsion}`$. Furthermore, the homotopy groups of a loop space modulo torsion give a graded Lie algebra where the Lie bracket is induced by the classical Samelson product, and the Hurewicz map is a morphism of graded Lie algebras. Thus the induced map $`\pi _{}\mathrm{\Omega }(M(𝒜^{k+1}))/\mathrm{Torsion}\mathrm{Prim}H_{}(\mathrm{\Omega }M(𝒜^{k+1});)`$ is an epimorphism of Lie algebras.
Since all spaces are simply connected, and are of finite type, the homotopy groups modulo torsion are finitely generated free abelian groups in any fixed degree. By a classical theorem of Milnor and Moore concerning rational homotopy groups, the induced map $`\pi _{}\mathrm{\Omega }(M(𝒜^{k+1}))/\mathrm{Torsion}\mathrm{Prim}H_{}(\mathrm{\Omega }M(𝒜^{k+1});)`$ is also a monomorphism. The result follows. ∎
By (5.1) above, there is an isomorphism of graded abelian groups
$$\mathrm{Prim}H_{}(\mathrm{\Omega }M(𝒜_j^{k+1}))\mathrm{Prim}H_{}(\mathrm{\Omega }M(𝒜_{j1}^{k+1}))L[d_j]_k.$$
Proceeding inductively, this implies that the Lie algebra $`\mathrm{Prim}H_{}(\mathrm{\Omega }M(𝒜^{k+1});)`$ is isomorphic to $`L[d_1]_k\mathrm{}L[d_{\mathrm{}}]_k`$ as a graded abelian group, where $`\{d_1,\mathrm{},d_{\mathrm{}}\}`$ are the exponents of $`𝒜`$. Thus the Lie algebras $`\mathrm{Prim}H_{}(\mathrm{\Omega }M(𝒜^{k+1});)`$ and $`E_0^{}(G(𝒜))_k`$ are additively isomorphic, see Theorem 4.3. To show that they are are isomorphic as Lie algebras, thereby completing the proof of Theorem 1.3, it remains to show that the Lie bracket structure of $`\mathrm{Prim}H_{}(\mathrm{\Omega }M(𝒜^{k+1});)`$ coincides with that of $`E_0^{}(G(𝒜))_k`$. This analysis parallels the determination of the Lie algebra structure of $`E_0^{}(G(𝒜))`$ in Section 4.
The fiber-type hyperplane arrangement $`𝒜=𝒜_{\mathrm{}}`$ is strictly linearly fibered over $`𝒜_\mathrm{}1`$, and $`|𝒜|=|𝒜_\mathrm{}1|+d_{\mathrm{}}`$. As before, write $`=𝒜_\mathrm{}1`$ and $`n=d_{\mathrm{}}`$. Recall the map $`g^{k+1}:M(^{k+1})F(^{k+1},n)`$ from (3.3). Recall also that the Lie algebra $`\mathrm{Prim}H_{}(\mathrm{\Omega }F(^{k+1},n);)`$ is denoted by $`(n)_k`$. Analogously, denote the Lie algebra $`\mathrm{Prim}H_{}(\mathrm{\Omega }M(𝒜^{k+1});)`$ by $`(𝒜)_k`$.
###### Theorem 5.4.
Let $`𝒜`$ and $``$ be fiber-type hyperplane arrangements with $`𝒜`$ strictly linearly fibered over $``$ and $`|𝒜|=||+n`$. Then the Lie algebra $`(𝒜)_k`$ is isomorphic to the semidirect product of $`()_k`$ by the free Lie algebra $`L[n]_k`$ determined by the Lie homomorphism $`\mathrm{\Theta }^k=\theta _n^k\gamma _{}^k:()_k\mathrm{Der}(L[n]_k)`$, where $`\gamma _{}^k:()_k(n)_k`$ is the map in loop space homology induced by $`g^{k+1}:M(^{k+1})F(^{k+1},n)`$, and $`\theta _n^k:(n)_k\mathrm{Der}(L[n]_k)`$ is given by $`\theta _n^k(B_{i,j})=\mathrm{ad}(B_{i,j})`$.
###### Proof.
The realization of the bundle $`p^{k+1}:M(𝒜^{k+1})M(^{k+1})`$ as the pullback of the bundle of configuration spaces $`p_{n+1}^{k+1}:F(^{k+1},n+1)F(^{k+1},n)`$ along the map $`g^{k+1}:M(^{k+1})F(^{k+1},n)`$ from Theorem 3.3 yields a commutative diagram of Hopf algebras
$$\begin{array}{ccccc}H_{}(\mathrm{\Omega }(^{k+1}\{n\text{points}\}))& & H_{}(\mathrm{\Omega }M(𝒜^{k+1}))& & H_{}(\mathrm{\Omega }M(^{k+1}))\\ \mathrm{id}& & & & \gamma _{}^k& & \\ H_{}(\mathrm{\Omega }(^{k+1}\{n\text{points}\}))& & H_{}(\mathrm{\Omega }F(^{k+1},n+1))& & H_{}(\mathrm{\Omega }F(^{k+1},n))\end{array}$$
with exact rows, and, on the level of primitives, a commutative diagram of Lie algebras
$$\begin{array}{ccccccccc}0& & L[n]_k& & (𝒜)_k& & ()_k& & 0\\ & & \mathrm{id}& & & & \gamma _{}^k& & \\ 0& & L[n]_k& & (n+1)_k& & (n)_k& & 0\end{array}$$
where $`\gamma _{}^k:()_k(n)_k`$ is induced by $`g^{k+1}:M(^{k+1})F(^{k+1},n)`$. Since the underlying bundles admit cross-sections, the rows in the above diagrams are split exact. Via these splittings, view $`()_k`$ and $`(n)_k`$ as Lie subalgebras of $`(𝒜)_k`$ and $`(n+1)_k`$ respectively.
From the above considerations, it follows that the Lie algebra $`(𝒜)_k`$ is isomorphic to the semidirect product of $`()_k`$ by $`L[n]_k`$ determined by the Lie homomorphism $`\mathrm{\Theta }^k:()_k\mathrm{Der}(L[n]_k)`$ given by $`\mathrm{\Theta }^k(b)=\mathrm{ad}_{L[n]_k}(b)`$ for $`b()_k`$. Moreover, for $`aL[n]_k`$, we have $`[b,a]=[\gamma _{}^k(b),a]`$ in $`L[n]_k`$. Thus $`\mathrm{ad}_{L[n]_k}(b)=\mathrm{ad}_{L[n]_k}(\gamma _{}^k(b))`$ in $`\mathrm{Der}(L[n]_k)`$ and $`\mathrm{\Theta }^k=\theta _n^k\gamma _{}^k`$. ∎
This result, together with Proposition 3.4, provides an inductive description of the Lie bracket structure of $`(𝒜)_k`$. The space $`M(^{k+1})`$ is $`2k`$-connected, and the cohomology algebra $`H^{}(M(^{k+1});)`$ is generated by classes $`a_H^{k+1}`$ in one-to-one correspondence with the hyperplanes $`H`$, see Corollary 2.3. These classes are of degree $`2k+1`$, and are dual to the elements of the basis $`\{C_H^{k+1}H\}`$ for $`H_{2k+1}(M(^{k+1});)`$ exhibited in Proposition 2.5. See also Remark 2.6.
The above observations imply that homology suspension induces an isomorphism
$$\sigma _{}:H_{2k}(\mathrm{\Omega }M(^{k+1});)H_{2k+1}(M(^{k+1});).$$
Let $`\beta _H^kH_{2k}(\mathrm{\Omega }M(^{k+1});)`$ be the unique class satisfying $`\sigma _{}(\beta _H^k)=C_H^{k+1}`$. Recall that the free Lie algebra $`L[n]_k`$ is generated by $`B_{1,n+1},\mathrm{},B_{n,n+1}`$.
###### Corollary 5.5.
For generators $`\beta _H^k`$ of $`()_k`$ and $`B_{m,n+1}`$ of $`L[n]_k`$, one has
$$\mathrm{\Theta }^k(\beta _H^k)(B_{m,n+1})=\underset{g(H)H_{i,j}}{}[B_{i,j},B_{m,n+1}].$$
###### Proof.
By Proposition 3.4, one has $`g_{}^{k+1}(C_H^{k+1})=A_{i,j}`$, where the sum is over all $`i`$ and $`j`$ for which $`g(H)H_{i,j}`$. Since the homology suspension $`\sigma _{}`$ is an isomorphism and $`\gamma _{}^k`$ is the map in loop space homology induced by $`g^{k+1}`$, one has $`\gamma _{}^k(\beta _H^k)=B_{i,j}`$, where the sum is over all $`i`$ and $`j`$ for which $`g(H)H_{i,j}`$. The result follows. ∎
To complete the proof of Theorem 1.3, assume inductively that the Lie algebras $`E_0^{}(G())_k`$ and $`()_k`$ are isomorphic. By Theorem 4.4, the Lie algebra $`E_0^{}(G(𝒜))`$ is the extension of $`E_0^{}(G())`$ by the free Lie algebra $`L[n]`$ (generated in degree one) determined by the Lie homomorphism $`\mathrm{\Theta }=\theta _ng_{}`$. Thus $`E_0^{}(G(𝒜))_k`$ may be realized as the extension of $`E_0^{}(G())_k`$ by the free Lie algebra $`L[n]_k`$ (generated in degree $`2k`$) determined by $`\mathrm{\Theta }`$ as specified in Definition 1.1. Similarly, by Theorem 5.4, the Lie algebra $`(𝒜)_k`$ is the extension of $`()_k`$ by the free Lie algebra $`L[n]_k`$ determined by the Lie homomorphism $`\mathrm{\Theta }^k=\theta _n^k\gamma _{}^k`$. A comparison of the results of Corollary 4.5 and Corollary 5.5 reveals that these extensions coincide. Therefore, the Lie algebras $`E_0^{}(G(𝒜))_k`$ and $`(𝒜)_k`$ are isomorphic.
Alternatively, Corollary 5.5 may be used to explicitly determine the Lie bracket structure in $`(𝒜)_k`$. As the argument is completely analogous to that which established Theorem 4.6, the result stated below without proof. The Lie algebra $`(𝒜)_k`$ is generated by $`\{\beta _H^kH𝒜\}`$. For a flat $`X𝖫(𝒜)`$, write $`\beta _X^k=_{XH}\beta _H^k`$.
###### Theorem 5.6.
Let $`𝒜`$ be a fiber-type hyperplane arrangement. Then, for each $`k1`$, the Lie bracket relations in $`(𝒜)_k`$ are given by
$$[\beta _X^k,\beta _H^k]=0,$$
for codimension two flats $`X𝖫(𝒜)`$ and hyerplanes $`H𝒜`$ containing $`X`$.
## 6. Homology of Iterated Loop Spaces
In this final section, the Poisson algebra structure on the homology of an iterated loop space of the complement of a redundant subspace arrangement associated to a fiber-type hyperplane arrangement is briefly analyzed.
For $`q>1`$, the homology of an $`q`$-fold loop space, $`\mathrm{\Omega }^qX`$, admits the structure of a graded Poisson algebra. Namely, there is a bilinear map given by the Browder operation
$$\lambda _{q1}:H_i(\mathrm{\Omega }^qX)H_j(\mathrm{\Omega }^qX)H_{i+j+q1}(\mathrm{\Omega }^qX)$$
which satisfies properties listed in \[4, pages 215–217\]. In particular, this pairing satisfies the axioms of a (graded) Poisson algebra, and is compatible with the Whitehead product structure for the classical Hurewicz homomorphism.
In the case where $`X=X_{\mathrm{}}(^{k+1})`$, $`k1`$, satisfies conditions (1)–(3) from the Introduction, these structures are analogues of classical constructions in homotopy theory. First, note that the single suspension $`\mathrm{\Sigma }X_{\mathrm{}}(^{k+1})`$ is homotopy equivalent to a bouquet of spheres. Thus there is an induced map $`\sigma ^2:\mathrm{\Sigma }^2X_{\mathrm{}}(^{k+1})X_{\mathrm{}}(^{k+2})`$ which induces an isomorphism on the first non-trivial homology group. The adjoint $`E^2:X_{\mathrm{}}(^{k+1})\mathrm{\Omega }^2X_{\mathrm{}}(^{k+2})`$ also induces an isomorphism on the first non-trivial homology group. This last map is an analogue of the classical Freudenthal double suspension where the spaces $`X_{\mathrm{}}(𝔼)`$ are replaced by single odd dimensional spheres. Looping $`E^2`$ is given by $`\mathrm{\Omega }(E^2):\mathrm{\Omega }X_{\mathrm{}}(^{k+1})\mathrm{\Omega }^3X_{\mathrm{}}(^{k+2})`$.
###### Theorem 6.1.
Let $`𝒜`$ be a fiber-type hyperplane arrangement in $`^{\mathrm{}}`$ with exponents $`\{d_1,\mathrm{},d_{\mathrm{}}\}`$. Then, for each $`k1`$,
1. The multiplicative map $`\mathrm{\Omega }(E^2):\mathrm{\Omega }M(𝒜^{k+1})\mathrm{\Omega }^3M(𝒜^{k+2})`$ induces an isomorphism on $`H_{2k}(;)`$, and is zero in degrees greater than $`2k`$.
2. If $`q>1`$, the homology of $`\mathrm{\Omega }^qM(𝒜^{k+1})`$, with any field coefficients, is a graded Poisson algebra with Poisson bracket given by the Browder operation for the homology of a $`q`$-fold loop space.
3. If $`q>1`$, then $`\mathrm{\Omega }^qM(𝒜^{k+1})`$ is homotopy equivalent to $`_{j=1}^{\mathrm{}}\mathrm{\Omega }^q(_{d_j}S^{2k+1})`$.
4. If $`1<q<2k+1`$, the homology of $`\mathrm{\Omega }^qM(𝒜^{k+1})`$, with coefficients in a field $`𝔽`$ of characteristic zero, is generated as a Poisson algebra by elements $`\beta _H`$ of degree $`2k+1q`$ for $`H𝒜`$. The Poisson bracket is given by the Browder operation $`\lambda _{q1}`$, and satisfies the relations
$$\lambda _{q1}[\beta _X,\beta _H],$$
for codimension two flats $`X𝖫(𝒜)`$ and hyerplanes $`H𝒜`$ containing $`X`$, where $`\beta _X=_{XH}\beta _H`$.
###### Sketch of Proof.
Part (a) follows from the fact that the homology of $`\mathrm{\Omega }^3M(𝒜^{k+2})`$ is abelian while the homology of $`\mathrm{\Omega }M(𝒜^{k+1})`$ is generated by Lie brackets of weight at least $`2`$ in homological degrees greater than $`2k`$.
Part (b) follows from the remarks at the beginning of this section.
Part (c) follows at once from the fact that the result holds in case $`q=1`$, which was established in Theorem 5.2.
In case $`q=1`$, the Browder operation $`\lambda _{q1}`$ is precisely the natural Lie bracket in the homology of a $`1`$-fold loop space, $`\mathrm{\Omega }M(𝒜^{k+1})`$. These Lie bracket relations are recorded in Theorem 5.6. As shown in \[4, pages 215–217\], a further property of the operation $`\lambda _{q1}`$ is that $`\sigma _{}\lambda _{q1}(x,y)=\lambda _{q2}(\sigma _{}x,\sigma _{}y)`$, where $`\sigma _{}`$ denotes the homology suspension. Thus by induction on $`q`$, the asserted Poisson bracket relations are satisfied modulo elements in the kernel of the suspension. Furthermore, $`\lambda _{q1}(x,y)`$ is primitive in case the classes $`x`$ and $`y`$ are primitive.
In characteristic zero, and in case $`q`$ is greater than $`1`$, the homology suspension induces an isomorphism on the module of primitives. Thus the asserted Poisson bracket relations are satisfied. ∎
###### Remark 6.2.
Let $`𝒜=𝒜_n`$ be the braid arrangement in $`^n`$. As noted in Example 2.1, one then has $`M(𝒜_n^{k+1})=F(^{k+1},n)`$ for all $`k`$. For the braid arrangement, the codimension two flats in $`𝖫(𝒜_n)`$ (the partition lattice) are of the forms
$$H_{i,j}H_{i,k}H_{j,k}\text{for}1i<j<kn,\text{and}H_{i,j}H_{k,l}\text{for}\{i,j\}\{k,l\}=\mathrm{}.$$
Thus by Theorem 6.1, for $`1<q<2k+1`$, the homology of $`\mathrm{\Omega }^qF(^{k+1},n)`$, with coefficients in a field $`𝔽`$ of characteristic zero, is generated as a Poisson algebra by elements $`B_{i,j}=\beta _{H_{i,j}}`$ of degree $`2k+1q`$ for $`1i<jn`$. Moreover, the Poisson bracket relations are given by the universal infinitesimal Poisson braid relations:
$`\lambda _{q1}[B_{i,j}+B_{i,k}+B_{j,k},B_{m,k}]`$ $`=0\text{for }m=i\text{ or }m=j\text{, and}`$
$`\lambda _{q1}[B_{i,j},B_{k,l}]`$ $`=0\text{for }\{i,j\}\{k,l\}=\mathrm{}\text{.}`$
As shown in , these are precisely the infinitesimal pure braid relations in case $`q=1`$, see also Examples 4.1 and 5.1.
It seems likely that, via the natural universal mapping property, one could define the “universal infinitesimal Poisson braid algebra,” and that the homology of $`\mathrm{\Omega }^qF(^{k+1},n)`$ with coefficients in a field $`𝔽`$ of characteristic zero is that algebra over $`𝔽`$.
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# 1 Introduction
## 1 Introduction
The persistent disagreement between solar neutrino data and theoretical expectations has been a long-standing problem in physics. Since the very first measurements , the Solar neutrino problem has remained as a puzzle, re-confirmed by new data on rates by GALLEX-SAGE as well as most recently published 825–day data collected by the Super-Kamiokande collaboration which goes beyond the simple rate measurements to include also rate–independent data such as the recoil electron spectra induced by solar neutrino interactions, as well as the zenith angle distributions . It has often been argued that these data can not be accounted for by astrophysics . Together with the atmospheric neutrino data these constitute the only present–day evidence in favour of physics beyond the Standard Model, providing a strong hint for neutrino conversion.
The most popular solutions of the solar neutrino anomalies are based on the idea of neutrino oscillations, either in vacuum or in the Sun due to the enhancement arising from matter effects .
Although these are the simplest neutrino conversion mechanisms there is considerable interest in alternative interpretations. For example it has long been noted that Majorana neutrinos may have non–zero transition magnetic moments which can generate spin–flavour conversions in the presence of a magnetic field. These are especially interesting for two reasons: (i) on general grounds neutrinos are expected to be Majorana particles and (ii) conversions induced by transition magnetic moments can be resonant in the Sun . There is also room for more exotic mechanisms such as flavour changing neutrino interactions which do not even require neutrino mass .
Here we will re-analyse the status of resonant spin–flavour solutions to the solar neutrino problem in the light of the most recent global set of solar neutrino data, including event rates as well as zenith angle distributions and recoil electron spectra induced by solar neutrino interactions in Superkamiokande which has attracted interest recently . In contrast to previous attempts we will adopt the general framework of self–consistent magneto–hydrodynamic (MHD) models of the Sun . A previous attempt in this direction is given in ref.. For definiteness we will concentrate in the recent proposal of Ref. where relatively simple analytic solutions have been given. We perform global fits of solar neutrino data for realistic solutions to the magneto-hydrodynamics equations inside the Sun. This requires adjusting both the neutrino parameters as well as optimizing the magnetic field profile. The arbitrariness associated to the latter is substantially reduced due to mathematics (they must be solutions of MHD equations) as well as reasonable physical requirements. This way and by neglecting neutrino mixing we obtain the simplest MHD-RSF solution to the solar neutrino problem, characterized by two effective parameters, $`\mathrm{\Delta }m^2`$ and $`\mu _\nu B_{max}`$, $`B_{max}`$ being the maximum magnitude of the magnetic field inside the convective region. Throughout this paper we have assumed that the neutrino transition magnetic moment $`\mu _\nu `$ is given in units of $`\mu _{11}\mu _\nu /10^{11}`$ $`\mu _B`$, where $`\mu _B`$ is the Bohr magneton and we set $`\mu _{11}1`$ everywhere. Our MHD-RSF solution can be meaningfully compared with the neutrino oscillation solutions to the solar neutrino problem. We find that our simplest two-parameter MHD-RSF fits to the solar neutrino data are slightly better than those for the oscillation solutions, but not in a statistically significant way. The required best fit points correspond to maximum magnetic field magnitudes in the convective zone smaller than 100 KG. We briefly discuss the prospects to distinguish our simplest MHD-RSF scenario from the neutrino oscillation solutions to the solar neutrino problem at future solar neutrino experiments, giving some predictions for the SNO experiment.
## 2 Static Magnetic Field Profiles in the Sun
In solar magneto-hydrodynamics (MHD, for short) one can explain the origin of solar magnetic fields from the dynamo mechanism at the bottom of the convective zone or, to be more specific, in the overshoot layer, where magnetic fields may be as strong as 300 kG . Such a picture is quite attractive and several MHD dynamo solutions has been known since long time ago (see for example ) However the corresponding magnetic field profiles are rather complicated and difficult to extract. For this reason there have been many attempts to mimic MHD properties through the use of ad hoc magnetic field profiles involving, for example, twisting fields .
Here we will follow an alternative approach using fully self-consistent solutions to the MHD equations inside the Sun. To achieve this we focus on the case of stationary solutions which are known analytically in terms of relatively simple functions . This way we obtain a simple and well-motivated magnetic field profile, without the full complexity that a dynamo model implies. In this section we will explain this model and discuss the limits on the shape parameters describing the field profile. We will also discuss how to relate this model with the dynamo picture of the solar interior.
### 2.1 Single-Mode Field Configurations
In this subsection we will describe the model we are using for the magnetic field profile. We consider only solutions to the equation for a static MHD plasma configuration in a gravitational field given by the equilibrium of the pressure force, the Lorentz force and the gravitational force
$$p\frac{1}{c}\stackrel{}{j}\times \stackrel{}{B}+\rho \mathrm{\Phi }=0,$$
(1)
where p is the pressure, $`j=(c/4\pi )rot\stackrel{}{B}`$ is the electric current, $`B`$ is the static magnetic field under consideration, $`\rho `$ is the matter density and $`\mathrm{\Phi }`$ is the gravitational potential.
This static MHD equations correspond to a quiet Sun and they admit axially symmetric solutions in the spherically symmetric gravitational field which can be simply expressed in terms of spherical Bessel functions and were first discussed in Ref. . For this model the magnetic field will be given by a family of solutions that depends on $`z_k`$, the roots of the spherical Bessel function $`f_{5/2}=\sqrt{z}J_{5/2}(z)`$, to ensure the boundary condition that $`\stackrel{}{B}`$ vanishes on the solar surface. Within the solar interior the magnetic field for any $`k`$ will be then given by
$`B_r^k(r,\theta )`$ $`=`$ $`2\widehat{B}^k\mathrm{cos}\theta \left[1{\displaystyle \frac{3}{r^2z_k\mathrm{sin}z_k}}\left({\displaystyle \frac{\mathrm{sin}(z_kr)}{z_kr}}\mathrm{cos}(z_kr)\right)\right],`$
$`B_\theta ^k(r,\theta )`$ $`=`$ $`\widehat{B}^k\mathrm{sin}\theta \left[2+{\displaystyle \frac{3}{r^2z_k\mathrm{sin}z_k}}\left({\displaystyle \frac{\mathrm{sin}(z_kr)}{z_kr}}\mathrm{cos}(z_kr)z_kr\mathrm{sin}(z_kr)\right)\right],`$
$`B_\varphi ^k(r,\theta )`$ $`=`$ $`\widehat{B}^kz_k\mathrm{sin}\theta \left[r{\displaystyle \frac{3}{rz_k\mathrm{sin}z_k}}\left({\displaystyle \frac{\mathrm{sin}(z_kr)}{z_kr}}\mathrm{cos}(z_kr)\right)\right],`$ (2)
where the coefficient $`\widehat{B}^k(B_{core})`$ is given by
$$\widehat{B}^k=\frac{B_{core}}{2(1z_k/\mathrm{sin}z_k)}.$$
(3)
Here $`\theta `$ is the polar angle and the distance $`r`$ has been normalized to $`R_{}=1`$. Taking into account the inclination of the solar equator to the ecliptics, where neutrinos propagate to the Earth, it follows that $`\theta `$ lies in the narrow range $`83^o97^o`$, depending on the season. In our calculations we have averaged over $`\theta `$ in the above range.
The modulus of the perpendicular component which is relevant to the neutrino spin-flavour takes the form
$$B_{}=\sqrt{B_\varphi ^2+B_\theta ^2}=B_{core}\frac{\mathrm{sin}\theta }{r}f(r),$$
(4)
where $`f(r)`$ is some known smooth function. Notice also that the behaviour of B at the solar center (r = 0)
$`B_r(0,\theta )`$ $`=`$ $`B_{core}\mathrm{cos}\theta ,`$
$`B_\theta (0,\theta )`$ $`=`$ $`B_{core}\mathrm{sin}\theta ,`$
$`B_\varphi (0,\theta )`$ $`=`$ $`B_{core}\mathrm{sin}\theta {\displaystyle \frac{z_k}{2}}{\displaystyle \frac{r}{R_{}}}0,`$ (5)
is completely regular, determined only by the parameter $`B_{core}`$. In Fig. 1 we display the perpendicular component of B for various $`k`$–values 1, 3 and 10, which correspond to the roots $`z_k=5.7`$, $`z_k=12.3`$ and $`z_k=34.5`$, respectively.
### 2.2 Astrophysical Constraints on Magnetic Fields
We now discuss the astrophysical restrictions on the free parameters $`B_{core}`$ and $`k`$ characterizing the model. We can see that the magnitude of a magnetic field at the center of the Sun is constrained by the Fermi-Chandrasekhar limit which implies
$$\eta =\frac{2}{15}\frac{5\gamma _06}{\gamma _01}\frac{(\widehat{B}^k)^2z_{2k}^2R_{}^4}{GM_{}^2}<\mathrm{\hspace{0.33em}1},$$
(6)
here $`\gamma _0`$ is polytropic index characterizing the equation of state (pressure $`P\rho ^{\gamma _0}`$), $`M_{}`$ is the solar mass and $`G`$ is Newton’s constant. This equation gives us an upper bound on $`B_{core}<\mathrm{\hspace{0.33em}2}`$ MGauss for $`k=1`$. For higher values of $`k`$ this constraint is even weaker.
Regarding with the values of $`k`$. These can be constrained by taking into account that in order to justify the use of a stationary solution, it is necessary that the diffusion time due to ohmic dissipation
$$t_{diss}(k)=\frac{4\pi \sigma _{cond}L^2(k)}{c^2}$$
(7)
must be bigger than the age of the Sun $`t_{}1.4\times 10^{17}s`$ . Here $`L(k)`$ denotes a characteristic spatial scale of the magnetic field which corresponds to the typical distance between subsequent nodes of the corresponding Bessel function. As we can see from Fig. 1 $`L(k=1)R_{}`$, while $`L(k)R_{}/k`$. In eq. (7) $`\sigma _{cond}=\omega _p^2/(4\pi \nu _{ep})`$ denotes the conductivity of the fully ionized hydrogen plasma. After substituting the plasma frequency
$$\omega _p=5.65\times 10^4\sqrt{\frac{n_e}{1cm^3}}s^1$$
and the $`e`$-$`p`$ collision frequency
$$\nu _{ep}=50\left(\frac{n_e}{1cm^3}\right)\left(\frac{T}{1K}\right)^{3/2}s^1$$
we obtain from eq. (7) an estimate of the magnetic field dissipation time $`t_{diss}(k)`$
$$t_{diss}(k)=\frac{6.4\times 10^7R_{}^21/s}{c^2}\times \frac{(T/K)^{3/2}}{k^2}>t_{}$$
(8)
Note that the dissipation time is shorter for higher $`k`$ values, long–lived field configurations being possible only if $`k`$ is small. For example, the dissipation time for the third mode is about an order magnitude less than that for the first mode. Moreover it depends on the value of the temperature that we take. Some typical values for the temperature are $`T_{\mathrm{min}}2.8\times 10^5K`$ for the bottom of the convective zone and $`T_{\mathrm{max}}1.6\times 10^7K`$ for the solar core . Thus, taking the optimistic estimate, $`T_{\mathrm{max}}`$, we obtain $`k<k_M=13`$, while, if we consider the average value $`T=4\times 10^6K`$ we will have $`k<k_M=5`$. In what follows we will consider values of $`k10`$.
### 2.3 Energy localization criterium
It is commonly accepted that magnetic fields measured at the surface of the Sun are weaker than within the convective zone interior where this field is supposed to be generated. It is known by observational data that the mean field value over the solar disk is of the order of 1 Gauss while in the solar spots magnetic field strength reaches 1 KG.
On the other hand the general knowledge of the solar magnetic field models is that the magnetic field increases at the overshoot layer, while being small at the solar interior, a picture rather opposite to the one we have seen in Fig. 1.
Although there is no direct information on the magnetic field magnitudes at the solar core, there are theoretical reasons which imply a central magnetic field less than 30 Gauss ; otherwise the present magnetic field in the convective zone would be too big, leading to a visible enhancement in sunspot activity.
This conflict can be avoided by taking advantage of the linear nature of the basic equilibrium MHD equation in eq. (1). This implies that any linear combination of solutions $`\stackrel{}{B}^k`$ ($`k=1,2,\mathrm{},k_M`$, for some fixed number $`k_M`$)
$$\stackrel{}{B}=c_1\stackrel{}{B_1}+c_2\stackrel{}{B_2}+\mathrm{}+c_M\stackrel{}{B_M}$$
(9)
is also a solution. As mentioned in section 2.2 we will adopt $`k_M10`$ in order to ensure that ohmic dissipation is acceptable and therefore justify the static approximation.
In order to ensure that the magnetic field energy is localized mainly within the convective region we will now supplement the constraints of section 2.2 by imposing that the magnetic field should vanish in the center of the Sun <sup>6</sup><sup>6</sup>6 Note that adopting a finite but small value for $`|\stackrel{}{B}(\stackrel{}{r}=0)|<\mathrm{\hspace{0.33em}30}`$ Gauss does not change significantly the profile of the total magnetic field eq. (9) which is the relevant quantity for describing neutrino propagation.
$$\stackrel{}{B}(\stackrel{}{r}=0)=0$$
(10)
The latter implies
$$c_1+c_2+\mathrm{}+c_{k_M}=0.$$
(11)
Therefore we will have, in principle, $`k_M1`$ free parameters.
We will require, in addition, that the magnetic field energy must be minimal in the region below the bottom of the convective zone, characterized by a certain value of $`r_0`$
$$E_B=_0^{r_0}d^3r\frac{\stackrel{}{B}^2(r)}{8\pi }.$$
(12)
This implies
$$\frac{}{c_i}E_B\mathrm{\Sigma }_{j=1}^{k_M1}c_ja_{ij}=0$$
(13)
where
$$a_{ij}=_0^{\stackrel{}{r_0}}d^3r(\stackrel{}{B^j}\stackrel{}{B^M})\stackrel{}{B^i}.$$
(14)
Without loss of generality we can assume that one of the coefficients is non-zero, which prevents us from having only the trivial solution $`c_i0`$. Taking $`c_{k_M1}0`$ we will have the linear non–homogeneous equation
$$\mathrm{\Sigma }_{j=1}^{M2}c_ja_{ij}=c_{k_M1}a_{_{ik_M1}}$$
(15)
which determines all the $`c_i`$ coefficients in terms of, say, $`c_{k_M1}`$. As expected on physical grounds, this last remaining parameter $`c_{k_M1}`$ corresponds to the maximum magnetic field magnitude in the convective region, i.e. $`B_{_{max}}`$ is proportional to $`c_{k_M1}`$. In general, all the coeficients $`c_i`$ will be different from zero for $`i<M`$ and alternate in sign. As an example for $`k_M=6`$ the coefficients are -0.338968, -0.261825, 1.29186, -1.77360, 1.86916, -0.786628.
The procedure sketched above provides a consistent method for combining individual mode solutions $`\stackrel{}{B_k}`$ of the static MHD equation, while fixing all of the coefficients of the linear combination, leaving as free parameters only the value of $`B_{_{max}}`$ inside the convective region and the value of $`k_M10`$. In Fig. (2) we show the resulting combined profiles for $`k_M=5,6,10`$. The parameter $`r_0`$ could also be taken as a free parameter but, on physical grounds, it should lie in a narrow range close to overshoot layer. We show explicitly that varying $`r_0`$ has little effect on our results.
## 3 Fitting the Solar Neutrino Data
The Majorana neutrino evolution Hamiltonian in a magnetic field is well–known to be four–dimensional . For definiteness and simplicity we will neglect neutrino mixing in what follows and consider first the case of active-active neutrino conversions. This will allow us to compare our $`\chi ^2`$-analysis with the previous ones . The $`\nu _e\overline{\nu }_{\mathrm{}}`$ conversions are described by the master Schrödinger evolution equation
$$i\left(\begin{array}{c}\dot{\nu }_e\hfill \\ \dot{\overline{\nu }_{\mathrm{}}}\hfill \end{array}\right)=\left(\begin{array}{cc}V_e\delta & \mu _\nu B_+\\ \mu _\nu B_{}& V_{\mathrm{}}+\delta \end{array}\right)\left(\begin{array}{c}\nu _e\\ \overline{\nu }_{\mathrm{}}\end{array}\right),$$
(16)
where $`\mu _\nu `$ denotes the neutrino transition magnetic moment in units of $`10^{11}`$ $`\mu _B`$, $`\mathrm{}`$ denoting either $`\mu `$ or $`\tau `$. Here $`B_\pm =B_x\pm iB_y`$ and $`\delta =\mathrm{\Delta }m^2/4E`$ is the neutrino mass parameter; $`V_e(t)=G_F\sqrt{2}(\rho (t)/m_p)(Y_eY_n/2)`$ and $`V_{\mathrm{}}(t)=G_F\sqrt{2}(\rho (t)/m_p)(Y_n/2)`$ are the neutrino vector potentials for $`\nu _e`$ and $`\nu _{\mathrm{}}`$ in the Sun given by the abundances of the electron ($`Y_e=m_pN_e(t)/\rho (t)`$) and neutron ($`Y_n=m_pN_n(t)/\rho (t)`$) components. In our numerical study of solar neutrino data we adopt the Standard Solar Model density profile of ref..
We solve Eq. (16) numerically by finding a solution of the Cauchy problem in the form of a set of wave functions $`\nu _a(t)=\nu _a(t)e^{i\mathrm{\Phi }_a(t)}`$ from which the neutrino survival probabilities $`P_{aa}(t)=\nu _a^{}\nu _a`$ are calculated. They obey the unitarity condition $`_aP_{aa}=1`$ where the subscript $`a`$ denotes $`a=e`$ for $`\nu _e`$ and $`a=\mathrm{}`$ for $`\overline{\nu }_{\mathrm{}}`$ respectively.
As an illustration we display in Fig. (3) the electron neutrino survival probablity $`P_{ee}`$ calculated in the MHD-RSF scheme from eq. (16) plotted versus $`E/\mathrm{\Delta }m^2`$. This is obtained with the magnetic field configurations given in Fig. (2).
We will now re-analyse the status of resonant spin–flavour solutions to the solar neutrino problem in the light of the most recent global set of solar neutrino data, including event rates as well as zenith angle distributions and recoil electron spectra induced by solar neutrino interactions in Superkamiokande which has attracted interest recently . It has been found that the quality of the fit to the solar neutrino data depends on the magnetic field profile. The best solutions have been obtained with a magnetic field around 100 KG in the convective zone and zero at the core (profile 3 in Ref , and profile 6 of Ref ) or an almost constant magnitude of the magnetic field, but with twisting direction using the profiles given in .
In contrast with previous work we will consider the fits obtained when we employ self–consistent solutions of MHD equations which obey the physical requirements we derived in sec. 2.2 using the procedure for combining magnetic field modes described in section 3. Our approach is *global* and allows us to compare quantitatively with other solutions to the solar neutrino problem within the same well calibrated theoretical calculation and fit procedure.
### 3.1 Rates
In order to determine the possible values of the parameters characterizing the MHD-RSF solution to the solar neutrino problem, we have first used the data on the total event rates measured at the Chlorine experiment in Homestake , at the two Gallium experiments GALLEX and SAGE and the 825-day Super–Kamiokande data sample, as given in table 1.
In our statistical treatment of the data for the combined fit we adopt the $`\chi ^2`$ definition given in ref. ,
$$\chi _R^2=\underset{i,j=1,3}{}(R_i^{th}R_i^{exp}),\sigma _{ij}^2(R_j^{th}R_j^{exp})$$
(17)
where $`R_i^{th}`$ is the theoretical prediction of the event rate in detector $`i`$ and $`R_i^{exp}`$ is the measured rate. The error matrix $`\sigma _{ij}`$ contains not only the theoretical uncertainties but also the experimental errors, both systematic and statistical.
The general expression of the expected event rate in the presence of oscillations in experiment $`i`$ is given by $`R_i^{th}`$
$`R_i^{th}`$ $`=`$ $`{\displaystyle \underset{k=1,8}{}}\varphi _k{\displaystyle }dE_\nu \lambda _k(E_\nu )\times [\sigma _{e,i}(E_\nu )P_{ee}(E_\nu ,t)`$
$`+\sigma _{x,i}(E_\nu )(1P_{ee}(E_\nu ,t))],`$
where $`E_\nu `$ is the neutrino energy, $`\varphi _k`$ is the total neutrino flux and $`\lambda _k`$ is the neutrino energy spectrum (normalized to 1) from the solar nuclear reaction $`k`$ with the normalization given in Ref. . Here $`\sigma _{e,i}`$ ($`\sigma _{x,i}`$) is the $`\nu _e`$ ($`\nu _x`$) (with $`x`$ being $`\overline{\mu }`$ or $`s`$ corresponding to active-active or active-sterile MHD-RSF conversions) cross section in the Standard Model with the target corresponding to experiment $`i`$, and $`P_{ee}(E_\nu ,t)`$ is the time–averaged $`\nu _e`$ survival probability.
For the Chlorine and Gallium experiments we use improved cross sections $`\sigma _i(E)`$ from Ref. . For the Super–Kamiokande experiment we calculate the expected signal with the corrected cross section given in the Appendix Sec. A.
The expected signal in the absence of oscillations, $`R_i^{\mathrm{BP98}}`$, can be obtained from Eq.(3.1) by substituting $`P_{ee}=1`$. In table 1 we also give the expected rates at the different experiments which we obtain using the fluxes of Ref. .
In Fig. (4) we display the region of MHD-RSF parameters allowed by the solar neutrino rates.
Our $`\chi ^2`$ analysis of the solar neutrino rates uses the magnetic field profiles discussed in Section 2. As we mentioned in that section, these profiles are characterized by $`k_M`$ and $`r_0`$. In table 2 we present the best fit points for $`k_M`$ from 4 to 8 and for $`r_0.6R`$ and for $`B_{_{max}}<300`$KGauss. In the same table we also show the best fit points for $`B_{_{max}}<100`$KGauss. We can see from this table that the $`\chi ^2`$ is pretty stable and does not depend significantly on the choice of $`k_M`$ and $`r_0`$ allowed by astrophysics. In Fig. (4) we display the region of MHD-RSF parameters allowed by the solar neutrino rates for the case $`M=6`$ and $`r_0=.6R`$. We can see that there are several allowed regions for different values of the magnetic field. As we already mentioned, in our analysis we have fixed the value of $`\mu _\nu `$ to be $`10^{11}\mu _B`$. Since the evolution equation depends on the product $`\mu _\nu B`$, for a smaller value of the neutrino magnetic moment the $`B_{_{max}}`$ in Fig. (4) would have to be correspondingly increased. In this sense, the local minima shown in table 2 for $`B_{_{max}}<100`$ KGauss allows a smaller $`\mu _\nu `$.
### 3.2 Zenith and Spectrum Fit
Apart from total event rates the water Cerenkov experiment also measures the zenith angle distribution of solar neutrino events as well as their electron recoil energy spectrum with their recent 825-day data sample .
The smallness of the $`\mathrm{\Delta }m^2`$ values indicated by the rates fit implies that no appreciable day–night variation of the counting rates is expected in our MHD-RSF solution. However the measured solar neutrino zenith angle data must be included in the analysis and we do that. This is necessary in order to enable us a meaningful comparison with vacuum and matter oscillations using the same statistical criteria , see definitions in the appendix. We obtain $`\chi _{zenith}^2=5.4`$ for the full range of parameters in the analysis, the same as for the no-oscillation case.
The recoil electron energy spectrum induced by solar neutrino interactions after 504 days of operation is given for energies above 6.5 MeV using the Low Energy (LE) analysis in which the recoil energy spectrum is divided into 16 bins, 15 bins of 0.5 MeV energy width and the last bin containing all events with energy in the range 14 MeV to 20 MeV. Below 6.5 MeV the background of the LE analysis increases very fast as the energy decreases. Super–Kamiokande has designed a new Super Low Energy (SLE) analysis in order to reject this background more efficiently so as to be able to lower their threshold down to 5.5 MeV. In their 825-day data they have used the SLE method and they present results for two additional bins with energies between 5.5 MeV and 6.5 MeV. In the appendix we present these data in table 7 as well as the details of our statistical analysis. Our results are almost independent of the choice of the parameters $`k_M`$ and $`r_0`$ in the physical range, thus establishing the robustness of the fit procedure. The predicted spectrum is essentially flat except for the upper part of the $`\mathrm{\Delta }m^2`$ region. As an example, we show in fig. 5 the excluded region at 99 % CL for the case $`k_M=6`$ and $`r_0=0.6`$.
### 3.3 Global Fit
As we have seen in the partial analysis, zenith and spectrum are essentially flat in the region of parameters which provide a good fit for the rates–only analysis. For this reason, the allowed regions are slightly modified by the inclusion of the zenith angular dependence and the energy spectrum data. As the results are statistically independent of the choice of $`k_M`$ and $`r_0`$ in the physical range, our analysis effectively involves only two parameters. It is therefore meaningful to compare it with the popular two–neutrino fits characterizing vacuum or matter–enhanced oscillations. In table 3, we show the best–fit points in the range of our study for different $`k_M`$ and $`r_0`$ values. Moreover, we show the local (global) minimum for $`B_{max}`$ less than 100 KG, which will be important to improve sensitivity on the transition magnetic moment of the neutrino. In fig. 5 we show the allowed region at 90% CL and 99% CL for the case $`r_0=0.6`$ and $`k_M=6`$. We have also investigated the effect of varying the hep flux, obtaining for the allowed regions results similar to the no–oscillation solution discussed previously in ref. , independently of the $`\mathrm{\Delta }m^2`$ and $`B_{max}`$ value, with a $`hep`$ normalization factor of 13.5.
We now move to the case of active-sterile MHD-RSF conversions. For this case one must make substitute $`\nu _s`$ for $`\nu _{\mathrm{}}`$ in eq. (16) and take into account that $`V_s=0`$. The results given above for active–active MHD-RSF conversions change when conversions involve sterile neutrinos. The best fit points ( and local ones ) are obtained with parameters slightly modified with respect to those obtained for the active–active case. In the rates only fit, the $`\chi _{rates}^2`$ is worse than for the active-active case, essentially due to the neutral current contribution in the Super–Kamiokande experiment. The zenith angle dependence and the recoil energy spectrum remains flat as before. The global fit for different $`r_0`$ and $`k_M`$, is shown in table 4.
## 4 MHD-RSF versus Oscillation Solutions
### 4.1 Present
From the results of the previous section it follows that our MHD-RSF solution to the solar neutrino problem provides a good description of the most recent solar neutrino data, including event rates as well as zenith angle distributions and recoil electron spectra induced by solar neutrino interactions in Superkamiokande. We have shown that our procedure is quite robust in the sense that the magnetic field profile has been determined in an essentially unique way. This effectively substitutes the neutrino mixing which characterizes the oscillation solutions by a single parameter $`B_{max}`$ characterizing the maximum magnitude of the magnetic field inside the convective region. The value of $`k_M`$ characterizing the maximum number of individual modes superimposed in order to obtain a realistic profile and the parameter $`r_0`$ characterizing the location of the convective region are severely restricted. The allowed $`k_M`$ values are restricted by ohmic dissipation arguments to be lower than 10 or so, while $`r_0`$ is close to $`0.6R_{}`$. We have found that our solar neutrino fits are pretty stable as long as $`k_M`$ exceeds 5 and $`r_0`$ lies in the relevant narrow range (see tables 2 and 3). Therefore our fits are effectively two–parameter fits ($`\mathrm{\Delta }m^2`$ and $`B_{max}`$) whose quality can be meaningfully compared with that of the fits obtained for the favored neutrino oscillation solutions to the solar neutrino problem. In table 5 we compare the various solutions of the solar neutrino problem with the MHD-RSF solutions for the lower magnetic field presented here.
Clearly the MHD-RSF fits seem somewhat better (though not in a statistically significant way) than those obtained for the MSW effect as well as just–so solutions . Notice that in table 5 we have used the same common calibrated theoretical procedures and statistical criteria. These results are for the case where the BP98 Standard Solar Model is adopted. We have also investigated the effect of varying the $`hep`$ flux, obtaining a $`hep`$ normalization factor of 13.5 to be compared with 12 for the SMA solution, 38 for the LMA solution and 15 for the VAC solution.
### 4.2 Future
Having performed our global analysis of the recent solar neutrino data within the framework of our MHD-RSF solution to the solar neutrino problem, we are in a position to calculate also the expected values of a number of observables to be measured by future solar neutrino experiments, such as SNO or Borexino. This task has been developed for the case of oscillation–type solutions to the solar neutrino problem in ref. . Here we will consider our alternative MHD-RSF solution described in sections 2 and 3, because of its theoretical elegance and the good quality of the global fits it provides. Again, the results of refs. will allow us to compare quantitatively our simplest MHD-RSF predictions with those associated with the vacuum (VAC) and MSW–type (SMA, LMA and LOW) solutions to the solar neutrino problem using the same well–calibrated theoretical calculation and fit procedure.
We determine the expected solar neutrino rates at SNO using the cross sections for the CC and NC $`\nu d`$ reactions given by ref. and the preliminary SNO collaboration estimates for the energy resolution, absolute energy scale and detection efficiencies. For definiteness we adopt the most optimistic threshold energy of 5 MeV which should be reached by the collaboration . We perform these calculations at the best–fit points which we have determined in the present paper, using 90 and 99 % CL error bars. For definiteness we have considered the global best fit points and local minima for $`B_{max}<100kG`$ given in table 3, for the case $`k_M=6`$ and $`r_0=0.6`$ and active-active MHD-RSF conversions.
We have calculated the neutral-to-charged-current event ratio (NC/CC for short) and our results are presented in Fig. (7).
Our predictions for the oscillation solutions agree relatively well with those of . The agreement is not perfect because we use the full zenith angle dependence in the analysis of the solar neutrino data instead of simply the day–night asymmetry employed in ref. . The size of the error bars displayed in Fig. (7) arises from the variation in the values of neutrino oscillation parameters, rather than from statistical and theoretical uncertainties, which are negligible .
Clearly from Fig. (7) we see that there is a substantial overlap between our MHD-RSF predictions and those found for each of the oscillation solutions (SMA, LMA, LOW, VAC). The overlap is especially large between the LMA and the MHD-RSF solutions. Taking into account the present theoretical uncertainties and a reasonable estimate of the experimental errors attainable, it follows that an unambiguous discrimination between our MHD-RSF solution and the neutrino oscillation–type solutions to the solar neutrino problem on the basis of the averaged event rates seems rather difficult. The expected features of the MHD-RSF recoil electron spectrum will be discussed elsewhere .
## 5 Discussion & Conclusions
We have re-analysed the status of resonant spin–flavour solutions to the solar neutrino problem in the framework of analytic solutions to the solar magneto-hydrodynamics equations, using the most recent global set of solar neutrino data. We have shown that our procedure is quite robust in the sense that the arbitrariness associated to the magnetic field profile has been almost eliminated due to both mathematical consistency and physical requirements. Effectively our analysis substitutes neutrino mixing by a single parameter $`B_{max}`$ characterizing the maximum magnitude of the magnetic field inside the convective region. The value of $`k_M`$ characterizing the maximum number of individual modes combined in a realistic profile and the parameter $`r_0`$ characterizing the location of the convective region are severely restricted. The allowed $`k_M`$ values are restricted by ohmic dissipation arguments to be lower than 10 or so, and we have found that our solar neutrino fits are pretty stable as long as $`k_M`$ exceeds 5. Moreover our fits are pretty stable within the relevant narrow range for $`r_0`$. This way we obtain effective two–parameter global fits of solar neutrino data for static MHD solutions characterized by $`\mathrm{\Delta }m^2`$ and $`B_{max}`$, since the magnetic field profile is essentially unique. This enables us to compare their quality with that of the fits obtained for the favored neutrino oscillation solutions to the solar neutrino problem. Adopting the Standard Solar Model we have found the MHD-RSF fits to be slightly better than the oscillation fits, though not in a statistically significant way. We have also analysed the prospects to distinguish our best MHD-RSF solution from the oscillation solutions (SMA, LMA, LOW, VAC) at future solar neutrino experiments. Both in the comparison of the present status of different solutions of the solar neutrino problem, as well as in their future predictions at SNO we have used a common well-calibrated theoretical procedure and statistical criteria. Taking into account the present theoretical uncertainties and the expected experimental errors attainable, an unambiguous discrimination between our MHD-RSF solution and the neutrino oscillation–type solutions to the solar neutrino problem at the SNO experiment seems rather difficult. On the other hand better measurements of rate–independent solar neutrino observables such as the day–night asymmetry and seasonality would be potentially useful, since our MHD-RSF predictions differ from the expectations of the oscillation schemes. For example, seasonality is expected to be smaller in our MHD-RSF solution than in MSW or just–so oscillations . On the other hand the day–night asymmetry of the MHD-RSF solution is negligible, in contrast with the MSW solutions .
Note, however, that the *complete* MHD-RSF solution is characterized also by a non–zero neutrino flavour mixing. This gives it the potential to be discriminated from the oscillation–type solutions . The most distinctive signal expected in this case consists of solar anti-neutrinos, which would provide a clear signal in water Cerenkov experiments . Moreover, for large enough neutrino mixing one expects also a sizeable suppression of the rates for $`pp`$ neutrinos, potentially testable at the GNO experiment. Last but not least, the possible time dependence of the charged current signal due to solar cycles still remains as a possible tool to discriminate the MHD-model from the oscillation schemes.
Note added: As we finished our paper there appeared the paper E. K. Akhmedov and J. Pulido, hep-ph/0005173, which also considers predictions for SNO observables in the RSF scheme employing the phenomenological magnetic field profiles they used previously in ref .
## Appendix A Zenith and Spectrum Data Samples and Fit Procedures
Here we summarize here the data used and the fit procedures adopted in this paper.
The zenith dependence data given by the Super–Kamiokande collaboration are shown in table 6.
The recoil electron spectrum data are given as
In table 7 $`\sigma _{i,stat}`$ is the statistical error, $`\sigma _{i,exp}`$ is the error due to correlated experimental errors, $`\sigma _{i,cal}`$ is the error due to the calculation of the expected spectrum, and $`\sigma _{i,uncorr}`$ is due to uncorrelated systematic errors.
In our study we use the experimental results from the Super–Kamiokande Collaboration on the recoil electron spectrum on the 18 energy bins including the results from the LE analysis for the 16 bins above 6.5 MeV and the results from the SLE analysis for the two low energy bins below 6.5 MeV, shown in table 7.
Notice that in table 7 we have symmetrized the errors to be included in our $`\chi ^2`$ analysis. We have explicitly checked that the exclusion region is very insensitive to this symmetrization. We define $`\chi ^2`$ for the spectrum as
$$\chi _S^2=\underset{i,j=1,18}{}(\alpha _{sp}\frac{R_i^{th}}{R_i^{\mathrm{BP98}}}R_i^{exp})\sigma _{ij}^2(\alpha _{sp}\frac{R_j^{th}}{R_j^{\mathrm{BP98}}}R_j^{exp})$$
(19)
where
$$\sigma _{ij}^2=\delta _{ij}(\sigma _{i,stat}^2+\sigma _{i,uncorr}^2)+\sigma _{i,exp}\sigma _{j,exp}+\sigma _{i,cal}\sigma _{j,cal}$$
(20)
Again, we introduce a normalization factor $`\alpha _{sp}`$ in order to avoid double-counting with the data on the total event rate which is already included in $`\chi _R^2`$. Notice that in our definition of $`\chi _S^2`$ we introduce the correlations amongst the different systematic errors in the form of a non-diagonal error matrix in analogy to our previous analysis of the total rates. These correlations take into account the systematic uncertainties related to the absolute energy scale and energy resolution.
The general expression of the expected rate in the presence of oscillations $`R^{th}`$ in a bin, is given from Eq.(3.1) but integrating within the corresponding electron recoil energy bin and taking into account that the finite energy resolution implies that the measured kinetic energy $`T`$ of the scattered electron is distributed around the true kinetic energy $`T^{}`$ according to a resolution function $`Res(T,T^{})`$ of the form
$$Res(T,T^{})=\frac{1}{\sqrt{2\pi }s}\mathrm{exp}\left[\frac{(TT^{})^2}{2s^2}\right],$$
(21)
where
$$s=s_0\sqrt{T^{}/\mathrm{MeV}},$$
(22)
and $`s_0=0.47`$ MeV for Super–Kamiokande . On the other hand, the distribution of the true kinetic energy $`T^{}`$ for an interacting neutrino of energy $`E_\nu `$ is dictated by the differential cross section $`d\sigma _\alpha (E_\nu ,T^{})/dT^{}`$, that we take from . The kinematic limits are
$$0T^{}\overline{T}^{}(E_\nu ),\overline{T}^{}(E_\nu )=\frac{E_\nu }{1+m_e/2E_\nu }.$$
(23)
For assigned values of $`s_0`$, $`T_{\mathrm{min}}`$, and $`T_{\mathrm{max}}`$, the corrected cross section $`\sigma _\alpha (E)`$ $`(\alpha =e,x)`$ is given as
$$\sigma _\alpha (E_\nu )=_{T_{\mathrm{min}}}^{T_{\mathrm{max}}}𝑑T_0^{\overline{T}^{}(E_\nu )}𝑑T^{}Res(T,T^{})\frac{d\sigma _\alpha (E_\nu ,T^{})}{dT^{}}.$$
(24)
Acknowledgements
We thank Alexei Bykov, Vladimir Kutvitsky, Dmitri Sokoloff and Victor Popov for useful discussions. This work was supported by DGICYT grant PB98-0693, by the European Commission under Intas Project 96-0659 and TMR contract ERBFMRX-CT96-0090, and by an Iberdrola research excellence grant. VBS and TIR were partially supported by the RFBR grant 00-02-16271, CPG was supported by the Generalitat Valenciana grant GV99-3-1-01 and OGM was supported by the CONACyT-Mexico grant J32220-E.
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# 1 Leading diagram and first subleading corrections.
CERN-TH/2000-139
LARGE-N EXPANSION, CONFORMAL FIELD THEORY
AND RENORMALIZATION-GROUP FLOWS
IN THREE DIMENSIONS
Damiano Anselmi
CERN, Theory Division, CH-1211, Geneva 23, Switzerland
Abstract
I study a class of interacting conformal field theories and conformal windows in three dimensions, formulated using the Parisi large-$`N`$ approach and a modified dimensional-regularization technique. Bosons are associated with composite operators and their propagators are dynamically generated by fermion bubbles. Renormalization-group flows between pairs of interacting fixed points satisfy a set of non-perturbative $`g1/g`$ dualities. There is an exact relation between the beta function and the anomalous dimension of the composite boson. Non-Abelian gauge fields have a non-renormalized and quantized gauge coupling, although no Chern–Simons term is present. A problem of the naive dimensional-regularization technique for these theories is uncovered and removed with a non-local, evanescent, non-renormalized kinetic term. The models are expected to be a fruitful arena for the study of odd-dimensional conformal field theory.
CERN-TH/2000-139 – May, 2000
Unexpected remnants of the renormalization algorithm in quantum field theory are the Adler–Bell–Jackiw anomalies , finite amplitudes arising from the quantum violation of classical conservation laws. Anomalies fall into two main classes: axial and trace. The axial anomalies obey all-order properties, such as the Adler–Bardeen theorem , and give important information about the low-energy physics, by means of the ’t Hooft anomaly-matching conditions . The trace anomaly is related to the beta function , by a formula $`\mathrm{\Theta }=\beta ^a𝒪^a`$, $`𝒪^a`$ being composite operators. Certain trace anomalies in external fields can be computed exactly in the IR limit of (supersymmetric) UV-free theories , where an exact beta function can also be derived . They reveal the intrinsic irreversibility of the renormalization-group flow, its relation to the invariant area of the graph of the beta function between the fixed points and the essential difference between marginal and relevant deformations .
Most of these powerful results apply only to even dimensions. Trace anomalies in external gravitational and flavour fields do not exist in odd dimensions. Nevertheless, an odd-dimensional formula for the irreversibility of the RG flow can in principle be written , because the relation $`\mathrm{\Theta }=\beta ^a𝒪^a`$ is completely general and so is the notion of invariant area of the graph of the beta function. It would be desirable to dispose of a web of non-trivial conformal field theories, conformal windows and RG flows in three dimensions to investigate these and related issues more closely. The purpose of this letter is to construct a large class of such theories and flows, and address the search for appropriate odd-dimensional generalizations of the properties mentioned above.
The beta functions of the most general power-counting renormalizable three-dimensional theory with a Chern–Simons vector field have been studied by Avdeev et al. in ref. . Three-dimensional quantum field theory is relevant for its possible applications in the domain of condensed-matter physics. However, the Chern–Simons models are parity-violating and this somewhat limits the range of their applicability. The three-dimensional $`\varphi ^6`$ theory is known to have a non-trivial fixed point in the large-$`N`$ expansion and a conformal window interpolating between the free limit and this point . Nevertheless, a large class of parity-preserving conformal windows is not known at present and will be constructed here.
The Chern–Simons coupling $`g_{\mathrm{cs}}`$ is not renormalized . The simplest argument to prove this fact proceeds as follows. Let us denote by $`\beta _{\mathrm{cs}}`$ the beta function of $`g_{\mathrm{cs}}.`$ The results of refs. , relating the trace anomaly to the beta functions, imply that, in our case, $`\mathrm{\Theta }`$ should contain a term proportional to the Chern–Simons form, multiplied by $`\beta _{\mathrm{cs}}`$. However, $`\mathrm{\Theta }`$ is gauge-invariant, while the Chern–Simons form is not. For this reason $`\beta _{\mathrm{cs}}`$ has to be identically zero. This kind of argument, essentially based on the properties of the trace anomaly, will be applied several times in this paper.
It was shown in ref. that the Chern–Simons coupling can be used to split the zeros of the beta function and generate a variety of non-trivial conformal windows. For example, the beta function of a $`\overline{\phi }\phi \overline{\psi }\psi `$-coupling with constant $`\eta `$ typically reads
$$\beta _\eta =a(\eta +bg_{\mathrm{cs}}^2)(\eta ^2g_{\mathrm{cs}}^4)$$
to the lowest order, $`a`$ and $`b`$ being some factors, possibily depending on the gauge group and the representation. The coupling $`g_{\mathrm{cs}}`$ is quantized in the non-Abelian case, $`g_{\mathrm{cs}}^2=1/N`$, and is arbitrarily small in the large-$`N`$ limit. Therefore, the existence of interacting fixed points at $`\eta =\pm g_{\mathrm{cs}}^2`$ and $`\eta =bg_{\mathrm{cs}}^2`$ is proved, in this limit, directly from perturbation theory. This construction is a three-dimensional analogue of the existence proof of a conformal window in QCD. There we have, to two loops,
$$\beta _{\mathrm{QCD}}=\beta _1\alpha ^2+\beta _2\alpha ^3+𝒪(\alpha ^5),\beta _1=\frac{1}{6\pi }(11N_c2N_f),\beta _2=\frac{25N_c^2}{(4\pi )^2},$$
where $`\beta _2`$ is written for $`\beta _1N_c`$ and $`N_c`$ large. The role of $`g_{\mathrm{cs}}`$ is here played by $`\beta _1/\beta _21/N_c`$. We see that all these constructions involve a large-$`N`$ limit of some sort. Our models will not be an exception in this respect.
The successful removal of divergences in quantum field theory is not restricted to the power-counting renormalizable theories. Non-renormalizable models in less than four dimensions were quantized long ago by Parisi, using a large-$`N`$ expansion . The four-fermion model has been studied in detail , and the technique has been applied to other cases, such as the $`S_{N1}`$ non-linear $`\sigma `$-model and the $`CP^{N1}`$ model . A challenging, open problem in quantum field theory is to classify the set of power-counting non-renormalizable theories that can be constructed in a perturbative sense, i.e. the appropriate generalization of the power-counting criterion .
For the purposes of this paper, the Parisi large-$`N`$ expansion is a powerful tool to construct non-trivial conformal field theories and conformal windows in three dimensions. The known four-fermion models are relevant perturbations of a certain subclass of these fixed points. Our models are power-counting non-renormalizable, because although they do not contain dimensionful parameters, certain bosonic fields do not have a propagator at the classical level. Such fields are associated with composite operators and can be scalars, but also Abelian and non-Abelian gauge vectors. The propagators are dynamically generated by fermion loops and the large-$`N`$ expansion is crucial to justify the resummation of fermion bubbles before the other diagrams, which are subleading.
To some extent, the construction presented here is a simple application of the general theory of Parisi, however formulated in a new way, which singles out the conformal properties and is more suitable to the research program that we have in mind. More importantly, I generate a whole class of RG flows (marginal deformations) interpolating between the conformal fixed points and show that they satisfy a remarkable set of non-perturbative strong–weak coupling dualities, also exhibited by an exact relation between the beta function and the anomalous dimension of the composite field. The non-Abelian gauge coupling is non-renormalized and has a discrete set of values. Observe that our theories do not contain a Chern–Simons term. I give a general argument proving the non-renormalization theorem, based on the trace anomaly.
I work in the Euclidean framework and use a modified dimensional-regularization technique. The naive dimensional technique is indeed not applicable to the theories studied here, nor to the more familiar four-fermion models, because the dynamically generated propagator does not regularize correctly. It is necessary to add a peculiar non-local term $`_{\mathrm{non}\mathrm{loc}}`$ to the classical lagrangian. This term does not generate new renormalization constants and is evanescent, therefore formally absent in $`D=3`$.
I start from the four-fermion model, written in terms of an auxiliary field $`\sigma `$:
$$_N=\underset{i=1}{\overset{N}{}}\overline{\psi }^i(/+\lambda \sigma )\psi ^i+\frac{1}{2}M\sigma ^2.$$
(1)
This theory was constructed rigorously in , where the existence of an interacting UV fixed point was established. A detailed study can be found in . There are two phases, and the chiral symmetry can be dynamically broken. The $`\sigma `$-field equation gives $`\sigma =\lambda \overline{\psi }\psi /M`$, whence the name “composite boson” for $`\sigma `$.
The theory is well-defined also if we set $`M=0`$. The model
$$=\underset{i=1}{\overset{N}{}}\overline{\psi }^i(/+\lambda \sigma )\psi ^i$$
(2)
is conformal both at the classical and quantum levels, as we now prove. We call it the $`\sigma _N`$ conformal field theory. At the classical level no scale is present. The renormalized lagrangian has the form
$$=Z_\psi \overline{\psi }/\psi +\lambda _\mathrm{B}Z_\sigma ^{1/2}Z_\psi \overline{\psi }\sigma \psi +_{\mathrm{non}\mathrm{loc}}.$$
$`_{\mathrm{non}\mathrm{loc}}`$ denotes the evanescent term to be discussed below. No $`\sigma ^3`$-term is generated by renormalization, because of the symmetry $`x_1x_1`$, $`\psi \gamma _1\psi `$, $`\sigma \sigma `$. The quadratic terms in $`\sigma `$ are also absent: i) the mass term $`M\sigma ^2`$ is not generated, because it is absent in the classical lagrangian and we can choose a subtraction scheme such that the cut-off appears only logarithmically in the quantum action; ii) no local kinetic term for $`\sigma `$ can be generated, since the field $`\sigma `$ has dimension 1 in $`D=3`$.
In general, the bare coupling can be written as $`\lambda _\mathrm{B}=\lambda Z_\lambda \mu ^{\epsilon /2}`$. However, the number of independent renormalization constants is equal to the number of independent fields and therefore we can interpret two $`Z`$’s as the wave-function renormalization constants of $`\psi `$ and $`\sigma `$, and set $`Z_\lambda 1`$. This ensures that $`\beta _\lambda 0`$ in $`D=3`$ and proves that the theory is conformal also at the quantum level. At the level of the trace anomaly, conformality (i.e. $`\mathrm{\Theta }0`$) follows from the fact that all local dimension-3 operators are proportional to the field equations.
The dynamical $`\sigma `$ kinetic term is generated by diagram (a), which, expanded around three dimensions, gives
$$(\mathrm{a})=\frac{N\lambda _\mathrm{B}^2}{(4\pi )^{D/2}}\frac{\mathrm{\Gamma }\left(2D/2\right)\mathrm{\Gamma }^2\left(D/21\right)}{\mathrm{\Gamma }\left(D2\right)}(k^2)^{D/21}=\frac{\lambda _\mathrm{B}^2N}{8}(k^2)^{(1\epsilon )/2}+𝒪(\epsilon ).$$
We fix the normalization with
$$\lambda ^2N=8+𝒪(1/N),$$
(3)
in $`D=3`$ and find, in momentum space,
$$\mathrm{\Gamma }_{\mathrm{kin}}[\sigma ]=\frac{1}{2}|\sigma (k)|^2\mu ^\epsilon (k^2)^{(1\epsilon )/2}+\frac{1}{2}M\sigma ^2.$$
(4)
From the diagrammatic point of view, the reader might find it easier to imagine that the mass $`M`$ is still non-zero, but small, and set it to zero at the end. In particular, at $`M0`$ it is immediate to resum the geometric series of the bubbles of type (a) (see Fig. 1). After inverting the $`\sigma `$ kinetic term and finding the propagator
$$\sigma (k)\sigma (k)=\frac{1}{M}\underset{L=0}{\overset{\mathrm{}}{}}(1)^L\frac{\mu ^{L\epsilon }(k^2)^{L(1\epsilon )/2}}{M^L}=\frac{1}{\mu ^\epsilon (k^2)^{(1\epsilon )/2}+M},$$
(5)
$`M`$ can be freely set to $`0`$, which we assume from now on. We see that the propagator of the $`\sigma `$-field is proportional to $`1/\sqrt{k^2}`$ in $`D=3`$. The propagator (5), however, does not regularize the theory properly, because it goes to zero too slowly at high energies. This fact becomes apparent in the calculations of the subleading corrections. Consider the example of diagram (b), where the dashed line is meant to be the $`\sigma `$-propagator (5). The integral
$$\frac{\mathrm{d}^{3\epsilon }p(p/+k/)}{(p+k)^2(p^2)^{(1\epsilon )/2}}$$
produces a $`\mathrm{\Gamma }(0)`$. The same holds for diagram (c). This phenomenon is very general and concerns theories of composite bosons in every dimension, and in particular the logarithmically trivial $`D=4`$ four-fermion models considered by Wilson in . We conclude that the naive dimensional-regularization procedure fails to regularize our theories.
The problem can be cured by giving a classical, but evanescent, kinetic term to the composite field $`\sigma `$, which at the leading order reads
$$_{\mathrm{non}\mathrm{loc}}=\frac{1}{2}|\sigma (k)|^2\sqrt{k^2}\left[1\frac{\lambda _\mathrm{B}^2N}{8}(k^2)^{\epsilon /2}\right].$$
(6)
$`_{\mathrm{non}\mathrm{loc}}`$ is renormalization-group invariant. This requirement is essential for an easier study of the theory. The new $`\mathrm{\Gamma }_{\mathrm{kin}}`$ is obtained by adding (6) to the old one, namely (4), henceforth producing the desired high-energy behaviour:
$$\mathrm{\Gamma }_{\mathrm{kin}}^{}[\phi ]=\frac{1}{2}|\sigma (k)|^2\sqrt{k^2},$$
(7)
which regularizes the theory correctly. It is easy to go through the usual proofs of renormalizability and locality of the counterterms with the improved dimensional technique.
In $`x`$-space we find, in $`D=3`$,
$$\sigma (x)\sigma (0)=\frac{1}{2\pi ^2|x|^2}.$$
(8)
This two-point function is intrinsically non-perturbative, since it equals the two-point function of an elementary scalar field with anomalous dimension $`+1/2`$.
The field $`\psi `$ has dimension $`\left(D1\right)/2`$ and (6) attributes exactly the same dimension to $`\sigma `$. Taking the $`\mu `$-derivative of the equation $`\lambda _\mathrm{B}=\lambda \mu ^{\epsilon /2},`$ we get
$$\beta (\lambda )=\frac{\mathrm{d}\lambda }{\mathrm{d}\mathrm{ln}\mu }=\frac{\epsilon }{2}\lambda .$$
Integrating the defining relation
$$\gamma _\sigma (\lambda )=\frac{1}{2}\frac{\mathrm{d}Z_\sigma (\lambda ,\epsilon )}{\mathrm{d}\mathrm{ln}\mu },$$
we get the $`\sigma `$-wave-function renormalization constant :
$$Z_\sigma (\lambda ,\epsilon )=\mathrm{exp}\left(\frac{4}{\epsilon }_0^\lambda \frac{\gamma _\sigma (\lambda ^{})}{\lambda ^{}}d\lambda ^{}\right).$$
We assume that we work in the minimal subtraction scheme. We want to find a closed expression for $`_{\mathrm{non}\mathrm{loc}}`$ that properly includes the subleading corrections. The requirements are that $`_{\mathrm{non}\mathrm{loc}}`$ be renormalization-group invariant and evanescent. An expression for $`_{\mathrm{non}\mathrm{loc}}`$ satisfying these properties reads, in momentum space,
$$_{\mathrm{non}\mathrm{loc}}=\frac{1}{2}|\sigma (k)|^2\sqrt{k^2}\left[1\frac{\lambda _\mathrm{B}^2N}{8}(k^2)^{\epsilon /2}\right]\mathrm{exp}\left(\frac{4}{\epsilon }_\lambda ^{\lambda _\mathrm{B}(k^2)^{\epsilon /4}}\frac{\gamma _\sigma (\lambda ^{})}{\lambda ^{}}d\lambda ^{}\right).$$
This formula is essentially unique, the alternatives differring by scheme redefinitions. Renormalization-group invariance is exhibited by rewriting $`_{\mathrm{non}\mathrm{loc}}`$ as
$$_{\mathrm{non}\mathrm{loc}}=\frac{1}{2}|\sigma _\mathrm{B}(k)|^2\sqrt{k^2}\left[1\frac{\lambda _\mathrm{B}^2N}{8}(k^2)^{\epsilon /2}\right]\mathrm{exp}\left(\frac{4}{\epsilon }_0^{\lambda _\mathrm{B}(k^2)^{\epsilon /4}}\frac{\gamma _\sigma (\lambda ^{})}{\lambda ^{}}d\lambda ^{}\right),$$
where $`\sigma _\mathrm{B}=\sigma Z_\sigma ^{1/2}`$. It is easy to prove that $`_{\mathrm{non}\mathrm{loc}}`$ is zero in $`D=3`$. Indeed, we have in the $`\epsilon 0`$ limit:
$$_{\mathrm{non}\mathrm{loc}}\frac{1}{2}|\sigma (k)|^2\sqrt{k^2}\left[1\frac{\lambda _\mathrm{B}^2N}{8}(k^2)^{\epsilon /2}\right]\left(\frac{\mu ^2}{k^2}\right)^{\gamma _\sigma (\lambda )}0.$$
A straightforward application of the Callan–Symanzik equations shows that the $`\sigma `$-two-point function has the form
$$\mathrm{\Gamma }_{\sigma \sigma }=A(\lambda )\sqrt{k^2}\left(\frac{\mu ^2}{k^2}\right)^{\gamma _\sigma (\lambda )},$$
(9)
or, in $`x`$-space,
$$\sigma (x)\sigma (0)=\frac{A^{}(\lambda )}{|x|^{2+2\gamma _\sigma (\lambda )}\mu ^{2\gamma _\sigma (\lambda )}}.$$
(10)
The numerical coefficients $`A(\lambda )`$ and $`A^{}(\lambda )`$ do not have here a direct physical meaning, because they are scheme-dependent and can be changed by redefining $`\mu `$.
Formulas (9) and (10) have the expected form for a conformal field theory. A non-vanishing anomalous dimension $`\gamma _\sigma (\lambda )`$ proves that the theory is interacting. We now calculate $`\gamma _\sigma (\lambda )`$ to the lowest order. We find, from diagrams (b) and (c),
$$Z_\psi =1\frac{\lambda ^2}{6\pi ^2\epsilon },Z_\sigma =1+\frac{4\lambda ^2}{3\pi ^2\epsilon }$$
(11)
and the anomalous dimensions are
$$\gamma _\psi =\frac{2}{3N\pi ^2},\gamma _\sigma =\frac{16}{3N\pi ^2}.$$
These values are in agreement with the calculations of (they can be checked using the formulas (2.35a-b) of that paper, after replacing $`N`$ with $`N/2,`$ since the authors of use doublets of complex spinors). Higher-order corrections have been studied by Gracey in refs. . It is important to remark that $`\gamma _\sigma `$ is negative. A negative anomalous dimension for the composite boson is not in contradiction with unitarity. We have already observed that the uncorrected $`\sigma `$-dimension is 1/2-larger than the minimum. The unitarity bound is therefore $`d_\sigma =1+\gamma _\sigma >1/2`$ or $`\gamma _\sigma >1/2`$, so that $`\gamma _\sigma `$ is allowed to have negative values in three dimensions. Observe that $`\gamma _\psi `$ is instead positive and could not be otherwise for a similar reason. To the first subleading order we have therefore the $`x`$-space correlator
$$\sigma (x)\sigma (0)=\frac{1}{2\pi ^2|x|^{232/(3N\pi ^2)}}.$$
Summarizing, we have formulated, via a large-$`N`$ expansion and an improved dimensional-regularization technique, a class of interacting conformal field theories in three dimensions. These theories are in general strongly coupled. They become weakly coupled for $`N`$ large, and free for $`N=\mathrm{}`$.
Now, we want to define renormalization-group flows interpolating between the $`\sigma _{N+M}`$ and the $`\sigma _N`$ conformal field theories. Let us consider the lagrangian
$$_{NM}=\underset{i=1}{\overset{M}{}}\overline{\chi }^i(/+g\sigma )\chi ^i+\underset{i=1}{\overset{N}{}}\overline{\psi }^i(/+\lambda \sigma )\psi ^i.$$
(12)
Here we expand perturbatively in $`g`$, or actually $`\overline{g}=g/\lambda `$. For $`\overline{g}=0`$ we have the $`\sigma _N`$ model plus $`M`$ free fermions. For $`\overline{g}=1`$ we have the $`\sigma _{N+M}`$ model. It is therefore natural to expect that the coupling $`\overline{g}`$ interpolates between the two fixed points. We can show that there is a non-trivial beta function by studying the first perturbative corrections. We combine the small-$`\overline{g}`$ expansion with the large-$`N`$ expansion. We also assume that $`\overline{g}^2M/N1`$. Since $`\overline{g}`$ varies from 0 to 1, this means that $`M`$ is much smaller than $`N`$. The renormalized lagrangian reads
$$_R=\underset{i=1}{\overset{M}{}}Z_\chi \overline{\chi }^i(/+gZ_{\overline{g}}Z_\sigma ^{1/2}\sigma )\chi ^i+\underset{i=1}{\overset{N}{}}Z_\psi \overline{\psi }^i(/+\lambda Z_\sigma ^{1/2}\sigma )\psi ^i+_{\mathrm{non}\mathrm{loc}}$$
and the evanescent, renormalization-group invariant, non-local kinetic term reads, in the general case:
$$_{\mathrm{non}\mathrm{loc}}=\frac{1}{2}|\sigma (k)|^2\sqrt{k^2}\left[1\frac{\lambda _\mathrm{B}^2N}{8}(k^2)^{\epsilon /2}\right]\mathrm{exp}\left(2_{\mathrm{ln}\mu }^{\mathrm{ln}\sqrt{k^2}}\gamma _\sigma (\mathrm{ln}\mu ^{})\mathrm{d}\mathrm{ln}\mu ^{}\right).$$
From the results (11) we easily get, to the lowest order,
$$Z_\psi =1\frac{4\mu ^\epsilon }{3N\pi ^2\epsilon },Z_\chi =1\frac{4\overline{g}^2\mu ^\epsilon }{3N\pi ^2\epsilon },Z_\sigma =1+\frac{32\mu ^\epsilon }{3N\pi ^2\epsilon },Z_{\overline{g}}=1+\frac{16\left(\overline{g}^21\right)\mu ^\epsilon }{3N\pi ^2\epsilon }$$
We therefore obtain
$$\beta _{\overline{g}}=\frac{16}{3N\pi ^2}\overline{g}\left(\overline{g}^21\right)+𝒪(\overline{g}/N^2,\overline{g}^5/N)$$
(13)
and conclude that the $`\sigma _N`$ model plus $`M`$ decoupled fermions is the UV limit of the flow and the $`\sigma _{N+M}`$ point is the IR limit. Remarkably, the first orders in $`\overline{g}`$ single out correctly both fixed points. This means that, presumably, every truncation of the perturbative expansion of $`\beta _{\overline{g}}`$ factorizes the expected $`\overline{g}\left(\overline{g}^21\right)`$. We show below that this is indeed the case. The theories with couplings $`\overline{g}`$ and $`\overline{g}`$ are clearly equivalent.
The flows (12) satisfy a natural strong–weak coupling duality, associated with the replacement $`\overline{g}1/\overline{g},`$ $`NM`$. The dual flow interpolates from the UV $`\sigma _M`$ model with $`N`$ free fermions to the IR $`\sigma ^{N+M}`$ model. Pairs of dual flows have the IR limits in common. Finally, the self-dual flow has $`N=M`$. We immediately realize that the $`\sigma _M`$ model plus $`N`$ free fermions is the fixed point at $`\overline{g}=\mathrm{}`$. It is natural to conjecture that the points $`\overline{g}=0,1,\mathrm{}`$ are all the fixed points of the exact beta function. The dual flows are plotted in Fig. 2.
The mentioned duality and fixed points are non-perturbative properties of the flows and are self-evident from the construction. We have already seen that, unexpectedly, the lowest order beta function (13), calculated for $`\overline{g}1`$, vanishes at $`\overline{g}=1`$. What is even more astonishing is that, with a little improvement, the beta function vanishes also at $`\overline{g}=\mathrm{}`$ and satisfies the mentioned duality exactly. To see this, let us relax the assumption $`\overline{g}^2M/N1`$, so that $`M`$ and $`N`$ can be of the same order. Diagram (a) is proportional to $`N+M\overline{g}^2.`$ The above formulas can be corrected replacing $`N`$ by $`N+M\overline{g}^2`$. In particular, the lowest-order beta function (13) becomes
$$\beta _{\overline{g}}=\frac{16}{3\pi ^2}\frac{\overline{g}\left(\overline{g}^21\right)}{\left(N+M\overline{g}^2\right)},$$
and does satisfy the $`\overline{g}1/\overline{g},`$ $`NM`$ duality, because
$$\beta _{1/\overline{g}}=\frac{16}{3\pi ^2}\frac{1/\overline{g}\left(1/\overline{g}^21\right)}{\left(M+N/\overline{g}^2\right)}.$$
The remarkable perturbative features that we have outlined are explained by an exact relation between the beta function and the anomalous dimension of $`\sigma `$, that we now derive. This formula is a sort of three-dimensional analogue of certain common formulas in four-dimensional supersymmetric theories, such as the NSVZ beta function , or the beta function of the superpotential coupling. We stress that in three dimensions we do not need supersymmetry for this.
We write the renormalized lagrangian in a manifestly dual form:
$`_R`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{M}{}}}V(g,M;\lambda ,N;\epsilon )\overline{\chi }^i(/+gU(g,M;\lambda ,N;\epsilon )\sigma )\chi ^i`$
$`+{\displaystyle \underset{i=1}{\overset{N}{}}}V(\lambda ,N;g,M;\epsilon )\overline{\psi }^i(/+\lambda U(\lambda ,N;g,M;\epsilon )\sigma )\psi ^i.`$
We have
$`Z_\chi `$ $`=`$ $`V(g,M;\lambda ,N;\epsilon ),Z_\psi =V(\lambda ,N;g,M;\epsilon ),`$
$`Z_\sigma `$ $`=`$ $`U^2(\lambda ,N;g,M;\epsilon ),Z_g={\displaystyle \frac{U(g,M;\lambda ,N;\epsilon )}{U(\lambda ,N;g,M;\epsilon )}},`$
and find
$$\beta _{\overline{g}}=\overline{g}\left(\gamma _\sigma \stackrel{~}{\gamma }_\sigma \right)\overline{g}\left[\gamma _\sigma (\overline{g},M;N)\gamma _\sigma (1/\overline{g},N;M)\right].$$
(14)
Observe that $`\gamma _\sigma (1,M;N)=\gamma _\sigma (M+N).`$ We can immediately check the duality of the exact beta function:
$$\beta _{1/\overline{g}}=\frac{1}{\overline{g}}\left[\gamma _\sigma (1/\overline{g},N;M)\gamma _\sigma (\overline{g},M;N)\right].$$
(15)
The beta function vanishes at the fixed points $`\overline{g}=0,\mathrm{}`$ and the solutions of
$$\gamma _\sigma (\overline{g},M;N)=\gamma _\sigma (1/\overline{g},N;M).$$
(16)
Using the fact that $`\gamma _\sigma (1,M;N)=\gamma _\sigma (M+N)`$ we know that $`\overline{g}=1`$ is a solution. We expect that this is the unique solution of the condition (16).
The trace anomaly reads
$$\mathrm{\Theta }=\beta _{\overline{g}}\sigma \underset{i=1}{\overset{M}{}}\overline{\chi }^i\chi ^i\beta _{\overline{g}}𝒪$$
and, correctly, does not vanish using the field equations.
More generally, we can consider the model
$$=\underset{i=1}{\overset{k}{}}\underset{j=1}{\overset{N_i}{}}V(\lambda _i,N_i;\lambda ,N;\epsilon )\overline{\psi }_{(i)}^j(/+\lambda _iU(\lambda _i,N_i;\lambda ,N;\epsilon )\sigma )\psi _{(i)}^j,$$
where the argument $`\lambda ,N`$ in $`(\lambda _i,N_i;\lambda ,N;\epsilon )`$ refers to the set of couples $`\lambda _j,N_j`$ with $`ji`$. Clearly, $`V(\lambda _i,N_i;\lambda ,N;\epsilon )`$ and $`U(\lambda _i,N_i;\lambda ,N;\epsilon )`$ are symmetric with respect to the exchanges $`\lambda _j,N_j\lambda _l,N_l`$ with $`j,li`$. Choosing $`\lambda _k`$ to be of order unity and all the other $`\lambda `$’s small, we have
$`\beta _i`$ $`=`$ $`\overline{\lambda }_i\left(\gamma _\sigma \stackrel{~}{\gamma }_\sigma ^{(i)}\right)`$
$``$ $`\overline{\lambda }_i[\gamma _\sigma (\overline{\lambda }_1,N_1;\mathrm{}\lambda _i,N_i\mathrm{};\overline{\lambda }_{k1},N_{k1};N_k)`$
$`\gamma _\sigma (\overline{\lambda }_1/\overline{\lambda }_i,N_1;\mathrm{}1/\overline{\lambda }_i,N_k\mathrm{};\overline{\lambda }_{k1}/\overline{\lambda }_i,N_{k1};N_i)],`$
$`\gamma _\sigma `$ $`=`$ $`\gamma _\sigma (\overline{\lambda }_1,N_1;\mathrm{};\overline{\lambda }_{k1},N_{k1};N_k)={\displaystyle \frac{\mathrm{d}U(\lambda _k,N_k;\lambda ,N;\epsilon )}{\mathrm{d}\mathrm{ln}\mu }},`$
where $`\overline{\lambda }_i=\lambda _i/\lambda _k`$ and $`i=1,\mathrm{},k1`$. The list of fixed points is obtained by assigning the values $`0,1,\mathrm{}`$ to $`\overline{\lambda }_1,\mathrm{},\overline{\lambda }_{k1}`$ in all possible ways, keeping in mind that when some $`\overline{\lambda }`$’s are infinite, it is immaterial whether the finite $`\overline{\lambda }`$’s are 0 or 1. In total, we have $`2^k1`$ fixed points, corresponding to the models $`\sigma _{_sN}`$ for all possible subsets $`s`$ of $`(N_1,\mathrm{},N_k)`$. The flows are naturally associated with the regular polyhedron having $`k`$ faces in $`k1`$ dimensions (triangle for $`k=3`$, tetrahedron for $`k=4`$, etc.) and the dualities are symmetries of this polyhedron. The RG patterns for $`k=2,3,4`$ are illustrated in Fig. 3.
The trace anomaly reads
$$\mathrm{\Theta }=\sigma \underset{i=1}{\overset{k1}{}}\beta _i\underset{j=1}{\overset{N_i}{}}\overline{\psi }_{(i)}^j\psi _{(i)}^j.$$
Flows interpolating between the UV $`\sigma _{N+M}`$ and IR $`\sigma _N`$ fixed points can be obtained by giving mass to $`M`$ fermions. For the general purposes mentioned in the introduction, these flows are less interesting than the pure RG flows, which preserve conformality at the classical level and run only due to the dynamical scale $`\mu `$ . In some cases, nevertheless, such as the vector models constructed below, giving masses to the fermions seems the only simple way to interpolate between pairs of fixed points, because a non-renormalization theorem forbids the running of the gauge coupling constant.
Vector four-fermion models were also considered in . Here I make a set of observations on the non-Abelian composite gauge bosons, and prove that their coupling constant is quantized and non-renormalized. The Abelian coupling, instead, is non-renormalized, but can take arbitrary values.
We start from the four-fermion model defined by the lagrangian
$$=\overline{\psi }^i/\psi ^i+\frac{\lambda ^2}{2M}\left[(\overline{\psi }^i\gamma _\mu \psi ^i)^2\right],$$
(17)
to which we associate the conformal field theory
$$=\overline{\psi }^i(/+i\lambda A/)\psi ^i.$$
The vector $`A_\mu `$ becomes dynamical at the quantum level and the resulting conformal theory is interacting. The diagram (a) generates the $`A_\mu `$-propagator at the leading order in the large-$`N`$ expansion and its kinetic term in the quantum action reads, in momentum space and coordinate space, respectively:
$$\mathrm{\Gamma }_{\mathrm{kin}}[A]=\frac{1}{2}\frac{\lambda ^2N}{16}A_\mu (k)A_\nu (k)\frac{k^2\delta _{\mu \nu }k_\mu k_\nu }{\sqrt{k^2}}=\frac{1}{4}\frac{\lambda ^2N}{16}F_{\mu \nu }\frac{1}{\sqrt{\mathrm{}}}F_{\mu \nu },$$
(18)
$`F_{\mu \nu }`$ denoting the field strength. At the leading order we set again
$$\frac{\lambda ^2N}{16}=1,\lambda _\mathrm{B}=\lambda \mu ^{\epsilon /2}.$$
Since the $`U(1)`$ currents are conserved, the subleading corrections can only change the coefficient of the quadratic term in (18), but cannot change the dimension of the vector. There is, nevertheless, a non-vanishing anomalous dimension for the fermion fields, calculable from diagram (b) or, alternatively, from (c). We find, using an analogue of the Feynman gauge,
$$\gamma _\psi =\frac{2}{3N\pi ^2}+𝒪\left(\frac{1}{N^2}\right),\gamma _A0.$$
Finally, the non-local lagrangian kinetic term of the vector reads
$$_{\mathrm{non}\mathrm{loc}}=\frac{1}{2}A_\mu (k)A_\nu (k)\frac{k^2\delta _{\mu \nu }k_\mu k_\nu }{\sqrt{k^2}}\left[1\frac{\lambda _\mathrm{B}^2N}{16}(k^2)^{\epsilon /2}\right]$$
and does not need subleading corrections, since $`\gamma _A=0`$.
It is straightforward to construct conformal field theories with non-Abelian gauge fields, using the same method. The lagrangian
$$=\overline{\psi }^i[\delta _{ij}/+i\lambda A/^aT_{ij}^a]\psi ^j$$
generates a gauge-field quantum action
$$\mathrm{\Gamma }_{\mathrm{kin}}[A]=\frac{1}{4}\frac{\lambda ^2C(T)N_f}{16}F_{\mu \nu }^a\frac{1}{\sqrt{\mathrm{}}}F_{\mu \nu }^a,$$
for the field strength
$$F_{\mu \nu }^a=_\mu A_\nu ^a_\nu A_\mu ^a\lambda f^{abc}A_\mu ^bA_\nu ^c.$$
With a natural normalization convention for the action we see that the gauge coupling is discretized and equals
$$g^2=\frac{16}{C(T)N_f}+𝒪\left(\frac{1}{N_f^2}\right).$$
We have the freedom to change the representation $`T`$ and $`N_f`$, but since the set of unitary representations is denumerable, the non-Abelian gauge coupling can take only discrete values. It remains to be seen whether we can interpolate between two models with different values of the gauge coupling by means of an RG flow. However, a non-renormalization theorem forbids this. Consider the trace anomaly $`\mathrm{\Theta }`$. The term responsible for the running of the gauge coupling should be proportional to the gauge beta function multiplied by a non-trivial dimension-3, local, gauge-invariant operator. However, there exists no such operator in the gauge sector, all candidates being proportional to the field equations. The usual term $`F_{\mu \nu }^2`$ has dimension 4, while terms proportional to $`D/\psi `$ are trivial.
Other interesting conformal field theories are given by the gauged $`\sigma _N`$ models
$$=\overline{\psi }^i[\delta _{ij}/+i\lambda A/^aT_{ij}^a+i\lambda ^{}\sigma \delta _{ij}]\psi ^j.$$
RG flows such as
$$=\underset{i=1}{\overset{N}{}}\overline{\psi }^i[\delta _{ij}/+i\lambda A/^aT_{ij}^a+i\lambda ^{}\sigma \delta _{ij}]\psi ^j+\underset{i=1}{\overset{M}{}}\overline{\chi }^i[\delta _{ij}/+i\lambda A/^aR_{ij}^a+ig\sigma \delta _{ij}]\chi ^j$$
do not change the gauge-coupling, by the non-renormalization theorem proved above, but only the $`\sigma `$ coupling. The patterns of their RG flows are similar to the RG patterns of the non-gauged $`\sigma _N`$ models, with the only difference that the duality symmetries involve also exchanges of the representations, such as $`RT`$, etc.
A non-Abelian coupling can take arbitrary values, but the non-renormalization theorem applies. Consider for example
$$=\underset{i=1}{\overset{N}{}}\overline{\psi }^i[/+i\lambda A/]\psi ^i+\underset{i=1}{\overset{M}{}}\overline{\chi }^i[/+i\lambda ^{}A/]\chi ^i.$$
Here the trace anomaly is still identically zero and the theory is conformal for arbitrary $`\lambda `$ and $`\lambda ^{}`$. We have
$$\mathrm{\Gamma }_{\mathrm{kin}}[A]=\frac{1}{4}\frac{\lambda ^2N+\lambda ^{}{}_{}{}^{2}M}{16}F_{\mu \nu }\frac{1}{\sqrt{\mathrm{}}}F_{\mu \nu }.$$
We can reduce to the original vector four-fermion model (17) by means of a relevant deformation. A mass perturbation, such as $`MA_\mu ^2/2`$, produces the propagator
$$A_\mu ^a(k)A_\nu ^b(k)=\frac{\delta ^{ab}\sqrt{k^2}}{k^2+M\sqrt{k^2}}\left(\delta _{\mu \nu }+\frac{k_\mu k_\nu }{M\sqrt{k^2}}\right).$$
The behaviour $`k_\mu k_\nu /(Mk^2)`$ at large momenta is not dangerous if the current is conserved , which happens for Abelian fields. The situation is similar to quantum electrodynamics in four dimensions, where the photon can be given a mass without spoiling the renormalizability. With non-Abelian gauge fields we have to advocate a symmetry-breaking mechanism. We consider
$$=\overline{\psi }^i(\delta _{ij}/+i\lambda A/^aT_{ij}^a)\psi ^j+|D_\mu \phi |^2+V(|\phi |)+\mathrm{\Lambda }\overline{\psi }\psi \overline{\phi }\phi +\mathrm{\Lambda }^{}\overline{\psi }T^a\psi \overline{\phi }R^a\phi \mathrm{}$$
and assume that the potential $`V(|\phi |)`$ is such that the scalar field has an expectation value $`|\phi |=M^{1/2}`$. We know that the theory is renormalizable in the large-$`N`$ expansion. We can integrate the vector field out by solving its field equation. For simplicity we write the formulas in the Abelian case, although the mechanism is not strictly necessary there. We have
$$A_\mu ^a=\frac{i}{2\lambda \left|\phi \right|^2}\left(\overline{\psi }\gamma _\mu \psi \overline{\phi }_\mu \phi +_\mu \phi \overline{\phi }\right)=\frac{i\overline{\psi }\gamma _\mu \psi }{\lambda |M^{1/2}+\eta |^2}\frac{1}{\lambda }_\mu \theta $$
having set $`\phi =\mathrm{e}^{i\theta }\left(M^{1/2}+\eta \right)/\sqrt{2}`$. The Goldstone boson $`\theta `$ is gauged away as usual and we remain with
$$=\overline{\psi }^i/\psi ^i\frac{\left(\overline{\psi }^i\gamma _\mu \psi ^i\right)^2}{2M\left|1+\eta /\sqrt{M}\right|^2}+\frac{1}{2}\left(_\mu \eta \right)^2+V(\eta )+\frac{\mathrm{\Lambda }}{2}\overline{\psi }\psi \left|1+\eta /\sqrt{M}\right|^2.$$
(19)
When the mass of $`\eta `$ is very large, (17) is recovered exactly. By construction, the theory (19) is renormalizable, although this is not evident in the final form. Since the limit of large $`\eta `$-mass can be taken at $`M`$ fixed, we see that (17) is also renormalizable.
A more direct way to get to (17) is by replacing $`V(|\phi |)`$ with $`i\alpha (\overline{\phi }\phi M)`$, such as in the $`S_{N1}`$ non-lineal $`\sigma `$-model . The field $`\alpha `$ is dynamical and acquires a propagator proportional to $`\sqrt{k^2}`$, which is however compatible with power counting. In this case, however, we need to take a large-$`N`$ limit also in the number of $`\phi `$ components.
Acknowledgements
I am grateful to F. Nogueira for drawing my attention to the properties of four-fermion theories in three dimensions and the large-$`N`$ expansion for non-renormalizable theories, for discussions, and for pointing out relevant references. I also thank C. De Calan for discussions and U. Aglietti for continuous assistance and critical remarks. Finally, I thank J. Gracey and A. Kapustin and for pointing out references and , respectively.
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# More on Optical Holonomic Quantum Computer
## 1 Introduction
This paper is a continuation of Fujii and Pachos and Chountasis and the aim is to give a mathematical reinforcement to .
Quantum Computer is a very attractive and challenging object for New Science.
After the breakthrough by P. Shor there has been remarkable progress in Quantum Computation (or Computer)(QC briefly). This discovery had a great influence on scientists. This drived not only theoreticians to finding other quantum algorithms, but also experimentalists to building quantum computers. See in outline. and are also very useful for non-experts.
On the other hand, Gauge Theories are widely recognized as the basis in quantum field theories. Therefore it is very natural to intend to include gauge theories in QC $`\mathrm{}`$ a construction of “gauge theoretical” quantum computation or of “geometric” quantum computation in our terminology. The merit of geometric method of QC is strong against the influence from the environment. See .
Zanardi and Rasetti in and proposed such an idea using non-abelian Berry phase (quantum holonomy), see also and . In their model a Hamiltonian (including some parameters) must be degenerated because an adiabatic connection is introduced using this degeneracy . In other words, a quantum computational bundle on some parameter space (see ) is introduced due to this degeneracy.
They gave a simple but interesting example to explain their idea. We believe that this example will become important in the near future. Therefore we treated it once more and gave an explicit form to the non-abelian Berry connections and curvatures, see Fujii and Pachos and Chountasis .
In a non-abelian Berry connection and curvature was calculated by making use of the product of unitary coherent operators based on Lie algebras $`𝐂`$ and $`su(1,1)`$. On the other hand in an another non-abelian Berry connection and curvature was calculated by making use of the product of unitary coherent operators based on Lie algebras $`su(2)`$ and $`su(1,1)`$. See also .
We want to generalize the results above. Namely we want to replace Lie algebras $`su(2)`$ and $`su(1,1)`$ with bigger Lie algebras $`su(n+1)`$ and $`su(n,1)`$ for $`n1`$. Fortunately we have many studies of coherent states based on $`su(n+1)`$ and $`su(n,1)`$, see, for examples, , and .
In conclusion it is not easy for us to generalize the result in up to this time because we meet some difficulty. But it is, in principle, possible to generalize the one in in spite of hard calculation. We show this point in this paper. We list full calculations in the case of $`n=2`$ and leave the remaining cases to interested readers.
It is not easy to predict the future of geometric quantum computation. However it is an arena worth challenging for mathematical physicists.
## 2 Mathematical Foundation of Quantum Holonomy
We start with mathematical preliminaries. Let $``$ be a separable Hilbert space over $`𝐂`$. For $`m𝐍`$, we set
$$St_m\left(\right)\left\{V=(v_1,\mathrm{},v_m)\times \mathrm{}\times \right|V^{}V=1_m\},$$
(1)
where $`1_m`$ is a unit matrix in $`M(m,𝐂)`$. This is called a (universal) Stiefel manifold. Note that the unitary group $`U(m)`$ acts on $`St_m\left(\right)`$ from the right:
$$St_m\left(\right)\times U\left(m\right)St_m\left(\right):(V,a)Va.$$
(2)
Next we define a (universal) Grassmann manifold
$$Gr_m\left(\right)\left\{XM\left(\right)\right|X^2=X,X^{}=X\mathrm{and}\mathrm{tr}X=m\},$$
(3)
where $`M()`$ denotes a space of all bounded linear operators on $``$. Then we have a projection
$$\pi :St_m\left(\right)Gr_m\left(\right),\pi \left(V\right)VV^{},$$
(4)
compatible with the action (2) ($`\pi \left(Va\right)=Va(Va)^{}=Vaa^{}V^{}=VV^{}=\pi \left(V\right)`$).
Now the set
$$\{U\left(m\right),St_m\left(\right),\pi ,Gr_m\left(\right)\},$$
(5)
is called a (universal) principal $`U(m)`$ bundle, see and . We set
$$E_m\left(\right)\left\{(X,v)Gr_m\left(\right)\times \right|Xv=v\}.$$
(6)
Then we have also a projection
$$\pi :E_m\left(\right)Gr_m\left(\right),\pi \left((X,v)\right)X.$$
(7)
The set
$$\{𝐂^m,E_m\left(\right),\pi ,Gr_m\left(\right)\},$$
(8)
is called a (universal) $`m`$-th vector bundle. This vector bundle is one associated with the principal $`U(m)`$ bundle (5) .
Next let $`M`$ be a finite or infinite dimensional differentiable manifold and the map $`P:MGr_m\left(\right)`$ be given (called a projector). Using this $`P`$ we can make the bundles (5) and (8) pullback over $`M`$ :
$`\{U\left(m\right),\stackrel{~}{St},\pi _{\stackrel{~}{St}},M\}P^{}\{U\left(m\right),St_m\left(\right),\pi ,Gr_m\left(\right)\},`$ (9)
$`\{𝐂^m,\stackrel{~}{E},\pi _{\stackrel{~}{E}},M\}P^{}\{𝐂^m,E_m\left(\right),\pi ,Gr_m\left(\right)\},`$ (10)
see . (10) is of course a vector bundle associated with (9).
Let $``$ be a parameter space and we denote by $`\lambda `$ its element. Let $`\lambda _\mathrm{𝟎}`$ be a fixed reference point of $``$. Let $`H_\lambda `$ be a family of Hamiltonians parametrized by $``$ which act on a Fock space $``$. We set $`H_0`$ = $`H_{\lambda _\mathrm{𝟎}}`$ for simplicity and assume that this has a $`m`$-fold degenerate vacuum :
$$H_0v_j=\mathrm{𝟎},j=1m.$$
(11)
These $`v_j`$’s form a $`m`$-dimensional vector space. We may assume that $`v_i|v_j=\delta _{ij}`$. Then $`(v_1,\mathrm{},v_m)St_m\left(\right)`$ and
$$F_0\left\{\underset{j=1}{\overset{m}{}}x_jv_j\right|x_j𝐂\}𝐂^m.$$
Namely, $`F_0`$ is a vector space associated with o.n.basis $`(v_1,\mathrm{},v_m)`$.
Next we assume for simplicity that a family of unitary operators parametrized by $``$
$$W:U(),W(\lambda _\mathrm{𝟎})=\mathrm{id}.$$
(12)
is given and $`H_\lambda `$ above is given by the following isospectral family
$$H_\lambda W(\lambda )H_0W(\lambda )^1.$$
(13)
In this case there is no level crossing of eigenvalues. Making use of $`W(\lambda )`$ we can define a projector
$$P:Gr_m\left(\right),P(\lambda )W(\lambda )\left(\underset{j=1}{\overset{m}{}}v_jv_j^{}\right)W(\lambda )^1$$
(14)
and have the pullback bundles over $``$
$$\{U\left(m\right),\stackrel{~}{St},\pi _{\stackrel{~}{St}},\},\{𝐂^m,\stackrel{~}{E},\pi _{\stackrel{~}{E}},\}.$$
(15)
For the later we set
$$|vac=(v_1,\mathrm{},v_m).$$
(16)
In this case a canonical connection form $`𝒜`$ of $`\{U\left(m\right),\stackrel{~}{St},\pi _{\stackrel{~}{St}},\}`$ is given by
$$𝒜=vac|W(\lambda )^1dW(\lambda )|vac,$$
(17)
where $`d`$ is a differential form on $``$, and its curvature form by
$$d𝒜+𝒜𝒜.$$
(18)
On the other hand the global form of the curvature is given by $`PdPdP`$, which is related to (18) by
$$PdPdP=WW^1=W\left(d𝒜+𝒜𝒜\right)W^1,$$
(19)
see , and .
Let $`\gamma `$ be a loop in $``$ at $`\lambda _\mathrm{𝟎}`$., $`\gamma :[0,1],\gamma (0)=\gamma (1)`$. For this $`\gamma `$ a holonomy operator $`\mathrm{\Gamma }_𝒜`$ is defined :
$$\mathrm{\Gamma }_𝒜(\gamma )=𝒫exp\left\{_\gamma 𝒜\right\}U\left(m\right),$$
(20)
where $`𝒫`$ means path-ordered. This acts on the fiber $`F_0`$ at $`\lambda _\mathrm{𝟎}`$ of the vector bundle $`\{𝐂^m,\stackrel{~}{E},\pi _{\stackrel{~}{E}},M\}`$ as follows : $`\text{x}\mathrm{\Gamma }_𝒜(\gamma )\text{x}`$. The holonomy group $`Hol(𝒜)`$ is in general subgroup of $`U\left(m\right)`$ . In the case of $`Hol(𝒜)=U\left(m\right)`$, $`𝒜`$ is called irreducible, see .
In the Holonomic Quantum Computer we take
$`\mathrm{Encoding}\mathrm{of}\mathrm{Information}\text{x}F_0,`$
$`\mathrm{Processing}\mathrm{of}\mathrm{Information}\mathrm{\Gamma }_𝒜(\gamma ):\text{x}\mathrm{\Gamma }_𝒜(\gamma )\text{x}.`$ (21)
## 3 Unitary Coherent Operators based on $`su(n+1)`$ and $`su(n,1)`$
We apply the results of last section to Quantum Optics and discuss about (optical) Holonomic Quantum Computation proposed by and .
Let $`a(a^{})`$ be the annihilation (creation) operator of the harmonic oscillator. If we set $`Na^{}a`$ (: number operator), then
$$[N,a^{}]=a^{},[N,a]=a,[a,a^{}]=1.$$
(22)
Let $``$ be a Fock space generated by $`a`$ and $`a^{}`$, and $`\{|n|n𝐍\{0\}\}`$ be its basis. The actions of $`a`$ and $`a^{}`$ on $``$ are given by
$$a|n=\sqrt{n}|n1,a^{}|n=\sqrt{n+1}|n+1,$$
(23)
where $`|0`$ is a vacuum ($`a|0=0`$).
Next we consider the system of $`n+1`$–harmonic oscillators. For $`1jn+1`$ we set
$`a_j`$ $`=`$ $`1\mathrm{}1a1\mathrm{}1(\mathrm{j}\mathrm{position}),`$
$`a_{j}^{}{}_{}{}^{}`$ $`=`$ $`1\mathrm{}1a^{}1\mathrm{}1(\mathrm{j}\mathrm{position}),`$ (24)
then it is easy to see
$$[a_i,a_j]=[a_i^{},a_j^{}]=0,[a_i,a_j^{}]=\delta _{ij},i,j=1,2,\mathrm{},n+1.$$
(25)
We also denote by $`N_j=a_j^{}a_j`$ its number operators.
The Fock space $`^{(n+1)}`$ fot the system of $`n+1`$–harmonic oscillators is the $`n+1`$–tensor product $`^{(n+1)}=\mathrm{}`$, and each $`a_j\mathrm{and}a_{j}^{}{}_{}{}^{}`$ acts on $`j`$–component of $`^{(n+1)}`$ like (23).
Now since we want to consider coherent states based on Lie algebras $`su(n+1)`$ and $`su(n,1)`$, we make use of Schwinger’s boson method, see , .
### 3.1 $`su(n+1)`$
If we set
$$E_{\alpha \beta }=a_{\alpha }^{}{}_{}{}^{}a_\beta ,1\alpha ,\beta n+1$$
(26)
then from (25) we find
$$[E_{\alpha \beta },E_{\gamma \delta }]=E_{\alpha \delta }\delta _{\beta \gamma }\delta _{\delta \alpha }E_{\gamma \beta },$$
(27)
where $`\delta =\mathrm{diag}.(1,\mathrm{},1)`$. That is, $`\left\{E_{\alpha \beta }\right|1\alpha ,\beta n+1\}`$ is a generator of Lie algebra $`u(n+1)`$. Then a set of generator $`\left\{E_{j,n+1}\right|1jn\}`$ plays an important role. For $`1jn`$ we set
$$J_{}^{j}{}_{+}{}^{}=a_j^{}a_{n+1},J_{}^{j}{}_{}{}^{}=a_{n+1}^{}{}_{}{}^{}a_j,J_{}^{j}{}_{3}{}^{}=\frac{1}{2}\left(a_j^{}a_ja_{n+1}^{}{}_{}{}^{}a_{n+1}\right),$$
(28)
then we have
$$[J_{}^{j}{}_{3}{}^{},J_{}^{j}{}_{+}{}^{}]=J_{}^{j}{}_{+}{}^{},[J_{}^{j}{}_{3}{}^{},J_{}^{j}{}_{}{}^{}]=J_{}^{j}{}_{}{}^{},[J_{}^{j}{}_{+}{}^{},J_{}^{j}{}_{}{}^{}]=2J_{}^{j}{}_{3}{}^{}.$$
(29)
Namely, $`\{J_{}^{j}{}_{+}{}^{},J_{}^{j}{}_{}{}^{},J_{}^{j}{}_{3}{}^{}\}`$ forms $`su(2)`$–algebra. From this we can construct unitary coherent operators based on $`su(2)`$–algebra : For $`1jn`$
$$U_j(\xi _j)=exp\left(\xi _ja_j^{}a_{n+1}\overline{\xi _j}a_{n+1}^{}{}_{}{}^{}a_j\right).$$
(30)
The disentangling formula for this operator is
$$U_j(\xi _j)=e^{\eta _ja_j^{}a_{n+1}}e^{\mathrm{log}\left(1+|\eta _j|^2\right)\frac{1}{2}\left(a_j^{}a_ja_{n+1}^{}a_{n+1}\right)}e^{\overline{\eta }_ja_{n+1}^{}{}_{}{}^{}a_j},\eta _j=\frac{\xi _j\mathrm{tan}\left(|\xi _j|\right)}{|\xi _j|}.$$
(31)
see . Therefore combining these operators we define a unitary coherent operator based on $`su(n+1)`$–algebra :
Definition 1 For $`\stackrel{}{\xi }=(\xi _1,\xi _2,\mathrm{},\xi _n)`$ we define
$$U\left(\stackrel{}{\xi }\right)\underset{j=1}{\overset{n}{}}U_j(\xi _j)(\mathrm{in}\mathrm{this}\mathrm{order}).$$
(32)
Fot simplicity we also set $`U\left(\stackrel{}{\xi }\right)=_{j=1}^nU_j`$. For the latter convenience let us calculate $`U\left(\stackrel{}{\xi }\right)^1\frac{}{\xi _j}U\left(\stackrel{}{\xi }\right)`$. It is easy to see
$$U\left(\stackrel{}{\xi }\right)^1\frac{}{\xi _j}U\left(\stackrel{}{\xi }\right)=U_n^1\mathrm{}U_{j+1}^1\left(U_j^1\frac{}{\xi _j}U_j\right)U_{j+1}\mathrm{}U_n.$$
On the other hand we have already calculated $`U_j^1\frac{}{\xi _j}U_j`$ :
Lemma 2 We have
$`U_j^1{\displaystyle \frac{}{\xi _j}}U_j=`$ $`{\displaystyle \frac{1}{2}}\left(1+{\displaystyle \frac{\mathrm{sin}(2|\xi _j|)}{2|\xi _j|}}\right)a_j^{}a_{n+1}`$ (33)
$`+`$ $`{\displaystyle \frac{\overline{\xi _j}}{2|\xi _j|^2}}\left(1\mathrm{cos}(2|\xi _j|)\right){\displaystyle \frac{1}{2}}\left(a_j^{}a_ja_{n+1}^{}a_{n+1}\right)`$
$`+`$ $`{\displaystyle \frac{\overline{\xi _j}^2}{2|\xi _j|^2}}\left(1+{\displaystyle \frac{\mathrm{sin}(2|\xi _j|)}{2|\xi _j|}}\right)a_{n+1}^{}a_j.`$
From this we easily obtain
$`U\left(\stackrel{}{\xi }\right)^1{\displaystyle \frac{}{\xi _j}}U\left(\stackrel{}{\xi }\right)=`$ $`{\displaystyle \frac{1}{2}}\left(1+{\displaystyle \frac{\mathrm{sin}(2|\xi _j|)}{2|\xi _j|}}\right)a_j^{}U_n^1\mathrm{}U_{j+1}^1a_{n+1}U_{j+1}\mathrm{}U_n`$ (34)
$`+`$ $`{\displaystyle \frac{\overline{\xi _j}}{2|\xi _j|^2}}\left(1\mathrm{cos}(2|\xi _j|)\right){\displaystyle \frac{1}{2}}\left(a_j^{}a_jU_n^1\mathrm{}U_{j+1}^1a_{n+1}^{}a_{n+1}U_{j+1}\mathrm{}U_n\right)`$
$`+`$ $`{\displaystyle \frac{\overline{\xi _j}^2}{2|\xi _j|^2}}\left(1+{\displaystyle \frac{\mathrm{sin}(2|\xi _j|)}{2|\xi _j|}}\right)U_n^1\mathrm{}U_{j+1}^1a_{n+1}^{}U_{j+1}\mathrm{}U_na_j`$
Therefore we have only to calculate the term $`U_n^1\mathrm{}U_{j+1}^1a_{n+1}U_{j+1}\mathrm{}U_n`$.
Lemma 3 We have
$`U_n^1\mathrm{}U_{j+1}^1`$ $`a_{n+1}`$ $`U_{j+1}\mathrm{}U_n=c_{n,j+1}a_{n+1}{\displaystyle \underset{l=j+1}{\overset{n}{}}}d_{l,j+1}a_l,\mathrm{where}`$
$`c_{n,j+1}`$ $``$ $`\mathrm{cos}(|\xi _n|)\mathrm{cos}(|\xi _{n1}|)\mathrm{}\mathrm{cos}(|\xi _{j+1}|),`$
$`d_{l,j+1}`$ $``$ $`{\displaystyle \frac{\overline{\xi _l}\mathrm{sin}(|\xi _l|)}{|\xi _l|}}\mathrm{cos}(|\xi _{l1}|)\mathrm{cos}(|\xi _{l2}|)\mathrm{}\mathrm{cos}(|\xi _{j+1}|),`$ (35)
$`\mathrm{for}j+1ln.`$
Fron these facts we obtain
Proposition 4 for $`1jn`$
$`U\left(\stackrel{}{\xi }\right)^1{\displaystyle \frac{}{\xi _j}}U\left(\stackrel{}{\xi }\right)`$ (36)
$`=`$ $`{\displaystyle \frac{1}{2}}\left(1+{\displaystyle \frac{\mathrm{sin}(2|\xi _j|)}{2|\xi _j|}}\right)\left\{c_{n,j+1}a_j^{}a_{n+1}{\displaystyle \underset{l=j+1}{\overset{n}{}}}d_{l,j+1}a_j^{}a_l\right\}`$
$`+{\displaystyle \frac{\overline{\xi _j}}{2|\xi _j|^2}}(1\mathrm{cos}(2|\xi _j|)){\displaystyle \frac{1}{2}}\{a_j^{}a_jc_{n,j+1}^2a_{n+1}^{}a_{n+1}+{\displaystyle \underset{l=j+1}{\overset{n}{}}}c_{n,j+1}d_{l,j+1}a_{n+1}^{}a_l`$
$`+{\displaystyle \underset{l=j+1}{\overset{n}{}}}c_{n,j+1}\overline{d}_{l,j+1}a_l^{}a_{n+1}{\displaystyle \underset{l,k=j+1}{\overset{n}{}}}\overline{d}_{l,j+1}d_{k,j+1}a_l^{}a_k\}`$
$`+{\displaystyle \frac{\overline{\xi _j}^2}{2|\xi _j|^2}}\left(1+{\displaystyle \frac{\mathrm{sin}(2|\xi _j|)}{2|\xi _j|}}\right)\left\{c_{n,j+1}a_{n+1}^{}a_j{\displaystyle \underset{l=j+1}{\overset{n}{}}}\overline{d}_{l,j+1}a_l^{}a_j\right\}.`$
### 3.2 $`su(n,1)`$
If we set
$`E_{\alpha \beta }=a_{\alpha }^{}{}_{}{}^{}a_\beta ,1\alpha ,\beta n,`$
$`E_{n+1,\alpha }=a_{n+1}a_\alpha ,E_{\alpha ,n+1}=a_{\alpha }^{}{}_{}{}^{}a_{n+1}^{}{}_{}{}^{},E_{n+1,n+1}=a_{n+1}^{}{}_{}{}^{}a_{n+1}+1,`$ (37)
then from (25) we find
$$[E_{\alpha \beta },E_{\gamma \delta }]=E_{\alpha \delta }\eta _{\beta \gamma }\eta _{\delta \alpha }E_{\gamma \beta },$$
(38)
where $`\eta =\mathrm{diag}.(1,\mathrm{},1,1)`$. That is, $`\left\{E_{\alpha \beta }\right|1\alpha ,\beta n+1\}`$ is a generator of Lie algebra $`u(n,1)`$. Then a set of generator $`\left\{E_{j,n+1}\right|1jn\}`$ plays an important role. For $`1jn`$ we set
$$K_{}^{j}{}_{+}{}^{}=a_j^{}a_{n+1}^{}{}_{}{}^{},K_{}^{j}{}_{}{}^{}=a_{n+1}a_j,K_{}^{j}{}_{3}{}^{}=\frac{1}{2}\left(a_j^{}a_j+a_{n+1}^{}{}_{}{}^{}a_{n+1}+1\right),$$
(39)
then we have
$$[K_{}^{j}{}_{3}{}^{},K_{}^{j}{}_{+}{}^{}]=K_{}^{j}{}_{+}{}^{},[K_{}^{j}{}_{3}{}^{},K_{}^{j}{}_{}{}^{}]=K_{}^{j}{}_{}{}^{},[K_{}^{j}{}_{+}{}^{},K_{}^{j}{}_{}{}^{}]=2K_{}^{j}{}_{3}{}^{}.$$
(40)
Namely, $`\{K_{}^{j}{}_{+}{}^{},K_{}^{j}{}_{}{}^{},K_{}^{j}{}_{3}{}^{}\}`$ forms $`su(1,1)`$–algebra. From this we can construct unitary coherent operators based on $`su(1,1)`$–algebra : For $`1jn`$
$$V_j(\zeta _j)=exp\left(\zeta _ja_j^{}a_{n+1}^{}{}_{}{}^{}\overline{\zeta _j}a_{n+1}a_j\right).$$
(41)
The disentangling formula for this operator is
$$V_j(\zeta _j)=e^{\kappa _ja_j^{}a_{n+1}^{}{}_{}{}^{}}e^{\mathrm{log}\left(1|\kappa _j|^2\right)\frac{1}{2}\left(a_j^{}a_j+a_{n+1}^{}a_{n+1}+1\right)}e^{\overline{\kappa }_ja_{n+1}a_j},\kappa _j=\frac{\zeta _j\mathrm{tanh}\left(|\zeta _j|\right)}{|\zeta _j|}.$$
(42)
see . Therefore combining these operators we define a unitary coherent operator based on $`su(n,1)`$–algebra :
Definition 5 For $`\stackrel{}{\zeta }=(\zeta _1,\zeta _2,\mathrm{},\zeta _n)`$ we define
$$V\left(\stackrel{}{\zeta }\right)\underset{j=1}{\overset{n}{}}V_j(\zeta _j)(\mathrm{in}\mathrm{this}\mathrm{order}).$$
(43)
Fot simplicity we also set $`V\left(\stackrel{}{\zeta }\right)=_{j=1}^nV_j`$. For the latter convenience let us calculate $`V\left(\stackrel{}{\zeta }\right)^1\frac{}{\zeta _j}V\left(\stackrel{}{\xi }\right)`$. It is easy to see
$$V\left(\stackrel{}{\zeta }\right)^1\frac{}{\zeta _j}V\left(\stackrel{}{\zeta }\right)=V_n^1\mathrm{}V_{j+1}^1\left(V_j^1\frac{}{\zeta _j}V_j\right)V_{j+1}\mathrm{}V_n.$$
On the other hand we have already calculated $`V_j^1\frac{}{\zeta _j}V_j`$ :
Lemma 6 We have
$`V_j^1{\displaystyle \frac{}{\zeta _j}}V_j=`$ $`{\displaystyle \frac{1}{2}}\left(1+{\displaystyle \frac{\mathrm{sinh}(2|\zeta _j|)}{2|\zeta _j|}}\right)a_j^{}a_{n+1}^{}`$ (44)
$`+`$ $`{\displaystyle \frac{\overline{\zeta _j}}{2|\zeta _j|^2}}\left(1+\mathrm{cosh}(2|\zeta _j|)\right){\displaystyle \frac{1}{2}}\left(a_j^{}a_j+a_{n+1}^{}a_{n+1}+1\right)`$
$`+`$ $`{\displaystyle \frac{\overline{\zeta _j}^2}{2|\zeta _j|^2}}\left(1+{\displaystyle \frac{\mathrm{sinh}(2|\zeta _j|)}{2|\zeta _j|}}\right)a_{n+1}a_j.`$
From this we easily obtain
$`V\left(\stackrel{}{\zeta }\right)^1{\displaystyle \frac{}{\zeta _j}}V\left(\stackrel{}{\zeta }\right)`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(1+{\displaystyle \frac{\mathrm{sinh}(2|\zeta _j|)}{2|\zeta _j|}}\right)a_j^{}V_n^1\mathrm{}V_{j+1}^1a_{n+1}^{}V_{j+1}\mathrm{}V_n`$ (45)
$`+`$ $`{\displaystyle \frac{\overline{\zeta _j}}{2|\zeta _j|^2}}\left(1+\mathrm{cosh}(2|\zeta _j|)\right){\displaystyle \frac{1}{2}}\left(a_j^{}a_j+1+V_n^1\mathrm{}V_{j+1}^1a_{n+1}^{}a_{n+1}V_{j+1}\mathrm{}V_n\right)`$
$`+`$ $`{\displaystyle \frac{\overline{\zeta _j}^2}{2|\zeta _j|^2}}\left(1+{\displaystyle \frac{\mathrm{sinh}(2|\zeta _j|)}{2|\zeta _j|}}\right)V_n^1\mathrm{}V_{j+1}^1a_{n+1}V_{j+1}\mathrm{}V_na_j`$
Therefore we have only to calculate the term $`V_n^1\mathrm{}V_{j+1}^1a_{n+1}^{}V_{j+1}\mathrm{}V_n`$.
Lemma 7 We have
$`V_n^1\mathrm{}V_{j+1}^1`$ $`a_{n+1}^{}`$ $`V_{j+1}\mathrm{}V_n=e_{n,j+1}a_{n+1}^{}+{\displaystyle \underset{l=j+1}{\overset{n}{}}}f_{l,j+1}a_l,\mathrm{where}`$
$`e_{n,j+1}`$ $``$ $`\mathrm{cosh}(|\zeta _n|)\mathrm{cosh}(|\zeta _{n1}|)\mathrm{}\mathrm{cosh}(|\zeta _{j+1}|),`$
$`f_{l,j+1}`$ $``$ $`{\displaystyle \frac{\overline{\zeta _l}\mathrm{sinh}(|\zeta _l|)}{|\zeta _l|}}\mathrm{cosh}(|\zeta _{l1}|)\mathrm{cosh}(|\zeta _{l2}|)\mathrm{}\mathrm{cosh}(|\zeta _{j+1}|),`$ (46)
$`\mathrm{for}j+1ln.`$
Fron these facts we obtain
Proposition 8 for $`1jn`$
$`V\left(\stackrel{}{\zeta }\right)^1{\displaystyle \frac{}{\zeta _j}}V\left(\stackrel{}{\zeta }\right)`$ (47)
$`=`$ $`{\displaystyle \frac{1}{2}}\left(1+{\displaystyle \frac{\mathrm{sinh}(2|\zeta _j|)}{2|\zeta _j|}}\right)\left\{e_{n,j+1}a_j^{}a_{n+1}^{}+{\displaystyle \underset{l=j+1}{\overset{n}{}}}f_{l,j+1}a_j^{}a_l\right\}`$
$`+{\displaystyle \frac{\overline{\zeta _j}}{2|\zeta _j|^2}}(1+\mathrm{cosh}(2|\zeta _j|)){\displaystyle \frac{1}{2}}\{a_j^{}a_j+e_{n,j+1}^2(a_{n+1}^{}a_{n+1}+1)+{\displaystyle \underset{l=j+1}{\overset{n}{}}}e_{n,j+1}\overline{f}_{l,j+1}a_l^{}a_{n+1}^{}`$
$`+{\displaystyle \underset{l=j+1}{\overset{n}{}}}e_{n,j+1}f_{l,j+1}a_{n+1}a_l+{\displaystyle \underset{l,k=j+1}{\overset{n}{}}}f_{l,j+1}\overline{f}_{k,j+1}a_l^{}a_k\}`$
$`+{\displaystyle \frac{\overline{\zeta _j}^2}{2|\zeta _j|^2}}\left(1+{\displaystyle \frac{\mathrm{sinh}(2|\zeta _j|)}{2|\zeta _j|}}\right)\left\{e_{n,j+1}a_{n+1}a_j+{\displaystyle \underset{l=j+1}{\overset{n}{}}}\overline{f}_{l,j+1}a_l^{}a_j\right\}.`$
## 4 Optical Holonomic Quantum Computer $`\mathrm{}`$ Generalization
Let $`H_0`$ be a Hamiltonian with nonlinear interaction produced by a Kerr medium., that is $`H_0=\mathrm{}\mathrm{X}N(N1)`$, where X is a certain constant, see . The eigenvectors of $`H_0`$ corresponding to $`0`$ is $`\{|0,|1\}`$, so its eigenspace is $`\mathrm{Vect}\{|0,|1\}𝐂^2`$. We correspond to $`0|0,1|1`$ for a generator of Boolean algebra $`\{0,1\}`$. The space $`\mathrm{Vect}\{|0,|1\}`$ is called 1-qubit (quantum bit) space, see or . Since we are considering the system of $`n+1`$ particles, the Hamiltonian that we treat in the following becomes
$$H_0=\underset{j=1}{\overset{n+1}{}}\mathrm{}\mathrm{X}N_j(N_j1).$$
(48)
The $`0`$–eigenspace of this Hamiltonian becomes therefore
$$F_0=\mathrm{Vect}\{|0,|1\}\mathrm{}\mathrm{Vect}\{|0,|1\}𝐂^2\mathrm{}𝐂^2𝐂^{2^{n+1}}.$$
(49)
We denote
$$\alpha _1,\mathrm{},\alpha _{n+1}|\beta _1,\mathrm{},\beta _{n+1}=\underset{j=1}{\overset{n+1}{}}\alpha _j|\beta _j=\underset{j=1}{\overset{n+1}{}}\delta _{\alpha _j\beta _j},$$
for $`|\alpha _1,\mathrm{},\alpha _{n+1},|\beta _1,\mathrm{},\beta _{n+1}F_0`$. We order the basis of $`F_0`$ as
$`|0`$ $`=`$ $`|0,0,\mathrm{},0,0,`$
$`|1`$ $`=`$ $`|0,0,\mathrm{},0,1,`$
$`\mathrm{}`$
$`|2^{n+1}2`$ $`=`$ $`|1,1,\mathrm{},1,0,`$
$`|2^{n+1}1`$ $`=`$ $`|1,1,\mathrm{},1,1.`$
and set
$$|vac=(|0,|1,\mathrm{},|2^{n+1}1).$$
(50)
Namely $`m=2^{n+1}1`$ in (16).
Here we consider the following isospectral family of $`H_0`$ above :
$`H_{(\stackrel{}{\xi },\stackrel{}{\zeta })}`$ $`=`$ $`W(\stackrel{}{\xi },\stackrel{}{\zeta })H_0W(\stackrel{}{\xi },\stackrel{}{\zeta })^1,`$ (51)
$`W(\stackrel{}{\xi },\stackrel{}{\zeta })`$ $`=`$ $`U(\stackrel{}{\xi })V(\stackrel{}{\zeta })U(\mathrm{})(n+1\mathrm{times}),W(\stackrel{}{0},\stackrel{}{0})=\mathrm{id}.`$ (52)
For this system we want to calculate a connection form (17) in the last section. For that we set : for $`1jn`$
$`A_{\xi _j}`$ $`=`$ $`vac|W(\stackrel{}{\xi },\stackrel{}{\zeta })^1{\displaystyle \frac{}{\xi _j}}W(\stackrel{}{\xi },\stackrel{}{\zeta })|vac,`$
$`A_{\zeta _j}`$ $`=`$ $`vac|W(\stackrel{}{\xi },\stackrel{}{\zeta })^1{\displaystyle \frac{}{\zeta _j}}W(\stackrel{}{\xi },\stackrel{}{\zeta })|vac.`$ (53)
Here we note
$`W(\stackrel{}{\xi },\stackrel{}{\zeta })^1{\displaystyle \frac{}{\xi _j}}W(\stackrel{}{\xi },\stackrel{}{\zeta })`$ $`=`$ $`V(\stackrel{}{\zeta })^1\left\{U(\stackrel{}{\xi })^1{\displaystyle \frac{}{\xi _j}}U(\stackrel{}{\xi })\right\}V(\stackrel{}{\zeta }),`$ (54)
$`W(\stackrel{}{\xi },\stackrel{}{\zeta })^1{\displaystyle \frac{}{\zeta _j}}W(\stackrel{}{\xi },\stackrel{}{\zeta })`$ $`=`$ $`V(\stackrel{}{\zeta })^1{\displaystyle \frac{}{\zeta _j}}V(\stackrel{}{\zeta }).`$
On the other hand in Proposition 4 and Proposition 8 we have already calculated
$`U(\stackrel{}{\xi })^1\frac{}{\xi _j}U(\stackrel{}{\xi })`$ and $`V(\stackrel{}{\zeta })^1\frac{}{\zeta _j}V(\stackrel{}{\zeta })`$.
From Proposition 4 we must calculate $`V^1a_\alpha ^{}a_\beta V=\left(V^1a_\alpha V\right)^{}\left(V^1a_\beta V\right),\mathrm{where}V=V_1V_2\mathrm{}V_n`$. Therefore let us calculate $`V^1a_\alpha V\mathrm{for}1\alpha n+1`$. But remarking that $`V_ja_k=a_kV_j\mathrm{for}1kj1`$ because $`V_jV_j(\zeta _j)=\mathrm{exp}\left(\zeta _ja_j^{}a_{n+1}^{}{}_{}{}^{}\overline{\zeta _j}a_{n+1}a_j\right)`$ we must calculate
$`V^1a_jV`$ $`=`$ $`V_n^1\mathrm{}V_j^1a_jV_j\mathrm{}V_n\mathrm{for}1jn,`$
$`V^1a_{n+1}V`$ $`=`$ $`V_n^1\mathrm{}V_1^1a_{n+1}V_1\mathrm{}V_n.`$
To calculate these is not so difficult. The result is
Lemma 9
$`V^1a_jV`$ $`=`$ $`\mathrm{cosh}(|\zeta _j|)a_j+{\displaystyle \frac{\zeta _j\mathrm{sinh}(|\zeta _j|)}{|\zeta _j|}}\{{\displaystyle \underset{l=j+1}{\overset{n}{}}}{\displaystyle \underset{k=j+1}{\overset{l1}{}}}\mathrm{cosh}(|\zeta _k|){\displaystyle \frac{\overline{\zeta _l}\mathrm{sinh}(|\zeta _l|)}{|\zeta _l|}}a_l`$ (55)
$`+`$ $`{\displaystyle \underset{k=j+1}{\overset{n}{}}}\mathrm{cosh}(|\zeta _k|)a_{n+1}^{}\}\mathrm{for}1jn,`$
$`V^1a_{n+1}V`$ $`=`$ $`{\displaystyle \underset{k=1}{\overset{n}{}}}\left\{{\displaystyle \underset{l=1}{\overset{k1}{}}}\mathrm{cosh}(|\zeta _l|)\right\}{\displaystyle \frac{\zeta _k\mathrm{sinh}(|\zeta _k|)}{|\zeta _k|}}a_k^{}+{\displaystyle \underset{l=1}{\overset{n}{}}}\mathrm{cosh}(|\zeta _l|)a_{n+1}.`$ (56)
Here we should understand that $`_k^l(\mathrm{})=1`$ if $`l<k`$.
Using this we can calculate $`V^1a_\alpha ^{}a_\beta V`$ and next calculate (54) in principle. But it is not easy for us to obtain a compact form for this up to this time. Therefore let us restrict to some special cases ($`n=1,2`$) and obtain complete forms.
### 4.1 Example $`\mathrm{}n=1`$
In this case we can obtain the connection form in a complete manner. See and also . For simplicity we set $`\xi _1=\xi `$ and $`\zeta _1=\zeta `$. The result is
Lemma 10 we have
$`W^1{\displaystyle \frac{}{\xi }}W`$
$`=`$ $`{\displaystyle \frac{1}{2}}\left(1+{\displaystyle \frac{\mathrm{sin}(2|\xi |)}{2|\xi |}}\right)\left\{\mathrm{cosh}(2|\zeta |)a_1^{}a_2+{\displaystyle \frac{\zeta \mathrm{sinh}(2|\zeta |)}{2|\zeta |}}\left(a_1^{}\right)^2+{\displaystyle \frac{\overline{\zeta }\mathrm{sinh}(2|\zeta |)}{2|\zeta |}}\left(a_2\right)^2\right\}`$
$`+`$ $`{\displaystyle \frac{\overline{\xi }}{2|\xi |^2}}\left(1\mathrm{cos}(2|\xi |)\right){\displaystyle \frac{1}{2}}\left(a_1^{}a_1a_2^{}a_2\right)`$
$`+`$ $`{\displaystyle \frac{\overline{\xi }^2}{2|\xi |^2}}\left(1+{\displaystyle \frac{\mathrm{sin}(2|\xi |)}{2|\xi |}}\right)\left\{\mathrm{cosh}(2|\zeta |)a_2^{}a_1+{\displaystyle \frac{\overline{\zeta }\mathrm{sinh}(2|\zeta |)}{2|\zeta |}}\left(a_1\right)^2+{\displaystyle \frac{\zeta \mathrm{sinh}(2|\zeta |)}{2|\zeta |}}\left(a_2^{}\right)^2\right\},`$
$`W^1{\displaystyle \frac{}{\zeta }}W`$
$`=`$ $`{\displaystyle \frac{1}{2}}\left(1+{\displaystyle \frac{\mathrm{sinh}(2|\zeta |)}{2|\zeta |}}\right)a_1^{}a_2^{}+{\displaystyle \frac{\overline{\zeta }}{2|\zeta |^2}}\left(1+\mathrm{cosh}(2|\zeta |)\right){\displaystyle \frac{1}{2}}\left(a_1^{}a_1+a_2^{}a_2+1\right)`$
$`+`$ $`{\displaystyle \frac{\overline{\zeta }^2}{2|\zeta |^2}}\left(1+{\displaystyle \frac{\mathrm{sinh}(2|\zeta |)}{2|\xi |}}\right)a_1a_2.`$ (58)
From this lemma it is easy to calculate $`A_\xi `$ and $`A_\zeta `$. Let us here remember
$$|vac=(|0,0,|0,1,|1,0,|1,1).$$
Before stating the result let us prepare some notations.
$$\widehat{E}=\left(\begin{array}{cccc}0& 0& 0& 0\\ 0& 0& 1& 0\\ 0& 0& 0& 0\\ 0& 0& 0& 0\end{array}\right),\widehat{F}=\left(\begin{array}{cccc}0& 0& 0& 0\\ 0& 0& 0& 0\\ 0& 1& 0& 0\\ 0& 0& 0& 0\end{array}\right),\widehat{H}=\left(\begin{array}{cccc}0& 0& 0& 0\\ 0& \frac{1}{2}& 0& 0\\ 0& 0& \frac{1}{2}& 0\\ 0& 0& 0& 0\end{array}\right).$$
(59)
$$\widehat{A}=\left(\begin{array}{cccc}0& 0& 0& 1\\ 0& 0& 0& 0\\ 0& 0& 0& 0\\ 0& 0& 0& 0\end{array}\right),\widehat{C}=\left(\begin{array}{cccc}0& 0& 0& 0\\ 0& 0& 0& 0\\ 0& 0& 0& 0\\ 1& 0& 0& 0\end{array}\right),\widehat{B}=\left(\begin{array}{cccc}\frac{1}{2}& 0& 0& 0\\ 0& 1& 0& 0\\ 0& 0& 1& 0\\ 0& 0& 0& \frac{3}{2}\end{array}\right).$$
(60)
Proposition 11 We have
$`A_\xi `$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(1+{\displaystyle \frac{\mathrm{sin}(2|\xi |)}{2|\xi |}}\right)\mathrm{cosh}(2|\zeta |)\widehat{F}{\displaystyle \frac{\overline{\xi }}{2|\xi |^2}}\left(1\mathrm{cos}(2|\xi |)\right)\widehat{H}`$ (61)
$`+{\displaystyle \frac{\overline{\xi }^2}{2|\xi |^2}}\left(1+{\displaystyle \frac{\mathrm{sin}(2|\xi |)}{2|\xi |}}\right)\mathrm{cosh}(2|\zeta |)\widehat{E},`$
$`A_\zeta `$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(1+{\displaystyle \frac{\mathrm{sinh}(2|\zeta |)}{2|\zeta |}}\right)\widehat{C}+{\displaystyle \frac{\overline{\zeta }}{2|\zeta |^2}}\left(1+\mathrm{cosh}(2|\zeta |)\right)\widehat{B}`$ (62)
$`+{\displaystyle \frac{\overline{\zeta }^2}{2|\zeta |^2}}\left(1+{\displaystyle \frac{\mathrm{sinh}(2|\zeta |)}{2|\zeta |}}\right)\widehat{A}.`$
Since the connection form $`𝒜`$ is anti-hermitian ($`𝒜^{}=𝒜`$), it can be written as
$$𝒜=A_\xi d\xi +A_\zeta d\zeta A_\xi ^{}d\overline{\xi }A_\zeta ^{}d\overline{\zeta },$$
(63)
so that it’s curvature form $`=d𝒜+𝒜𝒜`$ becomes
$``$ $`=`$ $`\left(_\xi A_\zeta _\zeta A_\xi +[A_\xi ,A_\zeta ]\right)d\xi d\zeta `$ (64)
$`\left(_\xi A_\xi ^{}+_{\overline{\xi }}A_\xi +[A_\xi ,A_\xi ^{}]\right)d\xi d\overline{\xi }`$
$`\left(_\xi A_\zeta ^{}+_{\overline{\zeta }}A_\xi +[A_\xi ,A_\zeta ^{}]\right)d\xi d\overline{\zeta }`$
$`\left(_\zeta A_\xi ^{}+_{\overline{\xi }}A_\zeta +[A_\zeta ,A_\xi ^{}]\right)d\zeta d\overline{\xi }`$
$`\left(_\zeta A_\zeta ^{}+_{\overline{\zeta }}A_\zeta +[A_\zeta ,A_\zeta ^{}]\right)d\zeta d\overline{\zeta }`$
$`\left(_{\overline{\xi }}A_\zeta ^{}_{\overline{\zeta }}A_\xi ^{}+[A_\zeta ^{},A_\xi ^{}]\right)d\overline{\xi }d\overline{\zeta }.`$
In this case we can calculate the curvature form completely. Now let us state our main result in this section.
Theorem 12
$`=`$
$`\left\{\left(1+{\displaystyle \frac{\mathrm{sin}(2|\xi |)}{2|\xi |}}\right){\displaystyle \frac{\overline{\zeta }\mathrm{sinh}(2|\zeta |)}{2|\zeta |}}\widehat{F}+{\displaystyle \frac{\overline{\xi }^2}{|\xi |^2}}\left(1+{\displaystyle \frac{\mathrm{sin}(2|\xi |)}{2|\xi |}}\right){\displaystyle \frac{\overline{\zeta }\mathrm{sinh}(2|\zeta |)}{2|\zeta |}}\widehat{E}\right\}d\xi d\zeta `$
$`\{{\displaystyle \frac{\xi }{|\xi |^2}}(1+\mathrm{cos}(2|\xi |))\mathrm{cosh}(2|\zeta |)\widehat{F}{\displaystyle \frac{\mathrm{sin}(2|\xi |)}{|\xi |}}(1+\mathrm{cosh}^2(2|\zeta |))\widehat{H}`$
$`+{\displaystyle \frac{\overline{\xi }}{|\xi |^2}}(1+\mathrm{cos}(2|\xi |))\mathrm{cosh}(2|\zeta |)\widehat{E}\}d\xi d\overline{\xi }`$
$`\left\{\left(1+{\displaystyle \frac{\mathrm{sin}(2|\xi |)}{2|\xi |}}\right){\displaystyle \frac{\zeta \mathrm{sinh}(2|\zeta |)}{2|\zeta |}}\widehat{F}+{\displaystyle \frac{\overline{\xi }^2}{|\xi |^2}}\left(1+{\displaystyle \frac{\mathrm{sin}(2|\xi |)}{2|\xi |}}\right){\displaystyle \frac{\zeta \mathrm{sinh}(2|\zeta |)}{2|\zeta |}}\widehat{E}\right\}d\xi d\overline{\zeta }`$
$`\left\{\left(1+{\displaystyle \frac{\mathrm{sin}(2|\xi |)}{2|\xi |}}\right){\displaystyle \frac{\overline{\zeta }\mathrm{sinh}(2|\zeta |)}{2|\zeta |}}\widehat{E}+{\displaystyle \frac{\xi ^2}{|\xi |^2}}\left(1+{\displaystyle \frac{\mathrm{sin}(2|\xi |)}{2|\xi |}}\right){\displaystyle \frac{\overline{\zeta }\mathrm{sinh}(2|\zeta |)}{2|\zeta |}}\widehat{F}\right\}d\zeta d\overline{\xi }`$
$`{\displaystyle \frac{\mathrm{sinh}(2|\zeta |)}{|\zeta |}}\left(2\widehat{B}\text{1}_4\right)d\zeta d\overline{\zeta }`$
$`+\left\{\left(1+{\displaystyle \frac{\mathrm{sin}(2|\xi |)}{2|\xi |}}\right){\displaystyle \frac{\zeta \mathrm{sinh}(2|\zeta |)}{2|\zeta |}}\widehat{E}+{\displaystyle \frac{\xi ^2}{|\xi |^2}}\left(1+{\displaystyle \frac{\mathrm{sin}(2|\xi |)}{2|\xi |}}\right){\displaystyle \frac{\zeta \mathrm{sinh}(2|\zeta |)}{2|\zeta |}}\widehat{F}\right\}d\overline{\xi }d\overline{\zeta }.`$ (65)
From this and the theorem of Ambrose–Singer (see ) it is easy to see that
Corollary
$$Hol(𝒜)=SU(2)\times U(1)U(4).$$
(66)
Therefore $`𝒜`$ is not irreducible.
### 4.2 Example $`\mathrm{}n=2`$
For this case we can also obtain the connection form in a complete manner. Let us perform.
Lemma 13 We have
$`W^1{\displaystyle \frac{}{\xi _1}}W={\displaystyle \frac{1}{2}}\left(1+{\displaystyle \frac{\mathrm{sin}(2|\xi _1|)}{2|\xi _1|}}\right)\left\{\mathrm{cos}(|\xi _2|)V^1a_1^{}a_3V{\displaystyle \frac{\overline{\xi _2}\mathrm{sin}(|\xi _2|)}{|\xi _2|}}V^1a_1^{}a_2V\right\}`$
$`+`$ $`{\displaystyle \frac{\overline{\xi _1}}{2|\xi _1|^2}}(1\mathrm{cos}(2|\xi _1|)){\displaystyle \frac{1}{2}}\{V^1a_1^{}a_1V{\displaystyle \frac{1+\mathrm{cos}(2|\xi _2|)}{2}}V^1a_3^{}a_3V`$
$`+`$ $`{\displaystyle \frac{\overline{\xi _2}\mathrm{sin}(2|\xi _2|)}{2|\xi _2|}}V^1a_3^{}a_2V+{\displaystyle \frac{\xi _2\mathrm{sin}(2|\xi _2|)}{2|\xi _2|}}V^1a_2^{}a_3V{\displaystyle \frac{1\mathrm{cos}(2|\xi _2|)}{2}}V^1a_2^{}a_2V\}`$
$`+`$ $`{\displaystyle \frac{\overline{\xi _1}^2}{2|\xi _1|^2}}\left(1+{\displaystyle \frac{\mathrm{sin}(2|\xi _1|)}{2|\xi _1|}}\right)\left\{\mathrm{cos}(|\xi _2|)V^1a_3^{}a_1V{\displaystyle \frac{\xi _2\mathrm{sin}(|\xi _2|)}{|\xi _2|}}V^1a_2^{}a_1V\right\},`$
$`W^1{\displaystyle \frac{}{\xi _2}}W={\displaystyle \frac{1}{2}}\left(1+{\displaystyle \frac{\mathrm{sin}(2|\xi _2|)}{2|\xi _2|}}\right)V^1a_2^{}a_3V`$
$`+`$ $`{\displaystyle \frac{\overline{\xi _2}}{2|\xi _2|^2}}\left(1\mathrm{cos}(2|\xi _2|)\right){\displaystyle \frac{1}{2}}\left(V^1a_2^{}a_2VV^1a_3^{}a_3V\right)`$
$`+`$ $`{\displaystyle \frac{\overline{\xi _2}^2}{2|\xi _2|^2}}\left(1+{\displaystyle \frac{\mathrm{sin}(2|\xi _2|)}{2|\xi _2|}}\right)V^1a_3^{}a_2V,`$
$`W^1{\displaystyle \frac{}{\zeta _1}}W={\displaystyle \frac{1}{2}}\left(1+{\displaystyle \frac{\mathrm{sinh}(2|\zeta _1|)}{2|\zeta _1|}}\right)\left\{\mathrm{cosh}(|\zeta _2|)a_1^{}a_3^{}+{\displaystyle \frac{\overline{\zeta _2}\mathrm{sinh}(|\zeta _2|)}{|\zeta _2|}}a_1^{}a_2\right\}`$
$`+`$ $`{\displaystyle \frac{\overline{\zeta _1}}{2|\zeta _1|^2}}(1+\mathrm{cosh}(2|\zeta _1|)){\displaystyle \frac{1}{2}}\{a_1^{}a_1+{\displaystyle \frac{1+\mathrm{cosh}(2|\zeta _2|)}{2}}(a_3^{}a_3+1)+{\displaystyle \frac{\overline{\zeta _2}\mathrm{sinh}(2|\zeta _2|)}{2|\zeta _2|}}a_2a_3`$
$`+`$ $`{\displaystyle \frac{\zeta _2\mathrm{sinh}(2|\zeta _2|)}{2|\zeta _2|}}a_2^{}a_3^{}+{\displaystyle \frac{1+\mathrm{cosh}(2|\zeta _2|)}{2}}a_2^{}a_2\}`$
$`+`$ $`{\displaystyle \frac{\overline{\zeta _1}^2}{2|\zeta _1|^2}}\left(1+{\displaystyle \frac{\mathrm{sinh}(2|\zeta _1|)}{2|\zeta _1|}}\right)\left\{\mathrm{cosh}(|\zeta _2|)a_1a_3+{\displaystyle \frac{\zeta _2\mathrm{sinh}(|\zeta _2|)}{|\zeta _2|}}a_2^{}a_1\right\},`$
$`W^1{\displaystyle \frac{}{\zeta _2}}W={\displaystyle \frac{1}{2}}\left(1+{\displaystyle \frac{\mathrm{sinh}(2|\zeta _2|)}{2|\zeta _2|}}\right)a_2^{}a_3^{}+{\displaystyle \frac{\overline{\zeta _2}}{2|\zeta _2|^2}}\left(1+\mathrm{cosh}(2|\zeta _2|)\right){\displaystyle \frac{1}{2}}\left(a_2^{}a_2+a_3^{}a_3+1\right)`$
$`+`$ $`{\displaystyle \frac{\overline{\zeta _2}^2}{2|\zeta _2|^2}}\left(1+{\displaystyle \frac{\mathrm{sinh}(2|\zeta _2|)}{2|\zeta _2|}}\right)a_2a_3,`$ (70)
where we remember $`V=V_1V_2=V_1(\zeta _1)V_2(\zeta _2)=\mathrm{exp}\left(\zeta _1a_1^{}a_{3}^{}{}_{}{}^{}\overline{\zeta _1}a_3a_1\right)\mathrm{exp}\left(\zeta _2a_2^{}a_{3}^{}{}_{}{}^{}\overline{\zeta _2}a_3a_2\right)`$.
Next let us calculate $`V^1a_i^{}a_jV\mathrm{for}1i,j3`$. From Lemma 9 we have
Corollary 14
$`V^1a_1V`$ $`=`$ $`\mathrm{cosh}(|\zeta _1|)a_1+{\displaystyle \frac{\zeta _1\mathrm{sinh}(|\zeta _1|)}{|\zeta _1|}}\left\{\mathrm{cosh}(|\zeta _2|)a_3^{}+{\displaystyle \frac{\overline{\zeta _2}\mathrm{sinh}(|\zeta _2|)}{|\zeta _2|}}a_2\right\},`$
$`V^1a_2V`$ $`=`$ $`\mathrm{cosh}(|\zeta _2|)a_2+{\displaystyle \frac{\zeta _2\mathrm{sinh}(|\zeta _2|)}{|\zeta _2|}}a_3^{},`$
$`V^1a_3V`$ $`=`$ $`{\displaystyle \frac{\zeta _1\mathrm{sinh}(|\zeta _1|)}{|\zeta _1|}}a_1^{}+\mathrm{cosh}(|\zeta _1|)\left\{\mathrm{cosh}(|\zeta _2|)a_3+{\displaystyle \frac{\zeta _2\mathrm{sinh}(|\zeta _2|)}{|\zeta _2|}}a_2^{}\right\}.`$
From this we obtain
Lemma 15
$`V^1a_1^{}a_1V={\displaystyle \frac{1+\mathrm{cosh}(2|\zeta _1|)}{2}}a_1^{}a_1+{\displaystyle \frac{\zeta _1\mathrm{sinh}(2|\zeta _1|)}{2|\zeta _1|}}\left\{\mathrm{cosh}(|\zeta _2|)a_1^{}a_3^{}+{\displaystyle \frac{\overline{\zeta _2}\mathrm{sinh}(|\zeta _2|)}{|\zeta _2|}}a_1^{}a_2\right\}`$
$`+`$ $`{\displaystyle \frac{\overline{\zeta _1}\mathrm{sinh}(2|\zeta _1|)}{2|\zeta _1|}}\left\{\mathrm{cosh}(|\zeta _2|)a_3a_1+{\displaystyle \frac{\zeta _2\mathrm{sinh}(|\zeta _2|)}{|\zeta _2|}}a_2^{}a_1\right\}`$
$`+`$ $`{\displaystyle \frac{1+\mathrm{cosh}(2|\zeta _1|)}{2}}\{{\displaystyle \frac{1+\mathrm{cosh}(2|\zeta _2|)}{2}}(a_3^{}a_3+1)+{\displaystyle \frac{\overline{\zeta _2}\mathrm{sinh}(2|\zeta _2|)}{2|\zeta _2|}}a_3a_2`$
$`+`$ $`{\displaystyle \frac{\zeta _2\mathrm{sinh}(2|\zeta _2|)}{2|\zeta _2|}}a_2^{}a_3^{}+{\displaystyle \frac{1+\mathrm{cosh}(2|\zeta _2|)}{2}}a_2^{}a_2\},`$
$`V^1a_1^{}a_2V=\mathrm{cosh}(|\zeta _1|)\left\{\mathrm{cosh}(|\zeta _2|)a_1^{}a_2+{\displaystyle \frac{\zeta _2\mathrm{sinh}(|\zeta _2|)}{|\zeta _2|}}a_1^{}a_3^{}\right\}`$
$`+`$ $`{\displaystyle \frac{\overline{\zeta _1}\mathrm{sinh}(|\zeta _1|)}{|\zeta _1|}}\left\{{\displaystyle \frac{1+\mathrm{cosh}(2|\zeta _2|)}{2}}a_3a_2+{\displaystyle \frac{\zeta _2\mathrm{sinh}(2|\zeta _2|)}{2|\zeta _2|}}\left(a_3^{}a_3+1\right)\right\}`$
$`+`$ $`{\displaystyle \frac{\zeta _1\mathrm{sinh}(|\zeta _1|)}{|\zeta _1|}}\left\{{\displaystyle \frac{\zeta _2\mathrm{sinh}(2|\zeta _2|)}{2|\zeta _2|}}a_2^{}a_2+{\displaystyle \frac{\zeta _{2}^{}{}_{}{}^{2}\left(1+\mathrm{cosh}(2|\zeta _2|)\right)}{2|\zeta _2|^2}}a_2^{}a_3^{}\right\},`$
$`V^1a_1^{}a_3V={\displaystyle \frac{\zeta _1\mathrm{sinh}(2|\zeta _1|)}{2|\zeta _1|}}a_1^{}a_1^{}+\mathrm{cosh}(2|\zeta _1|)\left\{\mathrm{cosh}(|\zeta _2|)a_1^{}a_3+{\displaystyle \frac{\zeta _2\mathrm{sinh}(|\zeta _2|)}{|\zeta _2|}}a_1^{}a_2^{}\right\}`$
$`+`$ $`{\displaystyle \frac{\overline{\zeta _1}\mathrm{sinh}(2|\zeta _1|)}{2|\zeta _1|}}\left\{{\displaystyle \frac{1+\mathrm{cosh}(2|\zeta _2|)}{2}}a_3^2+{\displaystyle \frac{2\mathrm{sinh}(2|\zeta _2|)}{2|\zeta _2|}}a_3a_2^{}+{\displaystyle \frac{\zeta _{2}^{}{}_{}{}^{2}\left(1+\mathrm{cosh}(2|\zeta _2|)\right)}{2|\zeta _2|^2}}a_2^{}a_2^{}\right\},`$
$`V^1a_2^{}a_2V={\displaystyle \frac{1+\mathrm{cosh}(2|\zeta _2|)}{2}}a_2^{}a_2+{\displaystyle \frac{\zeta _2\mathrm{sinh}(2|\zeta _2|)}{2|\zeta _2|}}a_2^{}a_3^{}+{\displaystyle \frac{1+\mathrm{cosh}(2|\zeta _2|)}{2}}\left(a_3^{}a_3+1\right)`$
$`+`$ $`{\displaystyle \frac{\overline{\zeta _2}\mathrm{sinh}(2|\zeta _2|)}{2|\zeta _2|}}a_3a_2,`$
$`V^1a_2^{}a_3V={\displaystyle \frac{\zeta _1\mathrm{sinh}(|\zeta _1|)}{|\zeta _1|}}\left\{\mathrm{cosh}(|\zeta _2|)a_2^{}a_1^{}+{\displaystyle \frac{\overline{\zeta _2}\mathrm{sinh}(|\zeta _2|)}{|\zeta _2|}}a_3a_1^{}\right\}`$
$`+`$ $`\mathrm{cosh}(|\zeta _1|)\left\{{\displaystyle \frac{1+\mathrm{cosh}(2|\zeta _2|)}{2}}a_2^{}a_3+{\displaystyle \frac{\zeta _2\mathrm{sinh}(2|\zeta _2|)}{2|\zeta _2|}}a_2^{}a_2^{}+{\displaystyle \frac{\overline{\zeta _2}\mathrm{sinh}(2|\zeta _2|)}{2|\zeta _2|}}a_3^2+{\displaystyle \frac{1+\mathrm{cosh}(2|\zeta _2|)}{2}}a_3a_2^{}\right\},`$
$`V^1a_3^{}a_3V={\displaystyle \frac{1+\mathrm{cosh}(2|\zeta _1|)}{2}}\left(a_1^{}a_1+1\right)+{\displaystyle \frac{\overline{\zeta _1}\mathrm{sinh}(2|\zeta _1|)}{2|\zeta _1|}}\left\{\mathrm{cosh}(|\zeta _2|)a_1a_3+{\displaystyle \frac{\zeta _2\mathrm{sinh}(|\zeta _2|)}{|\zeta _2|}}a_1a_2^{}\right\}`$
$`+`$ $`{\displaystyle \frac{\zeta _1\mathrm{sinh}(2|\zeta _1|)}{2|\zeta _1|}}\left\{\mathrm{cosh}(|\zeta _2|)a_3^{}a_1^{}+{\displaystyle \frac{\overline{\zeta _2}\mathrm{sinh}(|\zeta _2|)}{|\zeta _2|}}a_2a_1^{}\right\}`$
$`+`$ $`{\displaystyle \frac{1+\mathrm{cosh}(2|\zeta _1|)}{2}}\{{\displaystyle \frac{1+\mathrm{cosh}(2|\zeta _2|)}{2}}a_3^{}a_3+{\displaystyle \frac{\zeta _2\mathrm{sinh}(2|\zeta _2|)}{2|\zeta _2|}}a_3^{}a_2^{}+{\displaystyle \frac{\overline{\zeta _2}\mathrm{sinh}(2|\zeta _2|)}{2|\zeta _2|}}a_2a_3`$
$`+`$ $`{\displaystyle \frac{1+\mathrm{cosh}(2|\zeta _2|)}{2}}(a_2^{}a_2+1)\}.`$
Making use of this lemma it is not difficult to calculate $`A_{\xi _j}`$ and $`A_{\zeta _j}(j=1,2)`$. Let us again remember
$$|vac=(|0,0,0,|0,0,1,|0,1,0,|0,1,1,|1,0,0,|1,0,1,|1,1,0,|1,1,1).$$
Therefore we have only to know that for $`1i,j3`$
$$vac|a_i^{}a_j|vac,vac|a_ia_j|vac\mathrm{and}vac|V^1a_i^{}a_jV|vac.$$
First let us determine $`vac|a_i^{}a_j|vac`$ and $`vac|a_ia_j|vac`$.
Lemma 16
$$vac|a_1^{}a_1|vac=\left(\begin{array}{cccccccc}0& 0& 0& 0& 0& 0& 0& 0\\ 0& 0& 0& 0& 0& 0& 0& 0\\ 0& 0& 0& 0& 0& 0& 0& 0\\ 0& 0& 0& 0& 0& 0& 0& 0\\ 0& 0& 0& 0& 1& 0& 0& 0\\ 0& 0& 0& 0& 0& 1& 0& 0\\ 0& 0& 0& 0& 0& 0& 1& 0\\ 0& 0& 0& 0& 0& 0& 0& 1\end{array}\right),vac|a_1^{}a_2|vac=\left(\begin{array}{cccccccc}0& 0& 0& 0& 0& 0& 0& 0\\ 0& 0& 0& 0& 0& 0& 0& 0\\ 0& 0& 0& 0& 0& 0& 0& 0\\ 0& 0& 0& 0& 0& 0& 0& 0\\ 0& 0& 1& 0& 0& 0& 0& 0\\ 0& 0& 0& 1& 0& 0& 0& 0\\ 0& 0& 0& 0& 0& 0& 0& 0\\ 0& 0& 0& 0& 0& 0& 0& 0\end{array}\right),$$
$$vac|a_1^{}a_3|vac=\left(\begin{array}{cccccccc}0& 0& 0& 0& 0& 0& 0& 0\\ 0& 0& 0& 0& 0& 0& 0& 0\\ 0& 0& 0& 0& 0& 0& 0& 0\\ 0& 0& 0& 0& 0& 0& 0& 0\\ 0& 1& 0& 0& 0& 0& 0& 0\\ 0& 0& 0& 0& 0& 0& 0& 0\\ 0& 0& 0& 1& 0& 0& 0& 0\\ 0& 0& 0& 0& 0& 0& 0& 0\end{array}\right),vac|a_1a_2|vac=\left(\begin{array}{cccccccc}0& 0& 0& 0& 0& 0& 1& 0\\ 0& 0& 0& 0& 0& 0& 0& 1\\ 0& 0& 0& 0& 0& 0& 0& 0\\ 0& 0& 0& 0& 0& 0& 0& 0\\ 0& 0& 0& 0& 0& 0& 0& 0\\ 0& 0& 0& 0& 0& 0& 0& 0\\ 0& 0& 0& 0& 0& 0& 0& 0\\ 0& 0& 0& 0& 0& 0& 0& 0\end{array}\right),$$
$$vac|a_1a_3|vac=\left(\begin{array}{cccccccc}0& 0& 0& 0& 0& 1& 0& 0\\ 0& 0& 0& 0& 0& 0& 0& 0\\ 0& 0& 0& 0& 0& 0& 0& 1\\ 0& 0& 0& 0& 0& 0& 0& 0\\ 0& 0& 0& 0& 0& 0& 0& 0\\ 0& 0& 0& 0& 0& 0& 0& 0\\ 0& 0& 0& 0& 0& 0& 0& 0\\ 0& 0& 0& 0& 0& 0& 0& 0\end{array}\right),vac|a_2^{}a_2|vac=\left(\begin{array}{cccccccc}0& 0& 0& 0& 0& 0& 0& 0\\ 0& 0& 0& 0& 0& 0& 0& 0\\ 0& 0& 1& 0& 0& 0& 0& 0\\ 0& 0& 0& 1& 0& 0& 0& 0\\ 0& 0& 0& 0& 0& 0& 0& 0\\ 0& 0& 0& 0& 0& 0& 0& 0\\ 0& 0& 0& 0& 0& 0& 1& 0\\ 0& 0& 0& 0& 0& 0& 0& 1\end{array}\right),$$
$$vac|a_2^{}a_3|vac=\left(\begin{array}{cccccccc}0& 0& 0& 0& 0& 0& 0& 0\\ 0& 0& 0& 0& 0& 0& 0& 0\\ 0& 1& 0& 0& 0& 0& 0& 0\\ 0& 0& 0& 0& 0& 0& 0& 0\\ 0& 0& 0& 0& 0& 0& 0& 0\\ 0& 0& 0& 0& 0& 0& 0& 0\\ 0& 0& 0& 0& 0& 1& 0& 0\\ 0& 0& 0& 0& 0& 0& 0& 0\end{array}\right),vac|a_2a_3|vac=\left(\begin{array}{cccccccc}0& 0& 0& 1& 0& 0& 0& 0\\ 0& 0& 0& 0& 0& 0& 0& 0\\ 0& 0& 0& 0& 0& 0& 0& 0\\ 0& 0& 0& 0& 0& 0& 0& 0\\ 0& 0& 0& 0& 0& 0& 0& 1\\ 0& 0& 0& 0& 0& 0& 0& 0\\ 0& 0& 0& 0& 0& 0& 0& 0\\ 0& 0& 0& 0& 0& 0& 0& 0\end{array}\right),$$
$$vac|a_3^{}a_3|vac=\left(\begin{array}{cccccccc}0& 0& 0& 0& 0& 0& 0& 0\\ 0& 1& 0& 0& 0& 0& 0& 0\\ 0& 0& 0& 0& 0& 0& 0& 0\\ 0& 0& 0& 1& 0& 0& 0& 0\\ 0& 0& 0& 0& 0& 0& 0& 0\\ 0& 0& 0& 0& 0& 1& 0& 0\\ 0& 0& 0& 0& 0& 0& 0& 0\\ 0& 0& 0& 0& 0& 0& 0& 1\end{array}\right),\mathrm{and}vac|a_j^2|vac=\mathrm{𝟎}_8j=1,2,3.$$
Next let us determine $`vac|V^1a_i^{}a_jV|vac`$. For that we prepare some notations : for $`1i,j3`$ we set
$$M_{ij}=vac|a_i^{}a_j|vac,N_{ij}=vac|a_ia_j|vac.$$
Then both $`M_{ij}`$ and $`N_{ij}`$ are real matrices and moreover satisfy
$`M_{ij}^{}=M_{ji}\mathrm{and}M_{ii}\mathrm{𝟎}_8,`$
$`N_{ij}=N_{ji}\mathrm{and}N_{ii}=\mathrm{𝟎}_8.`$
Then we have
Lemma 17
$`vac|V^1a_1^{}a_1V|vac={\displaystyle \frac{1+\mathrm{cosh}(2|\zeta _1|)}{2}}M_{11}+{\displaystyle \frac{\zeta _1\mathrm{sinh}(2|\zeta _1|)}{2|\zeta _1|}}\left\{\mathrm{cosh}(|\zeta _2|)N_{13}^{}+{\displaystyle \frac{\overline{\zeta _2}\mathrm{sinh}(|\zeta _2|)}{|\zeta _2|}}M_{12}\right\}`$
$`+`$ $`{\displaystyle \frac{\overline{\zeta _1}\mathrm{sinh}(2|\zeta _1|)}{2|\zeta _1|}}\{\mathrm{cosh}(|\zeta _2|)N_{13}+{\displaystyle \frac{\zeta _2\mathrm{sinh}(|\zeta _2|)}{|\zeta _2|}}M_{12}^{}\}+{\displaystyle \frac{1+\mathrm{cosh}(2|\zeta _1|)}{2}}\{{\displaystyle \frac{1+\mathrm{cosh}(2|\zeta _2|)}{2}}(M_{33}+E)`$
$`+`$ $`{\displaystyle \frac{\overline{\zeta _2}\mathrm{sinh}(2|\zeta _2|)}{2|\zeta _2|}}N_{23}+{\displaystyle \frac{\zeta _2\mathrm{sinh}(2|\zeta _2|)}{2|\zeta _2|}}N_{23}^{}+{\displaystyle \frac{1+\mathrm{cosh}(2|\zeta _2|)}{2}}M_{22}\},`$
$`vac|V^1a_1^{}a_2V|vac=\mathrm{cosh}(|\zeta _1|)\left\{\mathrm{cosh}(|\zeta _2|)M_{12}+{\displaystyle \frac{\zeta _2\mathrm{sinh}(|\zeta _2|)}{|\zeta _2|}}N_{13}^{}\right\}`$
$`+`$ $`{\displaystyle \frac{\overline{\zeta _1}\mathrm{sinh}(|\zeta _1|)}{|\zeta _1|}}\left\{{\displaystyle \frac{1+\mathrm{cosh}(2|\zeta _2|)}{2}}N_{23}+{\displaystyle \frac{\zeta _2\mathrm{sinh}(2|\zeta _2|)}{2|\zeta _2|}}\left(M_{33}+E\right)\right\}`$
$`+`$ $`{\displaystyle \frac{\zeta _1\mathrm{sinh}(|\zeta _1|)}{|\zeta _1|}}\left\{{\displaystyle \frac{\zeta _2\mathrm{sinh}(2|\zeta _2|)}{2|\zeta _2|}}M_{22}+{\displaystyle \frac{\zeta _{2}^{}{}_{}{}^{2}\left(1+\mathrm{cosh}(2|\zeta _2|)\right)}{2|\zeta _2|^2}}N_{23}^{}\right\},`$
$`V^1a_1^{}a_3V=\mathrm{cosh}(2|\zeta _1|)\left\{\mathrm{cosh}(|\zeta _2|)M_{13}+{\displaystyle \frac{\zeta _2\mathrm{sinh}(|\zeta _2|)}{|\zeta _2|}}N_{12}^{}\right\}`$
$`+`$ $`{\displaystyle \frac{\overline{\zeta _1}\mathrm{sinh}(2|\zeta _1|)}{2|\zeta _1|}}{\displaystyle \frac{2\mathrm{sinh}(2|\zeta _2|)}{2|\zeta _2|}}M_{23},`$
$`V^1a_2^{}a_2V={\displaystyle \frac{1+\mathrm{cosh}(2|\zeta _2|)}{2}}M_{22}+{\displaystyle \frac{\zeta _2\mathrm{sinh}(2|\zeta _2|)}{2|\zeta _2|}}N_{23}^{}+{\displaystyle \frac{1+\mathrm{cosh}(2|\zeta _2|)}{2}}\left(M_{33}+E\right)`$
$`+`$ $`{\displaystyle \frac{\overline{\zeta _2}\mathrm{sinh}(2|\zeta _2|)}{2|\zeta _2|}}N_{23},`$
$`V^1a_2^{}a_3V={\displaystyle \frac{\zeta _1\mathrm{sinh}(|\zeta _1|)}{|\zeta _1|}}\left\{\mathrm{cosh}(|\zeta _2|)N_{12}^{}+{\displaystyle \frac{\overline{\zeta _2}\mathrm{sinh}(|\zeta _2|)}{|\zeta _2|}}M_{13}\right\}+\mathrm{cosh}(|\zeta _1|)\mathrm{cosh}(2|\zeta _2|)M_{23},`$
$`V^1a_3^{}a_3V={\displaystyle \frac{1+\mathrm{cosh}(2|\zeta _1|)}{2}}\left(M_{11}+E\right)+{\displaystyle \frac{\overline{\zeta _1}\mathrm{sinh}(2|\zeta _1|)}{2|\zeta _1|}}\left\{\mathrm{cosh}(|\zeta _2|)N_{13}+{\displaystyle \frac{\zeta _2\mathrm{sinh}(|\zeta _2|)}{|\zeta _2|}}M_{12}^{}\right\}`$
$`+`$ $`{\displaystyle \frac{\zeta _1\mathrm{sinh}(2|\zeta _1|)}{2|\zeta _1|}}\left\{\mathrm{cosh}(|\zeta _2|)N_{13}^{}+{\displaystyle \frac{\overline{\zeta _2}\mathrm{sinh}(|\zeta _2|)}{|\zeta _2|}}M_{12}\right\}`$
$`+`$ $`{\displaystyle \frac{1+\mathrm{cosh}(2|\zeta _1|)}{2}}\{{\displaystyle \frac{1+\mathrm{cosh}(2|\zeta _2|)}{2}}M_{33}+{\displaystyle \frac{\zeta _2\mathrm{sinh}(2|\zeta _2|)}{2|\zeta _2|}}N_{23}^{}+{\displaystyle \frac{\overline{\zeta _2}\mathrm{sinh}(2|\zeta _2|)}{2|\zeta _2|}}N_{23}`$
$`+`$ $`{\displaystyle \frac{1+\mathrm{cosh}(2|\zeta _2|)}{2}}(M_{22}+E)\}.`$ (76)
Here we have denoted by $`E`$ the unit matrix in $`M(8,𝐂)`$.
Using these lemmas we can obtain $`A_{\xi _j}`$ and $`A_{\zeta _j}(j=1,2)`$ completely. Next we must calculate the curvature form making use of the connection form, but it is too hard. We leave its calculation to interested readers.
## 5 Discussion
We in this paper defined unitary coherent operators based on Lie algebras $`su(n+1)`$ and $`su(n,1)`$ and, making use of these, calculated non–abelian Berry connections of quantum computational bundles proposed by Zanardi and Rasetti . For $`n=1`$ and $`2`$ we gave an explicit form to them. This ia a generalization of that of Pachos and Chountasis .
But for $`n3`$ we could not give explicit ones due to complexity. Therefore our paper is far from complete. As $`n`$ becomes large our culculation will become miserable. Moreover we didn’t perform the calculation of curvatures except for $`n=1`$.
We have a lot of problems to be performed. We expect that many young mathematical physicists with brute force will enter in this field.
Acknowledgment.
The author wishes to thank K. Funahashi and Y. Machida for their helpful comments and suggestions. I also thank J. Pachos and S. Chountasis for some useful suggestions.
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# The neutron resonance: modeling photoemission and tunneling data in the superconducting state of Bi2Sr2CaCu2O8+δ
## Abstract
Motivated by neutron scattering data, we develop a model of electrons interacting with a magnetic resonance and use it to analyze angle resolved photoemission (ARPES) and tunneling data in the superconducting state of Bi<sub>2</sub>Sr<sub>2</sub>CaCu<sub>2</sub>O<sub>8+δ</sub>. We not only can explain the peak-dip-hump structure observed near the $`(\pi ,0)`$ point, and its particle-hole asymmetry as seen in SIN tunneling spectra, but also its evolution throughout the Brillouin zone, including a velocity ‘kink’ near the d-wave node.
Recent advances in the momentum resolution of ARPES spectroscopy have led to a detailed mapping of the spectral lineshape in the high-T<sub>c</sub> superconductor Bi<sub>2</sub>Sr<sub>2</sub>CaCu<sub>2</sub>O<sub>8+δ</sub> (BSCCO) throughout the Brillouin zone. In contrast to normal state data where well defined excitations do not exist, quasiparticle peaks were identified below $`T_c`$ along the entire Fermi surface. Moreover, it has been known for some time that near the $`(\pi ,0)`$ ($`M`$) point of the zone, the spectral function shows an anomalous lineshape, the so called ‘peak-dip-hump’ structure. The new data indicate a) near the d-wave node of the superconducting gap, the dispersion shows a characteristic ‘kink’ feature: for $`|\omega |<\omega _{kink}`$, the spectra exhibit sharp peaks with a weaker dispersion; above this, broad peaks with a stronger dispersion; b) away from the node, the dispersion kink develops into a ‘break’; the two resulting branches are separated by an energy gap, and overlap in momentum space; c) towards $`M`$, the break evolves into a pronounced spectral ‘dip’ separating the almost dispersionless quasiparticle branch from the weakly dispersing high energy branch (the ‘hump’); d) the kink, break, and dip features all occur at the same energy, independent of position in the zone.
Features similar to the ARPES spectrum near the $`M`$ point were earlier observed in tunneling spectroscopy. Experimental SIN junctions on BSCCO show a characteristic asymmetry, with a more pronounced dip-hump structure on the occupied side. On the other hand, SIS junctions reveal a strong dip-hump feature on both bias sides.
There have been several theoretical treatments which assigned the anomalous ARPES lineshape near the $`M`$ point of the zone to the coupling between spin fluctuations and electrons. Here, we are able to explain features a)-d) of the ARPES data, as well as the SIN tunneling asymmetry, in terms of the combined effect of A) the flat electronic dispersion near the $`M`$ point of the zone and B) coupling of the fermionic degrees of freedom to a bosonic mode which is sharp in energy and peaked in momentum near $`\stackrel{}{Q}=(\pi ,\pi )`$. Our main result is that the anomalous features in the dispersion and lineshape for all points in the zone have the same origin.
A resonance mode with these characteristics is observed in inelastic neutron scattering experiments in bilayer cuprates in the superconducting state. The neutron resonance lies below a gapped continuum, the latter having a signal typically a factor of 30 less than the maximum at $`\stackrel{}{Q}`$ at the mode energy. In order to extract the essential physics, we concentrate on the mode part and neglect the continuum. The latter contributes mainly to additional damping at higher energies. We treat the mode in a semi-phenomenological way, taking the relevant parameters from experiment. We then calculate the resulting electronic self energy to second order in the coupling constant. From the self energy we directly obtain the spectral function measured by ARPES, which we then use to calculate the tunneling conductance.
The retarded self energy on the real energy axis is given by
$$\widehat{\mathrm{\Sigma }}^R=\frac{ig^2}{2\mu _B^2}\widehat{\tau }_3\left(\widehat{G}^K\chi ^R+\widehat{G}^R\chi ^K\right)\widehat{\tau }_3$$
(1)
with $`AB(\stackrel{}{k},ϵ)_\stackrel{}{q}_{\mathrm{}}^{\mathrm{}}\frac{d\omega }{2\pi }A(\stackrel{}{k}\stackrel{}{q},ϵ\omega )B(\stackrel{}{q},\omega )`$, $`\widehat{\tau }_i`$ Pauli matrices in particle-hole space, and $`g`$ the effective coupling constant. The Keldysh ($`K`$) components are given in terms of retarded ($`R`$) and advanced ($`A`$) functions by $`\widehat{G}^K=(\widehat{G}^R\widehat{G}^A)(12f)`$ and $`\chi ^K=(\chi ^R\chi ^A)(1+2b)`$, with the usual Bose ($`b`$) and Fermi ($`f`$) distribution functions.
The model for the mode is based on measurements of the spin susceptibility from inelastic neutron scattering experiments. The mode energy will be denoted by $`\mathrm{\Omega }`$ and its energy width is assumed to be irrelevant for the self energy. This assumption will be confirmed later. This leads to the following model for the mode part of the susceptibility
$$\chi ^{R/A}(\stackrel{}{q},\omega )=f(\stackrel{}{q})\left(\frac{1}{\omega \mathrm{\Omega }\pm i\delta }\frac{1}{\omega +\mathrm{\Omega }\pm i\delta }\right)$$
(2)
Here $`f(\stackrel{}{q})`$ describes the momentum dependence of the mode and is assumed to be enhanced at the $`(\pi ,\pi )`$ point. Using the correlation length $`\xi `$ we write it as
$$f(\stackrel{}{q})=\frac{\chi _\stackrel{}{Q}\xi ^2}{\xi ^2+4(\mathrm{cos}^2\frac{q_x}{2}+\mathrm{cos}^2\frac{q_y}{2})}$$
(3)
Experimentally the energy integrated susceptibility at the $`(\pi ,\pi )`$ wavevector, $`\pi \chi _\stackrel{}{Q}`$, was determined to be $`0.95\mu _B^2`$ per plane for BSCCO, leading to $`\chi _\stackrel{}{Q}=0.3\mu _B^2`$. For the correlation length, we take a conservative estimate of $`\xi =2a`$. This corresponds to a full width half maximum of $`0.26`$Å<sup>-1</sup>, as observed in YBCO, but somewhat smaller than that estimated for BSCCO ($`0.52`$Å<sup>-1</sup>) which we feel is somewhat broad. The mode energy was chosen to be $`\mathrm{\Omega }=39`$ meV, which represents the reported values between 35 and 43 meV.
Though $`\chi `$ is ‘renormalized’, we use a bare $`\widehat{G}`$ in Eq. 1. This is the same approximation as in Ref. , where it was shown that this is better than using renormalized $`\widehat{G}`$ with vertex corrections neglected. This is unlike the electron-phonon problem, where Migdal’s theorem applies. We take the success of explaining the experimental features as strong support of this approximation.
The bare Green’s functions with normal state dispersion $`\xi _\stackrel{}{k}`$, gap function $`\mathrm{\Delta }_\stackrel{}{k}`$, and excitation energy $`E_\stackrel{}{k}=\sqrt{\xi _\stackrel{}{k}^2+\mathrm{\Delta }_\stackrel{}{k}^2}`$ are
$$\widehat{G}^{R/A}(\stackrel{}{k},ϵ)=\frac{\widehat{\alpha }_\stackrel{}{k}}{ϵE_\stackrel{}{k}\pm i\mathrm{\Gamma }}+\frac{\widehat{\beta }_\stackrel{}{k}}{ϵ+E_\stackrel{}{k}\pm i\mathrm{\Gamma }}$$
(4)
where $`\alpha _{11}=(1+\xi _\stackrel{}{k}/E_\stackrel{}{k})/2`$, $`\beta _{11}=(1\xi _\stackrel{}{k}/E_\stackrel{}{k})/2`$, $`\alpha _{12}=\beta _{12}=\mathrm{\Delta }_\stackrel{}{k}/2E_\stackrel{}{k}`$, etc. For the normal state dispersion we use a six-parameter tight binding fit. We neglect bilayer splitting, as experiments suggest it is absent in BSCCO. A characteristic feature of this dispersion is a flat band with a saddlepoint at $`M`$ with energy $`\xi _M=`$34 meV. The superconducting gap function is taken to be the d-wave $`\mathrm{\Delta }_\stackrel{}{k}=\mathrm{\Delta }_0(\mathrm{cos}k_x\mathrm{cos}k_y)/2`$ with a maximal gap value of $`\mathrm{\Delta }_0=35`$ meV. We have chosen $`\mathrm{\Gamma }=5`$ meV as an intrinsic lifetime broadening. The coupling constant relevant for our model is $`g^2\chi _\stackrel{}{Q}`$, chosen to be 0.125 $`\text{eV}^2\mu _B^2`$. Given a value $`\chi _\stackrel{}{Q}=0.3\mu _B^2`$, this corresponds to $`g=0.65`$ eV, the same value as used in previous spin fluctuation work. We performed the $`\omega `$-integration in Eq. 1 analytically and the correlation product in momentum space via fast Fourier transform, using $`256\times 256`$ points in the Brillouin zone.
In Fig. 1 we show the renormalization function $`Z(ϵ)=1\mathrm{\Sigma }_0(ϵ)/ϵ`$, where $`\mathrm{\Sigma }_0`$ is the $`\widehat{\tau }_0`$ component of the $`\widehat{\mathrm{\Sigma }}`$ matrix. Since the $`\stackrel{}{q}`$ integral in Eq. 1 is dominated by the regions around the $`M`$ point where the band is flat and close to the chemical potential, there are features in the imaginary part of the self-energy connected with the two extremal energies $`\mathrm{\Delta }_0`$ and $`E_M=\sqrt{\xi _M^2+\mathrm{\Delta }_0^2}`$. These features do not show dispersion, but a change in magnitude with position in the zone which is determined by the momentum width of the mode. This is the central result of this paper. More generally, the imaginary part of the self energy is enhanced between the values
$`ϵ_1=\mathrm{\Delta }_0+\mathrm{\Omega }`$ and $`ϵ_2=E_M+\mathrm{\Omega }`$. For $`\xi _M`$ approaching the chemical potential, $`E_M`$ approaches $`\mathrm{\Delta }_0`$ resulting in a peak-like feature in the self energy. For our case, $`ϵ_1=74`$ meV and $`ϵ_2=88`$ meV. Because the spectral weight of the mode is maximal near $`\stackrel{}{Q}=(\pi ,\pi )`$, the $`M`$ points of the zone, which are connected by $`\stackrel{}{Q}`$, benefit mostly. This results in stronger features in the self energy near the $`M`$ points compared to e.g. the nodal points. The peaked structure in the imaginary part of the self energy results (via Kramers-Kronig relations) in an enhancement of the real part of the renormalization function for $`|ϵ|<ϵ_1`$, and a reduction of it for $`|ϵ|>ϵ_2`$, as shown in Fig. 1. This leads to a renormalization of the low-energy dispersion of the spectra compared to the high energy part. Since the experimental energy width of the neutron resonance is smaller than the variation in energy of typical features in the self energy, this confirms our assumption that the energy width of the mode is not relevant.
The spectral function is obtained by
$$A(\stackrel{}{k},ϵ)=2\text{Im}\left[\left(\widehat{G}^R(\stackrel{}{k},ϵ)^1\widehat{\mathrm{\Sigma }}^R(\stackrel{}{k},ϵ)\right)^1\right]_{11}$$
(5)
In Fig. 2 we show the spectral functions for momentum cuts through the $`M`$ point towards $`Y=(\pi ,\pi )`$ ($`MY`$ cut), and parallel to $`MY`$ through the nodal point. Due to particle-hole coherence factors, there are quasiparticle peaks at $`M`$ on both sides of the chemical potential. On the negative energy side, the peak is more pronounced since $`\xi _M`$ is negative, and a strong dip feature is present. The asymmetry of the dip feature is a combined effect of the $`\widehat{\tau }_3`$ component of $`\widehat{\mathrm{\Sigma }}`$, which introduces particle-hole asymmetry, and the inherent particle-hole asymmetry of the band structure near the $`M`$ point. Going from $`M`$ towards the Fermi surface (Fig. 2, bottom), the hump feature quickly loses weight as observed in ARPES.
In the top panel of Fig. 2 we show spectra for a cut parallel to $`MY`$ through the order parameter node at the Fermi surface. Near the node there is only one peak crossing the chemical potential. The dip-hump features are very weak near the node and are presumably overshadowed by the additional lifetime effects due to the continuum part of the spin susceptibility. Note the much broader peaks for higher energies, $`|ϵ|>80`$ meV, compared to the sharper peaks near the chemical potential, as observed in ARPES experiments.
In Fig. 3 we present our results for the dispersion obtained from the maxima of the occupied part of the spectral function, $`A(\stackrel{}{k},ϵ)f(ϵ)`$. Near the $`M`$ point we observe an almost dispersionless strong peak feature at roughly the gap energy $`\mathrm{\Delta }_0`$, and a weaker hump feature at slightly below $`ϵ_2`$, consistent with experimental finding c). The peak feature, which without interaction with the mode would be at $`E_M`$, is pushed towards the chemical potential, thus ending up close to $`\mathrm{\Delta }_0`$ for not too small coupling constants. The position of the hump feature is strongly dependent on the coupling constant. We adjusted $`g`$ to reproduce the experimental value of about -130 meV for the hump feature at $`M`$; this choice also results in the weak dispersion of the hump feature as observed in experiment. As one goes away from $`M`$, the dispersion of the hump extents further below $`ϵ_2`$ and the peak starts to show dispersion, until a characteristic break in the dispersion with a jump at $``$-80 meV develops. This is exactly the experimental finding in point b). Note the stability of the characteristic -80 meV energy value for the break/dip feature throughout the zone. This is a result of the dominance of the region near $`M`$ in the $`\stackrel{}{q}`$ sum in Eq. 1, which sets the energy scale. Thus we confirm point d) of the experimental findings.
As the nodal point is approached, the self energy becomes weaker due to the momentum dependence of the mode. The sudden change in the linewidth for a cut parallel to $`MY`$ through the node (panel 2 of Fig. 3), as discussed in Fig. 2, occurs around -80 meV, in accordance with point a). We still observe a weak break feature, which will be smeared out by additional lifetime broadening from the continuum part of the susceptibility. This weak break is also reduced for a smaller coupling (panel 1), or if the Lorentzian in Eq. 3 is replaced by a Gaussian. Note that in accordance with experiments, the velocity near the nodal point is reduced compared to that for higher binding energies, causing a velocity ‘kink’.
Knowing the spectral function throughout the zone, we are able to calculate the tunneling spectra given a tunneling matrix element $`T_{\stackrel{}{k}\stackrel{}{p}}`$. From the SIN tunneling current $`I(V)`$ one obtains the differential conductance, $`dI/dV`$. As usual we neglect the energy dependence of the SIN matrix element $`|M_\stackrel{}{k}|^2=2e_\stackrel{}{p}|T_{\stackrel{}{k}\stackrel{}{p}}|^2A_N(\stackrel{}{p},ϵ)`$, where $`A_N`$ is the spectral function of the normal metal. The SIN tunneling current is then given by
$`I(V)={\displaystyle \underset{\stackrel{}{k}}{}}|M_\stackrel{}{k}|^2{\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle \frac{dϵ}{2\pi }}A(\stackrel{}{k},ϵ)\left\{f(ϵ)f(ϵ+eV)\right\}`$ (6)
In the top panels of Fig. 4, we show results for the SIN $`dI/dV`$ for several coupling strengths. We model the tunneling matrix element for two extreme cases: for incoherent tunneling we assume a constant $`|M_\stackrel{}{k}|^2=M_0^2`$, whereas for coherent tunneling we use $`|M_\stackrel{}{k}|^2=\frac{1}{4}M_1^2(\mathrm{cos}k_x\mathrm{cos}k_y)^2`$.
In both cases, we observe a clear asymmetry with a dip-hump structure on the negative bias side and a very weak feature on the positive side of the spectrum, as in experiment . The pronounced asymmetry is a result of the shallow band near the $`M`$ point, $`\xi _M\mathrm{\Omega }`$, which enhances the coupling to the resonance mode for populated states. Note that the hump position is strongly dependent on the coupling constant in contrast to the position of the dip minimum. The asymmetry in the peak height on either side of the spectrum is sensitive to the coupling constant too. In weak coupling the negative bias peak is higher due to the Van Hove singularity at the $`M`$ point. For stronger coupling the pronounced dip at negative bias reduces the height of the coherence peak on this side and shifts the hump to higher energies. For $`g`$=$`0.65`$ eV (full lines in Fig. 4) the peaks at positive and negative bias have roughly the same height, as in experiment .
For an SIS junction, the single particle tunneling current is given in terms of the spectral functions by
$`I(V)=2e{\displaystyle \underset{\stackrel{}{k}\stackrel{}{p}}{}}|T_{\stackrel{}{k}\stackrel{}{p}}|^2{\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle \frac{dϵ}{2\pi }}A(\stackrel{}{k},ϵ)A(\stackrel{}{p},ϵ+eV)`$ (7)
$`\times \left\{f(ϵ)f(ϵ+eV)\right\}`$ (8)
Again we show results for incoherent tunneling ($`|T_{\stackrel{}{k}\stackrel{}{p}}|^2=T_0^2`$) and for coherent tunneling with conserved parallel momentum, $`|T_{\stackrel{}{k}\stackrel{}{p}}|^2=\frac{1}{16}T_1^2(\mathrm{cos}k_x\mathrm{cos}k_y)^4\delta (\stackrel{}{k}_{||}\stackrel{}{p}_{||})`$. We show the SIS $`dI/dV`$ in the bottom panels of Fig. 4. Our theoretical SIS curves for incoherent tunneling resemble very closely the experimental results for BSCCO, unlike for coherent tunneling which exhibits negative $`dI/dV`$ regions due to the strong anisotropy of $`T_{\stackrel{}{k}\stackrel{}{p}}`$. Note that the dip-hump feature is strong on both sides for an SIS junction in contrast to the SIN results.
In conclusion, we have shown that the momentum dispersion of the ARPES spectra, as detailed in recent experiments, can be explained by a simple model which has as components A) a flat band region near the chemical potential in the normal state dispersion near the $`(\pi ,0)`$ point of the zone; B) a nearly dispersionless bosonic mode which is peaked in momentum near the $`(\pi ,\pi )`$ point, and which interacts with the fermionic degrees of freedom. The theoretical tunneling spectra obtained with the same parameter set are consistent with the experimental findings of an asymmetry of the peak-dip-hump structure in SIN tunneling spectra.
The authors would like to thank A. Kaminski and J.C. Campuzano for discussions concerning their photoemission data, and J. Zasadzinski concerning tunneling data. This work was supported by the U.S. Dept. of Energy, Basic Energy Sciences, under Contract No. W-31-109-ENG-38.
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# Literatur
On the Quantum Quarter Plane and the Real Quantum Plane
Konrad Schmüdgen
Fakultät für Mathematik und Informatik
Universität Leipzig, Augustusplatz 10, 04109 Leipzig, Germany
E-mail: schmuedg@mathematik.uni-leipzig.de
## Zusammenfassung
Suppose $`q\pm 1`$ is a complex number of modulus one. Let $`𝒪(_q^2)`$ be the $``$-algebra with two hermitean generators $`x`$ and $`y`$ satisfying the relation $`xy=qyx`$. Using operator representations of the $``$-algebra $`𝒪(_q^2)`$ on Hilbert space and the Weyl calculus of pseudodifferential operators we construct $``$-algebras of “functions” on the quantum quarter plane $`_q^{++}`$ and on the real quantum plane $`_q^2`$ which are left module $``$-algebras for the Hopf $``$-algebra $`𝒰_q(gl_2())`$. We define a family $`h_k`$, $`k^2`$, of covariant positive linear functionals on these $``$-algebras and study the actions of the $``$-algebras $`𝒪(_q^2)`$ and $`𝒰_q(gl_2())`$ on the associated Hilbert spaces. Quantum analogs of the partial Fourier transforms and the Fourier transform are found. A differential calculus on the “function” $``$-algebras is also developed and investigated.
Mathematics Subject Classifications (1991), 17B37,81R50,47D40
0. Introduction
Suppose that $`q\pm 1`$ is a complex number of modulus one. Let $`𝒪(_q^2)`$ be the $``$-algebra with two hermitean generators $`x`$ and $`y`$ satisfying the relation
$$xy=qyx.$$
In quantum group theory $`𝒪(_q^2)`$ is usually called the coordinate $``$-algebra of the real quantum plane. It is well-known that $`𝒪(_q^2)`$ is a left module $``$-algebra of the Hopf $``$-algebra $`𝒰_q(gl_2())`$ with action given by formulas (75)–(77) below. However these structures are not sufficient in order to study analytic properties of the real quantum plane. In the undeformed case $`q=1`$ the $``$-algebra $`𝒪(_q^2)`$ is just the polynomial algebra $`[x,y]`$ in two hermitean indeterminates $`x`$ and $`y`$ equipped with the usual action of the Lie algebra $`gl_2()`$. In this situation we can replace the polynomial algebra by the larger $``$-algebra $`𝒜(^2):=[x,y]+C_0^{\mathrm{}}(^2)`$ of functions on $`^2`$ and extend the action of $`gl_2()`$ to $`𝒜(^2)`$. On the algebra $`C_0^{\mathrm{}}(^2)`$ we can study differential and integral calculus and so we can develop analysis on $`^2`$. Roughly speaking, for the real quantum plane we try to proceed in a similar way.
An interesting approach to the quantum space $`_q^d`$ has been developed by M. Rieffel \[R\] in the framework of his theory of deformation quantization. This approach is essentially based on $`C^{}`$-algebras.
In the undeformed case the points of $`^2`$ are in one-to-one correspondence to the well-behaved irreducible $``$-representations (see \[S1\]) of the polynomial algebra $`[x,y]`$ in Hilbert space. Thus we are lead to look for well-behaved irreducible $``$-representations of the $``$-algebra $`𝒪(_q^2)`$. This problem has been studied in \[S2\]. In this paper we consider four irreducible well-behaved $``$-representations of $`𝒪(_q^2)`$. They are defined as follows. We fix a real number $`\gamma `$ such that $`q=e^{2\pi i\gamma }`$ and two reals $`\alpha `$ and $`\beta `$ such that $`\alpha \beta =\gamma `$. Let $`𝒫`$ and $`𝒬`$ be the self-adjoint operators on the Hilbert space $`L^2()`$ given by $`(𝒫f)(x)=(2\pi i)^1f^{}(x)`$ and $`𝒬f=xf(x)`$. Then, for $`ϵ,ϵ^{}\{+,\}`$ there exists an irreducible $``$-representation $`\rho _{ϵϵ^{}}`$ of the $``$-algebra $`𝒪(_q^2)`$ on $`L^2(^2)`$ such that
$$\rho _{ϵϵ^{}}(x)=ϵe^{2\pi \alpha 𝒬},\rho _{ϵϵ^{}}(y)=ϵ^{}e^{2\pi \beta 𝒫}.$$
Because of the spectra of operators $`\rho _{ϵϵ^{}}(x)`$ and $`\rho _{ϵϵ^{}}(y)`$, we think of the $``$-representation $`\rho _{ϵϵ^{}}`$ as realization of the algebra $`𝒪(_q^2)`$ on the quantum quarter plane $`_q^{ϵϵ^{}}`$.
Let us sketch the main ideas of our investigations and begin with the $``$-representation $`\rho _{++}`$ corresponding to the (open) quantum quarter plane $`_q^{++}`$. We want to construct “functions on the quantum quarter plane which are vanishing at the boundary”. In order to do so, we define these “functions” as pseudodifferential operators by means of the Weyl calculus \[Fo\], \[St\], that is,
$$Op(a)=\gamma \widehat{a}(\alpha s,\beta t)e^{2\pi i(s\alpha 𝒬+t\beta 𝒫)}𝑑s𝑑t.$$
As symbol class we take the set of all functions $`a(x_1,x_2)`$ on $`^2`$ which are in the intersection of domains of operators $`e^{2\pi c_1𝒬_1}e^{2\pi d_1𝒫_1}e^{2\pi c_2𝒬_2}e^{2\pi d_2𝒫_2}`$, where $`c_1,c_2,d_1,d_2`$. This set is a $``$-algebra, denoted $`𝔄(^2)`$, with respect to the twisted product and involution of pseudodifferential operators. The first aim of our construction is to extend the algebraic structure and the left action of $`𝒰_q(gl_2)`$ to the direct sum
$$𝒜(_q^{++}):=𝒪(_q^2)+𝔄(^2)$$
such that $`𝒜(_q^{++})`$ becomes a left $`𝒰_q(gl_2())`$-module $``$-algebra. We think of $`𝒜(_q^{++})`$ as counterpart of the $``$-algebra $`𝒜(^{++}):=[x,y]+C_0^{\mathrm{}}(^{++})`$ of functions on the quarter plane $`^{++}`$ equipped with the action of Lie algebra $`gl_2()`$. For each $`k^2`$ there exists a faithful linear functional $`h_k`$ on the $``$-algebra $`𝔄(^2)`$ which is covariant with respect to the left action of $`𝒰_q(gl_2())`$. These functionals $`h_k`$ can be viewed as quantum anologs of the state on the $``$-algebra $`C_0^{\mathrm{}}(^{++})`$ given by the Lebesgue measure. Further, there are two $`U_q(gl_2)`$-covariant differential calculi on the algebra $`𝒪(_q^2)`$ invented in \[PW\] and \[WZ\]. We extend one of these calculi to a differential calculus on the larger algebra $`𝒜(_q^{++})`$. Thus the key ingredients for a differential and integral calculus on the quantum quarter plane $`_q^{++}`$ are developed. Similar considerations can be done for three other quantum quarter planes. We strongly believe that this approach, with some technical modifications, could serve as a guide-line for the constructions of other non-compact quantum spaces as well.
The next main step is the construction of the function $``$-algebra for the real quantum plane $`_q^2`$. The idea is to obtain the quantum plane by “gluing together the four quantum quarter planes at the coordinate axis”. In order to do so, it is natural to begin with the direct sum $`𝔄_0(^2)_4`$ of the four $``$-algebras $`𝔄(^2)`$ corresponding to the quantum quarter planes. The elements of $`𝔄_0(^2)_4`$ are interpreted as “functions on the quantum plane which are vanishing at the coordinate axis”. As in the case of the “ordinary” plane we consider $`𝔄_0(^2)_4`$ as subspace of the Hilbert space obtained from a covariant positive linear functional. Then the generators $`E^{}:=q^{1/2}(qq^1)E\mathrm{and}F^{}:=q^{1/2}(qq^1)F`$ of $`U_q(gl_2[))`$ act on $`𝔄_0(^2)_4`$ as symmetric operators which are not essentially self-adjoint. In order to remedy this defect, we extend $`𝔄_0(^2)_4`$ to a larger $``$-algebra $`𝔄(^2)_4`$ by allowing, roughly speaking, symbols having singularities. All hermitean generators $`E^{},F^{},q^{1/4}K_1,q^{1/4}K_2`$ of $`𝒰_q(gl_2())`$ and $`x,y`$ of $`𝒪(_q^2)`$ act on the larger function $``$-algebra $`𝔄(^2)_4`$ by essentially self-adjoint operators. Because of the singularities of symbols, we do not get actions of the whole $``$-algebras $`𝒰_q(gl_2())`$ and $`𝒪(_q^2)`$ on $`𝔄(^2)_4`$.
The aims and main steps of our construction are explained, at least to some extent, by the preceding discussion. However, the rigorous undertaking of this program requires a number of technical lemmas on unbounded operator theory, on quantum groups, and on the Weyl calculus. We have collected these results in a rather long preliminary Section 1. Moreover, notation and terminology are fixed in Section 1. The reader might start at Section 2.
Let us describe the content of this paper more in detail. In Section 2 we investigate the left $`𝒰_q(gl_2())`$-module $``$-algebra $`𝒪(_q^2)`$ and the covariant first order differential calculus on $`𝒪(_q^2)`$. In Section 3 the corresponding formulas and structures are extended to a larger auxilary $``$-algebra $`𝒲`$. This $``$-algebra contains in an operator representation on the Hilbert space $`L^2(^2)`$ the operators $`e^{2\pi i(s\alpha 𝒬+t\beta 𝒫)},s,t`$. In Section 4 the left action of $`𝒰_q(gl_2())`$ on these operators is used in order to derive a left action on operators $`Op(a)`$ and so on symbols $`a𝔄(^2)`$. In this manner, the symbol algebra $`𝔄(^2)`$ and the direct sum $`𝒜(_q^{++})=𝒪(_q^2)+𝔄(^2)`$ of vector spaces become left $`𝒰_q(gl_2())`$-module $``$-algebras. This is the function algebra of the quantum quarter plane $`_q^{++}`$ mentioned above. We also extend the differential calculus of $`𝒪(_q^2)`$ to the function algebra $`𝒜(_q^{++})`$. In Section 5 we define for each $`k^2`$ a $`𝒰_q(gl_2())`$-covariant faithful positive linear functional $`h_k`$ on the $``$-algebra $`𝔄(^2)`$ by a weighted integral over the symbol. The $``$-representations $`\psi _k`$ of the $``$-algebras $`𝒪(_q^2)`$ and $`𝒰_q^{tw}(gl_2())`$ on the associated unitary space $`𝔄_k:=(𝔄(^2),,_k)`$ are described on the generators. The $``$-representations $`\psi _k`$ and product and involution of the $``$-algebra $`𝒜(_q^{++})`$ are transformed by a unitary transformation to the Hilbert space $`L^2(^2)`$. These transformed structures are essentially used for the construction of the quantum plane in Section 6. In the last subsection of Section 5 a uniqueness theorem for the covariant functional $`h_k`$ is proved. Section 6 is devoted to the construction of the real quantum plane from four quantum quarter planes. The function algebra of $`^2`$ can be thought as direct sum of function algebras of the four quarter planes with boundary conditions $`f(+0,y)=f(0,y)`$ and $`f(x,+0)=f(x,0)`$. We first give an equivalent formulation of this picture and use then the corresponding formulas as motivation for the definitions of structures of the real quantum plane. We also define three unitary operators $`_x^q`$, $`_y^q`$ and $`^q`$ which interchange up to some powers of the generators $`K_1`$ and $`K_2`$ the coordinate operators $`x`$, $`y`$ and the corresponding $`q`$-deformed partial derivatives $`𝒟_x^q`$, $`𝒟_y^q`$, respectively. The unitaries $`_x^q`$, $`_y^q`$, and $`^q`$ can be considered as quantum analogs of the partial Fourier transforms and the Fourier transform, respectively.
1. Preliminaries
1.1 Algebraic Preliminaries
All algebras in this paper are over the complex field. All the notions and facts on Hopf algebras and quantum groups used in this paper can be found in \[Mo\] and \[KS\], see also \[FRT\].
Let $`𝒰`$ be a Hopf algebra. We use the Sweedler notation $`\mathrm{\Delta }(f)=f_{(1)}f_{(2)}`$ for the comultiplication $`\mathrm{\Delta }(f)`$ of $`f𝒰`$. A left $`𝒰`$-module algebra $`𝒵`$ is an algebra (without unit in general) which is a left $`𝒰`$-module with left action $``$ such that
$$f(zz^{})=(f_{(1)}z)(f_{(2)}z^{}),z,z^{}𝒵,f𝒰.$$
(1)
A dual pairing of two Hopf algebras $`𝒰`$ and $`𝒜`$ is a bilinear mapping $`,:𝒰\times 𝒜`$ such that
$$\mathrm{\Delta }(f),a_1a_2=f,a_1a_2,f_1f_2,a=f_1f_2,\mathrm{\Delta }(a),$$
$$f\mathrm{,1}=\epsilon (f),1,a=\epsilon (a),S(f),a=f,S(a)$$
for all $`f,f_1,f_2𝒰`$ and $`a,a_1,a_2𝒜`$. By a dual pairing of two Hopf $``$-algebras $`𝒰`$ and $`𝒜`$ we mean a dual pairing $`,`$ of the Hopf algebras $`𝒰`$ and $`𝒜`$ which has the additional property that
$$f^{},a=\overline{f,S(a)^{}}\text{and}f,a^{}=\overline{S(f)^{},a},f𝒰,a𝒜.$$
(2)
Let $`.,.`$ be a dual pairing of Hopf algebras $`𝒰`$ and $`𝔄`$. Any right $`𝔄`$-comodule algebra $`𝒵`$ is a left $`𝒰`$-module algebra with left action $``$ defined by
$$fz=f,z_{(1)}z_{(0)},f𝒰,z𝒵,$$
(3)
where $`\varphi (z)=z_{(0)}z_{(1)}`$ is the Sweedler notation for the right coaction $`\varphi `$.
Lemma 1. Suppose that $`.,.`$ is a dual pairing of Hopf $``$-algebras $`𝒰`$ and $`𝔄`$. If $`𝒵`$ is a right $`𝔄`$-comodule $``$-algebra with right coaction $`\varphi :𝒵𝒵𝔄`$, then the associated left action of $`𝒰`$ on $`𝒵`$ satisfies the condition
$$(fz)^{}=(S(f)^{})z^{},f𝒰,z𝒵.$$
(4)
Proof. Since $`𝒵`$ is a $`𝔄`$-comodule $``$-algebra, the coaction $`\varphi `$ preserves the involution, so that
$$\varphi (z^{})(z^{})_{(0)}(z^{})_{(1)}=(z_{(0)})^{}(z_{(1)})^{}\varphi (z)^{}.$$
Using this condition and the second relation of (2) we conclude that
$$\begin{array}{c}\hfill S(f)^{}z^{}=S(f)^{},(z^{})_{(1)}(z^{})_{(0)}=\overline{f,(z_{(1)})^{}}(z_{(0)})^{}=(fz)^{}.\mathrm{}\end{array}$$
From now on we suppose that $`𝒰`$ is a Hopf $``$-algebra. Equation (4) in Lemma 1 gives the motivation for the following definition: A $``$-algebra $`𝒵`$ is called a left $`𝒰`$-module $``$-algebra if $`𝒵`$ is a left $`𝒰`$-module algebra with left action $``$ such that equation (4) holds. Then Lemma 1 says any right $`𝔄`$-comodule $``$-algebra $`𝒵`$ is a left $`𝒰`$-module $``$-algebra with respect to the associated left action.
Suppose that $`𝒵`$ is a left $`𝒰`$-module $``$-algebra with left action $``$ and let $`\chi `$ be a linear functional on the Hopf $``$-algebra $`𝒰`$. We shall say that a linear functional $`h`$ on $`𝒵`$ is covariant with respect to $`\chi `$ if
$$h(fz)=\chi (f)h(z),f𝒰,z𝒵.$$
(5)
Suppose for a moment that $`h0`$ is covariant with respect to $`\chi `$. Then it follows from the conditions of a left action that $`\chi `$ is a character, that is,
$$\chi (fg)=\chi (f)\chi (g),f,g𝒰,\mathrm{and}\chi (1)=1.$$
(6)
If in addition $`h`$ is hermitian (that is, $`h(z^{})=\overline{h(z)}`$ for $`z𝒵`$), then we conclude from (2) and (4) that
$$\overline{\chi (f)}=\chi (S(f)^{}),f𝒰.$$
(7)
Note that a linear function $`h`$ on the $`𝒰`$-module algebra $`𝒵`$ is invariant if and only if $`h`$ is covariant with respect to the counit $`\epsilon `$.
Lemma 2. Let $`h`$ be a linear functional on the left $`𝒰`$-module $``$-algebra $`𝒵`$ and set
$$y,x:=h(x^{}y),x,y𝒵.$$
(8)
Consider the following three conditions:
(i) $`h`$ is covariant with respect to $`\chi `$.
(ii) $`\chi (f)y,x=f_{(2)}y,S(f_{(1)})^{}x`$ for $`f𝒰,x,y𝒵.`$
(iii) $`fy,x=\chi (f_{(2)})y,f_{(1)}^{}x`$ for $`f𝒰,x,y𝒵.`$
Then we have (i)$``$(ii)$``$(iii). If $`𝒵`$ has a unit, then (iii)$``$(i) and so all three conditions are equivalent.
Proof. (i)$``$(ii): Using the formulas (8), (4) and (1) and the fact that $`SS=\mathrm{id}`$ in any Hopf $``$-algebra we get
$`f_{(2)}y,S(f_{(1)})^{}x`$ $`=h((S(f_{(1)})^{}x)^{}(f_{(2)}y))=h((S(S(f_{(1)})^{})^{}x^{})(f_{(2)}y))`$
$`=h((f_{(1)}x^{})(f_{(2)}y))=h(fx^{}y)=\chi (f)y,x.`$
(ii)$``$(iii): Using once more the relation $`SS=\mathrm{id}`$ and condition (ii) we compute
$`fy,x`$ $`=f_{(2)}y,\epsilon ((f_{(1)}^{})x=f_{(3)}y,S^1(f_{(2)}^{})f_{(1)}^{}x`$
$`=f_{(3)}y,S(f_{(2)})^{}(f_{(1)}^{}x)=\chi (f_{(2)})y,f_{(1)}^{}x.`$
(iii)$``$(i): Applying condition (iii) with $`x=1`$ we obtain
$`h(fy)`$ $`=fy\mathrm{,1}=\chi (f_{(2)})y,f_{(1)}^{}1=\chi (f_{(2)})y,\epsilon (f_{(1)}^{})1`$
$`=\chi (f_{(2)})\overline{\epsilon (f_{(1)}^{})}h(1^{}y)=\chi (f)h(y).\mathrm{}`$
The special case where $`\chi `$ is the counit $`\epsilon `$ and $`𝒵`$ has a unit will be stated separately as
Corollary 3. A linear functional $`h`$ on the left $`𝒰`$-module $``$-algebra $`𝒵`$ with unit is invariant (that is, $`h(fz)=\epsilon (f)h(z)`$ for $`f𝒰`$ and $`z𝒵`$) if and only if $`fy,x=y,f^{}x`$ for all $`x,y𝒵`$ and $`f𝒰`$.
Suppose $`h`$ is an invariant linear functional on the left $`𝒰`$-module $``$-algebra $`𝒵`$ such that the form (8) is a scalar product. Then, by the implication (i)$``$(iii) of Lemma 2, the left action of $`𝒰`$ on $`𝒵`$ is a $``$-representation of the $``$-algebra on the unitary space $`(𝒵,,`$).
Let $`\sigma _1`$ and $`\sigma _2`$ be automorphisms of an algebra $`𝒵`$. Recall that a linear mapping $`𝒟`$ of $`𝒵`$ is called at $`(\sigma _1,\sigma _2)`$-derivation if
$$𝒟(z_1z_2)=\sigma _1(z_1)𝒟(z_2)+𝒟(z_1)\sigma _2(z_2),z_1,z_2𝒵.$$
A first order differential calculus (briefly, a FODC) over an algebra $`𝒵`$ is a $`𝒵`$-bimodule $`\mathrm{\Gamma }`$ equipped with a linear mapping $`\mathrm{d}:𝒵\mathrm{\Gamma }`$ such that $`\mathrm{\Gamma }`$ is the linear span of elements $`z_1\mathrm{d}z_2,z_1,z_2𝒵`$, and $`\mathrm{d}(z_1z_2)=z_1\mathrm{d}z_2+\mathrm{d}z_1z_2`$ for $`z_1,z_2𝒵`$.
Next we recall the definitions of the Hopf algebras $`𝒰_q(gl_2)`$, $`𝒰_q(sl_2)`$ and $`𝒪(GL_q(2))`$ as used in what follows. Let $`𝒰_q(gl_2)`$ be the complex unital algebra with generators $`E`$, $`F`$, $`K_1`$, $`K_2`$, $`K_1^1`$, $`K_2^1`$ and defining relations
$$K_1K_2=K_2K_1,K_jK_j^1=K_j^1K_j=1\mathrm{for}j=\mathrm{1,2},$$
$$K_1EK_1^1=q^{1/2}E,K_2EK_2^1=q^{1/2}E,K_1FK_1^1=q^{1/2}F,K_2FK_2^1=q^{1/2}F,$$
$$EFFE=\lambda ^1(K^2K^2),$$
where we set
$$K:=K_1K_2^1\mathrm{and}\lambda :=qq^1.$$
The algebra $`𝒰_q(\mathrm{gl}_2)`$ is a Hopf algebra with structure maps given on the generators by
$$\mathrm{\Delta }(K_j)=K_jK_j,\mathrm{\Delta }(E)=EK+K^1E,\mathrm{\Delta }(F)=FK+K^1F,$$
$$\epsilon (K_j)=1,\epsilon (E)=\epsilon (F)=0,S(K_j)=K_j^1,S(E)=qE,S(F)=q^1F$$
for $`j=\mathrm{1,2}.`$ Note that the element $`L:=K_1K_2`$ is group-like and central. Let $`𝒰_q(sl_2)`$ denote the subalgebra of the algebra $`𝒰(gl_2)`$ generated by the elements $`E,F,K,K^1`$. Clearly, $`𝒰_q(sl_2)`$ is a Hopf subalgebra of $`𝒰(gl_2)`$.
Let $`𝒪(GL_q(2))`$ be the coordinate Hopf algebras of the quantum group $`GL_q(2)`$ and let $`u_{jk},j,k=\mathrm{1,2},`$ be the entries of the corresponding fundamental matrix. There exists a dual pairing $`.,.`$ of the Hopf algebras $`𝒰_q(gl_2)`$ and $`𝒪(GL_q(2))`$. It is determined by the values on the generators $`K_1,K_2,E,F`$ and $`u_{11},u_{12},u_{21},u_{22},`$ respectively, which are given by
$$K_1,u_{11}=K_2,u_{22}=q^{1/2},K_1,u_{22}=K_2,u_{11}=E,u_{21}=F,u_{12}=1$$
(9)
and zero otherwise.
The algebra with generators $`x`$ and $`y`$ satisfying the relation $`xy=qyx`$ is called the coordinate algebra $`𝒪(_q^2)`$ of the quantum plane $`_q^2`$. It is a right $`𝒪(GL_q(2))`$-comodule algebra with right coaction $`\phi `$ given by
$$\phi (x)=xu_{11}+yu_{21},\phi (y)=xu_{12}+yu_{22}.$$
(10)
Let $`\stackrel{ˇ}{𝒪}(_q^2)`$ denote the algebra with geneators $`x,x^1,y,y^1`$ and relations
$$xy=qyx,xx^1=x^1x=1,yy^1=yy^1=1.$$
That is, $`\stackrel{ˇ}{𝒪}(_q^2)`$ is the localization of $`𝒪(_q^2)`$ with respect to the elements $`x`$ and $`y`$ and the algebra $`𝒪(_q^2)`$ can be considered as a subalgebra of $`\stackrel{ˇ}{𝒪}(_q^2)`$.
Assume now that the deformation parameter $`q`$ is of modulus one. Then there exists an involution $`ff^{}`$ on the algebra $`𝒰_q(gl_2)`$ determined by
$$E^{}:=qE,F^{}:=q^1F,K_1^{}:=K_1,K_2^{}:=K_2.$$
(11)
Equipped with this involution the Hopf algebra $`𝒰_q(gl_2)`$ is a Hopf $``$-algebra denoted by $`𝒰_q(gl_2())`$. We often work with the hermitean elements
$$E^{}:=q^{1/2}\lambda E\mathrm{and}F^{}:=q^{1/2}\lambda F$$
of the $``$-algebra $`𝒰_q(gl_2())`$. Further, there is an involution $`ff^{}`$ given by
$$E^{}:=q^1E,F^{}:=qF,K_1^{}:=q^{1/2}K_1,K_2^{}:=q^{1/2}K_2$$
(12)
such that $`𝒰_q(gl_2)`$ becomes a $``$-algebra. It will be denoted by $`𝒰_q^{tw}(gl_2())`$. In Section 5 we study covariant linear functionals with respect to the character $`\chi `$ on $`𝒰_q(gl_2())`$ defined by $`\chi (K_1)=\chi (K_2)=q^{1/2}`$ and $`\chi (E)=\chi (F)=0`$. Then, by Lemma 2, the corresponding left action of $`𝒰_q(gl_2)`$ is a $``$-representation of the $``$-algebra $`𝒰_q^{tw}(gl_2())`$.
The Hopf algebra $`𝒪(GL_q(2))`$ is a Hopf $``$-algebra, denoted $`𝒪(GL_q(2,))`$, with involution determined on the generators by $`u_{jk}^{}=u_{jk},j,k=\mathrm{1,2}`$. The dual pairing $`.,.`$ of the Hopf algebras $`𝒰_q(gl_2)`$ and $`𝒪(GL_q(2))`$ given by (9) is also a dual pairing of the Hopf $``$-algebras $`𝒰_q(gl_2())`$ and $`𝒪(GL_q(2,))`$.
Further, there exist an involution of the algebra $`𝒪(_q^2))`$ given by
$$x^{}=x,y^{}=y$$
(13)
such that this algebra is a $``$-algebra. It is denoted by $`𝒪(_q^2)`$ and called the coordinate $``$-algebra of the real quantum plane. From the preceding formulas we see at once that the right coaction $`\phi `$ of $`𝒪(GL_q(2))`$ on $`𝒪(_q^2)`$ preserves the corresponding involutions, that is, $`𝒪(_q^2)`$ is a right comodule $``$-algebra of the Hopf $``$-algebra $`𝒪(GL_q(2,))`$. Hence, by Lemma 1, $`𝒪(_q^2)`$ is a left module $``$-algebra for the Hopf $``$-algebra $`𝒰(gl_2())`$.
Remark 1. In the literature the involution of $`𝒰_q(gl_2())`$ is often defined by the requirements $`E^{}=E,F^{}=F,K_1^{}=K_1,K_2^{}=K_2`$. The latter defines indeed an involution which makes $`𝒰_q(gl_2)`$ into a Hopf $``$-algebra. However, with respect to this involution the dual pairing with $`𝒪(GL_q(2,))`$ is not a dual pairing of Hopf $``$-algebras and the $``$-algebra $`𝒪(_q^2)`$ is not a left module $``$-algebra.
1.2. Operator-theoretic Preliminaries
First we fix some notation. Let $`𝒥(a,b)`$ be the strip $`\{z:a<\mathrm{Im}z<b\}`$. The Fourier transform $``$ and its inverse are used in the form
$$(f)(x)=\widehat{f}(x)=e^{2\pi itx}f(t)𝑑t,(^1f)(x)=e^{2\pi itx}f(t)𝑑t.$$
(14)
Throughout, we denote the domain of an operator $`T`$ by $`𝒟(T)`$ and the scalar product of $`L^2(^n)`$ by $`(,)`$. Let $`𝒫`$ and $`𝒬`$ be the self-adjoint operators on the Hilbert space $`L^2()`$ defined by
$$(𝒫f)(x)=\frac{1}{2\pi i}f^{}(x)\text{and}(𝒬f)(x)=xf(x).$$
The operators $`𝒫`$ and $`𝒬`$ are unitarily equivalent by the Fourier transform
$$𝒬^1=𝒫,𝒫^1=𝒬.$$
(15)
The first assertion of the following lemma describes the domain and the action of the operators $`e^{2\pi \beta 𝒫}`$ for real $`\beta `$. We formulate the result for $`\beta >0`$; the case $`\beta <0`$ is completely similar.
Lemma 4. (i): Suppose that $`\beta >0`$. Let $`g(z)`$ be a holomorphic function on the strip $`𝒥(0,\beta )`$ such that
$$\underset{0<y<\beta }{sup}\underset{\mathrm{}}{\overset{\mathrm{}}{}}|g(x+iy)|^2𝑑x<\mathrm{}.$$
Then there exist functions $`g(x)L^2()`$ and $`g_\beta (x)L^2()`$ such that $`lim_{y0}g_y=g`$ and $`lim_{y\beta }g_y=g_\beta `$ in $`L^2()`$, where $`g_y(x):=g(x+iy)`$ for $`x`$ and $`y(0,\beta )`$. Setting $`g(x+i\beta ):=g_\beta (x),x`$, we have
$$\underset{n\mathrm{}}{lim}g(x+n^2i)=g(x)\mathrm{and}\underset{n\mathrm{}}{lim}g(x+(\beta n^2)i)=g(x+\beta i)\mathrm{a}.\mathrm{e}.\mathrm{on}.$$
The function $`g`$ belongs to the domain $`𝒟(e^{2\pi \beta 𝒫})`$ and we have
$$(e^{2\pi \beta 𝒫}g)(x)=g(x+\beta i).$$
(16)
Conversely, each function $`g`$ in the domain $`𝒟(e^{2\pi \beta 𝒫})`$ arises in this way.
(ii): For any $`\beta `$ and $`\delta >0`$, the vector space $`𝒟_\delta :=\mathrm{Lin}\{e^{\delta x^2+cx}:c\}`$ is a core for the self-adjoint operators $`e^{2\pi \beta 𝒫}`$ and $`e^{2\pi \beta 𝒬}`$.
Proof. \[S2\], Lemma 1–3. $`\mathrm{}`$
Throughout this paper we assume that the deformation parameter $`q\pm 1`$ of modulus one and that $`\gamma `$ is a fixed real number such that
$$q=e^{2\pi i\gamma }.$$
Further, let $`\alpha `$ and $`\beta `$ denote real numbers such that $`\alpha \beta =\gamma `$ and put
$$X=e^{2\pi \alpha 𝒬}\text{and}Y=e^{2\pi \beta 𝒫}.$$
(17)
From (16) it follows that the operators $`X`$ and $`Y`$ defined by (17) satisfy the relation $`XY\eta =qYX\eta `$ for each vector $`\eta 𝒟(XY)𝒟(YX)`$. Therefore, for each $`ϵ,ϵ^{}\{+,\}`$, there exist a unique faithful $``$-representation $`\rho _{ϵϵ^{}}`$ of the $``$-algebra $`𝒪(_q^2)`$ on the domain $`𝔄()`$ such that
$$\rho _{ϵϵ^{}}(x)=ϵX𝔄(),\rho _{ϵϵ^{}}(y)=ϵ^{}Y𝔄().$$
(18)
Let $`𝔄()`$ be the set of entire holomorphic functions $`a(x)`$ on the complex plane satisfying
$$\underset{\delta _1<y<\delta _2}{sup}\underset{\mathrm{}}{\overset{\mathrm{}}{}}|a(x+iy)|^2e^{2sx}𝑑x<\mathrm{}$$
(19)
for all $`s,\delta _1,\delta _2`$, $`\delta _1<\delta _2`$. From Lemma 4 we easily derive that
$$𝔄()=\underset{n,m=\mathrm{}}{\overset{+\mathrm{}}{}}𝒟(X^nY^m)=\underset{n,m=\mathrm{}}{\overset{+\mathrm{}}{}}𝒟(Y^nX^m).$$
Clearly, $`𝔄()`$ is invariant under the Fourier transform and its inverse.
Throughout, we denote by $`f_\alpha `$ the function
$$f_\alpha (x)=2\mathrm{sinh}\pi \beta (2x+\alpha i)$$
(20)
and by $`L_\alpha `$ and $`R_\alpha `$ the operators on the Hilbert space $`L^2()`$ given by
$$L_\alpha =\overline{f_\alpha }(𝒫)e^{2\pi \alpha Q},R_\alpha =e^{2\pi \alpha 𝒬}f_\alpha (𝒫).$$
(21)
Some properties of these operators are collected in the next lemma.
Lemma 5. (i) $`L_\alpha `$ is a closed symmetric operator.
(ii) $`R_\alpha `$ is the adjoint operator of $`L_\alpha `$.
(iii) $`𝔄()`$ is a core for the operator $`L_\alpha `$.
(iv) $`f_\alpha (𝒫)^1𝔄()`$ is a core for the operator $`R_\alpha `$.
Proof. By formula (15), we can replace $`𝒫`$ by $`𝒬`$ and $`𝒬`$ by $`𝒫`$. But then the assertions (i)–(iii) have been stated in \[S2\] and \[S4\].
It remains to prove assertion (iv). First note that $`f_\alpha (𝒫)^1`$ is a bounded normal operator on the Hilbert space $`L^2()`$, so $`_\alpha :=f_\alpha (𝒫)^1𝔄()`$ is a dense linear subspace of $`L^2()`$. We show that
$$(R_\alpha _\alpha )^{}L_\alpha .$$
(22)
Indeed, suppose that $`\zeta 𝒟((R_\alpha _\alpha )^{})`$. Then there exists a vector $`\xi L^2()`$ such that $`R_\alpha \eta ,\zeta =\eta ,\xi `$ for all $`\eta _\alpha `$. Writing $`\eta `$ as $`\eta =f_\alpha (𝒫)^1\eta ^{}`$ with $`\eta ^{}𝔄()`$, we obtain $`e^{2\pi \alpha 𝒬}\eta ^{},\zeta =\eta ^{},\overline{f_\alpha }(𝒫)^1\xi `$. Since $`𝔄()`$ is a core for $`e^{2\pi \alpha 𝒬}`$ by Lemma 4(ii), the latter implies that $`\zeta 𝒟(e^{2\pi \alpha 𝒬})`$ and $`e^{2\pi \alpha 𝒬}\zeta =\overline{f_\alpha }(𝒫)^1\xi `$. Thus, we have $`\zeta 𝒟(\overline{f_\alpha }(𝒫)e^{2\pi \alpha 𝒬})=𝒟(L_\alpha )`$ which proves (22).
By the assertion of (i), (22) implies that $`(R_\alpha _\alpha )^{}(L_\alpha )^{}=R_\alpha `$. But the latter means that $`_\alpha `$ is a core for $`R_\alpha `$. $`\mathrm{}`$
Next we essentially use some results obtained in \[S4\]. We restate them here using a slight different notation. For $`\delta _1,\delta _2`$, $`\delta _1>\delta _2`$, let $`(\delta _1,\delta _2)`$ denote the set of all holomorphic functions $`f`$ on the strip $`𝒥(\delta _1,\delta _2)`$ satisfying
$$\underset{\delta _1<y<\delta _2}{sup}\underset{\mathrm{}}{\overset{\mathrm{}}{}}|f(x+iy)|^2e^{sx^2}𝑑x<\mathrm{}$$
for all $`s>0`$. By Lemma 2 in \[S4\], each function $`f(\delta _1,\delta _2)`$ has a.e. boundary limits $`f(x+i\delta _1)`$ and $`f(x+i\delta _2)`$ on $``$. For notational simplicity we assume that $`\alpha >0`$. With some obvious modifications all results remain valid for $`\alpha <0`$.
We apply Theorem 1 in \[S4\] to the function $`f(x):=2\mathrm{sinh}2\pi \beta x`$ and with $`\alpha `$ replaced by $`\alpha /2`$. Note that $`f(xi\alpha /2)=f_\alpha (x)`$, where $`f_\alpha `$ is defined by $`(\text{20})`$. Then Theorem 1 in \[S4\] can be restated as follows.
Lemma 6. There exist holomorphic functions $`w_\alpha (\alpha \mathrm{,0})`$ and $`v_\alpha (\alpha ,\alpha )`$ such that
$`|w_\alpha (x)|=|v_\alpha (x)|=1a.e.\mathrm{on},`$ (23)
$`w_\alpha (x)=f_\alpha (x)v_\alpha (x\alpha i),v_\alpha (x)=f_\alpha (x)w_\alpha (x\alpha i)a.e.\mathrm{on}.`$ (24)
The functions $`w_\alpha ,v_\alpha `$ are uniquely determined up to a constant factor of modulus one by these properties.
Let $`W_\alpha `$ and $`A_\alpha `$ denote the operator matrices
$$W_\alpha (𝒫)=\left(\begin{array}{cc}w_\alpha (𝒫)& 0\\ 0& v_\alpha (𝒫)\end{array}\right),A_\alpha =\left(\begin{array}{cc}0& L_\alpha \\ R_\alpha & 0\end{array}\right),B_\alpha =\left(\begin{array}{cc}0& e^{2\pi \alpha Q}\\ e^{2\pi \alpha Q}& 0\end{array}\right).$$
(25)
Since $`|w_\alpha |=|v_\alpha |=1`$ a.e. on $``$ by (23) and $`L_\alpha ^{}=R_\alpha `$ by Lemma 5, $`W_\alpha (𝒫)`$ is a unitary operator and $`A_\alpha `$ and $`B_\alpha `$ are self-adjoint operators on the Hilbert space $`L^2()L^2()`$.
Lemma 7. $`W_\alpha (𝒫)^{}A_\alpha W_\alpha (𝒫)=B_\alpha andW_\alpha (𝒫)^{}B_\alpha W_\alpha (𝒫)=A_\alpha `$.
Proof. By (15), we have $`W_\alpha (𝒫)^1=W_\alpha (𝒬)`$, $`e^{2\pi \alpha 𝒬}^1=e^{2\pi \alpha 𝒫}`$ and $`L_\alpha ^1=2\mathrm{sin}\pi \beta (2x\alpha i)e^{2\pi \alpha 𝒫}`$, that is, $`L_\alpha ^1`$ is the operator $`L_f`$ with $`f(x)=2\mathrm{sinh}2\pi \beta x`$ and $`\alpha `$ replaced by $`\alpha /2`$ in the notation of \[S4\]. Thus, under the Fourier transform the assertion $`W_\alpha (𝒫)^{}A_\alpha W_\alpha (𝒫)=B_\alpha `$ is just equation (24) in \[S4\].
Next we prove that $`W_\alpha (𝒫)^{}B_\alpha W_\alpha (𝒫)=A_\alpha `$. Since the self-adjoint operator $`A_\alpha `$ has no proper self-adjoint extension in the same Hilbert space, it suffices to show that $`W_\alpha (𝒫)^{}B_\alpha W_\alpha (𝒫)A_\alpha `$ which means that
$`w_\alpha (𝒫)^{}e^{2\pi \alpha 𝒬}v_\alpha (𝒫)L_\alpha \overline{f_\alpha }(𝒫)e^{2\pi \alpha 𝒬},`$
$`v_\alpha (𝒫)^{}e^{2\pi \alpha 𝒬}w_\alpha (𝒫)R_\alpha e^{2\pi \alpha 𝒬}f_\alpha (𝒫).`$
Note that $`\overline{f_\alpha (x)}=f_\alpha (x)`$. Applying the unitary transformation $`^1`$ and using (15) it follows that the latter relations are equivalent to
$`\overline{w_\alpha (x)}e^{2\pi \alpha 𝒫}v_\alpha (x)f_\alpha (x)e^{2\pi \alpha 𝒫},`$ (26)
$`v_\alpha (x)e^{2\pi \alpha 𝒫}w_\alpha (x)e^{2\pi \alpha 𝒫}f_\alpha (x).`$ (27)
Recall that $`f(x\pm \frac{\alpha }{2}i)=f_{\pm \alpha }(x)`$. Therefore, formula (24) can be rewritten as
$`f_\alpha (x)e^{2\pi \alpha 𝒫}=w_\alpha e^{2\pi \alpha 𝒫}\overline{v_\alpha },`$ (28)
$`e^{2\pi \alpha 𝒫}f_\alpha (x)=v_\alpha e^{2\pi \alpha 𝒫}\overline{w_\alpha }.`$ (29)
Let $`\eta 𝔄()`$. Since $`\overline{\eta }𝒟(f_\alpha (x)e^{2\pi \alpha 𝒫})`$ and hence $`\overline{v_\alpha }\overline{\eta }𝒟(e^{2\pi \alpha 𝒫})`$ by (28), we have $`v_\alpha \eta 𝒟(e^{2\pi \alpha 𝒫})`$. The relations (26) combined with the facts that $`w_\alpha (\alpha \mathrm{,0})`$ and $`v_\alpha (\alpha ,\alpha )`$ imply that $`v_\alpha (x+\alpha i)=f_\alpha (x+\alpha i)w_\alpha (x)`$ a.e. on $``$ (see e.g. formula (7) in \[S4\]). Since $`f_\alpha (x+\alpha i)=f_\alpha (x)`$, we obtain
$$\overline{w_\alpha (x)}e^{\pi \alpha 𝒫}v_\alpha (x)\eta =\overline{w_\alpha (x)}v_\alpha (x+\alpha i)e^{2\pi \alpha 𝒫}\eta =f_\alpha (x)e^{2\pi \alpha 𝒫}\eta .$$
That is, the operators $`\overline{w_\alpha }e^{2\pi \alpha 𝒫}v_\alpha `$ and $`f_\alpha e^{2\pi \alpha 𝒫}`$ coincide on the domain $`𝔄()`$. By Lemma 5, $`𝔄()`$ is a core for the closed operator $`L_\alpha `$ and so for $`f_\alpha (x)e^{2\pi \alpha 𝒫}=`$
$`L_\alpha ^1`$. Thus we conclude that
$$\overline{w_\alpha }e^{2\pi \alpha 𝒫}v_\alpha f_\alpha e^{2\pi \alpha 𝒫}.$$
Next we verify the second relation (24). By (29), $`\overline{w_\alpha (x)}f_\alpha (x)^1\overline{\eta (x)}𝒟(e^{2\pi \alpha 𝒫})`$ and hence $`w_\alpha (x)f_\alpha (x)^1\eta (x)𝒟(e^{2\pi \alpha 𝒫})`$. From (26) we derive that $`w_\alpha (x+\alpha i)=f_\alpha (x+\alpha i)v_\alpha (x)`$ (see formula (6) in \[S4\]). Therefore, for $`\phi (x)=f_\alpha (x)^1\eta (x)`$ with $`\eta 𝔄()`$ we obtain
$`\overline{v_\alpha (x)}e^{2\pi \alpha 𝒫}w_\alpha (x)\phi `$ $`=\overline{v_\alpha (x)}e^{2\pi \alpha 𝒫}x_\alpha (x)f_\alpha (x)^1\eta `$
$`=\overline{v_\alpha (x)}w_\alpha (x+\alpha i)f_\alpha (x+\alpha i)^1e^{2\pi \alpha 𝒫}\eta `$
$`=\overline{v_\alpha (x)}v_\alpha (x)e^{2\pi \alpha 𝒫}\eta =e^{2\pi \alpha 𝒫}f_\alpha (x)\phi .`$
Thus, the operators $`\overline{v_\alpha }e^{2\pi \alpha 𝒫}w_\alpha `$ and $`e^{2\pi \alpha 𝒫}f_\alpha `$ coincide on the dense domain $`f_\alpha (x)^1𝔄()`$. Since $`f_\alpha (𝒫)^1𝔄()`$ is a core for $`R_\alpha `$ by Lemma 5 and $`R_\alpha ^1=e^{2\pi \alpha 𝒫}f_\alpha (x)`$, it follows that $`\overline{v_\alpha }e^{2\pi \alpha 𝒫}w_\alpha e^{2\pi \alpha 𝒫}f_\alpha (x)`$. This proves (24) and completes the proof of Lemma 7. $`\mathrm{}`$
Lemma 8. Let $`c,d`$ and $`\delta _0>0`$. Suppose that $`8|dc|<1`$. Then the vector space $`_{\delta _0}=\mathrm{Lin}\{e_{t,\delta }(x)=e^{2\pi (itx\delta x^2)};t\mathrm{,0}<\delta <\delta _0\}`$ is dense in $`𝔄()`$ with respect to the norm $`_{c,d}=(e^{2\pi c𝒬}+e^{2\pi c𝒬})(e^{2\pi d𝒫}+e^{2\pi d𝒫})`$.
Proof. Since both operators $`e^{2\pi c𝒬}+e^{2\pi c𝒬}`$ and $`e^{2\pi d𝒫}+e^{2\pi d𝒫}`$ on the Hilbert space $`L^2()`$ are self-adjoint and greater than the identity, the operator domain $`_{c,d}:=𝒟((e^{2\pi c𝒬}+e^{2\pi c𝒬})(e^{2\pi d𝒫}+e^{2\pi d𝒫}))`$ equipped with the norm $`_{c,d}`$ is a Hilbert space. Assume to the contrary that the assertion of the lemma is not true. Then there exists a non-zero vector $`\psi _0_{c,d}`$ such that
$$((e^{2\pi c𝒬}+e^{2\pi c𝒬})(e^{2\pi d𝒫}+e^{2\pi d𝒫})\psi _0,(e^{2\pi c𝒬}+e^{2\pi c𝒬})(e^{2\pi d𝒫}+e^{2d𝒫})e_{t,\delta })=0$$
(30)
for $`t`$ and $`0<\delta <\delta _0`$. In order to write this relation in another form, we set
$$\psi =(e^{2\pi c𝒬}+e^{2\pi c𝒬})(e^{2\pi d𝒫}+e^{2\pi d𝒫})\psi _0,\xi _\delta =e^{2\pi \delta x^2}(e^{2\pi cx}+e^{2\pi cx})\psi .$$
(31)
Since $`\psi L^2()`$, it follows that $`e^{n|x|}\xi _\delta (x)L^1()`$ for all $`n`$. Further, we compute
$$\left((e^{2\pi d𝒫}+e^{2\pi d𝒫})e_{t,\delta }\right)(x)=e^{2\pi \delta (x^2d^2)}\left(e^{2\pi ix(t+2d\delta )+\pi t\delta }+e^{2\pi ix(t2d\delta )2\pi t\delta }\right).$$
(32)
Therefore, condition (30) means that the Fourier transform $`\widehat{\xi }_\delta `$ of the $`L^1`$-function $`\xi _\delta `$ satisfies the relation
$$e^{2\pi td}\widehat{\xi }_\delta (t+2\delta d)+e^{2\pi td}\widehat{\xi }_\delta (t2\delta d)=0,t.$$
(33)
Setting
$$\eta _\delta (t):=e^{\pi (2\delta )^1t^2\pi (4\delta d)^1ti}\widehat{\xi }_\delta (t),$$
(34)
equation (33) is obviously equivalent to the relation
$$\eta _\delta (t+2\delta d)=\eta _\delta (t2\delta d),t.$$
(35)
Since $`e^{n|x|}\xi _\delta (x)L^1()`$ for $`n`$, $`\widehat{\xi }_\delta `$ is an entire holomorphic function on the complex plane and so is $`\eta _\delta `$ by (34). Therefore, by (35), the function $`\eta _\delta `$ is bounded on the real axis and hence $`e^{n|t|}\widehat{\xi }_\delta (t)L^1()L^2()`$ for $`n`$ by (34). Consequently, $`\xi _\delta (x)`$ is an entire holomorphic function and $`\xi _\delta (x)𝒟(e^{n|𝒫|})`$ for $`n`$. Hence, by (31),
$$\zeta (x):=(e^{2\pi cx}+e^{2\pi cx})\psi (x)=e^{2\pi \delta x^2}\xi _\delta (x)$$
(36)
is also an entire function.
Computing $`\xi _\delta (x)`$ by the inverse Fourier transform from $`\widehat{\xi }_\delta (t)`$ and using equation (33), we derive that
$$\xi _\delta (x)+e^{8\pi \delta d^2+8\pi i\delta x}\xi _\delta (x+2di)=0,x.$$
(37)
Inserting (36) into (37) we conclude that
$$\zeta (x)+\zeta (x+2di)=0.$$
(38)
Since $`8|dc|<1`$ by assumption, the function $`(e^{2\pi cx}+e^{2\pi cx})^1`$ is holomorphic on the strip $`𝒥_{d,\epsilon }:=\{x:|\mathrm{Im}x|<2d+\epsilon \}`$ for small $`\epsilon >0`$ and
$$\mathrm{inf}\{\left|(e^{2\pi cx}+e^{2\pi cx})^1\right|;x𝒥_{d,\epsilon }\}>0.$$
(39)
Therefore, since $`\xi _\delta (x)𝒟(e^{4\pi d𝒫})`$ as noted above, we conclude from Lemma 4 that the function $`e^{2\pi \delta x^2}\psi (x)=(e^{2\pi cx}+e^{2\pi cx})^1\xi _\delta (x)`$ belongs to the domain $`𝒟(e^{4\pi d𝒫})`$ and
$$e^{4\pi d𝒫}(e^{2\pi \delta x^2}\psi (x))=e^{2\pi \delta (x2di)^2}\psi (x+2di).$$
(40)
Note that $`\psi (x)L^2()`$ by construction and $`\psi (x+2di)L^2()`$ by (36), (38) and (39). Since $`e^{2\pi \delta x^2}\psi (x)\psi (x)`$ and $`e^{4\pi d𝒫}(e^{2\pi \delta x^2}\psi )\psi (x+2di)`$ as $`\delta 0`$ by (40) and the operator $`e^{4\pi d𝒫}`$ is closed, it follows that $`\psi 𝒟(e^{4\pi d𝒫})`$ and $`(e^{4\pi d𝒫}\psi )(x)=\psi (x+2di)`$. Applying this fact and formula (38) we obtain
$$\begin{array}{cc}\hfill (e^{4\pi d𝒫}& \psi ,\psi )=\psi (x+2di)\overline{\psi (x)}dx\hfill \\ & =(e^{2\pi cx}+e^{2\pi cx})(e^{2\pi c(x+2di)}+e^{2\pi c(x+2di)})^1|\psi (x)|^2𝑑x.\hfill \end{array}$$
(41)
Because of the assumption $`8|cd|<1`$, we have $`\mathrm{cos}4\pi cd>0`$. Hence the function under the integral sign in (41) is non-negative. On the other hand, since $`\psi 0`$ by construction, we have $`(e^{4\pi d𝒫}\psi ,\psi )>0`$. Thus we arrived at a contradiction and the assertion of Lemma 8 is proved. $`\mathrm{}`$
Next we turn to some facts on tensor products of certain operators. If $`T_1`$ and $`T_2`$ are closed operators on a Hilbert space $``$, then the symbol $`T_1T_2`$ means the closure of the linear operator on the domain $`𝒟(T_1)𝒟(T_2)`$ in the Hilbert space $``$ defined by $`(T_1T_2)(\eta _1\eta _2)=T_1\eta _1T_2\eta _2`$.
Let $`𝒫_j`$ and $`𝒬_j,j=\mathrm{1,2}`$, be the self-adjoint operators on the Hilbert space $`L^2(^2)`$ given by
$$(𝒫_jf)(x_1,x_2)=\frac{1}{2\pi i}\frac{f}{x_j}(x_1,x_2)\mathrm{and}(𝒬_jf)(x_1,x_2)=x_jf(x_1,x_2).$$
Let $`^{++}:=\{\mu =(\mu _1,\mu _2)^2:\mu _1>0,\mu _2>0\}`$. For $`\mu =(\mu _1,\mu _2)^2`$ we set
$`S(e^{2\pi \mu 𝒬}):=(e^{2\pi \mu _1𝒬_1}I+e^{2\pi \mu _1𝒬_1}I)(Ie^{2\pi \mu _2𝒬_2}+Ie^{2\pi \mu _2𝒬_2}),`$
$`e^{2\pi \mu 𝒬}:=e^{2\pi \mu _1𝒬_1}e^{2\pi \mu _2𝒬_2}.`$
The operators $`S(e^{2\pi \mu 𝒫}),e^{2\pi \mu 𝒫}`$, $`e^{2\pi \mu |𝒫|}`$ and $`e^{2\pi \mu |𝒬|}`$ are defined in a similar manner. Then we have
$$𝒟_{\mu ,\nu }:=𝒟(S(e^{2\pi \mu 𝒬})S(e^{2\pi \nu 𝒫}))=\underset{\epsilon ,\delta _2^2}{}𝒟(e^{2\pi \epsilon \mu 𝒬}e^{2\pi \delta \nu 𝒫})=\underset{\epsilon ,\delta _2^2}{}𝒟(e^{2\pi \epsilon \nu 𝒫}e^{2\pi \delta \mu 𝒬}),$$
(42)
where $`_2^2=\{\epsilon =(\epsilon _1,\epsilon _2):\epsilon _1,\epsilon _2\{1,1\}\}`$ and $`2\pi \epsilon \mu :=(2\pi \epsilon _1\mu _1\mathrm{,2}\pi \epsilon _2\mu _2)`$. If $`\nu ^{++}`$, then $`𝒟_{\mu ,\nu }`$ is the vector space of all holomorphic functions on $`\{(z_1,z_2)^2:|\mathrm{Im}z_j|<|\nu _j|,j=\mathrm{1,2}\}`$ satisfying
$$\underset{|y_j|<|\nu _j|}{sup}|a(x_1+iy_1,x_2+iy_2)|^2e^{4\pi (|\mu _1x_1|+|\mu _2x_2|)}𝑑x_1𝑑x_2<\mathrm{}.$$
(43)
The latter fact can be proved in a similar manner as Lemma 1.1 in \[S2\] using the Paley-Wiener Theorem.
Lemma 9. (i) Suppose that $`\mu =(\mu _1,\mu _2)^2`$ and $`\nu =(\nu _1,\nu _2)^{++}`$.
If $`f𝒟(S(e^{2\pi \nu 𝒫})e^{2\pi \mu 𝒬})𝒟(e^{2\pi \mu 𝒬}S(e^{2\pi \nu 𝒫}))`$, then
$$|f(x_1+iy_1,x_2+iy_2)|\frac{1}{2\pi }((\nu _1|y_1|)(\nu _2|y_2|))^{1/2}e^{2\pi (\mu _1x_1+\mu _2x_2)}e^{2\pi \nu |𝒫|}e^{2\pi \mu 𝒬}f$$
(44)
for $`x_1,x_2,y_1,y_2,|y_1|<\nu _1,|y_2|<\nu _2.`$
(ii) Let $`\mu =(\mu _1,\mu _2),\nu =(\nu _1,\nu _2)_{++}`$. If $`f𝒟_{\mu ,\nu }`$, then
$`|f(x_1+iy_1,x_2+iy_2)|`$
$`{\displaystyle \frac{1}{2\pi }}((\nu _1|y_1|)(\nu _2|y_2|))^{1/2}e^{2\pi (\mu _1|x_1|+\mu _2|x_2|)}{\displaystyle \underset{\epsilon ,\delta _2^2}{}}e^{2\pi \epsilon \nu 𝒫}e^{2\pi \delta \mu 𝒬}f`$ (45)
for $`x_1,x_2,y_1,y_2,|y_1|<\nu _1,|y_2|<\nu _2`$. The vector space $`𝒟_{\mu ,\nu }`$ is contained in the Schwartz space $`𝒮(^2)`$.
Proof. (i): Setting $`g=f`$ and $`\epsilon _j=\nu _j|y_j|,j=\mathrm{1,2}`$, and using formulas (15), we get
$`\left|e^{2\pi (\mu _1x_1+\mu _2x_2)}f(x_1+iy_1,x_2+iy_2)\right|=\left|(e^{2\pi y𝒫}e^{2\pi \mu 𝒬}^1g)(x_1,x_2)\right|`$
$`=\left|(^1e^{2\pi y𝒬}e^{2\pi \mu 𝒫}g)(x_1,x_2)\right|`$
$`=\left|{\displaystyle e^{2\pi i(x_1t_1+x_2t_2)}(e^{2\pi y𝒬}e^{2\pi \mu 𝒫}g)(t_1,t_2)𝑑t_1𝑑t_2}\right|`$
$`\left({\displaystyle e^{4\pi (\epsilon _1|t_1|+\epsilon _2|t_2|)}𝑑t_1𝑑t_2}\right)^{1/2}e^{2\pi \epsilon |𝒬|}e^{2\pi y𝒬}e^{2\pi \mu 𝒫}g_{L^2(^2)}`$
$`=((2\pi \epsilon _1)(2\pi \epsilon _2))^{1/2}e^{2\pi \epsilon |𝒬|2\pi y𝒬}e^{2\pi \mu 𝒫}g`$
$`={\displaystyle \frac{1}{2\pi }}(\epsilon _1\epsilon _2)^{1/2}e^{2\pi ((\nu _1|y_1|)|t_1|y_1t_1+(\nu _2|y_2|)|t_2|y_2t_2)}e^{2\pi \mu 𝒫}g`$
$`{\displaystyle \frac{1}{2\pi }}(\epsilon _1\epsilon _2)^{1/2}e^{2\pi (\nu _1|t_1|+\nu _2|t_2|)}e^{2\pi \mu 𝒫}f`$
$`={\displaystyle \frac{1}{2\pi }}((\nu _1|y_1|)(\nu _2|y_2|))^{1/2}e^{2\pi \nu |𝒫|}e^{2\pi \mu 𝒬}f,`$
which proves (44). Note that by the domain assumptions on $`f`$ the function $`g=f`$ belongs to the corresponding operator domains.
(ii): Since obviously $`e^{2\pi \nu |𝒫|}e^{2\pi \mu 𝒬}f_{\epsilon ,\delta }e^{2\pi \epsilon \nu 𝒫}e^{2\pi \delta \mu 𝒬}f`$, inequality (S0.Ex40) follows at once from (44) applied with $`\mu `$ replaced by $`\epsilon \mu `$.
Finally, we prove that $`𝒮(^2)𝒟_{\mu ,\nu }`$. Let $`a𝒟_{\mu ,\nu }`$. By (42), we have $`a𝒟(e^{\mu |𝒬|})`$ which implies that $`a𝒟(𝒬_1^n𝒬_2^m)`$ for all $`n,m_0`$. Similarly, since $`(a)𝒟_{\nu ,\mu }`$ by (15), we have $`(a)𝒟(𝒬_1^n𝒬_2^m)`$ and hence $`a𝒟(𝒫_1^n𝒫_2^m)`$ for $`n,m_0`$. Both conditions implies that $`a`$ belongs to the Schwartz space $`𝒮(^2)`$ (see, for instance, Example 10.2.14 in \[S1\] for this apparantly weaker characterization of the Schwartz space). $`\mathrm{}`$
Lemma 10. Let $`c=(c_1,c_2),d=(d_1,d_2)^2`$, $`\delta _1>0`$, $`\delta _2>0`$. Suppose that $`8|c_jd_j|<1`$ for $`j=\mathrm{1,2}`$. Then the vector space $`_{\delta _1}_{\delta _2}`$ is dense in $`𝔄(^2)`$ with respect to the norm
$$_{c,d}:=S(e^{2\pi c𝒬})S(e^{2\pi d𝒫}).$$
(46)
Proof. Assume the contrary. Then there exists a vector $`\psi 0`$ which is orthogonal in $`L^2(^2)`$ to $`S(e^{2\pi c𝒬})S(e^{2\pi d𝒫})(_{\delta _1}_{\delta _2})`$. From the assertion of Lemma 8 it follows that $`(e^{2\pi c_j𝒬_j}+e^{2\pi c_j𝒬_j})(e^{2\pi d_j𝒫_j}+e^{2\pi d_j𝒫_j})_{\delta _j}`$ is dense in $`L^2()`$ for $`j=\mathrm{1,2}`$. But this in turn implies that $`\psi =0`$. $`\mathrm{}`$
Let $`𝔄(^2)`$ be the intersection of all domains $`𝒟(S(e^{c𝒬})S(e^{d𝒫}))`$, $`c,d^2`$, or equivalently the vector space of all holomorphic functions on $`^2`$ satisfying condition (44) for all $`\mu ,\nu ^2`$. Let $`\tau `$ denote the locally convex topology on $`𝔄(^2)`$ defined by the family of norms (46), $`c,d,^2`$. Since it obviously suffices to take a countable subfamily of such norms, the topology $`\tau `$ is metrizable. Since $`𝔄(^2)`$ is the intersection of domains $`𝒟(e^{2\pi c𝒬}e^{2\pi d𝒫})`$, $`𝔄(^2)`$ is complete with respect to this topology. Thus $`𝔄(^2)[\tau ]`$ is a Frechet space. The space $`𝔄(^2)`$ will play a crucial role as symbol algebra for the Weyl calculus.
1.3. The Weyl Calculus
In this subsection we shall be concerned with pseudodifferential operators on the Hilbert space $`L^2()`$ defined by means of the Weyl calculus. Our standard references in this matter are the books \[Fo\] and \[St\], see also \[GV\] and \[H\]. The Weyl correspondence assigns an operator $`Op(a)`$ to any function $`a`$ on $`^2`$ such that $`\widehat{a}L^1(^2)`$ by
$$Op(a)=\gamma \widehat{a}(\alpha s,\beta t)e^{2\pi i(s\alpha 𝒬+t\beta 𝒫)}𝑑s𝑑t.$$
(47)
Recall that $`\widehat{a}`$ is the Fourier transform (14) of the function $`a`$. $`\alpha `$ and $`\beta `$ are real numbers such that $`\alpha \beta =\gamma `$ and $`q=e^{2\pi i\gamma }`$. Since $`\widehat{a}L^1(^2)`$, the integral (47) can be understood as a Bochner integral and it defines a bounded operator $`Op(a)`$ on the Hilbert space $`L^2()`$.
Let us restate some well-known facts on the Weyl calculus (see \[Fo\], Chapter 2). The operator $`Op(a)`$ acts by the formula
$$(Op(a)f)(x)=a(\frac{1}{2}(x+y),t)e^{2\pi i(xy)t}f(y)𝑑y𝑑t.$$
(48)
For the operator product $`Op(a)Op(b)`$ and the adjoint operator $`Op(a)^{}`$ we have
$$Op(a)Op(b)=Op(a\mathrm{\#}b)\mathrm{and}Op(a)^{}=Op(a^{}),$$
(49)
where the symbols $`a\mathrm{\#}b`$ and $`a^{}`$ are defined by
$`(a\mathrm{\#}b)(x_1,x_2):=`$
$`4{\displaystyle a(u_1,u_2)b(v_1,v_2)e^{4\pi i[(x_1u_1)(x_2v_2)(x_1v_1)(x_2u_2)]}𝑑u_1𝑑u_2𝑑v_1𝑑v_2},`$ (50)
$`a^{}(x_1,x_2):=\overline{a(x_1,x_2)},x_1,x_2.`$ (51)
Lemma 11. Let $`\mu =(\mu _1,\mu _2)`$, $`\nu =(\nu _1,\nu _2)`$, $`\mu ^{}=(\mu _1^{},\mu _2^{})`$, $`\nu ^{}=(\nu _1^{},\nu _2^{})^{++}`$, $`a𝒟_{\mu ,\nu },b𝒟_{\nu ^{},\mu ^{}}`$. For $`t`$, we have
$`e^{2\pi t𝒬_1}(a\mathrm{\#}b)=(e^{2\pi t𝒬_1}a)\mathrm{\#}(e^{\pi t𝒫_2}b)\mathrm{if}|\mathrm{Re}t|<\mu _1,|\mathrm{Re}t|<2\nu _2^{},`$ (52)
$`e^{2\pi t𝒬_1}(a\mathrm{\#}b)=(e^{\pi t𝒫_2}a)\mathrm{\#}(e^{2\pi t𝒬_1}b)\mathrm{if}|\mathrm{Re}t|<2\nu _2,|\mathrm{Re}t|<\mu _1^{},`$ (53)
$`e^{2\pi t𝒬_2}(a\mathrm{\#}b)=(e^{2\pi t𝒬_2}a)\mathrm{\#}(e^{\pi t𝒫_1}b)\mathrm{if}|\mathrm{Re}t|<\mu _1,|\mathrm{Re}t|<2\nu _1^{},`$ (54)
$`e^{2\pi t𝒬_2}(a\mathrm{\#}b)=(e^{\pi t𝒫_1}a)\mathrm{\#}(e^{2\pi t𝒬_2}b)\mathrm{if}|\mathrm{Re}t|<2\nu _1,|\mathrm{Re}t|<\mu _2^{},`$ (55)
$`e^{2\pi t𝒫_1}(a\mathrm{\#}b)=(e^{2\pi 𝒫_1}a)\mathrm{\#}(e^{2\pi t𝒫_1}b)\mathrm{if}|\mathrm{Re}t|<\nu _1,|\mathrm{Re}t|<\nu _1^{}.`$ (56)
$`e^{2\pi t𝒫_1}(a\mathrm{\#}b)=(e^{4\pi t𝒬_2}a)\mathrm{\#}(e^{4\pi t𝒬_2}b)\mathrm{if}|\mathrm{Re}t|<\nu _2/2,|\mathrm{Re}t|<\nu _2^{}/2.`$ (57)
$`e^{2\pi t𝒫_2}(a\mathrm{\#}b)=(e^{2\pi t𝒫_2}a)\mathrm{\#}(e^{2\pi t𝒫_2}b)\mathrm{if}|\mathrm{Re}t|<\nu _2,|\mathrm{Re}t|<\nu _2^{}.`$ (58)
$`e^{2\pi t𝒫_2}(a\mathrm{\#}b)=(e^{4\pi t𝒬_1}a)\mathrm{\#}(e^{4\pi t𝒬_1}b)\mathrm{if}|\mathrm{Re}t|<\nu _1/2,|\mathrm{Re}t|<\nu _1^{}/2.`$ (59)
Proof. As samples, we carry out the proofs of formulas (52) and (57). The other equations are proved by a similar reasoning.
First we prove formula (57) for real $`t`$. It is well-known (see \[Fo\], p. 104) that the Fourier transform of the product $`a\mathrm{\#}b`$ is the twisted convolution of the Fourier transform $`(a)`$ and $`(b)`$, that is $`(a\mathrm{\#}b)=(a)_t(b)`$, where
$$(c_td)(x_1,x_2)=c(u_1,u_2)d(x_1u_1,x_2u_2)e^{\pi i(x_1u_2x_2u_1)}𝑑u_1𝑑u_2.$$
Using the preceding fact and formula (15) we compute
$`(e^{2\pi t𝒫_1}(a\mathrm{\#}b))(x_1,x_2)=\left(e^{2\pi t𝒬_1}(a\mathrm{\#}b)\right)(x_1,x_2)`$
$`=e^{2\pi tx_1}((a)_t(b))(x_1,x_2)`$
$`={\displaystyle (a)(u_1,u_2)(b)(x_1u_1,x_2u_2)e^{\pi i(x_1(u_22ti)x_2u_1)}𝑑u_1𝑑u_2}`$
$`={\displaystyle (a)(u_1,u_2+2ti)(b)(x_1u_1,x_2u_22ti)e^{\pi i(x_1u_2x_2u_1)}𝑑u_1𝑑u_2}`$
$`={\displaystyle \left(e^{4\pi t𝒫_2}(a)\right)(u_1,u_2)(e^{4\pi t𝒫_2}(b))(x_1u_1,x_2u_2)e^{\pi i(x_1u_2x_2u_1)}𝑑u_1𝑑u_2}`$
$`=\left(e^{4\pi t𝒫_2}(a)_te^{4\pi t𝒫_2}(b)\right)(x_1,x_2)`$
$`=(\left(e^{4\pi t𝒬_2}a\right)_t(e^{4\pi t𝒬_2}b))(x_1,x_2)`$
$`=\left(e^{4\pi tQ_2}a\mathrm{\#}e^{4\pi t𝒬_2}b\right)(x_1,x_2)`$
which in turn implies (57). It remains to justify the fourth equality sign which follows by the formal substitution $`u_2u_2+2ti`$. First we note that the assumptions $`a𝒟_{\mu ,\nu }`$ and $`b𝒟_{\mu ^{},\nu ^{}}`$ imply that $`(a)𝒟_{\nu ,\mu }`$ and $`(b)𝒟_{\nu ^{},\mu ^{}}`$, so that $`(a)𝒟(e^{\pm 2\pi \nu _2𝒫_2})`$ and $`(b)𝒟(e^{\pm 2\pi \nu _2^{}𝒫_2})`$. Therefore, since $`2|t|<\nu _2`$ and $`2|t|<\nu _2^{}`$, the function
$$(a)(u_1,u_2)(b)(x_1u_1,x_2u_2)e^{\pi i(x_1(u_22ti)x_2u_1)}$$
of $`u_2`$ is holomorph on a strip $`\epsilon <\mathrm{Im}u_2<2|t|+\epsilon `$ of the complex $`u_2`$-plane for some small $`\epsilon >0`$. Hence the integral of this function along the boundary of the rectangle with corners $`R,R,R+2ti,R+2ti`$ vanishes. In order to justify the substitution $`u_2u_2+2ti`$, it is sufficient to show that the corresponding integrals from $`\pm R`$ to $`\pm R+2ti`$ tend to zero as $`R+\mathrm{}`$. Using formula (S0.Ex40) we estimate
$`\left|{\displaystyle \underset{0}{\overset{2t}{}}}{\displaystyle \underset{\mathrm{}}{\overset{\mathrm{}}{}}}(a)(u_1,\pm R+si)(b)(x_1u_1,x_2(\pm R+si))e^{\pi i(x_1(\pm R+si2ti)x_2u_1)}𝑑s𝑑u_1\right|`$
$`C\left|{\displaystyle \underset{0}{\overset{2t}{}}}{\displaystyle \underset{\mathrm{}}{\overset{\mathrm{}}{}}}e^{2\pi (\nu _1|u_1|+\nu _2R)2\pi (\nu _1^{}|x_1u_1|+\nu _2^{}|x_2R|)+\pi x_1(2ts)}𝑑s𝑑u_1\right|`$
$`C_{x_1,x_2}e^{2\pi \nu _2R},`$
where $`C`$ and $`C_{x_1,x_2}`$ are not depending on $`R`$. Since $`\nu _2>0`$, the integral goes to zero if $`R+\mathrm{}`$. This proves formula (57) for real $`t`$.
Next we prove (52) for real $`t`$. From the definition (50) of the product $`\mathrm{\#}`$ we obtain
$`(e^{2\pi tQ_1}(a\mathrm{\#}b))(x_1,x_2)=`$
$`{\displaystyle e^{2\pi tu_1}a(u_1,u_2)b(v_1,v_2)e^{4\pi i[(x_1u_1)(x_2v_2ti/2)(x_1v_1)(x_2u_2)]}𝑑u_1𝑑u_2𝑑v_1𝑑v_2}.`$
Recall that $`a,b𝒟_{\mu ,\nu }`$ by assumption. Hence we have $`a𝒟(e^{\pm 2\pi \mu _1𝒬_1})`$ and $`b𝒟(e^{\pm 2\pi \nu _2^{}𝒫_2})`$. Since $`|t|<\mu _1`$ and $`|t|<2\nu _2^{}`$, the latter implies that $`a𝒟(e^{2\pi t𝒬_1})`$ and $`b𝒟(e^{\pi t𝒫_2})`$. In fact, we even have that $`e^{2\pi t𝒬_1}a,e^{\pi t𝒫_2}a𝒟_{\stackrel{~}{\mu },\stackrel{~}{\nu }}`$ for certain $`\stackrel{~}{\mu },\stackrel{~}{\nu }^{++}`$. Equation (52) follows from the preceding formula by the formal substitution $`v_2v_2+it/2`$. In order to show that this formal replacement is justified we integrate in the complex $`v_2`$-plane along the boundary of the rectangle with corners $`R,R,R+it/2,R+it/2`$, where $`R>0`$. To complete the proof, it suffices to show that the integrals from $`\pm R`$ to $`\pm R+it/2`$ tend to zero as $`R+\mathrm{}`$. Indeed, using formula (S0.Ex40) and the assumptions $`|t|<\mu _1`$ and $`|t|<2\nu _2^{}`$, we estimate
$`|{\displaystyle \underset{0}{\overset{t/2}{}}}{\displaystyle }dsdu_1du_2dv_1e^{2\pi tu_1}a(u_1,u_2)b(v_1,\pm R+si)`$
$`e^{4\pi i[(x_1u_1)(x_2(\pm R+si)ti/2)(x_1v_1)(x_2u_2)]}|`$
$`C\left|{\displaystyle \underset{0}{\overset{t/2}{}}}{\displaystyle }e^{2\pi tu_1}e^{2\pi (\mu _1|u_1|+\mu _2|u_2|+\mu _1^{}|v_1|+\mu _2^{}R)}e^{4\pi (x_1u_1)(t/2s)}dsdu_1du_2dv_1\right|`$
$`C_{x_1}^{}e^{2\pi \mu _2^{}R}\left|{\displaystyle \underset{\mathrm{}}{\overset{\mathrm{}}{}}}{\displaystyle \underset{0}{\overset{t/2}{}}}e^{2\pi (2su_1\mu _1|u_1|}du_1ds\right|`$
$`C_{x_1}^{\prime \prime }e^{2\pi \mu _2^{}R}{\displaystyle \underset{\mathrm{}}{\overset{\mathrm{}}{}}}e^{2\pi |t||u_1|2\pi \mu _1|u_1|}𝑑u_1C_{x_1}^{\prime \prime \prime }e^{2\pi \mu _2^{}R},`$
where $`C,C_{x_1}^{},C_{x_1}^{\prime \prime },C_{x_1}^{\prime \prime \prime }`$ are numbers not depending on $`R`$. Since $`\mu _2^{}>0`$, the integral goes to zero if $`R+\mathrm{}`$. This completes the proof of (52) for real $`t`$.
For imaginary $`t`$ the above reasoning works as well. In this case we are lead to real translations of $`u_2`$ and $`v_2`$, respectively, which are possible by the translation invariance of the Lebesgue measure. The case of general $`t`$ follows by combining the real and the imaginary cases. $`\mathrm{}`$
For $`\mu ,\nu ^{++}`$, let $`𝒟^{\mu ,\nu }`$ denote the intersection of domains $`𝒟_{\mu ^{},\nu ^{}}`$ (see (42)), where $`\mu ^{},\nu ^{}^{++}`$, $`\mu _j^{}<\mu _j`$, $`\nu _j^{}<\nu _j`$ for $`j=\mathrm{1,2}`$.
Corollary 12. Let $`\mu ,\nu ^{++}`$. If $`\mu _1<2\nu _2`$ and $`\mu _2<2\nu _1`$, then $`𝒟^{\mu ,\nu }`$ is a $``$-algebra with product $`\mathrm{\#}`$ and involution $``$ defined (50) and (51), respectively. In particular, $`𝔄(^2)`$ is a $``$-algebra.
Proof. Since $`\mu _1<2\nu _2`$ and $`\mu _2<2\nu _1`$, we conclude from formulas (52), (54), (56), (58) and (42) that $`a,b𝒟^{\mu ,\nu }`$ implies that $`a\mathrm{\#}b𝒟^{\mu ,\nu }`$. By (42) it is obvious that $`a^{}𝒟^{\mu ,\nu }`$ for $`a𝒟^{\mu ,\nu }`$. Thus $`𝒟^{\mu ,\nu }`$ is a $``$-algebra. Since $`𝔄(^2)`$ is the intersection of all domains $`𝒟^{\mu ,\nu }`$, $`𝔄(^2)`$ is a $``$-algebra as well. $`\mathrm{}`$
Lemma 13. Suppose that $`\mu ,\nu ^{++}`$. Let $``$ denote the norm of $`L^2(^2)`$. If $`a,b𝒟_{\mu ,\nu }=𝒟(S(e^{2\pi \mu 𝒬})S(e^{2\pi \nu 𝒫}))`$, then $`a,b,\overline{b}\mathrm{\#}a𝒮()`$ and we have
$`{\displaystyle a(x_1,x_2)\overline{b}(x_1,x_2)𝑑x_1𝑑x_2}`$ $`={\displaystyle (\overline{b}\mathrm{\#}a)(x_1,x_2)𝑑x_1𝑑x_2},`$ (60)
$`a\mathrm{\#}b`$ $`ab.`$ (61)
Proof. By Lemma 9(ii) and Corollary 12, we have $`a,b𝒮(^2)`$ and so $`\overline{b}\mathrm{\#}a𝒮(^2)`$. From Proposition 2 in \[St\], p. 555, it follows that $`Op(a)`$ and $`Op(b)`$ are Hilbert-Schmidt operators and
$$Op(a),Op(b)_{HS}=\mathrm{Tr}Op(b)^{}Op(a)=(a,b)\mathrm{and}Op(a)_{HS}=a,$$
(62)
where $`\mathrm{Tr}`$ is the trace and $`,_{HS}`$ and $`_{HS}`$ denote scalar product and Hilbert-Schmidt norm of Hilbert-Schmidt operators, respectively. Using formulas (62) and (49) and the submultiplicativity of the Hilbert-Schmidt norm we obtain
$`a\mathrm{\#}b`$ $`=Op(a\mathrm{\#}b)_{HS}=Op(a)Op(b)_{HS}`$
$`Op(a)_{HS}Op(b)_{HS}=ab.`$
This proves (61).
Put $`c:=b^{}\mathrm{\#}a`$. By (48), the operator $`Op(c)`$ is an integral operator with kernel
$$K_c(x_1,x_2)=c(\frac{1}{2}(x_1+x_2),t)e^{2\pi i(x_1x_2)t}𝑑t.$$
(63)
Since $`c𝒮(^2)`$, the function $`d`$ defined by $`d(y_1,y_2)=c(y_1,t)e^{2\pi iy_2t}𝑑t`$ is in $`𝒮(^2)`$ and so is the function $`K_c(x_1,x_2)=d(\frac{1}{2}(x_1+x_2),x_1x_2)`$. It is well-known that any integral operator with kernel in the Schwartz space $`𝒮(^2)`$ is a trace class operator on $`L^2()`$ and that its trace is given by the integral over the diagonal. Using this fact and formulas (49) and (63) we get
$`\mathrm{Tr}Op(b)^{}Op(a)=\mathrm{Tr}Op(b^{}\mathrm{\#}a)=\mathrm{Tr}Op(c)`$
$`=`$ $`{\displaystyle K_c(x_1,x_1)𝑑x_1}={\displaystyle c(x_1,x_2)𝑑x_1𝑑x_2}`$
$`=`$ $`{\displaystyle a(x_1,x_2)\overline{b}(x_1,x_2)𝑑x_1𝑑x_2}.`$
Comparing the latter with (62), formula (60) follows. $`\mathrm{}`$
Our next proposition says that $`𝔄(^2)[\tau ]`$ is a Frechet $``$-algebra with approximate identity.
Proposition 14. (i) Provided with the product $`\mathrm{\#}`$, the involution $``$ and the locally convex topology $`\tau ,𝔄(^2)`$ is a Frechet topological $``$-algebra.
(ii) Set $`f_\epsilon (x_1,x_2):=e^{\pi \epsilon (x_1^2+x_2^2)}.`$ For each $`a𝔄(^2)`$, we have
$$\underset{\epsilon +0}{lim}f_\epsilon \mathrm{\#}a=\underset{\epsilon +0}{lim}a\mathrm{\#}f_\epsilon =a$$
(64)
in the locally convex space $`𝔄(^2)[\tau ]`$.
Proof. (i): Recall that by definition the topology $`\tau `$ is generated by the family of norms $`e^{2\pi c𝒬}e^{2\pi d𝒫}`$, $`c,d^2`$, where $``$ is the norm of $`L^2(^2)`$. Fix $`c,d^2`$ and put
$$\mu =(\mu _1,\mu _2),\mu _1:=d_1+c_1/2,\mu _2:=d_2c_2/2.$$
(65)
By (53), (55), (56), (58) and (61), we obtain
$$e^{2\pi c𝒬}e^{2\pi d𝒫}(a\mathrm{\#}b)=(e^{2\pi \mu 𝒫}a)\mathrm{\#}(e^{2\pi c𝒬}e^{2\pi d𝒫}b)e^{2\pi \mu 𝒫}ae^{2\pi c𝒬}e^{2\pi d𝒫}b$$
(66)
for $`a,b𝔄(^2)`$. Since $`e^{2\pi c𝒬}e^{2\pi d𝒫}a^{}=e^{2\pi c𝒬}e^{2\pi d𝒫}a`$, product and involution are $`\tau `$-continuous, so $`𝔄(^2)[\tau ]`$ is indeed a topological $``$-algebra. Since $`𝔄(^2)[\tau ]`$ is a Frechet space as noted above, it is a Frechet topological $``$-algebra.
(ii): Let $`b𝔄(^2)`$ and $`\mu =(\mu _1,\mu _2)^2`$. Our aim is to prove by explicit estimations that
$$\underset{\epsilon +0}{lim}(e^{2\pi \mu 𝒫}f_\epsilon )\mathrm{\#}b=b$$
(67)
in $`L^2(^2)`$. Note first that from the well-known equation
$$e^{\pi (sc)^2\epsilon 2\pi its}𝑑s=\epsilon ^{1/2}e^{\pi \epsilon ^1t^2}e^{2\pi itc},c,$$
(68)
we obtain that
$$(e^{2\pi \mu 𝒫}f_\epsilon )(x_1,x_2)=\epsilon ^1e^{\pi \epsilon ^1(x_1^2+x_2^2)}e^{2\pi (\mu _1x_1+\mu _2x_2)}.$$
Using the latter formulas and the definition (50) of the product $`\mathrm{\#}`$ we compute
$`|b(x_1,x_2)((e^{2\pi \mu 𝒫}f_\epsilon )\mathrm{\#}b)(x_1,x_2)|`$
$`=\left|b(x_1,x_2)4{\displaystyle (e^{2\pi \mu 𝒫}f_\epsilon )(2x_22v_2\mathrm{,2}v_12x_1)b(v_1,v_2)e^{4\pi i(v_1x_2x_1v_2)}𝑑v_1𝑑v_2}\right|`$
$`=|\frac{4}{\epsilon }{\displaystyle }(b(x_1,x_2)b(v_1,v_2))dv_1dv_2`$ (69)
$`e^{4\pi \epsilon ^1((x_2v_2)^2+(v_1x_1)^2)+4\pi ((x_2v_2)\mu _1+(v_1x_1)\mu _2)+4\pi i(v_1x_2x_1v_2)}|`$
$`\frac{4}{\epsilon }{\displaystyle e^{4\pi \epsilon ^1((x_1v_1)^2+(x_2v_2)^2)+4\pi ((x_2v_2)\mu _1(x_1v_1)\mu _2)}|b(x_1,x_2)b(v_1,v_2)|𝑑v_1𝑑v_2}.`$ (70)
Fix a number $`\delta >0`$. Since $`b𝔄(^2)`$,(S0.Ex40) holds for arbitrary $`\nu ,\mu ^{++}`$. Hence there exists $`M`$ such that for $`(x_1,x_2)^2`$,
$$e^{(4\pi |\mu _1|+1)|x_1|+(4\pi |\mu _2|+1)|x_2|}|b(x_1,x_2)|M.$$
(71)
Further, since $`e^{|x_1|+|x_2|}b(x_1,x_2)0`$ as $`|x_1|+|x_2|+\mathrm{}`$ by (71), the function $`c(x_1,x_2):=e^{|x_1|+|x_2|}b(x_1,x_2)`$ is uniformly continuous on $`^2`$. Thus there exists $`\delta _1`$ such that $`1>\delta _1>0`$ and
$$(1+eM)e^{4\pi (|\mu _1|+|\mu _2|)}|c(x_1,x_2)c(v_1,v_2)|<\delta $$
for $`(x_1v_1)^2+(x_2v_2)^2<\delta _1^2`$. From this and (71) we easily derive that
$$e^{|x_1|+|x_2|}e^{4\pi ((x_2v_2)\mu _2(x_1v_1)\mu _1)}|b(x_1,x_2)b(v_1,v_2)|<\delta $$
(72)
when $`(x_1v_1)^2+(x_2v_2)^2<\delta _1^2`$.
Next we turn to the domain where $`(x_1v_1)^2+(x_2v_2)^2\delta _1^2`$. Obviously, there exists $`K`$ such that
$$\pi \epsilon ^1(t_1^2+t_2^2)+4\pi (|\mu _1|+1)|t_1|+(|\mu _2|+1)|t_2|K$$
(73)
for all $`(t_1,t_2)^2`$ and $`1>\epsilon >0`$. If $`(x_1v_1)^2+(x_2v_2)^2\delta _1^2`$, then by (71) and (73) we obtain that
$`e^{|x_1|+|x_2|}e^{2\pi \epsilon ^1((x_1v_1)^2+(x_2v_2)^2)}e^{4\pi ((x_2v_2)\mu _1(x_1v_1)\mu _2)}|b(x_1,x_2)b(v_1,v_2)|`$
$`e^{\pi \epsilon ^1((x_1v_1)^2+(x_2v_2)^2)}e^K2Me^{\pi \epsilon ^1\delta _1^2}e^K2M<\delta `$ (74)
for sufficiently small $`\epsilon >0`$.
Using the relation
$$\frac{4}{\epsilon }e^{2\pi \epsilon ^1((x_1v_1)^2+(x_2v_2)^2)}𝑑v_1𝑑v_2=2$$
by (68), it follows from estimates (S0.Ex76), (72) and (S0.Ex80) that
$$|b(x_1,x_2)((e^{2\pi \mu 𝒫}f_\epsilon )\mathrm{\#}b)(x_1,x_2)|3\delta e^{|x_1||x_2|}$$
and so $`b(e^{2\pi \mu 𝒫}f_\epsilon )\mathrm{\#}b12\delta `$ for small $`\epsilon >0`$. This proves (67).
Now let $`a𝔄(^2)`$ and $`c,d^2`$. Let $`\mu `$ be as in (65). Applying (67) with $`b=e^{2\pi c𝒬}e^{2\pi d𝒫}a`$, we get
$$e^{2\pi c𝒬}e^{2\pi d𝒫}(f_\epsilon \mathrm{\#}aa)=(e^{2\pi \mu 𝒫}f_\epsilon )\mathrm{\#}(e^{2\pi c𝒬}e^{2\pi d𝒫}a)e^{2\pi c𝒬}e^{2\pi d𝒫}a0$$
as $`\epsilon +0`$. This proves that $`\underset{\epsilon +0}{lim}f_\epsilon \mathrm{\#}a=a`$ in $`𝔄(^2)[\tau ]`$. Applying the involution we obtain the second equality in (64). $`\mathrm{}`$
Remark 2. Upon scaling and multiplying by parameters, the operators $`Op(f_\epsilon )`$, $`\epsilon >0`$, form the so-called Hermite semigroup $`e^{2\pi t(𝒫^2+𝒬^2)},t>0`$, acting on the Hilbert space $`L^2(^2)`$, see \[Fo\], pp. 236–238.
In this paper we shall mainly use the symbol algebra $`𝔄(^2)`$. However, for most considerations it suffices to work with the smaller symbol algebras
$`𝔄_{ex}(^2):=\mathrm{Lin}\{e^{\epsilon _1x_1^2\epsilon _2x_2^2+c_1x_1+c_2x_2};\epsilon _1>0,\epsilon _2>0,c_1,c_2\},`$
$`𝔄_{pex}(^2):=\mathrm{Lin}\{x_1^{n_1}x_2^{n_2}e^{\epsilon _1x_1^2\epsilon _2x_2^2+c_1x_1+c_2x_2};\epsilon _j>0,c_j,n_j_0\}.`$
Both $`𝔄_{ex}(^2)`$ and $`𝔄_{pex}(^2)`$ are $``$-algebras with multiplication (50) and involution (51). In order to prove this assertion it is sufficient to show that $`a\mathrm{\#}b`$ is in $`𝔄_{ex}(^2)`$ resp. $`𝔄_{pex}(^2)`$ when $`a`$ and $`b`$ are so. In the case of $`𝔄_{ex}(^2)`$ this can be verified by direct computation of the twisted product $`a\mathrm{\#}b`$ using formula (68). From formula (3) in \[GV\] it follows at once that $`a\mathrm{\#}b𝔄_{pex}(^2)`$ for $`a,b𝔄_{pex}(^2)`$.
2. The Coordinate Algebra $`𝒪(_𝐪^\mathrm{𝟐})`$ of the Quantum Plane
2.1 $`𝒪(_q^2)`$ as a left module algebra of $`𝒰_q(gl_2)`$
Let $``$ be the left action of $`𝒰_g(gl_2)`$ on $`𝒪(_q^2)`$ associated with the right coaction of $`𝒪(GL_q(2))`$ defined by (10). From (10) and (9) we then obtain
$`K_1x=q^{1/2}x,K_1y=y,K_2x=x,K_2y=q^{1/2}y,`$ (75)
$`Ex=y,Ey=0,Fx=0,Fy=x.`$ (76)
Moreover, since $`\epsilon (K_1)=\epsilon (K_2)=1`$ and $`\epsilon (E)=\epsilon (F)=0`$, we also have
$$K_11=1,K_21=1,E1=0,F1=0.$$
(77)
The following proposition derives the action of the generators $`K_1,K_2,E,F`$ of $`𝒰_q(gl_2)`$ on general elements of the algebra $`𝒪(_q^2)`$. We set
$$D_{q^2}(f)(x):=\frac{f(x)f(q^2x)}{(1q^2)x}.$$
Proposition 15. If $`g`$ and $`h`$ are complex polynomials in a single variable, then we have
$`K_1(g(x)h(y))`$ $`=g(q^{1/2}x)h(y),K_2(g(x)h(y))=g(x)h(q^{1/2}y),`$ (78)
$`E(g(x)h(y))`$ $`=q^{1/2}D_{q^2}(g(q^{1/2}))(x)yh(q^{1/2}y),`$ (79)
$`F(g(x)h(y))`$ $`=q^{1/2}g(q^{1/2}x)xD_{q^2}(h(q^{1/2}))(y).`$ (80)
Proof. Since $`𝒪(_q^2)`$ is a $`𝒰_q(gl_2)`$-comodule algebra, equation (1) holds. The assertion follows from this equation combined with formulas (75) and (77). For the generators $`K_1`$ and $`K_2`$ this is obvious. We carry out the proof of formula (80). The proof of formula (79) is similar.
Since $`\mathrm{\Delta }(F)=FK+K^1F`$, it follows from (1) that
$$F(zz^{})=(Fz)(Kz^{})+(K^1z)(Fz^{}).$$
(81)
Recall that $`Fx=0`$ and $`F1=0`$ by (75) and (77). Using these facts it follows from (81) by induction on $`n`$ that $`Fx^n=0`$ for $`n_0`$. Thus we have $`Fg(x)=0`$.
Since $`K^1g(x)=K_1^1K_2g(x)=K_1^1g(x)=g(q^{1/2}x)`$, we derive from (81), applied to $`z=g(x)`$ and $`z^{}=h(y)`$, that
$$F(g(x)h(y))=g(q^{1/2}x)(Fh(y)).$$
(82)
Therefore, in order to prove (80) it suffices to show that
$$Fy^n=q^{1/2}q^{n/2}xD_{q^2}(y^n),n.$$
(83)
We prove (83) by induction on $`n`$. If $`n=1`$, then (83) is true by (75). If (83) is valid for a, then it follows from (81) and (75) that
$`Fy^{n+1}`$ $`=(Fy)(Ky^n)+(K^1y)(Fy^n)`$
$`=x(q^{n/2}y^n)+(q^{1/2}y)(q^{1/2}(1q^2)^1(1q^{2n})q^{n/2}xy^{n1})`$
$`=q^{n/2}(1+q^2(1q^2)^1(1q^{2n}))xy^n`$
$`=q^{1/2}(1q^2)^1(1q^{2(n+1)})q^{(n+1)/2}xy^n`$
$`=q^{1/2}q^{(n+1)/2}xD_{q^2}(y^{n+1}),`$
which proves (83) in the case of $`n+1`$. $`\mathrm{}`$
For $`zO(_q^2)`$, we define
$$𝒟_x^q(z)=Ky^1E^{}z,𝒟_y^q(z)=Kx^1F^{}z,$$
(84)
where the elements $`y^1`$ and $`x^1`$ of $`\widehat{𝒪}(_q^2)`$ act by left multiplication on $`𝒪(_q^2)`$. From (79) and (80) we obtain
$`𝒟_x^q(g(x)h(y))=q^{1/2}\lambda D_{q^2}(g)(qx)h(qy),`$ (85)
$`𝒟_y^q(g(x)h(y))=q^{1/2}\lambda g(x)D_{q^2}(h)(qy).`$ (86)
for polynomials $`g`$ and $`h`$.
For $`r,s_0`$, let $`\sigma _{rs}`$ denote the automorphism of the algebra $`𝒪(_q^2)`$ defined by $`\sigma _{rs}(z)=K_1^rK_2^sz,z𝒪(_q^2)`$. From formulas (78)-(80) or (85) we easily derive that $`𝒟_x^q`$ is a $`(\sigma _{\mathrm{2,0}},\sigma _{2,2})`$-derivation and $`𝒟_y^q`$ is a $`(\sigma _{0,2},\sigma _{\mathrm{2,2}})`$-derivation of the algebra $`𝒪(_q^2)`$, that is, for $`z_1,z_2𝒪(_q^2)`$ we have
$`𝒟_x^q(z_1z_2)=(K_1^2z_1)𝒟_x^q(z_2)+𝒟_x^q(z_1)(K_1^2K_2^2z_2),`$
$`𝒟_y^q(z_1z_2)=(K_2^2z_1)𝒟_y^q(z_2)+𝒟_y^q(z_1)(K_1^2K_2^2z_2).`$
In the limit $`q1`$ the preceding equations go into the Leibniz rule. We shall consider the linear mappings $`𝒟_x^q`$ and $`𝒟_y^q`$ as $`q`$-deformed partial derivatives of the algebra $`𝒪(_q^2)`$.
2.2 Covariant Differential Calculus on $`𝒪(_q^2)`$
As shown in \[PW\] and \[WZ\], there are two distinguished first order differential calculi $`\mathrm{\Gamma }_+`$ and $`\mathrm{\Gamma }_{}`$ on $`𝒪(_q^2)`$. For both calculi, the set of differentials $`\{dx,dy\}`$ is a basis for the right (and for the left) $`𝒪(_q^2)`$-module of first order forms. Therefore, for any $`z𝒪(_q^2)`$ there exist uniquely determined elements $`_x(z)`$ and $`_y(z)`$ of $`𝒪(_q^2)`$, called partial derivatives of $`z`$, such that
$$dz=dx_x(z)+dy_y(z).$$
(87)
The bimodule structures of the calculi $`\mathrm{\Gamma }_+`$ and $`\mathrm{\Gamma }_{}`$ are described by the following commutation relations:
$`\mathrm{\Gamma }_+:`$ $`x\mathrm{d}y=q\mathrm{d}yx+(q^21)\mathrm{d}xy,y\mathrm{d}x=q\mathrm{d}xy,`$ (88)
$`x\mathrm{d}x=q^2\mathrm{d}xx,y\mathrm{d}y=q^2\mathrm{d}yy.`$ (89)
$`\mathrm{\Gamma }_{}:`$ $`y\mathrm{d}x=q^1\mathrm{d}xy+(q^21)\mathrm{d}yx,x\mathrm{d}y=q^1\mathrm{d}yx,`$ (90)
$`x\mathrm{d}x=q^2\mathrm{d}xx,y\mathrm{d}y=q^2\mathrm{d}yy.`$ (91)
From these relations we see that $`\eta _+:=y^2xdx`$ and $`\eta _{}:=x^2ydy`$ are non-zero central elements of the bimodules $`\mathrm{\Gamma }_+`$ and $`\mathrm{\Gamma }_{}`$, respectively. Recall that an element $`\eta `$ of a bimodule over an algebra $`𝒵`$ is called central if $`\eta z=z\eta `$ for all $`z𝒵`$.
Note that the relations for $`\mathrm{\Gamma }_+`$ go into the relations of $`\mathrm{\Gamma }_{}`$ if we interchange the coordinates $`x`$ and $`y`$ and the numbers $`q`$ and $`q^1`$. The partial derivatives $`_x`$ and $`_y`$, considered as linear mappings of $`𝒪(_q^2)`$, and the coordinate functions $`x`$ and $`y`$, acting on $`𝒪(_q^2)`$ by left multiplication, satisfy the relations:
$$\begin{array}{cc}\hfill \mathrm{\Gamma }_+:& _xy=qy_x,_yx=qx_y,\hfill \\ & _xxq^2x_x=1+(q^21)y_y,_yyq^2y_y=1.\hfill \\ \hfill \mathrm{\Gamma }_{}:& _xy=q^1y_x,_yx=q^1x_y,\hfill \\ & _xxq^2x_x=1,_yyq^2y_y=1+(q^21)x_x.\hfill \end{array}$$
From these formulas one derives by induction the expressions for the actions of $`_x`$ and $`_y`$ on general elements of $`𝒪(_q^2)`$. If $`g`$ and $`h`$ are complex polynomials in a single variable, then we have:
$`\mathrm{\Gamma }_+:_x(g(y)h(x))=g(qy)D_{q^2}(h)(x),_y(g(y)h(x))=D_{q^2}(g)(y)h(x),`$ (92)
$`\mathrm{\Gamma }_{}:_x(g(x)h(y))=D_{q^2}(g)(x)h(y),_y(g(x)h(y))=g(q^1x)D_{q^2}(h)(y).`$ (93)
All these facts and formulas are well-known. We now give another description of these calculi. Let $`\mathrm{\Omega }`$ be the free bimodule of the localization algebra $`\stackrel{ˇ}{𝒪}(_q^2)`$ generated by a central vector space $`V`$. That is, $`\mathrm{\Omega }`$ is the vector space $`\stackrel{ˇ}{𝒪}(_q^2)V`$ with bimodule structure given by
$$u\left(_jz_je_j\right)v:=_juz_jve_j,$$
(94)
where $`u,v,z_j\stackrel{ˇ}{𝒪}(_q^2),e_jV`$. For notational simplicity, we write $`ze`$ instead of $`ze`$, where $`z\stackrel{ˇ}{𝒪}(_q^2)`$ and $`eE`$. Fix two elements $`e_1,e_2V`$ and put
$$\begin{array}{cc}& \omega _+:=q^2y^2x^2e_1+x^2e_2,\omega _{}=q^2x^2y^2e_1+y^2e_2,\hfill \\ & d_+z:=\omega _+zz\omega _+,d_{}z:=\omega _{}zz\omega _{},z𝒪(_q^2).\hfill \end{array}$$
Let us abbreviate $`𝒵:=𝒪(_q^2)`$. Obviously, $`\stackrel{~}{\mathrm{\Gamma }}_\epsilon :=𝒵d_\epsilon 𝒵𝒵`$ is a $`𝒵`$-bimodule and the mapping $`d_\epsilon :𝒵\stackrel{~}{\mathrm{\Gamma }}_\epsilon `$ satisfies the Leibniz rule for $`\epsilon =+,`$. Thus, the pair $`(\stackrel{~}{\mathrm{\Gamma }}_\epsilon ,d_\epsilon )`$ is a first order differential calculus over the algebra $`𝒵=𝒪(_q^2)`$. For the differentials of the coordinate functions we obtain
$`d_+x=(q^21)y^2x^1e_1,d_+y=(q^21)q^2y^3x^2e_1+(q^21)yx^2e_2,`$ (95)
$`d_{}x=(q^21)q^2x^3y^2e_1+(q^21)xy^2e_2,d_{}y=(q^21)x^2y^1e_1.`$ (96)
Lemma 16. Suppose that the elements $`e_1`$ and $`e_2`$ are linearly independent. Then the first order differential calculi $`\mathrm{\Gamma }_\epsilon `$ and $`\stackrel{~}{\mathrm{\Gamma }}_\epsilon ,\epsilon =+,`$, are isomorphic.
Proof. Since $`\{dx,dy\}`$ is a free left $`𝒪(_q^2)`$-module basis of $`\mathrm{\Gamma }_\epsilon `$, there is a well-defined left $`𝒪(_q^2)`$-module homomorphism $`\psi _\epsilon :\mathrm{\Gamma }_\epsilon \stackrel{~}{\mathrm{\Gamma }}_\epsilon `$ such that
$$\psi _\epsilon (udx+vdy)=ud_\epsilon x+vd_\epsilon y,u,v𝒪(_q^2).$$
In order to prove that $`\psi _\epsilon `$ is an $`𝒪(_q^2)`$-bimodule homomorphism, it suffices to show that the relations (88) and (89) resp. (90) and (91) hold also in $`\mathrm{\Omega }_+`$ and $`\mathrm{\Omega }_{}`$. As a sample, we verify the first relation of (90). The other relations follow by similar straightforward computations. Using formulas (96) and the commutation rules in the algebra $`\stackrel{ˇ}{𝒪}(_q^2)`$, we obtain
$`q^1d_{}xy+(q^21)d_{}yx=`$
$`(q^21)(qx^3y^2e_1y+q^1xy^2e_2y+(q^21)x^2y^1e_1x)=`$
$`(q^21)((qx^3y^1+(q^21)qx^3y^1)e_1+yxy^2e_2)=yd_{}x.`$
From the construction it is clear that $`\psi _\epsilon `$ is a surjective FODC homomorphism. We show that $`\psi _{}`$ is injective and suppose that $`ud_{}x+vd_{}y=0`$ for some elements $`u,v𝒪(_q^2)`$. Inserting the expressions from (96) and using the assumption that $`e_1`$ and $`e_2`$ are linearly independent, we get
$$uq^2x^3y^2+vx^2y^1=0,vxy^2=0$$
which in turn implies that $`u=v=0`$. The proof for $`\psi _+`$ is similar. $`\mathrm{}`$
We shall identify the isomorphic calculi $`\mathrm{\Gamma }_\epsilon `$ and $`\stackrel{~}{\mathrm{\Gamma }}_\epsilon `$. The above approach to the calculi $`\mathrm{\Gamma }_\epsilon `$ is convenient for many purposes. Among others, it allows us easily to extend these calculi to larger algebras.
The partial derivatives $`_x`$ and $`_y`$ can be also expressed in terms of the action of the generators of $`𝒰_q(gl_2)`$. Combining the formulas (79), (80) and (92) we obtain for the calculus $`\mathrm{\Gamma }_{}`$ the relations
$$_x(z)=q^{\frac{3}{2}}y^1EK_1^3K_2z,_y(z)=q^{\frac{1}{2}}x^1FK_1^3K_2z$$
for $`z𝒪(_q^2)`$, where the action of $`E,F,K_1,K_2`$ is given by Proposition 15 and the elements $`y^1`$ and $`x^1`$ act by left multiplication on $`𝒪(_q^2)`$.
3. An Auxilary $``$-algebra $`𝒲`$
3.1 The $`𝒰_q(gl_2())`$-module $``$-algebra $`𝒲`$
Let $`𝒲`$ denote the $``$-algebra generated by the operators
$$W(s,t):=e^{2\pi i(s\alpha 𝒬+t\beta 𝒫)},s,t.$$
(97)
These operators satisfy the relations
$`W(s_1,t_1)W(s_2,t_2)`$ $`=e^{\pi i\gamma (s_2t_1s_1t_2)}W(s_1+s_2,t_1+t_2),`$ (98)
$`W(s,t)^{}`$ $`=W(\overline{s},\overline{t})`$ (99)
for $`s_1,t_1,s_2,t_2,s,t`$. By Lemma 4, the operator $`W(s,t)`$ acts as
$$(W(s,t)f)(x)=e^{2\pi is\alpha x+\pi ist\gamma }f(x+t).$$
(100)
Equations (98) and (100) hold for vectors contained in the corresponding operator domains. For instance, they hold on the domains $`𝒟_\delta `$, where $`\delta >0`$, and $`𝔄(^2)`$ in the Hilbert space $`L^2(^2)`$. Each of the dense subspaces is an invariant dense core for all operators $`W(s,t)`$. From (17) and (100) we see that
$$W(i\mathrm{,0})=X,W(0,i)=Y,W(s\mathrm{,0})=X^{is},W(0,t)=Y^{is},s.$$
(101)
Our next aim is to define a left action of the Hopf algebra $`𝒰_q(gl_2)`$ on $`𝒲`$. Let us identify the generators $`x`$ with $`X=W(i\mathrm{,0})`$ and $`y`$ with $`Y=W(0,i)`$. Then $`𝒪(_q^2)`$ becomes a $``$-subalgebra of $`𝒲`$. We now use formulas (78)-(80) (which have been proved only for polynomials $`g`$ and $`h`$!) as a motivation and extend them formally to the functions $`g(X)=X^{is}`$ and $`h(Y)=Y^{it},s,t`$, of the positive self-adjoint operators $`X`$ and $`Y`$ defined by (17). Throughout we interpret expressions $`(q^{k/2}X)^{is}`$ and $`(q^{k/2}Y)^{it}`$ as $`e^{\pi k\gamma s}e^{2\pi is\alpha 𝒬}`$ and $`e^{\pi k\gamma t}e^{2\pi it\beta 𝒫}`$, respectively, for $`k=0,\pm \mathrm{1,3}`$ and $`s,t`$. Recall that
$$W(s,t)=e^{\pi i\gamma st}W(s\mathrm{,0})W(0,t)=e^{\pi i\gamma st}X^{is}Y^{it}$$
(102)
by (97). Applying now (78) and (80) formally (!) and using (102) we derive
$`\lambda FW(s,t)`$ $`=\lambda e^{\pi i\gamma st}((FX^{is})(KY^{it})+(K^1X^{is})(FYX^{it})`$
$`=\lambda e^{\pi i\gamma st}(0+(q^{1/2}X)^{is}q^{1/2}XD_{q^2}((q^{1/2}Y)^{it}))`$
$`=q^{1/2}(1q^2)e^{\pi i\gamma st}e^{\pi \gamma s}X^{is}X\left(\frac{(q^{1/2}Y)^{it}(q^{1/2}q^2Y)^{it}}{(1q^2)Y}\right)`$
$`=q^{1/2}e^{\pi \gamma s}e^{\pi i\gamma st}X^{is}X\left(e^{\pi \gamma t}e^{3\pi \gamma t}\right)Y^{it}Y^1`$
$`=q^{1/2}e^{\pi \gamma s}\left(e^{\pi \gamma t}e^{3\pi \gamma t}\right)e^{\pi i\gamma st}X^{is+1}Y^{it1}`$
$`=q^{1/2}e^{\pi \gamma s}\left(e^{\pi \gamma t}e^{3\pi \gamma t}\right)e^{\pi i\gamma st}W(si\mathrm{,0})W(0,t+i)`$
$`=q^{1/2}e^{\pi \gamma s}\left(e^{\pi \gamma t}e^{3\pi \gamma t}\right)e^{\pi i\gamma st}e^{\pi i\gamma (si)(t+i)}W(si,t+i)`$
$`=\left(e^{2\pi \gamma t}e^{2\pi \gamma t}\right)W(si,t+i).`$
The formulas for the actions of the other generators $`E,K_1`$ and $`K_2`$ are derived by a similar formal reasoning. Replacing (80) by (79) and (78) we obtain
$`\lambda EW(s,t)=(e^{2\pi \gamma s}e^{2\pi \gamma s})W(s+i,ti),`$
$`K_1W(s,t)=e^{\pi \gamma s}W(s,t),K_2W(s,t)=e^{\pi \gamma t}W(s,t).`$
We now take the above formulas which have been obtained by formal algebraic manipulations as the starting point for the rigorous definition of a left action of $`𝒰_q(gl_2)`$ on the $``$-algebra $`𝒲`$. That is, for $`s,t`$ we define
$`EW(s,t)=\lambda ^1(e^{2\pi \gamma s}e^{2\pi \gamma s})W(s+i,ti),`$ (103)
$`FW(s,t)=\lambda ^1(e^{2\pi \gamma t}e^{2\pi \gamma t})W(si,t+i),`$ (104)
$`K_1W(s,t)=e^{\pi \gamma s}W(s,t),K_2W(s,t)=e^{\pi \gamma t}W(s,t).`$ (105)
Proposition 17. With definitions (103)–(105), $`𝒲`$ is a left $`𝒰_q(gl_2())`$-module $``$-algebra.
Proof. Since the set of operators $`W(s,t),s,t`$, is linearly independent as easily shown, the preceding definitions extend uniquely to well-defined linear mappings of $`𝒲`$ into itself. It is straightforward to check that the terms $`K_1Eq^{1/2}EK_1,K_2Eq^{1/2}EK_2,K_1Fq^{1/2}FK_1,K_2Fq^{1/2}FK_2`$ and $`\lambda EF\lambda FEK^2+K^2`$ applied to an arbitrary basis element $`W(s,t)`$ of $`𝒲`$ vanish. Thus, formulas (103)–(105) define indeed a left action of the algebra $`U_q(gl_2)`$ on $`𝒲`$. That the left module $`𝒲`$ is a $`𝒰_q(gl_2)`$-module algebra means that (1) is satisfied. It suffices to check this condition for the generators $`f=E,F,K_1,K_2`$, $`K_1^1`$, $`K_2^1`$ and $`z=W(s,t),z^{}=W(s^{},t^{}),s,t,s^{},t^{}`$. As a sample, we carry out this for the generator $`f=E`$. Using (103), (105) and (100) we compute
$`\lambda (EW(s,t))(KW(s^{},t^{}))+\lambda (K^1W(s,t))(EW(s^{},t^{}))`$
$`=(e^{2\pi \gamma s}e^{2\pi \gamma s})e^{\pi \gamma (s^{}t^{})}W(s+i,ti)W(s^{},t^{})`$
$`+e^{\pi \gamma (ts})(e^{2\pi \gamma s^{}}e^{2\pi \gamma s^{}})W(s,t)W(s^{}+it^{}i)`$
$`=(e^{2\pi \gamma s}e^{2\pi \gamma s})e^{\pi \gamma (s^{}t^{})}e^{\pi i\gamma (s^{}(ti)(s+i)t^{})}W(s+s^{}+i,t+t^{}i)`$
$`+e^{\pi \gamma (ts})(e^{2\pi \gamma s^{}}e^{2\pi \gamma s^{}})e^{\pi i\gamma ((s^{}+i)ts(t^{}i))}W(s+s^{}+i,t+t^{}i)`$
$`=(e^{2\pi \gamma (s+s^{})}e^{2\pi \gamma (s+s^{})})e^{\pi i\gamma (s^{}tst^{})}W(s+s^{}+i,t+t^{}i)`$
$`=\lambda e^{\pi i\gamma (s^{}tst^{})}EW(s+s^{},t+t^{})`$
$`=\lambda E(W(s,t)W(s^{},t^{})).`$
This proves (1) in the case $`f=E`$.
Finally, it remains to check that (4) holds. Since $`𝒲`$ is a left $`𝒰_q(gl_2)`$-module algebra as just shown, it suffices to do this for the generators $`f`$ of $`𝒰_q(gl_2)`$. Again we restrict ourselves to the case $`f=\lambda E,z=W(s,t)`$. Since $`S(\lambda E)^{}=\overline{\lambda }(q)E^{}=\lambda E`$ by (11) and $`W(s,t)^{}=W(\overline{s},\overline{t})`$, we have
$`(\lambda EW(s,t))^{}=(e^{2\pi \gamma s}e^{2\pi \gamma s})W(s+i,ti))^{}`$
$`=(e^{2\pi \gamma \overline{s}}e^{2\pi \gamma \overline{s}})W(\overline{s}+i,\overline{t}i)=\lambda EW(\overline{s},\overline{t})`$
$`=S(\lambda E)^{}W(s,t)^{}.\mathrm{}`$
The $``$-algebra $`𝒲`$ consists of Hilbert space operators and formulas (103)–(105) have been derived by using formal operator calculus. However, the content of Proposition 17 is purely algebraic: It is obvious that the complex vector space $`𝒲`$ with basis $`W(s,t),s,t`$ is a $``$-algebra with multiplication and involution defined by (98) and (99). Proposition 17 says that $`𝒲`$ is a $`𝒰_q(gl_2())`$-module $``$-algebra with respect to the left action defined by (103)–(105).
Recall that $`𝒪(_q^2)`$ is a $`𝒰_q(gl_2())`$-module $``$-subalgebra of $`𝒲`$. By definition, the products $`xw`$ and $`yw`$ for $`w𝒲`$ are the operator products $`Xw`$ and $`Yw`$, respectively, in the Hilbert space $`L^2()`$. From (97) and (101) we obtain
$`xW(s,t)=e^{\pi \gamma t}W(si,t),yW(s,t)=e^{\pi \gamma s}W(s,ti),`$ (106)
$`W(s,t)x=e^{\pi \gamma t}W(si,t),W(s,t)y=e^{\pi \gamma s}W(s,ti).`$ (107)
3.2 Covariant differential calculus on $`𝒲`$
In this subsection we extend the differential calculus $`\mathrm{\Gamma }_{}`$ of $`𝒪(_q^2)`$ to $`𝒲`$. In order to do so, we use the approach given in 2.2 with $`𝒵=𝒲`$ and write $`\mathrm{\Gamma },d,\omega `$ for $`\mathrm{\Gamma }_{},d_{},\omega _{}`$, respectively.
As in 2.2, we set $`\omega =q^2x^2y^2e_1+y^2e_2`$ and define
$$dz=\omega zz\omega ,z𝒲.$$
Obviously, $`\mathrm{\Gamma }:=𝒲d𝒲𝒲`$ is a first order differential calculus over $`𝒲`$ with differentiation $`d`$ such that the differentials $`dx,dy`$ form a free left $`𝒲`$-module basis of $`\mathrm{\Gamma }`$. Because of this property, the partial derivatives $`_x(z)`$ and $`_y(z)`$ are well-defined by (87). In order to compute the latter for $`z=W(s,t)`$, we use the commutation rules $`xW(s,t)=e^{2\pi \gamma t}W(s,t)x`$ and $`yW(s,t)=e^{2\pi \gamma s}W(s,t)y`$ (by (106) and (107)) and the expressions (96) for $`dx`$ and $`dy`$. Comparing coefficients in (87), we obtain for $`s,t`$,
$`_x(W(s,t))={\displaystyle \frac{1e^{4\pi \gamma s}}{1q^2}}e^{\pi \gamma t}W(s+i,t),_y(W(s,t))={\displaystyle \frac{1e^{4\pi \gamma t}}{1q^2}}e^{3\pi \gamma s}W(s,t+i.)`$
4. The $`𝒰_𝐪(\mathrm{𝐠𝐥}_\mathrm{𝟐}())`$-module $``$-algebra $`𝒜(_𝐪^{++})`$
4.1 In the preceding section we extended the action of the Hopf $``$-algebra $`𝒰_q(gl_2())`$ on $`𝒪(_q^2)`$ to the larger $``$-algebra $`𝒲`$ such that $`𝒲`$ is a left module $``$-algebra of $`𝒰_q(gl_2())`$. We now go one step further and make the $``$-algebra $`𝔄(^2)`$ into a left $`𝒰_q(gl_2())`$-module $``$-algebra. In order to do so we use the formulas (103)–(105) in order to derive the corresponding formulas for the action of the generators $`E,F,K_1,K_2`$ on $`Op(a)`$. Suppose that $`a𝔄(^2)`$. For the generator $`E`$ we obtain
$$\begin{array}{cc}& \lambda EOp(a)=\gamma \widehat{a}(\alpha s,\beta t)(\lambda EW(s,t))𝑑s𝑑t\hfill \\ & =\gamma \widehat{a}(\alpha s,\beta t)(e^{2\pi \gamma s}e^{2\pi \gamma s})W(s+i,ti)𝑑s𝑑t\hfill \\ & =\gamma \widehat{a}(\alpha (si),\beta (t+i))(e^{2\pi \gamma (si)}e^{2\pi \gamma (si)})W(s,t)𝑑s𝑑t.\hfill \end{array}$$
Let us explain the steps of this computation. The first equality is only a formal interchanging of integrals and left action, while the second follows from formula (103). The third equality is obtained by the formal replacements $`ss+i`$ and $`tti`$. These substitutions are justified by a standard argument from complex analysis which has been used already in the proof of Lemma 11: The integral of the holomorphic operator-valued function
$$s\widehat{a}(\alpha s,\beta t)(e^{2\pi \gamma s}e^{2\pi \gamma s})W(s+i,ti)$$
along the boundary of the rectangle $`R,R,Ri,Ri`$ for fixed $`t`$ and $`R>0`$ is zero. By Lemma 9, the integrals from $`R`$ to $`Ri`$ and from $`Ri`$ to $`R`$ tend to zero as $`R+\mathrm{}`$. Arguing similarly for the variable $`t`$, the third equality is obtained. In order to complete this reasoning, we note that the function
$$\widehat{a}(\alpha (si),\beta (t+i))(e^{2\pi \gamma (si)}e^{2\pi \gamma (si)})$$
is the Fourier transform of the function $`a_E𝔄(^2)`$ defined by
$$a_E(x_1,x_2):=e^{2\pi (\beta x_2\alpha x_1)}(a(x_1+\beta i,x_2)a(x_1\beta i,x_2)).$$
(108)
Thus, we have seen that $`\lambda EOp(a)=Op(a_E)`$. Using (104) and (105) instead of (103) a similar reasoning shows that $`\lambda FOp(a)=Op(a_F)`$ and $`K_jOp(a)=Op(a_{K_j})`$, where the symbol $`a_F,a_{K_j}𝔄(^2)`$ are given by
$$\begin{array}{cc}\hfill a_F(x_1,x_2)& =e^{2\pi (\alpha x_1\beta x_2)}(a(x_1,x_2+\alpha i)a(x_1,x_2\alpha i)),\hfill \\ \hfill a_{K_1}(x_1,x_2)& =a(x_1\frac{\beta }{2}i,x_2),\hfill \\ \hfill a_{K_2}(x_1,x_2)& =a(x_1,x_2\frac{\alpha }{2}i).\hfill \end{array}$$
Summarizing, in terms of the symbol we have derived the following formulas for the actions of the generators $`E,F,K_1,K_2`$ of $`𝒰_q(gl_2())`$:
$`(Ea)(x_1,x_2)`$ $`=\lambda ^1e^{2\pi (\beta x_2\alpha x_1)}(a(x_1+\beta i,x_2)a(x_1\beta i,x_2)),`$ (109)
$`(Fa)(x_1,x_2)`$ $`=\lambda ^1e^{2\pi (\alpha x_1\beta x_2)}(a(x_1,x_2+\alpha i)a(x_1,x_2\alpha i)),`$ (110)
$`(K_1a)(x_1,x_2)`$ $`=a(x_1\frac{\beta }{2}i,x_2),(K_2a)(x_1,x_2)=a(x_1,x_2\frac{\alpha }{2}i).`$ (111)
The derivation of these formulas is rigorous except for the justification of the interchanging of integrals and actions. This could be made rigorous by introducing appropriate locally convex topologies. We shall not proceed this way, because we shall use formulas (109)–(111) only as definitions of the action of $`U_q(gl_2())`$ on $`𝔄(^2)`$ and prove the corresponding properties directly in 4.3.
Note that formulas (109)–(111) and also formulas (112) and (113) below are meaningful for larger classes of symbols rather than $`𝔄(^2)`$. For instance, for the function $`a(x_1,x_2)=e^{2\pi i(\alpha sx_1+\beta tx_2)}`$ (which is of course not in $`𝔄(^2)`$ ) we have $`Op(a)=e^{2\pi i(\alpha s𝒬+\beta t𝒫)}=W(s,t)`$. In this case formulas (109)–(111) reduces to the equations (103)–(105) derived in the preceding section. If we allow the symbols to be distributions, then we recover also formulas (78)–(80).
In a similar manner the product of the operators $`Op(a)`$ with operators $`X`$ and $`Y`$ can be computed by using formulas (106) and (107) (or (48)). We then obtain $`XOp(a)=Op({}_{x}{}^{}a)`$, $`Op(a)X=Op(a_x)`$, $`YOp(a)=Op({}_{y}{}^{}a)`$ and $`Op(a)Y=Op(a_y)`$, where the symbol $`{}_{x}{}^{}a,a_x,{}_{y}{}^{}a,a_y𝔄(^2)`$ are given by
$`{}_{x}{}^{}a(x_1,x_2)=e^{2\pi \alpha x_1}a(x_1,x_2+\frac{\alpha }{2}i),a_x(x_1,x_2)=e^{2\pi \alpha x_1}a(x_1,x_2\frac{\alpha }{2}i),`$
$`{}_{y}{}^{}a(x_1,x_2)=e^{2\pi \beta x_2}(a(x_1\frac{\beta }{2}i),x_2),a_y(x_1,x_2)=e^{2\pi \beta x_2}a(x_1+\frac{\beta }{2}i,x_2).`$
Let $`𝒜(_q^{++})`$ denote the direct sum of vector spaces $`𝒪(_q^2)`$ and $`𝔄(^2)`$.
Lemma 18. There is a unique structure of a $``$-algebra on $`𝒜(_q^{++})`$ such that $`𝒪(_q^2)`$ and $`𝔄(^2)`$ are $``$-subalgebras of $`𝒜(_q^{++})`$ and the products of the generators $`x,y`$ of $`𝒪(_q^2)`$ and symbols $`a𝔄(^2)`$ are given by
$`xa(x_1,x_2)=e^{2\pi \alpha x_1}a(x_1,x_2+\frac{\alpha }{2}i),ax(x_1,x_2)=e^{2\pi \alpha x_1}a(x_1,x_2\frac{\alpha }{2}i),`$ (112)
$`ya(x_1,x_2)=e^{2\pi \beta x_2}a(x_1\frac{\beta }{2}i,x_2),ay(x_1,x_2)=e^{2\pi \beta x_2}a(x_1+\frac{\beta }{2}i,x_2).`$ (113)
Proof. We first note that the maps $`z\rho _{++}(z)`$ (see (18)) and $`aOp(a)`$ (see (49)) are faithful $``$-representations of the $``$-algebras $`𝒪(_q^2)`$ and $`𝔄(^2)`$ on the domain $`𝔄()`$ on the Hilbert space $`L^2()`$. Since $`\rho _{++}(z)Op(a)=Op({}_{z}{}^{}a)`$ and $`Op(a)\rho _{++}(z)=Op(a_z)`$ for $`z=x,y`$ by the above definitions, it is clear that the sum $`\rho _{++}(𝒪(_q^2))+Op(𝔄(^2))`$ of the images of these $``$-algebras is a $``$-algebra of unbounded operators on the domain $`𝔄()`$. It is easily seen that the only bounded operators in $`\rho _{++}(𝒪(_q^2))`$ are the multiples of the identity operator and that no operator in $`Op(𝔄(^2))`$ is a multiple of the identity. Thus, $`\rho _{++}(𝒪(_q^2))Op(𝔄(^2))=\{0\}`$. Hence the map $`𝒥:(z,a)\rho _{++}(z)+Op(a)`$ of $`𝒜(_q^{++})`$ to $`\rho _{++}(𝒪(_q^2))+Op(𝔄(^2))`$ is bijective. The unique $``$-algebra structure on $`𝒜(_q^{++})`$ for which $`𝒥`$ is a $``$-homomorphism has obviously the desired properties. $`\mathrm{}`$
We shall show by Theorem 21 below that $`𝒜(_q^{++})=𝒪(_q^2)+𝔄(^2)`$ is even a left $`𝒰(gl_2())`$-module $``$\- algebra. We call this left $`𝒰(gl_2())`$-module $``$\- algebra $`𝒜(_q^{++})`$ the $``$-algebra of functions on the quantum quarter plane. Obviously, the $``$-subalgebra $`𝒪(_q^2)`$ is considered as the algebra generated by the two coordinate functions $`x`$ and $`y`$ of the quantum quarter plane. The elements of $`𝔄(^2)`$ can be interpreted as “functions on the quantum quarter plane which go rapidly to zero at the boundary of the quantum quarter plane”. Note that $`𝔄(^2)`$ is a two-sided $``$-ideal of the $``$-algebra $`𝒜(_q^{++})`$.
4.2 In this subsection we introduce two useful algebra homomorphisms in order to understand the algebraic content behind formulas (109)–(111). Let $`_q`$ denote the complex unital algebra with generators $`x_1,x_1^1,y_1,y_1^1`$, $`x_2`$, $`x_2^1`$,$`y_2,y_2^1`$ and defining relations
$$x_jy_j=q^{1/8}y_jx_j,x_jx_j^1=x_j^1x_j=1,y_jy_j^1=y_j^1y_j=1\text{for}j=\mathrm{1,2},$$
(114)
$$x_1x_2=x_2x_1,y_1y_2=y_2y_1,x_1y_2=y_2x_1,x_2y_1=y_1x_2,$$
(115)
where we set $`q^{1/8}:=e^{\pi \gamma i/4}`$. The subalgebra $`_j`$, $`j=\mathrm{1,2}`$, generated by $`x_j,x_j^1,y_j`$, $`y_j^1`$ is nothing but the localization of the algebra $`𝒪(_{q^{1/8}}^2)`$ at the elements $`x_j`$ and $`y_j`$, and $`_q`$ is just the tensor product of the algebras $`_1`$ and $`_2`$.
Lemma 19. There are injective algebra homomorphisms $`\psi :𝒰_q(gl_2)_q`$ and $`\psi :𝒪(_q^2)_q`$ such that
$`\psi (E)`$ $`=\lambda ^1x_2^2x_1^2(y_1^4y_1^4),`$ (116)
$`\psi (F)`$ $`=\lambda ^1x_1^2x_2^2(y_2^4y_2^4),`$ (117)
$`\psi (K_1)`$ $`=y_1^2,\psi (K_2)=y_2^2,`$ (118)
$`\psi (x)`$ $`=x_1^2y_2^2,\psi (y)=x_2^2y_1^2.`$ (119)
Proof. In oder to prove the assertion for $`𝒰_q(gl_2)`$ it suffices to check that the operators $`\psi (E),\psi (F),\psi (K_1)`$ and $`\psi (K_2)`$ satisfy the defining relations of the algebra $`𝒰_q(gl_2)`$. Using the relations (114)–(115) of the algebra $`_q`$ we obtain
$`\lambda \psi (E)\psi (K_1)`$ $`=x_2^2x_1^2(y_1^4y_1^4)y_1^2=x_2^2(x_1^2y_1^2)(y_1^4y_1^4)`$
$`=x_2^2(q^{1/4})^2y_1^2x_1^2(y_1^4y_1^4)=q^{1/2}y_1^2x_2^2x_1^2(y_1^4y_1^4)`$
$`=q^{1/2}\lambda \psi (K_1)\psi (E).`$
The relations $`\psi (E)\psi (K_2)=q^{1/2}\psi (K_2)\psi (E),\psi (F)\psi (K_1)=q^{1/2}\psi (K_1)\psi (F),`$
$`\psi (F)\psi (K_2)=q^{1/2}\psi (K_2)\psi (F)`$ and $`\lambda (\psi (E)\psi (F)\psi (F)\psi (E))=\psi (K)^2\psi (K)^2`$ are verified by similar computations. Obviously we have $`\psi (x)\psi (y)=q\psi (y)\psi (x)`$. Hence the above formulas define indeed algebra homomorphisms of $`𝒰_q(gl_2)`$ and $`𝒪(_q^2)`$ into $`_q`$. Since the sets $`\{E^kK_1^nK_2^mF^l;k,l_0,n,m\}`$, $`\{x^ky^n;k,n_0\}`$ and $`\{x_1^kx_2^ly_1^ny_2^m;k,l,n,m_0\}`$ are vector space bases of $`𝒰_q(gl_2)`$, $`𝒪(_q^2)`$ and $`_q`$, respectively, it follows easily from formulas (109)-(119) that the mappings $`\psi :𝒰_q(gl_2)_q`$ and $`\psi :𝒪(_q^2)_q`$ are injective. $`\mathrm{}`$
Since $`|q|=1`$, $`_q`$ is a $``$-algebra with involution determined by $`x_j^{}:=x_j`$ and $`y_j^{}:=y_j`$, $`j=\mathrm{1,2}.`$ The algebra homomorphism $`\psi :𝒰_q^{tw}(gl_2())_q`$ does not preserve the involution. The next lemma shows that $`\psi `$ is similar to a $``$-homomorphism.
Lemma 20. For $`z𝒰_q(gl_2())`$ and $`z𝒪(_q^2)`$, define
$$\phi (z)=x_1x_2y_1^1y_2\psi (z)(x_1x_2y_1^1y_2)^1.$$
Then $`\phi :𝒰_q^{tw}(gl_2())_q`$ and $`\phi :𝒪(_q^2)_q`$ are injective $``$-homomorphisms of the corresponding $``$-algebras. In fact, we have
$`\phi (E^{})`$ $`=x_2^2x_1^1(y_1^4y_1^4)x_1^1=x_2^2x_1^2(q^{1/2}y_1^4q^{1/2}y_1^4)`$
$`=x_2^2(q^{1/2}y_1^4q^{1/2}y_1^4)x_1^2,`$ (120)
$`\phi (F^{})`$ $`=x_1^2x_2^1(y_2^4y_2^4)x_2^1=x_1^2x_2^2(q^{1/2}y_2^4q^{1/2}y_2^4),`$
$`=x_1^2(q^{1/2}y_2^4q^{1/2}y_2^4)x_2^2,`$ (121)
$`\phi (q^{1/4}K_1)`$ $`=\psi (K_1)=y_1^2,\phi (q^{1/4}K_2)=\psi (K_2)=y_2^2,`$ (122)
$`\phi (x)`$ $`=\psi (x)=x_1^2y_2^2,\phi (y)=\psi (y)=x_2^2y_1^2.`$ (123)
Proof. Clearly, $`\phi :𝒰_q(gl_2())_q`$ and $`\phi :𝒪(_q^2)_q`$ are injective homomorphisms, because $`\psi `$ are by Lemma 19. Therefore, it is sufficient to prove that $`\phi (z^{})=\phi (z)^{}`$ for the four generators $`z=E^{},F^{},q^{1/4}K_1,q^{1/4}K_2`$ of $`𝒰_q^{tw}(gl_2())`$ and the two generators $`z=x,y`$ of $`𝒪(_q^2)`$. Since all these generators $`z`$ and their images $`\phi (z)`$ are hermitean, it suffices to check formulas (S0.Ex135)–(123). The latter formulas follow by straightforward computations from (109)-(119) combined with the relations (114)-(115) of the algebra $`_q`$. As a sample we verify (S0.Ex135) and compute
$`\phi (\lambda E)`$ $`=x_1x_2y_1^1y_2\psi (\lambda E)(x_1x_2y_1^1y_2)^1`$
$`=x_1x_2y_1^1y_2x_2^2x_1^2(y_1^4y_1^4)y_2^1y_1x_2^1x_1^1`$
$`=x_1x_2q^{1/4}x_1^2y_1^1q^{1/4}x_2^2y_2(y_1^4y_1^4)y_2^1y_1x_2^1x_1^1`$
$`=q^{1/2}x_2^2x_1^1(y^4y_1^4)x_1^1`$
which gives the first formula of(S0.Ex135). The second and third formulas of (S0.Ex135) follow by applying once more the commutation rules (114) and (115). $`\mathrm{}`$
The $``$-homomorphisms $`\phi `$ of $`𝒰_q^{tw}(gl_2())`$ and $`𝒪(_q^2)`$ are crucial in what follows.
4.3 Let us return to the left action of the Hopf $``$-algebra $`𝒰(gl_2())`$ on $`𝔄(^2)`$ given by the formulas (109)–(111). We define a $``$-representation $`\rho _0`$ of the $``$-algebra $`_q`$ on the invariant dense domain $`𝔄(^2)`$ of the Hilbert space $`L^2(^2)`$ by
$$\rho _0(x_1)=e^{\pi \alpha 𝒬_1},\rho _0(y_1)=e^{{\scriptscriptstyle \frac{\pi }{2}}\beta 𝒫_1},\rho _0(x_2)=e^{\pi \beta 𝒬_2},\rho _0(y_2)=e^{{\scriptscriptstyle \frac{\pi }{2}}\alpha 𝒫_2},$$
(124)
where $`q^{1/4}=e^{\pi \gamma i/2}`$ and as always $`\alpha \beta =\gamma `$. It is obvious that these operators satisfy the relations of the $``$-algebra $`_q`$, so (124) defines indeed a $``$-representation of $`_q`$. Inserting (124) into (116)–(118) we see that equations (109)–(111) can be expressed as
$$fa=\rho _0(\psi (f))a,a𝔄(^2),$$
(125)
for the generators $`f=E,F,K_1,K_2`$ of $`𝒰_q(gl_2)`$. We now take this equation as a definition for arbitrary elements $`f𝒰_q(gl_2)`$. Since $`\rho _0\psi `$ is an algebra homomorphism, (125) gives a well-defined left action of the algebra $`𝒰_q(gl_2)`$ on $`𝔄(^2)`$. Recall from 2.1 that we have also a left action $``$ of $`𝒰_q(gl_2)`$ on $`𝒪(_q^2)`$. Hence the equation
$$f(z+a):=fz+fa,f𝒰_q(gl_2),z𝒪(_q^2),a𝔄(^2),$$
(126)
defines a left action of $`𝒰_q(gl_2)`$ on the direct sum $`𝒜(_q^{++})=𝒪(_q^2)+𝔄(^2)`$. In terms of the $``$-representations $`\rho _0`$ formulas (112) and (113) can be written as
$`xa=\rho _0\psi (x)a=\rho _0(x_1^2y_2^2)a,ax=\rho _0(x_1^2y_2^2)a,`$ (127)
$`ya=\rho _0\psi (y)a=\rho _0(x_2^2y_1^2)a,ay=\rho _0(x_2^2y_1^2)a.`$ (128)
The main result of this section is the following theorem.
Theorem 21. With the preceding definitions, the $``$-algebra $`𝒜(_q^{++})`$ of functions on the quantum quarter plane is a left $`𝒰_q(gl_2())`$-module $``$-algebra.
Proof. We already noticed that $``$ is a left action of the algebra $`𝒰_q(gl_2)`$ on $`𝒜(_q^{++})`$. It remains to show that conditions (1) and (4) are fulfilled for arbitrary elements $`z,z^{}𝒜(_q^{++})`$ and $`f𝒰_q(gl_2)`$.
We first prove that $`𝔄(^2)`$ is a $`𝒰_q(gl_2)`$-left module algebra. Since $``$ is a left action of $`𝒰_q(gl_2)`$, it suffices to prove (1) for the generators $`f=\lambda E,\lambda F,K_1,K_2`$. These verifications are lengthy but straigthforward. We restrict ourselves to the case $`f=\lambda E`$. Then we compute
$`(\lambda Ea)\mathrm{\#}(Kb)(x_1,x_2)+(K^1a)\mathrm{\#}(\lambda Eb)(x_1,x_2)`$
$`=`$ $`4{\displaystyle 𝑑u_1𝑑u_2𝑑v_1𝑑v_2e^{4\pi i[(x_1u_1)(x_2v_2)(x_1v_1)(x_2v_2)]}}`$
$`\{e^{2\pi (\beta u_2\alpha u_1)}(a(u_1+\beta i,u_2)a(u_1\beta i,u_2))b(v_1\frac{\beta }{2}i,v_2+\frac{\alpha }{2}i)`$
$`+a(u_1+\frac{\beta }{2}i,u_2\frac{\alpha }{2}i)e^{2\pi (\beta v_2\alpha v_1)}(b(v_1+\beta i,v_2)b(v_1\beta i,v_2))\}`$
$`=`$ $`4{\displaystyle 𝑑u_1𝑑u_2𝑑v_1𝑑v_2a(u_1,u_2)b(v_1,v_2)}`$
$`\{e^{2\pi (\beta u_2\alpha (u_1\beta i))+4\pi i[(x_1u_1+\beta i)(x_2v_2i\alpha /2)(x_1v_1i\beta /2)(x_2u_2)]}`$
$`e^{2\pi (\beta u_2\alpha (u_1+\beta i))+4\pi i[(x_1u_1+\beta i)(x_2v_2+i\alpha /2)(x_1v_1i\beta /2)(x_2u_2)]}`$
$`+e^{2\pi (\beta v_2\alpha (v_1\beta i))+4\pi i[(x_1u_1+i\beta /2)(x_2v_2)(x_1v_1+\beta i)(x_2u_2i\alpha /2)]}`$
$`e^{2\pi (\beta v_2\alpha (v_1+\beta i))+4\pi i[(x_1u_1+i\beta /2)(x_2v_2)(x_1v_1\beta i)(x_2u_2i\alpha /2)]}\}`$
$`=`$ $`4{\displaystyle 𝑑u_1𝑑u_2𝑑v_1𝑑v_2a(u_1,u_2)b(v_1,v_2)}`$
$`\{e^{2\pi (\beta x_2\alpha x_1)+4\pi i[(x_1\beta iu_1)(x_2v_2)(x_1\beta iv_1)(x_2u_2)]}`$
$`+e^{2\pi (\beta v_2\alpha v_1)+4\pi i[(x_1+\beta iu_1)(x_2v_2)(x_1+\beta iv_1)(x_2u_2)]}\}`$
$`=`$ $`(\lambda E(a\mathrm{\#}b))(x_1,x_2).`$
The first equality is obtained by inserting the formulas (109) and (111) for the actions of $`E`$ and $`K`$ and (50) for the product $`\mathrm{\#}`$ of the algebra $`𝔄(^2)`$. The second equality follows by the substitution $`u_1u_1+\beta i,v_1v_1\frac{\beta }{2}i,v_2v_2+\frac{\alpha }{2}i`$ of the first summand and similar replacements of the other three summands. As noted in the considerations preceding (51), these substitutions are justified because of Lemma 9. Next let us consider the expressions in the four exponentials after the second equality sign. By regrouping these terms we see that the first and the fourth exponentials cancel, while the second and third ones can be reexpressed as the exponentials after the third equality sign. The fourth equality follows by applying once more formulas (50) and (109). By a similar reasoning condition (1) can be checked for the other generators $`f=\lambda F,K_1,K_2`$. Thus, $`𝔄(^2)`$ is a $`𝒰_q(gl_2)`$-left module algebra.
Recall from 2.1 that $`𝒪(_q^2)`$ is also a $`𝒰_q(gl_2)`$-left module algebra. Therefore, in order to prove that the sum $`𝒜(_q^{++})=𝒪(_q^2)+𝔄(^2)`$ is a $`𝒰_q(gl_2)`$-left module algebra, it remains to show that
$`f(wa)`$ $`=(f_{(1)}w)(f_{(2)}a),`$ (129)
$`f(aw)`$ $`=(f_{(1)}a)(f_{(2)}w)`$ (130)
for $`f𝒰_q(gl_2),w𝒪(_q^2)`$ and $`a𝔄(^2)`$. It is easily seen that equation (129) holds for the product $`fg`$ and arbitrary $`w`$ and $`a`$ provided that (129) holds for $`f`$ and arbitrary $`w`$ and $`a`$ and also for $`g`$ and arbitrary $`w`$ and $`a`$. Hence it suffices to check condition (129) for elements $`f`$ from a set $`M`$ of generators of the algebra $`𝒰_q(gl_2)`$ and for arbitrary $`w`$ and $`a`$. Suppose in addition that $`M`$ is a vector space such that $`\mathrm{\Delta }(M)𝒰_q(gl_2)M`$. Let $`w,w^{}𝒪(_q^2)`$ such that (129) holds for $`w`$ and all $`fM`$ and $`a`$ and also for $`w^{}`$ and all $`fM`$ and $`a`$. We show that then (129) holds for the product $`ww^{}`$ and arbitrary $`fM`$ and $`a`$ by computing
$`(f_{(1)}(ww^{}))(f_{(2)}a)=(f_{(1)}w)(f_{(2)}w^{})(f_{(3)}a)=(f_{(1)}w)(f_{(2)}(w^{}a))=f(ww^{}a).`$
Note that for the second equality we used that $`\mathrm{\Delta }(f)𝒰_q(gl_2)M`$ by assumption and so (129) is valid for the elements in the second tensor factor of $`\mathrm{\Delta }(f)`$. Applying the preceding with $`M=\mathrm{Lin}\{E,F,K_1,K_2\}`$ we conclude that condition (129) is fulfilled provided that it holds for $`f=E,F,K_1,K_2`$, $`w=x,y`$ and arbitrary $`a𝔄(^2)`$. Arguing in a similar manner with condition (130) it follows that it is sufficient to verify (130) for the generators $`f=E,F,K_1,K_2`$ and $`w=x,y`$. As a sample, we prove equation (129) for $`f=E`$ and $`w=x`$. Using the formulas (75), (116), (118) and (127) we obtain
$`(\lambda Ex)(Ka)+`$ $`(K^1x)(\lambda Ea)`$
$`=\lambda y(Ka)+q^{1/2}x(\lambda Ea)`$
$`=\lambda \rho _0(x_2^2y_1^2)\rho _0(y_1^2y_2^2)a+q^{1/2}\rho _0(x_1^2y_2^2)\rho _0(x_1^2x_2^2(y_1^4y_1^4))a`$
$`=\rho _0(\lambda y_1^4x_2^2y_2^2+q^{1/2}y_2^2x_2^2(y_1^4y_1^4))a`$
$`=\rho _0((qy_1^4q^1y_1^4)x_2^2y_2^2)a`$
$`=\rho _0(x_1^2x_2^2(y_1^4y_1^4))\rho _0(x_1^2y_2^2)a`$
$`=\lambda E(xa).`$
The other verifications are carried out in a similar manner. Thus we have shown that $`𝒜(_q^{++})`$ is a $`𝒰_q(gl_2)`$-left module algebra.
Finally, we turn to condition (4). Because $`𝒜(_q^{++})`$ is a left $`𝒰_q(gl_2)`$-module algebra, it is enough to prove (4) for the generators $`f=\lambda E,\lambda F,K_1,K_2`$. We verify (4) for $`f=\lambda E`$ and $`z=a𝔄(^2)`$. Since $`S(\lambda E)^{}=\lambda E`$, we obtain
$$\begin{array}{cc}\hfill (\lambda Ea)^{}(x_1,x_2)& =e^{2\pi (\beta x_2\alpha x_1)}(\overline{a(x_1+\beta i,x_2)}\overline{a(x_1\beta i,x_2)})\hfill \\ & =e^{2\pi (\beta x_2\alpha x_1)}(\overline{a}(x_1\beta i,x_2)\overline{a}(x_1+\beta i,x_2))\hfill \\ & =(S(\lambda E)^{}\overline{a})(x_1,x_2).\hfill \end{array}$$
By similar computations we check that condition (4) is satisfied for $`f=\lambda F,K_1,K_2`$. Hence it follows that (4) holds for $`f𝒰_q(gl_2())`$ and $`z=a𝔄(^2)`$. Since $`𝒪(_q^2)`$ is a left $`𝒰_q(gl_2())`$-module $``$-algebra, (4) holds also for $`z=a𝒪(_q^2)`$. Thus it remains to show that (4) is satisfied for products $`z=wa`$ and $`z=aw`$, where $`a𝔄(^2)`$ and $`w𝒪(_q^2)`$. We carry out this for $`z=wa`$ and compute
$`S(f)^{}(wa)^{}`$ $`=S(f_{(2)}^{}a^{})S(f_{(1)}^{}w^{})`$
$`=(f_{(2)}a)^{}(f_{(1)}w)^{}=((f_{(1)}w)(f_{(2})a))^{}=(f(wa))^{}.`$
This completes the proof of Theorem 21. $`\mathrm{}`$
We close this subsection by collecting some additional useful relations. From the above formulas and the defining relations of the algebra $`_q`$ we derive the following cross commutation relations for elements $`\psi (E),\psi (F),\psi (K)`$ and $`\psi (x),\psi (y)`$ in the algebra $`_q`$:
$`\psi (E)\psi (x)q^{1/2}\psi (x)\psi (E)=\psi (y)\psi (K),\psi (E)\psi (y)=q^{1/2}\psi (y)\psi (E),`$
$`\psi (F)\psi (x)=q^{1/2}\psi (x)\psi (F),\psi (F)\psi (y)q^{1/2}\psi (y)\psi (F)=\psi (x)\psi (K),`$
$`\psi (K)\psi (x)=q^{1/2}\psi (x)\psi (K),\psi (K)\psi (y)=q^{1/2}\psi (y)\psi (K).`$
The images of the algebras $`𝒰_q(sl_2())`$ and $`𝒪(_q^2)`$ under the algebra homomorphism $`\psi `$ do not generate the whole algebra $`_q`$. But they are large enough such that fourth powers of the generators $`x_1,x_2,y_1^1,y_2`$ can be expressed as
$`y_1^4`$ $`=q^1+q^{1/2}\lambda \psi (EK^1)\psi (x)\psi (y)^1,y_2^4=qq^{1/2}\lambda \psi (FK^1)\psi (y)\psi (x)^1,`$
$`x_1^4`$ $`=\psi (x^4)(qq^{1/2}\lambda \psi (FK^1)\psi (y)\psi (x)^1),`$
$`x_2^4`$ $`=\psi (y^4)(q^1+q^{1/2}\lambda \psi (EK^1\psi (x)\psi (y)^1).`$
4.4. Covariant differential calculus on the quantum quarter plane
In this subsection we extend the differential calculus $`\mathrm{\Gamma }=\mathrm{\Gamma }_{}`$ of $`𝒪(_q^2)`$ to the larger algebra $`𝒜(_q^{++})`$. As in 3.2, we use the approach developed in 2.2, but now with the algebra $`𝒵=𝒜(_q^{++})`$. We briefly repeat the construction from 2.2. Let $`V`$ be a two-dimensional vector space with basis $`\{e_1,e_2\}`$. The vector space $`\mathrm{\Omega }=𝒜(_q^{++})V`$ becomes a $`𝒜(_q^{++})`$-bimodule with bimodule structure defined by (94). The differentiation $`d`$ is defined by the commutator with the element $`\omega =q^2x^2y^2e_1+y^2e_2`$ of the $`𝒜(_q^{++})`$-bimodule $`\mathrm{\Omega }`$, that is,
$$dz=\omega zz\omega ,z𝒜(_q^{++}).$$
It is clear that $`\mathrm{\Gamma }:=𝒵d𝒵𝒵`$ is a first order differential calculus over $`𝒵=𝒜(_q^{++})`$ with differentiation $`d`$ such that $`\{dx,dy\}`$ is a free left $`𝒵`$-module basis of $`\mathrm{\Gamma }`$.
Let us compute the partial derivatives $`_x(a)`$ and $`_y(a)`$ for $`a𝔄(^2)`$. Using formulas (112), (113) and (124) we obtain from the definition of $`d`$ that
$`da`$ $`=q^2x^2y^2e_1a+y^2e_2aq^2ax^2y^2e_1ay^2e_2`$
$`=q^2x^2y^2\rho _0(1y_1^8y_2^8)ae_1+y^2\rho _0(1y_1^8)ae_2.`$ (131)
On the other hand, by (96) and (87) we have
$$da=(q^21)(q^2x^3y^2_x(a)e_1+xy^2_x(a)e_2+x^2y^1_y(a)e_1).$$
(132)
Comparing the coefficient of $`e_1`$ and $`e_2`$ in (S0.Ex172) and (132) we derive
$$_x(a)=(1q^2)^1x^1\rho _0(1y_1^8)a,_y(a)=(1q^2)^1y^1\rho _0(y_1^8(1y_2^8))a$$
or equivalently
$`_x(a)=(1q^2)^1e^{2\pi \alpha x_1}(a(x_1,x_2\frac{\alpha }{2}i)a(x_12\beta i,x_2\frac{\alpha }{2}i)),`$
$`_y(a)=(1q^2)^1e^{2\pi \beta x_2}(a(x_1\frac{3}{2}\beta i,x_2)a(x_1\frac{3}{2}\beta i,x_22\alpha i)).`$
In terms of the $``$-representation $`\rho _0`$ defined by (124) and the actions of generators $`E,F,K_1,K_2,x,y`$ given by (109)–(111), these formulas can be written as
$`_x(a)=(1q^2)^1\rho _0(x_1^2y_2^2(1y_1^8))a=q^{\frac{3}{2}}y^1EK_1^3K_2a,`$
$`_y(a)=(1q^2)^1\rho _0(x_2^2y_1^6(1y_2^8))a=q^{\frac{1}{2}}x^1FK_1^3K_2a`$
for $`a𝔄(^2)`$. In 2.2 we have shown that the two latter expressions of $`_x(a)`$ and $`_y(a)`$ hold also for elements of $`𝒪(_q^2)`$. Therefore, we have proved that
$$_x(a)=q^{3/2}y^1EK_1^3K_2a,_y(a)=q^{1/2}x^1FK_1^3K_2a$$
for all elements $`a`$ in the $``$-algebra $`𝒜(_q^{++})=𝒪(_q^2)+𝔄(^2)`$.
5. Covariant linear functionals on the quantum quarter plane
5.1 In the first subsection we construct for any $`k=(k_1,k_2)^2`$ a $`𝒰_q(gl_2())`$-covariant linear functional $`h_k`$ on $`𝔄(^2)`$ and show that it defines a scalar product $`a,b_k`$ on the $``$-algebra $`𝔄(^2)`$. This functional and the associated scalar product can be considered as $`q`$-analogs of the state given by Lebesgue measure and the $`L^2`$-scalar product on the classical quarter plane.
From the defining relations of the algebra $`𝒰_q(gl_2)`$ it is clear that there exists a unique character $`\tau `$ on $`𝒰_q(gl_2)`$ such that
$$\chi (K_1)=\chi (K_2)=q^{1/2}\mathrm{and}\chi (E)=\chi (F)=0.$$
Since the restriction of $`\chi `$ to the Hopf subalgebra $`𝒰_q(sl_2)`$ is the counit, any $`𝒰_q(sl_2)`$-covariant linear functional with respect to $`\chi `$ is $`𝒰_q(sl_2)`$-invariant.
Proposition 22. For $`k=(k_1,k_2)^2`$ and $`a𝔄(^2)`$ we define
$$h_k(a)=e^{2\pi (\alpha _kx_1+\beta _kx_2)}a(x_1,x_2)𝑑x_1𝑑x_2,$$
(133)
where
$$\alpha _k:=\alpha +2\beta ^1k_1,\beta _k:=\beta +2\alpha ^1k_2.$$
(i) The linear functional $`h_k`$ on the $`𝒰_q(gl_2())`$-module $``$-algebra $`𝔄(^2)`$ is covariant with respect to the character $`\chi `$.
(ii) For $`s,t`$ and a $`𝔄(^2)`$,
$$h_k(a(x_1+\beta s,x_2+\alpha t))=e^{2\pi \gamma (s+t)4\pi (k_1s+k_2t)}h_k(a).$$
(134)
(iii) $`h_k`$ is continuous on the Frechet space $`𝔄(^2)[\tau ]`$. More precisely, for any $`\epsilon =(\epsilon _1,\epsilon _2),\epsilon _1>0,\epsilon _2>0`$, we have
$$|h_k(a)|\frac{1}{2\pi \sqrt{\epsilon _1\epsilon _2}}S(e^{\epsilon 𝒬})e^{2\pi (\alpha _k𝒬_1+\beta _k𝒬_2)}a,a𝔄(^2).$$
Proof. (i): It suffices to verify condition (5) for the generators $`f=K_1,K_2,E,F`$. That is, we have to show that
$$h_k(K_1a)=h_k(K_2a)=q^{1/2}h_k(a)\mathrm{and}h_k(Ea)=h_k(Fa)=0$$
for $`a𝔄(_q^2)`$. By formulas (109)-(111), the latter conditions are equivalent to the relations
$`{\displaystyle e^{2\pi (\alpha _kx_1+\beta _kx_2)}}`$ $`a(x_1i\beta /2,x_2)dx_1dx_2`$
$`={\displaystyle e^{2\pi (\alpha _kx_1+\beta _kx_2)}a(x_1,x_2i\alpha /2)𝑑x_1𝑑x_2}`$
$`=e^{\pi \alpha \beta i}{\displaystyle e^{4\pi (\beta ^1k_1x_1+\alpha ^1k_2x_2)}a(x_1,x_2)𝑑x_1𝑑x_2},`$
$`{\displaystyle e^{2\pi (\alpha _kx_1+(\beta _kx_2)}}`$ $`a(x_1+\beta i,x_2)dx_1dx_2`$
$`={\displaystyle e^{4\pi (\beta ^1k_1x_1+(\alpha ^1k_2+\beta )x_2)}a(x_1\beta i,x_2)𝑑x_1𝑑x_2},`$
$`{\displaystyle e^{4\pi ((\alpha +\beta ^1k_1)x_1+\alpha ^1k_2x_2)}}`$ $`a(x_1,x_2+\alpha i)dx_1dx_2`$
$`={\displaystyle e^{4\pi ((\alpha +\beta ^1k_1)x_1+\alpha ^1k_2x_2)}a(x_1,x_2\alpha i)𝑑x_1𝑑x_2}.`$
These identities follow by the formal replacements $`(x_1,x_2)(x_1+i\beta /2,x_2)`$, $`(x_1,x_2)(x_1,x_2i\alpha /2)`$ and $`(x_1,x_2)(x_1,x_22\alpha i)`$, respectively. Similarly as above, these substitutions are justified by integrating in the complex plane and using the asymptotic estimate of Lemma 9.
(ii): The formula follows by the substitution $`(x_1,x_2)(x_1\beta s,x_2\alpha t)`$.
(iii) follows from (133) and the Cauchy-Schwarz inequality. $`\mathrm{}`$
Let $`,_k`$ be the sesquilinear form defined by means of the functional $`h_k`$ on the $``$-algebra $`𝔄(^2)`$, that is,
$$a,b_k=h_k(b^{}\mathrm{\#}a),a,b𝔄(^2).$$
(135)
Recall that $`(,)`$ denotes the scalar product of the Hilbert space $`L^2(^2)`$. For $`k^2`$, we abbreviate
$$T_k:=e^{\pi (\alpha _k𝒬_1\frac{\beta _k}{2}𝒫_1)}e^{\pi (\beta _k𝒬_2+\frac{\alpha _k}{2}𝒫_2)}e^{\pi \alpha _k𝒬_1}e^{{\scriptscriptstyle \frac{\pi }{2}}\beta _k𝒫_1}e^{\pi \beta _k𝒬_2}e^{{\scriptscriptstyle \frac{\pi }{2}}\alpha _k𝒫_2}.$$
(136)
Proposition 23. The sesquilinear form $`,_k`$ is a scalar product and $`T_k`$ is an isometric linear isomorphism of the unitary space $`(𝔄(^2),,_k)`$ on the unitary space $`(𝔄(^2),(,))`$. For $`a,b𝔄(^2)`$, we have
$$a,b_k=e^{2\pi (\alpha _kx_1+\beta _kx_2)}a(x_1+i\beta _k/4,x_2i\alpha _k/4)\overline{b}(x_1i\beta _k/4,x_2+i\alpha _k/4)𝑑x_1𝑑x_2.$$
(137)
Proof. For $`a,b𝔄(^2)`$, we compute
$`a,b_k`$ $`={\displaystyle e^{2\pi (\alpha _kx_1+\beta _kx_2)}(b^{}\mathrm{\#}a)(x_1,x_2)𝑑x_1𝑑x_2}`$
$`={\displaystyle \left(e^{2\pi (\alpha _k𝒬_1+\beta _k𝒬_2)}b\right)^{}\mathrm{\#}\left(e^{\pi (\alpha _k𝒫_2\beta _k𝒫_1)}a\right)(x_1,x_2)𝑑x_1𝑑x_2}`$
$`={\displaystyle \left(e^{2\pi (\alpha _k𝒬_1+\beta _k𝒬_2)}b\right)^{}(x_1,x_2)\left(e^{\pi (\alpha _k𝒫_2\beta _k𝒫_1)}a\right)(x_1,x_2)𝑑x_1𝑑x_2}`$
$`=(e^{\pi (\alpha _k𝒫_2\beta _k𝒫_1)}a,e^{2\pi (\alpha _k𝒬_1+\beta _k𝒬_2)}b)`$
$`=(T_k^2a,b)=(T_ka,T_kb).`$ (138)
Here the first equality combines the definitions (133) and (135) of $`h_k`$ and $`,_k`$, respectively. The second equality follows from (52) and (54), while the third one follows from formula (60). The fourth equality is just the definition of the scalar product ($`,`$.), and the fifth and the sixth are easily derived from (16). Since the operator $`T_k`$ is a bijective linear mapping of $`𝔄(^2)`$, we conclude from formula (138) that$`,_k`$ is indeed a scalar product on the vector space $`𝔄(^2)`$. The expression in (137) is obtained from (138) by inserting the actions of the operator $`T_k`$. Further, it follows from (138) that $`T_k`$ is an isometric linear isomorphism of the unitary space $`𝔄_k:=(𝔄(^2),,_k)`$ onto the unitary space $`𝔄(,):=(𝔄(^2),(,))`$. $`\mathrm{}`$
5.2 In this subsection we investigate the left actions of the algebras $`𝒰_q(gl_2)`$ and $`𝒪(_q^2)`$ on the unitary space $`𝔄_k:=(𝔄(^2),,_k)`$. Among others, we shall transform these actions to the domain $`𝔄(^2)`$ in the Hilbert space $`L^2(^2)`$ by means of the unitary operator $`T_k`$.
Recall from 4.2 that the map $`\rho _0\psi `$ defines left actions of the algebras $`𝒰_q(gl_2)`$ and $`𝒪(_q^2)`$ on $`𝔄(^2)`$. For the generators $`E,F,K_1,K_2`$ of $`𝒰_q(gl_2)`$ this action has been also given by formulas (109)–(111), see also (125). For the algebra $`𝒪(_q^2)`$ the action $`\rho _0\psi `$ is just the left multiplication in the larger algebra $`𝒜(_q^{++})`$, see Lemma 18 and formulas (127) and (128). Let $`\psi _k`$ denote the action $`\rho _0\psi `$ of $`𝒰_q(gl_2)`$ and $`𝒪(_q^2)`$ considered as representation on the unitary space $`𝔄_k`$. Since $`T_k`$ is a unitary transformation of $`𝔄_k(𝔄(^2),,_k)`$ on $`(𝔄(^2),(,))`$ by Proposition 23, $`\psi _k`$ is unitarily equivalent to the representation
$$\mathrm{\Psi }_k():=T_k\psi _k()T_k^1$$
(139)
on the domain $`𝔄(^2)`$ in the Hilbert space $`L^2(^2)`$. Further, the compositions $`\mathrm{\Phi }:=\rho _0\phi `$ of the $``$-homomorphisms $`\phi `$ (defined in Lemma 20) and the $``$-representation $`\rho _0`$ of $`_q`$ (defined by (124)) are also $``$-representations of the $``$-algebras $`𝒰_q^{tw}(gl_2())`$ and $`𝒪(_q^2)`$, respectively, on the domain $`𝔄(^2)`$ in $`L^2(^2)`$.
Let $`T`$ denote the operator $`T_k`$ defined by (136) for $`k=(\mathrm{0,0})`$, that is,
$$T=e^{\pi \alpha 𝒬_1}e^{{\scriptscriptstyle \frac{\pi }{2}}\beta 𝒫_1}e^{\pi \beta 𝒬_2}e^{{\scriptscriptstyle \frac{\pi }{2}}\alpha 𝒫_2}=e^{\pi \alpha 𝒬_1{\scriptscriptstyle \frac{\pi }{2}}\beta 𝒫_1}e^{\pi \beta 𝒬_2+{\scriptscriptstyle \frac{\pi }{2}}\alpha 𝒫_2}.$$
(140)
Using the operator
$`\begin{array}{cc}\hfill C_k& :=e^{2\pi k_1\beta ^1𝒬_1}e^{\pi k_2\alpha ^1𝒫_1}e^{2\pi k_2\alpha ^1𝒬_2}e^{\pi k_1\beta ^1𝒫_2}\hfill \\ & =e^{\pi (2k_1\beta ^1𝒬_1k_2\alpha ^1𝒫_1)}e^{\pi (2k_2\alpha ^1𝒬_2+k_1\beta ^1𝒫_2)}\hfill \end{array}`$ (141)
acting on the Hilbert space $`L^2(^2)`$, we can write the operator $`T_k`$, as
$$T_k=i^{k_1k_2}C_kT.$$
(142)
Comparing formulas (140) and (124) we see that $`T=\rho _0(x_1x_2y_1^1y_2)`$. Therefore, by Lemma 13, we get
$`\mathrm{\Psi }_0(z)=T\psi _0(z)T_0^1=\rho _0((x_1x_2y_1^1y_2)\psi (z)(x_1x_2y_1^1y_2)^1)=\rho _0\phi (z)=\mathrm{\Phi }(z)`$
for $`z𝒰_q(gl_2)`$ and $`z𝒪(_q^2)`$. That is, we have $`\mathrm{\Psi }_0=\mathrm{\Phi }`$ for both $``$-algebras $`𝒰_q^{tw}(gl_2())`$ and $`𝒪(_q^2)`$.
Next we relate the representations $`\mathrm{\Psi }_k`$ and $`\mathrm{\Psi }_0=\mathrm{\Phi }`$. Using the definitions of the operator $`C_k`$ and $`\mathrm{\Phi }(f)=\rho _0\phi (f),f=E^{},F^{},K_1,K_2,x,y,`$ we compute
$`C_k\mathrm{\Phi }(f)C_k^1`$ $`=(1)^{k_1+k_2}\mathrm{\Phi }(f)\mathrm{for}f=E^{},F^{},K,`$ (143)
$`C_k\mathrm{\Phi }(K_j)C_k^1`$ $`=(1)^{k_j}\mathrm{\Phi }(K_j)\mathrm{for}j=\mathrm{1,2},`$ (144)
$`C_k\mathrm{\Phi }(x)C_k^1`$ $`=\mathrm{\Phi }(x),C_k\mathrm{\Phi }(y)C_k^1=\mathrm{\Phi }(y)`$ (145)
Because of the formulas (142) and (139) we therefore have
$`\mathrm{\Psi }_k(f)`$ $`=(1)^{k_1+k_2}\mathrm{\Phi }(f)\mathrm{for}f=E^{},F^{},K,`$ (146)
$`\mathrm{\Psi }_k(K_j)`$ $`=(1)^{k_j}\mathrm{\Phi }(K_j)\mathrm{for}j=\mathrm{1,2},`$ (147)
$`\mathrm{\Psi }_k(x)`$ $`=\mathrm{\Phi }(x),\mathrm{\Psi }_k(y)=\mathrm{\Phi }(y).`$ (148)
In order to complete the picture we collect the formulas for the operators $`\mathrm{\Phi }(f)`$, where $`f=E^{},F^{},K_1,K_2,x,y`$. Recall from (20) and (21) that $`L_\alpha `$ denotes the operator $`L_\alpha =\overline{f_\alpha }(𝒫)e^{2\pi \alpha Q}`$, where $`f_\alpha (x)=2\mathrm{sinh}\pi \beta (2x+\alpha i)`$. Combining (S0.Ex135)-(123) and (124) we obtain
$`\mathrm{\Phi }(E^{})`$ $`=L_\alpha e^{2\pi \beta 𝒬_2},\mathrm{\Phi }(F^{})=e^{2\pi \alpha 𝒬_1}L_\beta ,`$ (149)
$`\mathrm{\Phi }(q^{1/4}K_1)`$ $`=e^{\pi \beta 𝒫_1}I,\mathrm{\Phi }(q^{1/4}K_2)=Ie^{\pi \alpha 𝒫_2},\mathrm{\Phi }(K)=e^{\pi \beta 𝒫_1}e^{\pi \alpha 𝒫_2},`$ (150)
$`\mathrm{\Phi }(x)`$ $`=e^{2\pi \alpha 𝒬_1}e^{\pi \alpha 𝒫_2},\mathrm{\Phi }(y)=e^{\pi \beta 𝒫_1}e^{2\pi \beta 𝒬_2}.`$ (151)
Let us briefly discuss the outcome of these considerations. Since the functional $`h_k`$ on the left $`𝒰_q(gl_2())`$-module $``$-algebra $`𝔄(^2)`$ is covariant with respect to $`\chi `$, it follows from Lemma 2,(i)$``$(ii), and the definition of the involution of $`𝒰_q^{tw}(gl_2())`$ that $`\psi _k=\rho \psi `$ is a $``$-representation of the $``$-algebra $`𝒰_q^{tw}(gl_2())`$ on the unitary space $`𝔄_k=(𝔄(^2),,_k)`$. The $``$-representation $`\psi _k`$ is unitarily equivalent to the $``$-representation $`\mathrm{\Psi }_k`$ of $`𝒰_q^{tw}(gl_2())`$ on the domain $`𝔄(^2)`$ in the Hilbert space $`L^2(^2)`$. The actions of the operators $`\mathrm{\Psi }_k(f)`$ for the generators $`f=E^{},F^{},K_1,K_2`$ are explicitly given by the formulas (146)–(147) and (149)–(150). Note that the dependence of the operators $`\mathrm{\Psi }_k(f)`$ on $`k^2`$ appears only in the signs in (146)-(147). In particular, if $`k_1`$ and $`k_2`$ are both even, then the $``$-representation $`\mathrm{\Psi }_k`$ of $`𝒰_q^{tw}(gl_2())`$ on $`𝔄_k`$ is unitarily equivalent to the fixed $``$-representation $`\mathrm{\Phi }`$ on the domain $`𝔄(^2)`$ in the Hilbert space $`L^2(^2)`$. From (148) and (151) we see that $`\mathrm{\Psi }_k`$ is a $``$-representation of $`𝒪(_q^2)`$ on the unitary space $`𝔄_k`$. For any $`k^2`$, the $``$-representation $`\mathrm{\Psi }_k`$ of $`𝒪(_q^2)`$ on $`𝔄_k`$ is unitarily equivalent to the $``$-representation $`\mathrm{\Phi }`$ of $`𝒪(_q^2)`$ on the domain $`𝔄(^2)`$ in $`L^2(^2)`$.
Remark 3. The preceding derivation shows the reason for the non-uniqueness of covariant functionals on the left $`𝒰_q(gl_2())`$-module $``$-algebra $`𝔄(^2)`$ from the technical side: For even numbers $`k_1`$ and $`k_2`$ the unbounded positive self-adjoint operator $`C_k`$ commutes with all representation operator $`\mathrm{\Phi }(z)`$, $`z𝒰_q^{tw}(gl_2())`$, so that the unbounded commutant of $`\mathrm{\Phi }(𝒰_q(gl_2))`$ is non-trivial. However, it can be shown that the $``$-representation $`\mathrm{\Phi }`$ of $`𝒰_q^{tw}(gl_2())`$ on $`L^2(^2)`$ is irreducibel. Hence the bounded commutant (more precisely, the strong bounded commutant, see \[S1\]) is trivial. The existence of examples of that kind is a well-known phenomena for unbounded operator algebras.
By the preceding we have expressed the actions $`\psi _k`$ of the algebras $`𝒰_q(gl_2)`$ and $`𝒪(_q^2)`$ on the unitary space $`𝔄_k=(𝔄(^2),,_k)`$ by means of the $``$-representations $`\mathrm{\Phi }`$ on the domain $`𝔄(^2)`$ in $`L^2(^2)`$. Since the representation $`\mathrm{\Phi }=T\psi _0T^1`$ is obtained from $`\psi _0`$ by the unitary operator $`T`$, it is natural to transform also the structure of the $``$-algebra $`𝒜(_q^{++})`$ and the covariant functional $`h:=h_0`$ under the bijective linear mapping $`T`$ of $`𝔄(^2)`$. That is, for $`f,g𝒜(_q^{++})`$ and $`a𝔄(^2)`$ we define
$$f\mathrm{}g=T(T^1fT^1g),f^{}=T(T^1(f)^{})\mathrm{and}\stackrel{~}{h}(a)=h(T^1a).$$
(152)
Since $`𝒜(_q^{++})`$ is a $``$-algebra with product $``$ and involution $`ff^{}`$, the vector space $`𝒜(_q^{++})`$ is a $``$-algebra, denoted $`\stackrel{~}{𝒜}(_q^{++})`$, with product $`\mathrm{}`$ and involution $`ff^{}`$. Further, since $`𝒜(_q^{++})`$ is a left $`𝒰_q(gl_2())`$-module $``$-algebra (by Theorem 21) and $`h`$ is covariant with respect to the left action $`\psi _0`$ (by Proposition 22), it is clear that $`\stackrel{~}{𝒜}(_q^{++})`$ is a left $`𝒰_q(gl_2())`$-module $``$-algebra and the linear functional $`\stackrel{~}{h}`$ is covariant with respect to the left action $`\mathrm{\Phi }=T\psi _0T^1`$.
Let us make the transformed structures more explicit. Suppose that $`f=z+a`$ and $`g=w+b`$, where $`z,w𝒪(_q^2)`$ and $`a,b𝔄(^2)`$. From (112), (113) and (140) it follows that $`xTc=Txc`$ and $`yTc=Tyc`$ and so $`vTc=Tvc`$ for all $`v𝒪(_q^2)`$ and $`c𝔄(^2)`$ . Hence we get
$$(z+a)\mathrm{}(w+b)=zw+zb+aw+a\mathrm{}b\mathrm{and}(z+a)^{}=z^{}+a^{}.$$
That is, product and involution of the $``$-algebra $`𝒪(_q^2)`$ remain unchanged and also the products of elements of $`𝒪(_q^2)`$ and $`𝔄(^2)`$. It remains to describe the transformed product and involution of the $``$-subalgebra $`𝔄(^2)`$. From the definition (140) of the operator $`T`$ and the formulas (52)–(59) we obtain
$`a\mathrm{}b`$ $`=a\mathrm{\#}(e^{{\scriptscriptstyle \frac{\pi }{2}}\beta 𝒫_1}e^{\pi \alpha 𝒬_1}e^{{\scriptscriptstyle \frac{\pi }{2}}\alpha 𝒫_2}e^{\pi \beta 𝒬_2}b)=(e^{{\scriptscriptstyle \frac{\pi }{2}}\beta 𝒫_1}e^{\pi \alpha 𝒬_1}e^{{\scriptscriptstyle \frac{\pi }{2}}\alpha 𝒫_2}e^{\pi \beta 𝒬_2}a)\mathrm{\#}b`$
$`=(e^{{\scriptscriptstyle \frac{\pi }{4}}\beta 𝒫_1}e^{{\scriptscriptstyle \frac{\pi }{2}}\alpha 𝒬_1}e^{{\scriptscriptstyle \frac{\pi }{4}}\alpha 𝒫_2}e^{{\scriptscriptstyle \frac{\pi }{2}}\beta 𝒬_2})a\mathrm{\#}(e^{{\scriptscriptstyle \frac{\pi }{4}}\beta 𝒫_1}e^{{\scriptscriptstyle \frac{\pi }{2}}\alpha 𝒬_1}e^{{\scriptscriptstyle \frac{\pi }{4}}\alpha 𝒫_2}e^{{\scriptscriptstyle \frac{\pi }{2}}\beta 𝒬_2})b,`$ (153)
$`a^{}(x_1,x_2)=(e^{\pi (\beta 𝒫_1\alpha 𝒫_2)}a)(x_1,x_2)=\overline{a}(x_1\frac{\beta }{2}i,x_2+\frac{\alpha }{2}i).`$ (154)
The vector space $`𝔄(^2)`$ equipped with this transformed $``$-algebra structure is denoted by $`\stackrel{~}{𝔄}(^2)`$. Inserting the definition (133) of the functional $`h=h_0`$ and substituting $`(x_1,x_2)(x_1+\frac{\beta }{4}i,x_2\frac{\alpha }{4}i)`$ we get
$$\stackrel{~}{h}(a)=e^{\pi (\alpha x_1+\beta x_2)}a(x_1,x_2)𝑑x_1𝑑x_2,a𝔄(^2).$$
(155)
Summarizing, we have transformed all structures obtained so far of the left $`𝒰_q(gl_2())`$-module $``$-algebra $`𝒜(_q^{++})`$ of “functions on the quantum quarter plane” under the the mapping $`T`$. The main advantage of this new picture is that the scalar product
$$(a,b)_{\stackrel{~}{h}}:=\stackrel{~}{h}(b^{}\mathrm{}a),a,b𝔄(^2),$$
derived from the transformed covariant functional $`\stackrel{~}{h}`$ is just the $`L^2`$-scalar product $`(,)`$ on $`^2`$. (This follows at once from the construction. It can be also verified directly by using formulas (S0.Ex196) and (155).) The latter fact will be crucial when we are looking for self-adjoint extensions of the representation operators $`\mathrm{\Phi }(E^{})`$ and $`\mathrm{\Phi }(F^{})`$ in a larger Hilbert space.
5.3 In this last subsection of Section 5 we prove two uniqueness result for covariant linear functionals. They say that under some technical assumptions equation (134) essentially characterizes the covariant functional $`h_k`$.
Proposition 24. Let $`k=(k_1,k_2)^2`$ and let $`h`$ be a faithful positive linear functional on the $``$-algebra $`𝔄_{pex}(^2)`$ which satisfies relation (134) for all $`s,t`$ and $`a𝔄_{pex}(^2)`$. Suppose that $`h`$ is continuous relative to the topology $`\tau `$. Then there exists a positive number $`\nu `$ such that $`h(a)=\nu h_k(a)`$ for all $`a𝔄_{pex}(^2)`$.
Proof. The proof mimics some of the preceding considerations in reversed order. Since $`h`$ is a faithful positive linear functional and $`T_k`$ is a bijective linear mapping of $`𝔄_{pex}(^2)`$, the equation
$$(a,b)_h=h(T_k^1b)^{}\mathrm{\#}(T_k^1a)),a,b𝔄_{pex}(^2),$$
(156)
defines a scalar product on $`𝔄_{pex}(^2)`$. (In the case when $`h=h_k`$ this scalar product is just the $`L^2`$-scalar product.) Consider the following four one-parameter groups of operators acting on the unitary space $`(𝔄_{pex}(^2),(,)_h)`$:
$`(U_1(t)f)(x_1,x_2)=e^{2\pi it\alpha x_1}f(x_1,x_2),(U_2(t)f)(x_1,x_2)=e^{2\pi it\beta x_2}f(x_1,x_2),`$
$`(V_1(t)f)(x_1,x_2)=f(x_1+\beta t,x_2),(V_2(t)f)(x_1,x_2)=f(x_1,x_2+\alpha t).`$
It is straightforward to verify the commutation relations
$`T_k^1U_1(t)=e^{\frac{\pi }{2}t\gamma +\pi tk_2}U_1(t)T_k^1,T_k^1U_2(t)=e^{\frac{\pi }{2}t\gamma \pi tk_1}U_2(t)T_k^1,`$
$`T_k^1V_1(t)=e^{\pi t\gamma +2\pi tk_1}V_1(t)T_k^1,T_k^1V_2(t)=e^{\pi t\gamma +2\pi tk_2}V_2(t)T_k^1`$
for arbitrary $`t`$. From these commutation rules and the assumption (134) we derive that the operators $`U_1(t),U_2(t),V_1(t),V_2(t)`$ preserve the scalar product $`(,)_h`$ for real $`t`$. Let us carry out this computation for $`U_1(t)`$. Using the definition of the scalar product $`(,)_h`$, the above relation for $`T_k^1U_1(t)`$, formula (136) and finally the assumption (134), we obtain
$`(U_1(t)a,U_1(t)b)_h`$ $`=h((T_k^1U_1(t)b)^{}\mathrm{\#}(T^1U(t)a))`$
$`=e^{\pi t\gamma +2\pi tk_2}h((e^{2\pi it\alpha 𝒬_1}T_k^1b)^{}\mathrm{\#}(e^{2\pi it\alpha 𝒬_1}T_k^1a))`$
$`=e^{\pi t\gamma +2\pi tk_2}h((e^{2\pi it\alpha 𝒬_1}(T_k^1b)^{})\mathrm{\#}(e^{2\pi it\alpha 𝒬_1}T_k^1a))`$
$`=e^{\pi t\gamma +2\pi tk_2}h(e^{\pi it\alpha 𝒫_2}(T_k^1b)^{}\mathrm{\#}T_k^1a))`$
$`=e^{\pi t\gamma +2\pi tk_2}e^{\pi t\gamma 2\pi tk_2}h((T_k^1b)^{}\mathrm{\#}(T_k^1a))`$
$`=(a,b)_h.`$ (157)
The corresponding proofs for the other operators are similar. The continuous extensions of the unitary operators $`U_j(t)`$ and $`V_j(t)`$ to the Hilbert space completion $`𝒢`$ of the unitary space $`(𝔄_{pex},(,)_h)`$ are denoted by the same symbols.
Before we continue let us note an estimate for the Hilbert space norm $`_h`$. Since the functional $`h`$ is $`\tau `$-continuous by assumption, there is a $`\tau `$-continuous seminorm $`𝔯`$ on $`𝔄_{pex}(^2)`$ such that
$$a_h^2=h((T_k^1b)^{}\mathrm{\#}(T_k^1a))𝔯((T_k^1b)^{}\mathrm{\#}(T_k^1a)),a𝔄_{pex}(^2).$$
By the definition of the topology $`\tau `$, we can take $`𝔯`$ to be a sum of norms
$$_{c,d}:=e^{2\pi (c_1𝒫_1+c_2𝒫_2)}e^{2\pi (d_1𝒬_1+d_2𝒬_2)}.,$$
(158)
where $`c=(c_1,c_2),d=(d_1,d_2)^2`$. Using the formulas (52)– (59) we conclude that the Hilbert space norm $`_h`$ is $`\tau `$-continuous on $`𝔄_{pex}(^2)`$.
Now let $`s`$ and let $`a𝔄_{pex}(^2)`$. For $`c,d^2`$, we then have
$`(U_1(t)Is2\pi i\alpha 𝒬_1)a_{c,d}`$
$`=`$ $`e^{2\pi (c_1𝒫_1+c_2𝒫_2)}e^{2\pi (d_1𝒬_1+d_2𝒬_2)}(U_1(t)Is2\pi i\alpha 𝒬_1)a`$
$`=`$ $`(e^{2\pi t\alpha c_1}U_1(t)Is2\pi i\alpha 𝒬_1s2\pi \alpha t)e^{2\pi (c_1𝒫_1+c_2𝒫_2)}e^{2\pi (d_1𝒬_1+d_2𝒬_2)}a`$
Since the norm $`_h`$ can be estimated by sums of norms $`_{c,d},c,d^2`$, it follows from the preceding equality, applied with $`s=0`$, that $`U_1(t)aa`$ in the Hilbert space $`𝒢`$ as $`t0`$. Therefore, the one-parameter unitary group $`tU_1(t)`$ on the Hilbert space $`𝒢`$ is strongly continuous. Next we set $`s=t`$. Then we conclude from the preceding that $`t^1(U_1(t)I)a2\pi i\alpha 𝒬_1`$ in the Hilbert space $`𝒢`$ as $`t0`$. (These facts are obvious in case of the Hilbert space $`L^2(^2)`$, but we do not yet know that $`(,)_h`$ is a multiple of the $`L^2`$-scalar product $`(,)`$.) The corresponding assertions for the other unitary groups follow by a similar reasoning.
It is obvious that the operators $`U_j(t)`$ and $`V_j(t)`$ satisfy the commutation relations
$`V_j(t)U_j(s)=e^{2\pi i\gamma ts}U_j(s)V_j(t),V_1(t)U_2(s)=U_2(s)V_1(t),V_2(t)U_1(s)=U_1(s)V_2(t),`$
$`U_1(t)U_2(s)=U_2(s)U_1(t),V_1(t)V_2(s)=V_2(s)V_1(t)`$
for $`j=\mathrm{1,2}`$ and $`s,t`$. That is, the unitary one-parameter groups $`U_j(t)`$ and $`V_j(t),j=\mathrm{1,2}`$ and $`t`$, fulfill the Weyl relation. Therefore, by the Stone-von Neumann uniqueness theorem \[Pu\] that there exist a Hilbert space $`𝒦`$ and a unitary transformation $`𝒥`$ of $`𝒢`$ on the Hilbert space $`\stackrel{~}{}:=L^2(^2)𝒦`$ such that
$$𝒥V_j(t)U_l(s)𝒥^1=\stackrel{~}{V}_j(t)\stackrel{~}{U}_l(s)\mathrm{for}j,l=\mathrm{1,2},s,t,$$
(159)
where $`\stackrel{~}{V}_j(t)`$ and $`\stackrel{~}{U}_l(s)`$ denote the unitary groups on the Hilbert space $`\stackrel{~}{}`$ defined by the same formulas as $`V_j(t)`$ and $`U_l(s)`$, respectively.
For $`\delta =(\delta _1,\delta _2)_{++}^2`$ and $`c=(c_1,c_2)^2`$, let $`e_{c,\delta }`$ denote the function $`e^{2\pi (c_1x_1+c_2x_2\delta _1x_1^2\delta _2x_2^2)}`$ on $`^2`$. Let
$$A_\delta :=(𝒫_1+2i\delta _1𝒬_1)(𝒫_12i\delta _1𝒬_1)+(𝒫_2+2i\delta _2𝒬_2)(𝒫_22i\delta _2𝒬_2)$$
be the operator on unitary space $`(𝔄_{pex}(^2),(,)_h)`$ and let $`\stackrel{~}{A}_\delta `$ denote the operator on $`\stackrel{~}{}`$ given by the same formula. Since $`(𝒫_j2i\delta _j𝒬_j)e_{0,\delta }=0`$ for $`j=\mathrm{1,2}`$, it is clear that $`A_\delta e_{0,\delta }=0`$. From (158) it follows that the unitary transformation $`𝒥`$ intertwines the generators of the unitary groups $`V_j(t),U_l(s)`$ and $`\stackrel{~}{V}_j(t),\stackrel{~}{U}_l(s)`$, respectively. Hence we also have $`\stackrel{~}{A}_\delta 𝒥(e_{0,\delta })=0`$. It is not diffult to show that $`\mathrm{ker}\stackrel{~}{A}_\delta =e_{0,\delta }𝒦`$. Hence we conclude that for arbitrary $`\delta _{++}^2`$ there exist a vector $`x_\delta 𝒦`$ such that $`𝒥(e_{0,\delta })=e_{0,\delta }x_\delta `$.
We next show that the vector $`x_\delta `$ does not depend on $`\delta _{++}^2`$. For $`\epsilon `$, let $`\epsilon _1:=(\epsilon \mathrm{,0})`$ and $`\epsilon _2:=(0,\epsilon )`$. It is not difficult to verify that
$$\epsilon ^1(e_{0,\delta +\epsilon _j}e_{0,\delta })2\pi 𝒬_j^2e_{0,\delta }0$$
(160)
as $`\epsilon 0`$ in the topology $`\tau `$ on $`𝔄_{pex}(^2)`$. Since the norm $`_h`$ is $`\tau `$-continuous as noted above, (160) holds also in the Hilbert space $`𝒢`$. Therefore, the image under $`𝒥`$ of the expression in (160) tends to zero in the Hilbert space $``$ as $`\epsilon 0`$. This means that
$$\left[\epsilon ^1(e_{0,\delta +\epsilon _j}e_{0,\delta })2\pi 𝒬_j^2e_{0,\delta }\right]x_\delta +e_{0,\delta +\epsilon _j}\epsilon ^1(x_{\delta +\epsilon _j}x_\delta )0$$
as $`\epsilon 0`$ in $``$. Using once more the fact that (160) holds in $`L^2(^2)`$, we conclude that $`\epsilon ^1(x_{\delta +\epsilon _j}x_\delta )0`$ in $`𝒦`$. That is, the partial derivatives of the $`𝒦`$-valued function $`\delta x_\delta `$ vanish on $`_{++}^2`$. Hence this function is constant with respect to $`\delta _{++}^2`$, say $`x_\delta =x`$.
Recall that $`𝔄_{pex}(^2)`$ is defined as the linear span of function $`x_1^nx_2^me_{c,\delta }`$, where $`n,m_0,c^2`$ and $`\delta _{++}^2`$. From (160) and the fact that $`𝒥(e_{0,\delta })=e_{0,\delta }x`$ it follows that $`𝒥(a)=ax`$ for all $`a𝔄_{pex}(^2)`$. Since $`𝒥`$ is unitary, we obtain
$$(a,b)_h=(𝒥(a),𝒥(b))_\stackrel{~}{}=(a,b)_{L(^2)}(x,x)_𝒦.$$
Thus the scalar product $`(,)_h`$ is a positive multiple of the $`L^2`$-scalar product $`(,)`$ on $`𝔄_{pex}(^2)`$, that is, $`(,)_h=\nu (,)`$, where $`\nu :=(x,x)`$.
By the latter and the definition (156) of the scalar product $`(,)_h`$, we have
$`h(b^{}\mathrm{\#}a)`$ $`=(T_ka,T_kb)_h=\nu (T_ka,T_kb)=\nu (T_k^2a,b)`$
$`=\nu (e^{2\pi \alpha _k𝒬_1\pi \beta _k𝒫_1}e^{2\pi \beta _k𝒬_2+\pi \alpha _k𝒫_2}a,b)`$
for $`a,b𝔄(^2)`$. Now we set $`b=f_\epsilon `$, where $`f_\epsilon `$ is the approximate identity from Proposition 14. Taking the limit $`\epsilon +0`$ in the topology $`\tau `$ and using (64), we obtain
$$h(a)=\nu e^{2\pi \alpha _kx_1+2\pi \beta _kx_2}a(x_1+\frac{\beta _k}{2}i,x_2\frac{\alpha _k}{2}i)𝑑x_1𝑑x_2.$$
Arguing as in the proof of Proposition 22, we substitute $`(x_1,x_2)(x_1i\beta _k/2,x_2+i\alpha _k/2)`$ in the latter formula and obtain $`h(a)=\nu h_k(a),a𝔄_{pex}(^2)`$. $`\mathrm{}`$
Theorem 25. Let $`k=(k_1,k_2)^2`$ and let $`c=(c_1,c_2),d=(d_1,d_2)^2`$ be such that $`8|c_1d_1|<1`$ and $`8|c_2d_2|<1`$. Suppose that $`h`$ is a faithful positive linear functional on the $``$-algebra $`𝔄(^2)`$ such that
$`h(a(x_1+\beta s,x_2+\alpha s))=e^{2\pi \gamma (s+t)4\pi (k_1s+k_2t)}h(a),`$
$`|h(a)|CS(e^{c𝒬})S(e^{d𝒫})a`$ (161)
for some positive constant $`C`$ and for all $`s,t`$ and $`a𝔄(^2)`$.
Then there is a positive constant $`\nu `$ such that $`h=\nu h_k`$.
Proof. By Proposition 24, we have $`h(a)=\nu h_k(a)`$ for all $`a𝔄_{pex}(^2)`$. By Lemma 10, $`𝔄_{pex}(^2)`$ is dense in $`𝔄(^2)`$ relative to the norm $`S(e^{c𝒬})S(e^{d𝒫})`$. Since $`h`$ and $`h_k`$ are both continuous with respect to this norm, it follows that $`h=\nu h_k`$ on $`𝔄(^2)`$. $`\mathrm{}`$
Remark 4. It is likely that the assertion of Theorem 25 remains valid if we assume only the $`\tau `$-continuity of the functional $`h`$ instead of inequality (S0.Ex219) with $`8|c_jd_j|<1,j=\mathrm{1,2}`$. For this it would be sufficient to know that $`𝔄_{pex}(^2)`$ is $`\tau `$-dense in $`𝔄(^2)`$.
6. The real quantum plane
6.1 In order to motivate the definitions given below, we first recall the following well-known representation of the Lie algebra $`gl_2()`$ on the $`C^{\mathrm{}}`$-functions of $`^2`$: The action of the generators $`e,f,h_1,h_2`$ of $`gl_2()`$ satisfying the relations
$$[h_1,e]=e,[h_1,f]=f,[h_2,e]=e,[h_2,f]=f,[e,f]=h_1h_2$$
is given by the formulas
$$e=y\frac{}{x},f=x\frac{}{y},h_1=x\frac{}{x},h_2=y\frac{}{y}.$$
(162)
The groups of relations (78)–(80) and (109)–(111) can be interpreted as quantum versions of the formulas (162). Likewise, we see that the linear mappings $`𝒟_x^q`$ and $`𝒟_y^q`$ can be thought as $`q`$-deformations of the partial derivatives $`\frac{}{x}`$ and $`\frac{}{y}`$, respectively. We shall not make this more explicit because we shall not need these details.
Let us explain the underlying idea of our construction of the real quantum plane in the classical situation. We consider the plane $`^2`$ as the direct sum of the four quarter planes $`^{++},^+,^+,^{}`$ that are glued together along the two coordinate axis. On the level of functions this means that a continuous function $`f`$ on $`^2\backslash \{(\mathrm{0,0})\}`$ is given by a 4-tuple $`(f_{++},f_+,f_+,f_{})`$ of continuous functions on the quarter planes satisfying the boundary conditions
$`f_{++}(+0,y)=f_+(0,y),f_+(+0,y)=f_{}(0,y),`$
$`f_{++}(x,+0)=f_+(x,0),f_+(x,+0)=f_{}(x,0).`$
We now turn to the quantum case and consider the direct sum
$$𝔅(^2)_4:=\stackrel{~}{𝔄}(^2)\stackrel{~}{𝔄}(^2)\stackrel{~}{𝔄}(^2)\stackrel{~}{𝔄}(^2)$$
(163)
of the four $``$-algebras $`\stackrel{~}{𝔄}(^2)`$ corresponding to the four quarter planes. Recall that product and involution of the $``$-algebra $`\stackrel{~}{𝔄}(^2)`$ are given by formulas (S0.Ex196) and (154). We could have also taken here the $``$-algebra $`𝔄(^2)`$. The reason why we prefer to work with $`\stackrel{~}{𝔄}(^2)`$ is that on direct sums of Hilbert spaces $`L^2(^2)`$ it is easier to describe self-adjoint extensions of the operators $`\mathrm{\Phi }(E^{})`$ and $`\mathrm{\Phi }(F^{})`$. Since the elements of $`\stackrel{~}{𝔄}(^2)`$ correspond to “functions on the quantum quarter plane which vanish at the boundaries”, no boundary condition occurs and we can just take the direct sum of the four $``$-algebras.
Next we want to make the direct sum $`(_q^2)=𝒪(_q^2)+𝔅(^2)_4`$ of the vector spaces $`𝒪(_q^2)`$ and $`𝔅(^2)_4`$ into a left $`𝒰_q(gl_2())`$-module $``$-algebra. This means that we have to define the products of the generators $`x`$ and $`y`$ of $`𝒪(_q^2)`$ by a 4-tuple $`a=(a_1,a_2,a_3,a_4)`$, $`a_j\stackrel{~}{𝔄}(^2)`$, and left actions of the generators $`E,F,K_1,K_2`$ of $`𝒰_q(gl_2)`$ on $`a`$. Let us look for a moment at the classical case and consider a quarter plane $`^{ϵϵ^{}}`$, where $`ϵ,ϵ^{}\{+,\}`$. From the formulas (162) we see that if we pass from $`^{++}`$ to $`^{ϵϵ^{}}`$ then $`x`$ has to be replaced by $`ϵx`$, $`y`$ by $`ϵ^{}y`$, $`e`$ by $`ϵϵ^{}e`$ and $`f`$ by $`ϵϵ^{}f`$, while $`h_1`$ and $`h_2`$ remain unchanged. This suggests to take the following definitions in the quantum case:
$`xa`$ $`=(xa_1,xa_2,xa_3,xa_4),ya=(ya_1,ya_2,ya_3,ya_4),`$
$`\mathrm{\Phi }(E^{})a`$ $`=(\mathrm{\Phi }(E^{})a_1,\mathrm{\Phi }(E^{})a_2,\mathrm{\Phi }(E^{})a_3,\mathrm{\Phi }(E^{})a_4),`$
$`\mathrm{\Phi }(F^{})a`$ $`=(\mathrm{\Phi }(F^{})a_1,\mathrm{\Phi }(F^{})a_2,\mathrm{\Phi }(F^{})a_3,\mathrm{\Phi }(F^{})a_4),`$
$`\mathrm{\Phi }(K_j)a`$ $`=(\mathrm{\Phi }(K_j)a_1,\mathrm{\Phi }(K_j)a_2,\mathrm{\Phi }(K_j)a_3,\mathrm{\Phi }(K_j)a_4),j=\mathrm{1,2},`$
where $`a=(a_1,a_2,a_3,a_4)`$, $`a_j\stackrel{~}{𝔄}(^2)`$. Reall that the action of the operators $`x\mathrm{\Phi }(x)`$, $`y\mathrm{\Phi }(y)`$ and $`\mathrm{\Phi }(f)`$, $`f=E^{},F^{},K_1,K_2`$, are described by formulas (149)–(151). Since $`\stackrel{~}{𝒜}(_q^{++})=𝒪(_q^2)+𝔄(^2)`$ is a left $`𝒰_q(gl_2())`$-module $``$-algebra with left action $`\mathrm{\Phi }`$, it is easily verified that the preceding formulas define indeed a unique $``$-algebra structure on $`(_q^2)_4`$ and that $`(_q^2)_4`$ becomes a left $`𝒰_q(gl_2())`$-module $``$-algebra in this manner.
From the preceding it is clear that the functional $`\stackrel{~}{h}`$ on $`𝔅(^2)_4`$ defined by
$$\stackrel{~}{h}(a)=\stackrel{~}{h}(a_1)+\stackrel{~}{h}(a_2)+\stackrel{~}{h}(a_3)+\stackrel{~}{h}(a_4),a=(a_1,a_2,a_3,a_4)𝔅(^2)_4,$$
(164)
is $`𝒰_q(gl_2())`$-covariant with respect to $`\chi `$, where $`\stackrel{~}{h}(a_j)`$ is given by (155).
6.2 The vector space $`𝔅(^2)_4`$ is also a dense domain of the Hilbert space
$$=L^2(^2)L^2(^2)L^2(^2)L^2(^2).$$
(165)
From the formulas (150) and (151) it is clear that the operators $`x`$, $`y`$ and $`\mathrm{\Phi }(q^{1/4}K_j)`$, $`j=\mathrm{1,2}`$, are essentially self-adjoint on the domain $`𝔅(^2)_4`$. From (149) and Lemma 5(ii) it follows that the operators $`\mathrm{\Phi }(E^{})`$ and $`\mathrm{\Phi }(F^{})`$ are symmetric but not essentially self-adjoint on the domain $`𝔅(^2)_4`$ in $``$. (Indeed, by Lemma 5(ii) the adjoint of $`L_\alpha `$ is the operator $`R_\alpha `$ and $`R_\alpha `$ is easily seen to be a proper extension of $`L_\alpha `$.) That the operators $`\mathrm{\Phi }(E^{})`$ and $`\mathrm{\Phi }(F^{})`$ are not essentially self-adjoint on the domain $`𝔅(^2)_4`$ is not surprising, because $`𝔅(^2)_4`$ contains only “functions which vanish at the boundaries of the four quantum quarter planes”. In the classical case the corresponding symmetric operators $`ie=yi\frac{}{x}`$ and $`if=xi\frac{}{y}`$ are also not essentially self-adjoint when the functions in the domain have boundary values zero at the $`x`$\- and $`y`$-axis. Our next step is to “glue together the four quarter quantum planes” and to obtain self-adjoint extensions of the symmetric operators $`\mathrm{\Phi }(E^{})`$ and $`\mathrm{\Phi }(F^{})`$ in this manner.
In order to explain the glueing procedure let us return to the classical situation and consider the direct sum of four quarter planes. Since the corresponding functions have the boundary values zero at the $`x`$\- and $`y`$-axis, the operators $`ie=yi\frac{}{x}`$ and $`if=xi\frac{}{y}`$ are not essentially self-adjoint. The particular self-adjoint extensions of the symmetric operators $`ie`$ and $`if`$ we are interested in are determined by the boundary conditions
$$f(+0,y)=f(0,y)\mathrm{and}f(x,+0)=f(x,0).$$
(166)
For simplicity, let us first look how the two upper quarter planes $`^+`$ and $`^{++}`$ are glued together along the positive $`y`$-axis. We identify a function $`f_+`$ on $`^+`$ with the function $`f_{++}`$ on $`^{++}`$ given by $`f_{++}(x,y):=f_+(x,y)`$, $`y>0,x`$. Let $`D_x`$ denote the symmetric operator $`i\frac{}{x}`$ on $`L^2(^{++})`$ with boundary condition $`f(+0,y)=0`$, $`y>0`$. Then the operator $`T_0=i\frac{}{x}`$ on $`^+^{++}`$ with boundary condition zero at the positive $`y`$-axis acts on the Hilbert space $`L^2(^{++})L^2(^{++})`$ and has the form
$$T_0=\left(\begin{array}{cc}D_x& 0\\ 0& D_x\end{array}\right)$$
Let $`T`$ denote the self-adjoint extension of $`T_0`$ with boundary condition $`f(+0,y)=f(0,y)`$, $`y>0`$, and let $`J_0`$ be the symmetry (that is, self-adjoint unitary)
$$J_0=\frac{1}{\sqrt{2}}\left(\begin{array}{cc}1& 1\\ 1& 1\end{array}\right).$$
One easily verifies that
$$J_0TJ_0=\left(\begin{array}{cc}0& D_x\\ D_x^{}& 0\end{array}\right),$$
where $`D_x^{}`$ is the adjoint of the closed symmetric operator $`D_x`$ on the Hilbert space $`L^2(^{++})`$.
In a similar manner we proceed with the four quarter planes $`^{++}`$, $`^+`$, $`^+`$, $`^{}`$. As above, we identify functions on $`^+,^+,^{}`$ with the corresponding functions on $`^{++}`$, and let $`D_y=i\frac{}{y}`$ on $`L^2(^{++})`$ with boundary condition $`f(x,+0)=0,x>0`$. Further, let $`T_{10}=i\frac{}{x}`$ and $`T_{20}=i\frac{}{y}`$ be the symmetric operators with boundary values zero at the positive $`y`$\- resp. $`x`$-axis acting on the Hilbert space
$$𝒢=L^2(^{++})L^2(^{++})L^2(^{++})L^2(^{++}).$$
Consider the symmetries
$$J_1=\frac{1}{\sqrt{2}}\left(\begin{array}{cccc}1& 1& 0& 0\\ 1& 1& 0& 0\\ 0& 0& 1& 1\\ 0& 0& 1& 1\end{array}\right),J_2=\frac{1}{\sqrt{2}}\left(\begin{array}{cccc}1& 0& 1& 0\\ 0& 1& 0& 1\\ 1& 0& 1& 0\\ 0& 1& 0& 1\end{array}\right)$$
on the Hilbert space $`𝒢`$. Since $`J_1`$ and $`J_2`$ commute, their product
$$J:=J_1J_2=\frac{1}{2}\left(\begin{array}{cccc}1& 1& 1& 1\\ 1& 1& 1& 1\\ 1& 1& 1& 1\\ 1& 1& 1& 1\end{array}\right)$$
(167)
is again a symmetry. We shall transform all structures by means of the symmetry $`J`$.
In order to save space, let us introduce abbreviations for some operator matrices. If $`z_1`$, $`z_2`$ and $`z`$ is an operator on a Hilbert space $`𝒦_0`$, we set
$`\sigma _1(z_1,z_2)=\left(\begin{array}{cccc}z_1& 0& 0& 0\\ 0& z_2& 0& 0\\ 0& 0& z_1& 0\\ 0& 0& 0& z_2\end{array}\right),\sigma _2(z_1,z_2)=\left(\begin{array}{cccc}z_1& 0& 0& 0\\ 0& z_1& 0& 0\\ 0& 0& z_2& 0\\ 0& 0& 0& z_2\end{array}\right),`$
$`\theta _1(z)=\left(\begin{array}{cccc}0& z& 0& 0\\ z^{}& 0& 0& 0\\ 0& 0& 0& z\\ 0& 0& z^{}& 0\end{array}\right),\theta _2(z)=\left(\begin{array}{cccc}0& 0& z& 0\\ 0& 0& 0& z\\ z^{}& 0& 0& 0\\ 0& z^{}& 0& 0\end{array}\right),`$
$`\kappa _1(z)=\left(\begin{array}{cccc}0& 0& 0& z\\ 0& 0& z^{}& 0\\ 0& z& 0& 0\\ z^{}& 0& 0& 0\end{array}\right),\kappa _2(z)=\left(\begin{array}{cccc}0& 0& 0& z\\ 0& 0& z& 0\\ 0& z^{}& 0& 0\\ z^{}& 0& 0& 0\end{array}\right)`$
and $`\sigma _j(z):=\sigma _j(z,z)`$, $`j=\mathrm{1,2}`$. The matrices $`\sigma _j(z_1,z_2),\theta _j(z)`$ and $`\kappa _j(z)`$ act as operators on the Hilbert space $`𝒦_0𝒦_0K_0𝒦_0`$.
Let us now continue the above discussion. In terms of the preceding notation we have $`T_{10}=\sigma _1(D_x)`$ and $`T_{20}=\sigma _2(D_y)`$. It is not difficult to check that the self-adjoint extensions $`T_1`$ and $`T_2`$ of $`T_{10}`$ and $`T_{20}`$, respectively, with boundary condition (166) satisfy the relations
$`JT_1J=J_2\kappa _1(D_x)J_2=\theta _1(D_x),JT_2J=J_1\kappa _2(D_y)J_1=\theta _2(D_y).`$
That is, the particular self-adjoint extensions $`T_1`$ and $`T_2`$ are characterized in a simple manner by means of the symmetry transformation $`J`$.
Let $`x`$ and $`y`$ denote the multiplication operators by the coordinate functions $`x`$ and $`y`$, respectively, on $`L^2(^{++})`$. Then the multiplication operators $`M_x`$ and $`M_y`$ by the coordinate functions on $`L^2(^2)`$ can be expressed as
$$JM_xJ=J\sigma _1(x)J=\theta _1(x),JM_yJ=J\sigma _2(y)J=\theta _2(y).$$
(168)
Further, the self-adjoint operators $`i\widehat{e}=iy\frac{}{x}`$ and $`i\widehat{f}=ix\frac{}{y}`$ with boundary condition (166) can be written as
$$Ji\widehat{e}J=JM_xT_2J=\kappa _2(xD_y),Ji\widehat{f}J=JM_yT_1J=\kappa _1(yD_x).$$
(169)
Finally, we also transform the structure of the $``$-algebra of functions on $`^2`$ under the symmetry $`J`$. If we consider functions $`f,g`$ on $`^2`$ as 4-tuples of functions on the quarter planes, then the product $`fg`$ is transformed under $`J`$ as
$$fg:=J(J(f)J(g)).$$
(170)
More explicitly, for $`f=(f_1,f_2,f_3,f_4)`$ and $`g=(g_1,g_2,g_3,g_4)`$ we compute
$`fg=1/4(f_1g_1+f_2g_2+f_3g_3+f_4g_4,f_1g_2+f_2g_1+f_3g_4+f_4g_3,`$
$`f_1g_3+f_2g_4+f_3g_1+f_4g_2,f_1g_4+f_2g_3+f_3g_2+f_4g_1),`$ (171)
where $`f_jg_k`$ and $`g_kf_j`$ mean the usual pointwise products of functions on the quarter plane. Obviously, the involution of functions is invariant under $`J`$, that is, we have $`J(f^{})=J(f)^{}`$.
Thus, we have described the operators and the algebra of functions on the plane by means of the symmetry $`J`$ and the corresponding operators and algebras of the quarter planes. The advantage of this approach is that it gives a convenient algebraic form for the particular self-adjoint extensions of the first order differential operators. There is the disadvantage that the algebra product has to be changed too. The preceding formulas and considerations will serve as guiding motivation for the construction of the real quantum plane in the next subsection.
6.3 In this subsection we develop the basics of the construction of the real quantum plane. Our starting point is the description of the quantum quarter plane by means of the left $`𝒰_q(gl_2())`$-module $``$-algebra $`\stackrel{~}{𝒜}(_q^{++})`$ developed in 4.2. In what follows we assume that
$$0<|\gamma |<1/3.$$
(172)
Let us begin by defining the left action of $`𝒰_q(gl_2)`$. Remembering the formulas (169), (149) and (150), we consider the self-adjoint operators on the Hilbert space $``$ (see (165)) defined by the $`4\times 4`$-operator matrices
$`𝖤:=\kappa _1(L_\alpha e^{2\pi \beta 𝒬_2}),𝖥:=\kappa _2(e^{2\pi \alpha 𝒬_1}L_{\beta ,\alpha }),𝖪_1:=I(e^{\pi \beta 𝒫_1}),𝖪_2:=I(e^{\pi \alpha 𝒫_2}).`$
Here $`I(z)`$ denotes the diagonal matrix with diagonal entries equal to $`z`$. Note that the entries of these matrices are just the operators occuring in formulas (149)–(150). Let $`𝔄(^2)_4`$ be the domain
$$𝔄(^2)_4:=𝔄_{12}(^2)𝔄_2(^2)𝔄_1(^2)𝔄(^2)$$
(173)
in the Hilbert space $``$, where
$`𝔄_1(^2)`$ $`=f_\alpha (𝒫_1)^1𝔄(^2),𝔄_2(^2)=f_\beta (𝒫_2)^1𝔄(^2),`$
$`𝔄_{12}(^2)`$ $`=f_\alpha (𝒫_1)^1𝔄(^2)+f_\beta (𝒫_2)^1𝔄(^2).`$
It is clear that $`𝔄(^2)𝔄_1(^2)𝔄_2(^2)𝔄_{12}(^2)`$. Since $`f_\alpha (𝒫_j)𝔄(^2)𝔄(^2)`$ for $`j=\mathrm{1,2}`$, $`𝔄(^2)_4`$ contains the domain $`𝔅(_q^2)`$ defined by (163). In particular, $`𝔄(^2)_4`$ is dense in $``$. Using the fact that the operator $`L_\alpha `$ in $`L^2()`$ has the adjoint operator $`R_\alpha =e^{2\pi \alpha 𝒬}f_\alpha (𝒫)`$ by Lemma 5(ii), it is clear that $`𝔄(^2)_4`$ is contained in the domains of the operators $`𝖤,𝖥,𝖪_j`$, $`j=\mathrm{1,2}`$. Define
$$\mathrm{\Phi }(E^{})=𝖤𝔄(^2)_4,\mathrm{\Phi }(F^{})=𝖥𝔄(^2)_4,\mathrm{\Phi }(q^{1/4}K_j)=𝖪_j𝔄(^2)_4.$$
(174)
Proposition 26. (i) $`𝔄(^2)_4`$ is contained in the domains of the operator products $`\mathrm{\Phi }(E^{})\mathrm{\Phi }(F^{})`$, $`\mathrm{\Phi }(F^{})\mathrm{\Phi }(E^{})`$, $`\mathrm{\Phi }(E^{})\mathrm{\Phi }(K_j)`$, $`\mathrm{\Phi }(K_j)\mathrm{\Phi }(E^{})`$, $`\mathrm{\Phi }(F^{})\mathrm{\Phi }(K_j)`$, and $`\mathrm{\Phi }(K_j)\mathrm{\Phi }(F^{})`$ for $`j=\mathrm{1,2}`$.
(ii) The operators $`\mathrm{\Phi }(E^{}),\mathrm{\Phi }(F^{})`$ and $`\mathrm{\Phi }(K_j)`$ satisfy the defining relations in terms of the generators $`E^{},F^{}`$ and $`K_j`$ of the algebra $`𝒰_q(gl_2)`$.
(iii) The operators $`\mathrm{\Phi }(E^{}),\mathrm{\Phi }(F^{})`$ and $`\mathrm{\Phi }(q^{1/4}K_j)`$ are essentially self-adjoint on the domain $`𝔄(^2)_4`$.
Proof. (i): We show that $`𝔄(^2)_4`$ is contained in the domain of the product $`\mathrm{\Phi }(E^{})\mathrm{\Phi }(F^{})`$. By the definition of operators $`\mathrm{\Phi }(E^{})`$ and $`\mathrm{\Phi }(F^{})`$ this means that
$`𝔄_{12}(^2)𝒟((L_\alpha e^{2\pi \beta 𝒬_2})(e^{2\pi \alpha Q_1}R_\beta )),𝔄_2(^2)𝒟((R_\alpha e^{2\pi \beta 𝒬_2})(e^{2\pi \alpha 𝒬_1}R_\beta )),`$
$`𝔄_1(^2)𝒟((L_\alpha e^{2\pi \beta 𝒬_2})(e^{2\pi \alpha Q_1}L_\beta )),𝔄(^2)𝒟((R_\alpha e^{2\pi \beta 𝒬_2})(e^{2\pi \alpha 𝒬_1}L_\beta )).`$
The last inclusion for the fourth component is obvious. Let us verify the assertion for the third component of $`\mathrm{\Phi }(E^{})\mathrm{\Phi }(F^{})`$. That is, we have to show that each vector $`\eta =f_\alpha (𝒫_1)^1\zeta 𝔄_1(^2)`$, where $`\zeta 𝔄(^2)`$, belongs to the domain of the product $`(L_\alpha e^{2\pi \beta 𝒬_2})(e^{2\pi \alpha Q_1}L_\beta )`$. In order to prove this, it suffices to check that for all $`\xi 𝔄()`$ the vector $`f_\alpha (𝒫)^1\xi `$ belongs to the domain of the operator $`e^{2\pi \alpha 𝒬}`$. Applying the Fourier transform we see that the latter is equivalent to the fact that for arbitrary $`\xi 𝔄()`$ the function $`f_\alpha (x)^1\xi (x)`$ belongs to the domain of the operator $`e^{2\pi \alpha 𝒫}`$. The assumption $`0<|\gamma |<1/3`$ implies that the infimum of the holomorphic function $`f_\alpha (x)`$ on the strip $`0<\mathrm{Im}x<\alpha `$ when $`\alpha >0`$ resp. $`0>\mathrm{Im}x>\alpha `$ when $`\alpha <0`$ is positive. It follows from the characterization of the operator $`e^{2\pi \alpha 𝒫}`$ given in Lemma 4 that the function $`f_\alpha (x)^1\xi `$ belongs to the domain of this operator. The corresponding proofs for the first and the seond components of $`\mathrm{\Phi }(E^{})\mathrm{\Phi }(F^{})`$ and for the other operator products are similar.
(ii): We carry out the proof for the relation $`E^{}F^{}F^{}E^{}=\lambda (K_1^2K_2^2K_1^2K_2^2)`$ of the algebra $`𝒰_q(gl_2)`$. For the other defining relations these verifications are straigthforward and will be omitted. Computing the commutator of the operator matrices $`𝖤`$ and $`𝖥`$, we obtain a diagonal operator having the operators
$$\begin{array}{cc}\hfill A_1& :=(L_\alpha e^{2\pi \beta 𝒬_2})(e^{2\pi \alpha Q_1}R_\beta )(e^{2\pi \alpha 𝒬_1}L_\beta )(R_\alpha e^{2\pi \beta 𝒬_2}),\hfill \\ \hfill A_2& :=(R_\alpha e^{2\pi \beta 𝒬_2})(e^{2\pi \alpha 𝒬_1}R_\beta )(e^{2\pi \alpha 𝒬_1}L_\beta )(L_\alpha e^{2\pi \beta 𝒬_2}),\hfill \\ \hfill A_3& :=(L_\alpha e^{2\pi \beta 𝒬_2})(e^{2\pi \alpha Q_1}L_\beta )(e^{2\pi \alpha 𝒬_1}R_\beta )(R_\alpha e^{2\pi \beta 𝒬_2}),\hfill \\ \hfill A_4& :=(R_\alpha e^{2\pi \beta 𝒬_2})(e^{2\pi \alpha 𝒬_1}L_\beta )(e^{2\pi \alpha 𝒬_1}R_\beta )(L_\alpha e^{2\pi \beta 𝒬_2})\hfill \end{array}$$
as diagonal entries. By (i), for any vector $`\eta =(\eta _1,\eta _2,\eta _3,\eta _4)𝔄(^2)`$ the $`j`$-th component $`\eta _j`$ belongs to the domain of each of the two summands of the operator $`A_j`$, $`j=\mathrm{1,2,3,4}`$. Using this fact we compute the terms $`A_j\eta _j`$ and obtain
$$A_j\eta _j=\lambda (e^{2\pi \beta 𝒫_1}e^{2\pi \alpha 𝒫_2}e^{2\pi \beta 𝒫_1}e^{2\pi \alpha 𝒫_2})\eta _j.$$
Thus, $`\mathrm{\Phi }(E^{})\mathrm{\Phi }(F^{})\eta \mathrm{\Phi }(F^{})\mathrm{\Phi }(E^{})\eta =\lambda (\mathrm{\Phi }(K_1)^2\mathrm{\Phi }(K_2)^2\mathrm{\Phi }(K_1)^2\mathrm{\Phi }(K_2)^2)\eta `$.
(iii): By Lemma 4(ii), the domain $`𝒟_\delta `$ is a core for the self-adjoint operator $`e^{c𝒫}`$, $`c`$. Since $`𝒟_\delta 𝒟_\delta 𝔄(^2)`$, the operator $`\mathrm{\Phi }(K_j)`$ is essentially self-adjoint even on the smaller domain $`𝔅(^2)`$. For the operators $`\mathrm{\Phi }(E^{})`$ and $`\mathrm{\Phi }(F^{})`$ the assertion follows from statements (iii) and (iv) of Lemma 5. $`\mathrm{}`$
It is easily seen that the domain $`𝔄(^2)_4`$ is not invariant under the operators $`\mathrm{\Phi }(E^{})`$ and $`\mathrm{\Phi }(F^{})`$. Therefore, we do not get a $``$-representation of the whole $``$-algebra $`𝒰_q^{tw}(gl_2())`$ on the domain $`𝔄(^2)_4`$ of the Hilbert space $`𝒢`$. However, the actions of the generators $`f=E^{},F^{},K_1,K_2`$ satisfy the defining relations and they have the hermiticity properties of the $``$-algebra $`𝒰_q^{tw}(gl_2())`$.
Our next step is to define a $``$-algebra structure on $`𝔄(^2)_4`$. Recall that we used the $``$-algebra $`\stackrel{~}{𝔄}(_q^{++})`$ with product $`\mathrm{}`$ and involution $``$ (see 4.2) as $``$-algebra of functions on the quantum quarter plane. As in the case of functions on $`^2`$ we transform the componentwise product of 4-tuples under the symmetry $`J`$ (see (170) and (S0.Ex238)). That is, for $`a=(a_1,a_2,a_3,a_4),b=(b_1,b_2,b_3,b_4)𝔄(^2)_4`$ we define
$`\begin{array}{cc}\hfill ab=1/4(& a_1\mathrm{}b_1+a_2\mathrm{}b_2+a_3\mathrm{}b_3+a_4\mathrm{}b_4,a_1\mathrm{}b_2+a_2\mathrm{}b_1+a_3\mathrm{}b_4+a_4\mathrm{}b_3,\hfill \\ & a_1\mathrm{}b_4+a_2\mathrm{}b_3+a_3\mathrm{}b_2+a_4\mathrm{}b_1,a_1\mathrm{}b_3+a_2\mathrm{}b_4+a_3\mathrm{}b_1+a_4\mathrm{}b_2),\hfill \end{array}`$ (175)
$`a^{}=(`$ $`a_1^{},a_2^{},a_3^{},a_4^{}).`$ (176)
For the product $`a_j\mathrm{}b_k`$ we take the symmetrized version in formula (S0.Ex196):
$`a_j\mathrm{}b_k=(e^{{\scriptscriptstyle \frac{\pi }{4}}\beta 𝒫_1}e^{{\scriptscriptstyle \frac{\pi }{2}}\alpha 𝒬_1}e^{{\scriptscriptstyle \frac{\pi }{4}}\alpha 𝒫_2}e^{{\scriptscriptstyle \frac{\pi }{2}}\beta 𝒬_2})a_j\mathrm{\#}(e^{{\scriptscriptstyle \frac{\pi }{4}}\beta 𝒫_1}e^{{\scriptscriptstyle \frac{\pi }{2}}\alpha 𝒬_1}e^{{\scriptscriptstyle \frac{\pi }{4}}\alpha 𝒫_2}e^{{\scriptscriptstyle \frac{\pi }{2}}\beta 𝒬_2})b_k.`$ (177)
The involution $`a_j^{}`$ of a $`a_j`$ is defined as in (154) by
$$a_j^{}(x_1,x_2)=(e^{\pi (\beta 𝒫_1\alpha 𝒫_2)}\overline{a}_j)(x_1,x_2).$$
(178)
Proposition 27. The vector space $`𝔄(^2)_4`$ is a $``$-algebra with product $``$ and involution $``$ given by the formulas (175), (177), (176) and (178).
Proof. By arguing as in the proof of assertion (i) of Proposition 26 it follows from assumption (172) that the components $`a_j,b_k`$ in formula (177) are in the corresponding operator domains and in a domain $`𝒟_{\mu ,\nu }`$ for certain $`\mu ,\nu ^2`$. Thus, the product (177) is indeed well-defined. That is the reason why we have chosen the symmetrized version in (S0.Ex196) rather than the two other formulas in (S0.Ex196). Note that for $`a,b𝔄(^2)`$ all three formulas in (S0.Ex196) coincide, but we dealing now with larger classes of symbols.
Now we prove that for $`a,b𝔄(^2)_4`$ the components of the product $`a\mathrm{}b`$ belong again to the corresponding component space in (173). As a sample, we show that $`𝔄_{12}(^2)\mathrm{}𝔄_{12}(^2)𝔄_{12}(^2)`$. Let $`a,bf_\alpha (𝒫_1)^1𝔄(^2)`$. From formulas (56) and (177) we get
$$f_\alpha (𝒫_1)(a\mathrm{}b)=q^{1/2}(f_\alpha (𝒫_1)a)\mathrm{}e^{2\pi \beta 𝒫_1}b+q^{1/2}e^{2\pi \beta 𝒫_1}a\mathrm{}f_\alpha (𝒫_1)b.$$
(179)
Since $`f_\alpha (𝒫_1)a`$ and $`f_\alpha (𝒫_1)b`$ are in $`𝔄_2(^2)`$, it follows from formulas (52)–(59) that $`f_\alpha (𝒫_1)(a\mathrm{}b)`$ is in the domain of all operator $`e^{2\pi c𝒬}e^{2\pi d𝒫}`$, $`c,d^2`$. Thus, we have $`f_\alpha (c𝒫_1)(a\mathrm{}b)𝔄_2(^2)`$. If $`a,bf_\beta (𝒫_2)^1𝔄(^2)`$, then we use formula (58) rather than (56) and obtain the identity
$$f_\beta (𝒫_2)(a\mathrm{}b)=q^{1/2}(f_\beta (𝒫_2)a)\mathrm{}e^{2\pi \alpha 𝒫_2}b+q^{1/2}e^{2\pi \alpha 𝒫_2}a\mathrm{}f_\beta (𝒫_2)b$$
which implies that $`f_\beta (𝒫_2)(a\mathrm{}b)𝔄_2(^2)`$. Similar verifications can be done for the other cases and components. Thus, we have shown that $`ab𝔄(^2)_4`$ for $`a,b𝔄(^2)_4`$.
Recall that by construction the product $``$ and the involution $``$ have been transformed from the products $`\mathrm{\#}`$ and the involution $``$, respectively, under the bijective mapping $`J`$. Hence $`𝔄(^2)_4`$ is a $``$-algebra, because the products $`\mathrm{\#}`$ and the involution $``$ satisfy all axioms of a $``$-algebra. $`\mathrm{}`$
Next we define the product of elements of the coordinate algebra $`𝒪(_q^2)`$ and the algebra $`𝔄(^2)_4`$. Because of (168) and (151), we consider the self-adjoint operators $`𝗑`$ and $`𝗒`$ on the Hilbert space $`𝒢`$ given by the the operator matrices
$$𝗑=\theta _1(e^{2\pi \alpha 𝒬_1}e^{\pi \alpha 𝒫_2}),𝗒=\theta _1(e^{\pi \beta 𝒫_1}e^{2\pi \beta 𝒬_2}).$$
and define
$$xa\mathrm{\Phi }(x)a=𝗑a,ya\mathrm{\Phi }(y)a=𝗒a$$
(180)
for $`a𝔄(^2)_4`$. Using once more assumption (172) it follows that each $`a𝔄(^2)_4`$ is contained in the domains of the self-adjoint operators $`𝗑`$ and $`𝗒`$, so that (180) is well-defined.
An arbitrary element $`a𝔄(^2)_4`$ is in general not in the domains of the powers $`𝗑^n`$ and $`𝗒^n`$ for $`n`$ and the expressions $`xa`$ and $`ya`$ defined by (180) are in general not in $`𝔄(^2)_4`$. Hence the direct sum $`𝒜(_q^2)=𝒪(_q^2)+𝔄(^2)_4`$ of $``$-algebras $`𝒪(_q^2)`$ and $`𝔄(^2)_4`$ is not an algebra with respect to the product (180). However, for certain elements $`z𝒪(_q^2)`$ and $`a𝔄(^2)_4`$ the above definition (180) gives indeed a well-defined product $`za`$ and $`𝒜(_q^2)`$ becomes a partial $``$-algebra with product (180) in this manner. We shall not pursue this further.
We now replace $`𝔄(^2)_4`$ by its $``$-subalgebra
$$𝔄_0(^2)_4:=𝔄(^2)𝔄(^2)𝔄(^2)𝔄(^2).$$
According to our general picture, the elements of this subalgebra can be considered as functions on the real quantum plane which are vanishing on the coordinate axis. For $`a𝔄_0(_q^2)`$, the elements $`\mathrm{\Phi }(x)a`$ and $`\mathrm{\Phi }(y)a`$ are obviously again in $`𝔄_0(^2)_4`$, so equation (180) defines a $``$-representation $`\mathrm{\Phi }`$ of the $``$-algebra $`𝒪(_q^2)`$ on the invariant dense domain $`𝔄_0(^2)_4`$ of the Hilbert space $``$. By Lemma 4(ii), the operators $`\mathrm{\Phi }(x^n)`$ and $`\mathrm{\Phi }(y^n)`$ are essentially self-adjoint on $`𝔄_0(^2)_4`$. For $`z𝒪(_q^2)`$ and $`a𝔄_0(^2)_4`$ we define
$$za:=\mathrm{\Phi }(z)a.$$
Then the direct sum $`𝒜_0(_q^2):=𝒪(_q^2)+𝔄_0(^2)_4`$ becomes a $``$-algebra with product (180) such that $`𝒪(_q^2)`$ and $`𝔄_0(^2)_4`$ are $``$-subalgebras. Indeed, $`𝒜_0(_q^2)`$ is merely the transformation of the left $`𝒰_q(gl_2())`$-module $``$-algebra $`𝔅(^2)_4`$ defined in 6.1 (see (163)) under the symmetry $`J`$. Therefore, $`𝒜_0(_q^2)`$ is a $``$-algebra with product (180) and there is a left action of $`𝒰_q(gl_2)`$ given by formulas (174) such that $`𝒜_0(_q^2)`$ is a left $`𝒰_q(gl_2())`$-module $``$-algebra.
For the study of the coordinate functions $`x`$ and $`y`$ the $``$-algebra $`𝔄_0(^2)_4`$ is “large enough”, since the operators $`\mathrm{\Phi }(x)`$ and $`\mathrm{\Phi }(y)`$ are essentially self-adjoint on $`𝔄_0(^2)_4`$. However, for the action of the generators of $`𝒰_q(gl_2())`$ it is not, because $`𝔄_0(^2)_4`$ is not a core for the essentially self-adjoint operators $`\mathrm{\Phi }(E^{})`$ and $`\mathrm{\Phi }(F^{})`$. The larger $``$-algebra $`𝔄(^2)_4`$ is needed in this case. Since we have only an action of the generators and not of the whole algebra $`𝒰_q(gl_2))`$, $`𝔄(^2)`$ cannot be a left $`𝒰_q(gl_2())`$-module algebra. But for the generators $`f=E^{},F^{},K_1,K_2`$ the two conditions (1) and (4) are valid on the $``$-algebra $`𝔄(^2)_4`$.
The proof of the latter assertion requires a number of computations. We carry out this verification only for the generator $`f=E^{}`$ and for elements of the form $`a=(a_1\mathrm{,0,0,0}),b=(b_1\mathrm{,0,0,0})𝔄(^2)_4`$, where $`a_1,b_1f_\alpha (𝒫_1)^1𝔄(^2)`$. First we note that from formulas (52)–(55) and (177) we derive that
$`e^{2\pi \alpha 𝒬_1}(c\mathrm{}d)`$ $`=q^{1/4}(e^{2\pi \alpha 𝒬_1}c)\mathrm{}(e^{\pi \alpha 𝒫_2}d)=q^{1/4}(e^{\pi \alpha 𝒫_2}c)\mathrm{}(e^{2\pi \alpha 𝒬_1}d),`$ (181)
$`e^{2\pi \beta 𝒬_2}(c\mathrm{}d)`$ $`=q^{1/4}(e^{2\pi \beta 𝒬_2}c)\mathrm{}(e^{\pi \beta 𝒫_1}d)=q^{1/4}(e^{\pi \beta 𝒫_1}c)\mathrm{}(e^{2\pi \beta 𝒬_2}d)`$ (182)
for arbitrary elements $`c,df_\alpha (𝒫_1)^1𝔄(^2)`$. Using the definitions of the product $``$ and of the operators $`\mathrm{\Phi }(E^{})`$ and $`\mathrm{\Phi }(K)`$ and formulas (179), (181) and (182) we compute
$`(\mathrm{\Phi }(E^{})(a\mathrm{}b))_1`$ $`=R_\alpha e^{2\pi \beta 𝒬_2}(a_1\mathrm{}b_1)`$
$`=e^{2\pi \alpha 𝒬_1}e^{2\pi \beta 𝒬_2}f_\alpha (𝒫_1)(a_1\mathrm{}b_1)`$
$`=q^{1/2}e^{2\pi \alpha 𝒬_1}e^{2\pi \beta 𝒬_2}((f_\alpha (𝒫_1)a_1)\mathrm{}e^{2\pi \beta 𝒫_1}b_1)`$
$`+q^{1/2}e^{2\pi \alpha 𝒬_1}e^{2\pi \beta 𝒬_2}(e^{2\pi \beta 𝒫_1}a_1\mathrm{}f_\alpha (𝒫_1)b_1)`$
$`=(e^{2\pi \alpha 𝒬_1}e^{2\pi \beta 𝒬_2}f_\alpha (𝒫_1)a_1)\mathrm{}(e^{\pi \beta 𝒫_1}e^{\pi \alpha 𝒫_2}b_1)`$
$`+(e^{\pi \beta 𝒫_1}e^{\pi \alpha 𝒫_2}a_1)\mathrm{}(e^{2\pi \alpha 𝒬_1}e^{2\pi \beta 𝒬_2}f_\alpha (𝒫_1)b_1)`$
$`=(\mathrm{\Phi }(E^{})a)_1\mathrm{}(\mathrm{\Phi }(K)b)_1+(\mathrm{\Phi }(K^1)a)_1\mathrm{}(\mathrm{\Phi }(E^{})b)_1,`$
where the lower index $`1`$ always refers to the first component. This proves condition (1) for $`f=E^{}`$ and the particular elements $`a`$ and $`b`$. The other cases can be treated in a similar manner.
Next we want to have a counterpart on $`𝔄(^2)_4`$ of the covariant linear functional $`hh_0`$ on $`𝔄(^2)`$. Recall that this counterpart on the $``$-algebra $`(^2)_4`$ is the functional $`\stackrel{~}{h}`$ defined by (164). We have to transform this functional $`\stackrel{~}{h}`$ under the symmetry $`J`$ by setting $`h(a):=\stackrel{~}{h}(Ja)`$. By the definitions (167) of $`J`$ and (164) of $`\stackrel{~}{h}`$ we obtain $`h(a)=2\stackrel{~}{h}(a_1)`$, where $`\stackrel{~}{h}(a_1)`$ is given by (155). Inserting formula (155) we are lead to define
$$h(a)=2e^{\pi (\alpha x_1+\beta x_2)}a_1(x_1,x_2)𝑑x_1𝑑x_2,a=(a_1,a_2,a_3,a_4)𝔄(^2)_4.$$
(183)
Let us check first that $`h(a)`$ is well-defined for arbitrary elements $`a𝔄(^2)_4`$. Indeed, if $`a𝔄(^2)_4`$, then we have $`a_1f_\alpha (𝒫_1)^1𝔄(^2)+f_\beta (𝒫_2)^1𝔄(^2)`$. From assumption (172) we conclude that $`a_1`$ is in the domain $`𝒟(e^{\pi \alpha 𝒬_1}e^{\pi \beta 𝒬_2})`$ and that $`(e^{\pi \alpha 𝒬_1}e^{\pi \beta 𝒬_2})a_1𝒟_{\nu ,\mu }`$ for some $`\nu ,\mu ^{++}`$. By Lemma 9(ii), the latter implies that the function $`e^{\pi (\alpha x_1+\beta x_2)}a_1(x_1,x_2)`$ is in the Schwartz space $`𝒮(^2)`$. Hence the integral in (183) exists.
From the construction it follows that the sesquilinear form $`,_h`$ associated with $`h`$ by (8) is the scalar product of the Hilbert space $``$ (see (165)). This can be also verified directly by using formulas (175), (176), (S0.Ex196), (154) and (52)-(59). That is, we have
$$a,b_h=h(b^{}a)=_^2(a_1\overline{b_1}+a_2\overline{b_2}+a_3\overline{b_3}+a_4\overline{b_4})𝑑x_1𝑑x_2$$
(184)
for $`a=(a_1,a_2,a_3,a_4),b=(b_1,b_2,b_3,b_4)𝔄(^2)_4`$. Thus we have reduced the scalar product $`,_h`$ to $`L^2`$-scalar products on $`^2`$. To achieve formula (184) was the main aim of our considerations and it is the reason why we have transformed all structures by means of the operator $`T`$ and the symmetry $`J`$.
Finally, let us turn to the differential calculus on the quantum plane. The construction from 3.2 carries over almost verbatim to the algebra $`𝒜_0(_q^2):=𝒪(_q^2)+𝔄_0(^2)_4`$ and yields a first order differential calculus $`\mathrm{\Gamma }`$ over the algebra $`𝒜_0(_q^2)`$ such that $`\{dx,dy\}`$ is a free left module basis of $`\mathrm{\Gamma }`$. The corresponding partial derivatives $`_x`$ and $`_y`$ are of the form
$`_x(a)=\left(\begin{array}{cccc}0& _x(a_1)& 0& 0\\ _x(a_2)& 0& 0& 0\\ 0& 0& 0& _x(a_3)\\ 0& 0& _x(a_4)& 0\end{array}\right).`$
$`_y(a)=\left(\begin{array}{cccc}0& 0& _y(a_1)& 0\\ 0& 0& 0& _y(a_2)\\ _y(a_3)& 0& 0& 0\\ 0& _y(a_4)& 0& 0\end{array}\right)`$
for $`a=(a_1,a_2,a_3,a_4)𝔄_0(^2)_4`$, where $`_x(a_j)`$ and $`_y(a_j)`$, $`j=\mathrm{1,2,3,4}`$, are as in 3.2.
Recall that $`𝒜(_q^2)=𝒪(_q^2)+𝔄(^2)_4`$ is not an algebra, because not all $`a𝔄(^2)_4`$ are in the domains of $`\mathrm{\Phi }(x)^n\mathrm{\Phi }(y)^m`$. Likewise, the differential $`da=\omega aa\omega `$ and the partial derivatives $`_x(a)`$ and $`_y(a)`$ are well-defined only for those elements $`a𝔄(^2)_4`$ belonging to the corresponding operator domains.
We close this subsection by listing the formulas of the operator matrices for the generators $`E^{},F^{},K_1,K_2,x`$ and $`y`$. Recall that
$`L_\alpha =\overline{f_\alpha }(𝒫)e^{2\pi \alpha Q},R_\alpha =e^{2\pi \alpha 𝒬}f_\alpha (𝒫),`$
$`f_\alpha (𝒫)=2\mathrm{sinh}\pi \beta (2𝒫+\alpha i)=q^{1/2}e^{2\pi \beta 𝒫}+q^{1/2}e^{2\pi \beta 𝒫}.`$
$`\overline{\mathrm{\Phi }(E^{})}=𝖤=\left(\begin{array}{cccc}0& 0& 0& L_\alpha e^{2\pi \beta 𝒬_2}\\ 0& 0& R_\alpha e^{2\pi \beta 𝒬_2}& 0\\ 0& L_\alpha e^{2\pi \beta 𝒬_2}& 0& 0\\ R_\alpha e^{2\pi \beta 𝒬_2}& 0& 0& 0\end{array}\right),`$
$`\overline{\mathrm{\Phi }(F^{})}=𝖥=\left(\begin{array}{cccc}0& 0& 0& e^{2\pi \alpha 𝒬_1}L_\beta \\ 0& 0& e^{2\pi \alpha 𝒬_1}L_\beta & 0\\ 0& e^{2\pi \alpha 𝒬_1}R_\beta & 0& 0\\ e^{2\pi \alpha 𝒬_1}R_\beta & 0& 0& 0\end{array}\right),`$
$`\overline{\mathrm{\Phi }(q^{1/4}K_1)}=𝖪_1=\left(\begin{array}{cccc}e^{\pi \beta 𝒫_1}I& 0& 0& 0\\ 0& e^{\pi \beta 𝒫_1}I& 0& 0\\ 0& 0& e^{\pi \beta 𝒫_1}I& 0\\ 0& 0& 0& e^{\pi \beta 𝒫_1}I\end{array}\right),`$
$`\overline{\mathrm{\Phi }(q^{1/4}K_2)}=𝖪_2=\left(\begin{array}{cccc}Ie^{\pi \alpha 𝒫_2}& 0& 0& 0\\ 0& Ie^{\pi \alpha 𝒫_2}& 0& 0\\ 0& 0& Ie^{\pi \alpha 𝒫_2}& 0\\ 0& 0& 0& Ie^{\pi \alpha 𝒫_2}\end{array}\right),`$
$`\overline{\mathrm{\Phi }(x)}=𝗑=\left(\begin{array}{cccc}0& e^{2\pi \alpha 𝒬_1}e^{\pi \alpha 𝒫_2}& 0& 0\\ e^{2\pi \alpha 𝒬_1}e^{\pi \alpha 𝒫_2}& 0& 0& 0\\ 0& 0& 0& e^{2\pi \alpha 𝒬_1}e^{\pi \alpha 𝒫_2}\\ 0& 0& e^{2\pi \alpha 𝒬_1}e^{\pi \alpha 𝒫_2}& 0\end{array}\right),`$
$`\overline{\mathrm{\Phi }(y)}=𝗒=\left(\begin{array}{cccc}0& 0& e^{\pi \beta 𝒫_1}e^{2\pi \beta 𝒬_2}& 0\\ 0& 0& 0& e^{\pi \beta 𝒫_1}e^{2\pi \beta 𝒬_2}\\ e^{\pi \beta 𝒫_1}e^{2\pi \beta 𝒬_2}& 0& 0& 0\\ 0& e^{\pi \beta 𝒫_1}e^{2\pi \beta 𝒬_2}& 0& 0\end{array}\right).`$
6.4 In this last subsection we introduce three unitary operators $`_x^q,_y^q`$ and $`^q`$ on the Hilbert space $``$ which can be considered as quantum analogs of the partial Fourier transforms and the Fourier transform on $`^2`$, respectively.
First let us note that the counterparts of the $`q`$-deformed partial derivatives $`𝒟_x^q`$ and $`𝒟_y^q`$ (see (84)) on the algebra $`𝒜(_q^2)`$ are defined by
$`𝒟_x^q:=\mathrm{𝖪𝗒}^1𝖤=\left(\begin{array}{cccc}0& L_\alpha e^{\pi \alpha 𝒫_2}& 0& 0\\ R_\alpha e^{\pi \alpha 𝒫_2}& 0& 0& 0\\ 0& 0& 0& L_\alpha e^{\pi \alpha 𝒫_2}\\ 0& 0& R_\alpha e^{\pi \alpha 𝒫_2}& 0\end{array}\right),`$
$`𝒟_y^q:=\mathrm{𝖪𝗑}^1𝖥=\left(\begin{array}{cccc}0& 0& e^{\pi \beta 𝒫_1}L_\beta & 0\\ 0& 0& 0& e^{\pi \beta 𝒫_1}L_\beta \\ e^{\pi \beta 𝒫_1}R_\beta & 0& 0& 0\\ 0& e^{\pi \beta 𝒫_1}R_\beta & 0& 0\end{array}\right)`$
Clearly, $`𝒟_x^q`$ and $`𝒟_y^q`$ are self-adjoint operators on the Hilbert space $``$.
Let $`u`$ be the unitary operator on $`L^2()`$ given by $`(uf)(x)=f(x)`$ and let $`w_\alpha ,v_\alpha `$ resp. $`w_\beta ,v_\beta `$ be the holomorphic functions from Lemma 6. Then
$`_x^q:=\sigma _1(u\overline{w}_\alpha (𝒫_1)I,u\overline{v}_\alpha (𝒫_1)I),_y^q:=\sigma _2(Iu\overline{w}_\beta (𝒫_2),Iu\overline{w}_\beta (𝒫_2)).`$
are commuting unitaries on the Hilbert space $``$. Set $`^q:=_x^q_y^q`$, that is,
$$^q=\left(\begin{array}{cccc}u\overline{w}_\alpha (𝒫_1)u\overline{w}_\beta (𝒫_2)& 0& 0& 0\\ 0& u\overline{v}_\alpha (𝒫_1)u\overline{w}_\beta (𝒫_2)& 0& 0\\ 0& 0& u\overline{w}_\alpha (𝒫_1)u\overline{v}_\beta (𝒫_2)& 0\\ 0& 0& 0& u\overline{v}_\alpha (𝒫_1)u\overline{v}_\beta (𝒫_2)\end{array}\right).$$
We call the unitaries $`_x^q`$ and $`_y^q`$ quantum partial Fourier transforms and $`^q`$ quantum Fourier transform of the real quantum plane. The reason for this terminology stems from the fact that, roughly speaking, these unitaries interchange the coordinate functions $`x,y`$ and the $`q`$-deformed partial derivatives $`𝒟_x^q,𝒟_y^q`$, respectively. More precisely, we have the following relations.
Proposition 27. (i) $`_x^q𝗑(_x^q)^1=𝒟_x^q,_x^q𝒟_x^q(_x^q)^1=𝗑,_x^q𝗒(_x^q)^1=𝖪_1^2𝗒,_x^q𝒟_y^q(_x^q)^1=𝖪_1^2𝒟_y^q.`$
(ii) $`_y^q𝗒(_y^q)^1=𝒟_y^q,_y^q𝒟_y^q(_x^q)^1=𝗒,_y^q𝗑(_y^q)^1𝖪_2^2𝗑,F_y^q𝒟_x^q(_y^q)^1=𝖪_2^2𝒟_x^q`$.
(iii) $`^q𝗑(^q)^1=𝖪_2^2𝒟_x^q,^q𝗒(^q)^1=𝖪_1^2𝒟_y^q,^q𝒟_x^q(^q)^1=𝖪_2^2𝗑,^q𝒟_y^q(^q)^1=𝖪_2^2𝗒.`$
Proof. (i): By Lemma 7, written in terms of matrix entries, we have
$`u\overline{w}_\alpha (𝒫)L_\alpha v_\alpha (𝒫)u=e^{2\pi \alpha 𝒬},u\overline{v}_\alpha (𝒫)R_\alpha w_\alpha (𝒫)u=e^{2\pi \alpha 𝒬},`$ (185)
$`u\overline{w}_\alpha (𝒫)e^{2\pi \alpha 𝒬}v_\alpha (𝒫)u=uL_\alpha u=L_\alpha ,u\overline{v}_\alpha (𝒫)e^{2\pi \alpha 𝒬}w_\alpha (𝒫)u=uR_\alpha u=R_\alpha .`$ (186)
Because of the two relations (185) the matrix entries of $`_x^q𝒟_x^q(_x^q)^1`$ and $`𝗑`$ coincide, while (186) implies that $`_x^q𝗑(F_x^q)^1`$ and $`𝒟_x^q`$ have the same matrix entries. This proves the first two relations of (i). The two other relations follow at once from the corresponding definitions combined with the fact that $`ue^{\pi \beta 𝒫}u=e^{\pi \beta 𝒫}`$.
The proof of (ii) is similar to the proof of (i). The relations of (iii) follow easily from (i) and (ii). $`\mathrm{}`$
Remark 5. The holomorphic functions $`w_\alpha `$ and $`v_\alpha `$ coincide with the functions $`w_1`$ and $`w_2`$ in \[S4\], where a closely related quantum Fourier transform of a $`q`$-deformed Heisenberg algebra appeared. Holomorphic functions of similar kind have been used by S.L. Woronowicz \[W\] in another context as quantum exponential functions for the quantum $`ax+b`$–group.
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# Properties of the Langevin and Fokker-Planck equations for scalar fields and their application to the dynamics of second order phase transitions
## I The Langevin and Fokker-Planck equations for scalar fields
We start by formulating the problem. The second order Langevin equation for a scalar fields $`\varphi (x,t)`$ with an arbitrary interaction potential $`V(\varphi )`$ is given by
$`\left(_t^2^2\right)\varphi (x)+{\displaystyle \frac{\delta V(\varphi )}{\delta \varphi (x)}}+\eta \dot{\varphi }(x)=\xi (𝐱,t).`$ (1)
The stochastic fields $`\xi (x,t)`$ are taken to be Gaussian and white, characterized by
$`\xi (𝐱,t)=0,\xi (𝐱,t)\xi (𝐱^{},t^{})=\mathrm{\Omega }\delta ^3(𝐱𝐱^{})\delta (tt^{}).`$ (2)
Although we have written this equation for a single scalar field $`\varphi (x,t)`$ the generalization of Eq. (1-2) and of what follows, to a $`O(N)`$ symmetric theory, with N the number of flavors, is straightforward.
We can proceed to reduce the order of this differential system, by introducing field generalized momenta conjugate to $`\varphi `$:
$`_t\varphi (x)`$ $`=`$ $`\pi (x),`$ (3)
$`_t\pi (x)`$ $`=`$ $`\eta \pi (x)+^2\varphi (x){\displaystyle \frac{\delta V(\varphi )}{\delta \varphi (x)}}+\xi (x).`$ (4)
An interesting limit of Eqs. (3-4) arises when $`\eta `$ is large, or more precisely when $`_t\pi (x)<<\eta \pi (x)`$. Then we can write
$`\eta _t\varphi (x)`$ $`=`$ $`^2\varphi (x){\displaystyle \frac{\delta V(\varphi )}{\delta \varphi (x)}}+\xi (x).`$ (5)
Although this is the more conventional form of the Langevin field equation I will refer to it as the overdamped limit of Eqs. (3-4).
The objective of the Langevin analysis is to measure the expectation value (over the stochastic fields) of any functional $`\rho `$ of the fields $`\pi (x),\varphi (x)`$, generated by the evolution of Eq (1-2). With the choice Eq. (2) this average is a Gaussian functional integral of the form
$`{\displaystyle D\xi \rho [\pi _\xi ,\varphi _\xi ]e^{\frac{1}{2\mathrm{\Omega }}{\scriptscriptstyle d^4x\xi ^2}}}`$ (6)
where the $`\xi `$ subscripts in the field and its momentum denote their functional dependence on the stochastic fields $`\xi (x,t)`$, through the evolution of Eq. (3-4). Relation (6) and (2) can be generalized to fields $`\xi `$ with more complicated Gaussian distributions.
In the presence of the stochastic fields $`\xi (x,t)`$ the initial conditions for the fields $`\pi (x),\varphi (x)`$ may only be know statistically themselves. They can then be expressed in terms of a functional probability distribution $`P[\pi ,\varphi ]`$, at the initial time.
Throughout the evolution we assume that a time dependent probability distribution $`P_{\mathrm{FP}}[\varphi (x),\pi (x),t]`$, a functional of the time independent fields and a function of time, exists. This is the Fokker-Planck probability distribution, which we require, as usual, to be positive definite and normalizable in the sense:
$`𝒩{\displaystyle D\pi D\varphi P_{FP}[\varphi ,\pi ,t]}=1.`$ (7)
Note that in general $`𝒩`$ as well as other expectation values of the fields may be formally infinite as the ultraviolet cutoff of the theory is taken to zero.
The time dependent expectation value over the stochastic fields $`\xi `$ of any quantity $`\rho [\varphi ,\pi ]`$ is then given by
$`\rho (t)=𝒩{\displaystyle D\pi D\varphi P_{FP}[\varphi ,\pi ,t]\rho [\varphi ,\pi ]}.`$ (8)
In order to be useful this picture requires the explicit knowledge of $`P_{FP}[\varphi ,\pi ,t]`$. The Langevin equation generates a time-dependent (Fokker-Planck) probability distribution $`P_{FP}[\pi ,\varphi ,t]`$, which we can write formally
$`P_{FP}[\pi ,\varphi ,t]=`$ (9)
$`{\displaystyle d^Dx\delta \left[\widehat{\pi }(x,t)\pi (x)\right]\delta \left[\widehat{\varphi }(x,t)\varphi (x)\right]},`$ (10)
where the brackets denote, as before, average over the stochastic fields. The static fields are the arguments of $`P_{FP}[\pi ,\varphi ,t]`$, whereas $`\widehat{\pi }(x,t),\widehat{\varphi }(x,t)`$ obey the Langevin Eqs. (3\- 4).
The probability $`P_{FP}[\pi ,\varphi ,t]`$, obeys a functional Fokker-Planck evolution equation that can be computed directly by differentiating Eq. (10). Following Zinn-Justin
$`_tP_{FP}=[{\displaystyle }d^Dx_t\widehat{\varphi }(x,t){\displaystyle \frac{\delta }{\delta \widehat{\varphi }(x,t)}}+_t\widehat{\pi }(x,t){\displaystyle \frac{\delta }{\delta \widehat{\pi }(x,t)}}]`$ (11)
$`\times \delta [\widehat{\pi }(x,t)\pi (x)]\delta [\widehat{\varphi }(x,t)\varphi (x)].`$ (12)
The properties of the $`\delta `$ allow us to trade the functional derivatives of $`\widehat{\pi }(x,t),\widehat{\varphi }(x,t)`$ for derivatives relative to $`\pi (x),\varphi (x)`$. These can be taken out of the stochastic average $`\mathrm{}`$. The only remaining term inside this average results from the appearance of the stochastic fields $`\xi `$ in the equation of motion for $`\widehat{\pi }(x,t)`$. This term is of the form
$`{\displaystyle d^Dx\xi (x,t)\delta \left[\widehat{\pi }(x,t)\pi (x)\right]}.`$ (13)
The Gaussianity of the stochastic fields enforces the identity
$`\xi (x,t)\delta \left[\widehat{\pi }(x,t)\pi (x)\right]=\mathrm{\Omega }{\displaystyle \frac{\delta }{\delta \xi (x,t)}}\delta \left[\widehat{\pi }(x,t)\pi (x)\right]`$ (14)
$`=\mathrm{\Omega }{\displaystyle \frac{\delta \widehat{\pi }(x,t)}{\delta \xi (x,t)}}{\displaystyle \frac{\delta }{\widehat{\pi }(x,t)}}\delta \left[\widehat{\pi }(x,t)\pi (x)\right].`$ (15)
Once again we can trade the functional derivative relative to $`\widehat{\pi }(x,t)`$ for another relative to $`\pi (x)`$. The expectation value of $`\frac{\delta \widehat{\pi }(x,t)}{\delta \xi (x,t)}`$ can in turn be obtained from the formal integration of the equation of motion by using a regularized delta function in time. We must effectively take half of this delta function to obtain
$`{\displaystyle \frac{\delta \widehat{\pi }(x,t)}{\delta \xi (x,t)}}={\displaystyle \frac{1}{2}}.`$ (16)
Bringing together all the terms results finally in
$`_tP_{FP}[\pi ,\varphi ,t]=_{\mathrm{FP}}P_{FP}[\pi ,\varphi ,t].`$ (17)
where
$`_{\mathrm{FP}}`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Omega }}{2}}{\displaystyle \frac{\delta ^2}{\delta \pi ^2}}+\pi {\displaystyle \frac{\delta }{\delta \varphi }}`$ (19)
$`{\displaystyle \frac{\delta }{\delta \pi }}\left(\eta \pi ^2\varphi +{\displaystyle \frac{\delta V(\varphi )}{\delta \varphi }}\right)`$
It is useful to think in terms of the functional operator $`_{\mathrm{FP}}`$ as the generator of infinitesimal time displacements of the probability functional $`P_{\mathrm{FP}}`$. If, as in most applications, the potential $`V(\varphi )`$ is explicitly time independent we can invoke a separation ansatz for $`P_{FP}`$ such that
$`P_{FP}[\pi ,\varphi ,t]=𝒫[\pi ,\varphi ]T(t)`$ (20)
Thus we can regard Eq. (17) as an analog of a functional Schrödinger equation, in imaginary time. Then we can write the time independent and dependent equations
$`_{\mathrm{FP}}𝒫_n=E_n𝒫_n,_tT(t)=E_nT(t).`$ (21)
The functional dependence on the fields is now limited to the static probability eigenfunctionals $`𝒫_n`$. The time evolution of the Fokker-Planck distribution is completely characterized by the spectrum of eigenvalues of $`_{\mathrm{FP}}`$, $`E_n`$. An orthogonal complete basis of functionals $`B_n`$ can in general be constructed from the set $`𝒫_n`$<sup>*</sup><sup>*</sup>*The set of functionals $`𝒫_n`$ will only be guaranteed to be an orthogonal basis if $`_{\mathrm{FP}}`$ is a Hermitian functional operator, which is not true in general. The relation between the $`B_n`$ and $`P_n`$ is usually very simple involving a power of the canonical distribution and it follows that $`E_n`$ are also the eigenvalues of the set $`B_n`$. See eg. for some examples in few degree of freedom problems. Note that it is not possible in general to separate Eq. (20) further into two equations, one in $`\pi `$ and another in $`\varphi `$ via an ansatz like $`P_{FP}=R[\pi ]S[\varphi ]`$.
On general grounds we expect a steady state corresponding to thermal equilibrium to be reached for long times. Formally, we can then project the evolution of $`P_{FP}`$ in terms of the eigenvalues $`E_n`$ and functionals $`B_n`$ as:
$`P_{FP}[\pi ,\varphi ,t]={\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}C_nB_n[\pi ,\varphi ]e^{E_nt}.`$ (22)
where the $`C_i`$’s are the projections of $`P_{FP}`$ at the initial time onto the basis of eigenfunctionals $`B_n`$:
$`C_n={\displaystyle D\pi D\varphi B_n[\pi ,\varphi ]P_{FP}[\pi ,\varphi ,t=0]},`$ (23)
where both $`B_n`$ and $`P_{FP}[\pi ,\varphi ,t]`$ are taken to be normalized.
If equilibrium is approached for long times the corresponding distribution must be associated with the zero mode of $`_{\mathrm{FP}}`$, $`E_0=0`$ i.e. it is stable. Contributions from eigenstates $`𝒫_n`$, $`n0`$, vanish exponentially, in a characteristic time $`t_{\mathrm{eq}}E_n^1`$, provided that the real part of $`E_n`$ is positive. This follows from the requirement that $`V(\varphi )`$ is bounded from below. For the $`E_n`$, solutions of the second time derivative eigenproblem there may be in addition an imaginary part. Thus the excited states decay away in time as an exponentially damped oscillator.
The eigenfunctional $`𝒫_0`$, characterizing the long time field probability is then interpreted as the equilibrium distribution. Taking $`E_0=0`$ in Eq. (21) we find
$`𝒫_0[\pi ,\varphi ]=𝒩\mathrm{exp}\left[\beta {\displaystyle d^Dx\left(\frac{\pi (x)^2}{2}+S(\varphi )\right)}\right]`$ (24)
with $`S(\varphi )=\left(\frac{1}{2}(\varphi (x))^2+V(\varphi )\right)`$. Here we took $`\mathrm{\Omega }=2\eta /\beta `$, which is the analog of Einstein’s relation for Brownian motion, and ensures the balance between fluctuations (sourced by the fields $`\xi (x,t)`$) and dissipation (through the $`\eta \pi (x,t)`$ in Eq. (4)) at long times.
It is interesting to note that it can be read directly from Eq. (24) that the momentum equilibrium variance
$`\pi (x)^2=1/\beta =T,`$ (25)
which expresses equipartition in a relativistic classical field theory. This quantity is also an excellent thermometer for the dynamical evolution.
Now, because the canonical distribution is Gaussian in $`\pi `$ we can proceed to obtain a reduced distribution in terms of $`\varphi `$ alone. Performing the Gaussian integral over all $`\pi `$ we obtain
$`P_{eq}[\varphi ]=𝒩^{}\mathrm{exp}\left[\beta {\displaystyle d^DxS(\varphi )}\right],`$ (26)
which is the canonical Boltzmann distribution. $`P_{eq}[\varphi ]`$ could have been equally obtained from the static solution of the Fokker-Planck equation associated with the perhaps more usual first derivative Langevin equation (see below). In this respect we see that in equilibrium it makes no difference to perform the evolution, with or without the second time derivative in Eq. (1). Away from equilibrium, or even if we simply wish to study the fluctuations around it, we should take the appropriate form of the evolution, choosing to keep or neglect the second time derivative according to the physical picture under consideration. The difference between both of these evolutions, is expressed in terms of the higher eigenvalues and eigenfunctionals associated with the different Hamiltonians, which I discuss below.
For a harmonic potential $`V[\varphi ]=\frac{m^2}{2}\varphi ^2`$ we can also read the equilibrium thermal propagator directly from Eq. (26). It is, in momentum space
$`\varphi _p\varphi _p={\displaystyle \frac{T}{p^2+m^2}}.`$ (27)
This is the canonical free Boltzmann propagator, which is the correct description for ideal classical waves at finite temperature.
It is illuminating to compare the above framework directly with that for the overdamped Langevin evolution. Following the same, but somewhat simpler, procedure it can be shown that the Fokker-Planck Hamiltonian takes the form
$`_{FP}={\displaystyle \frac{\delta }{\delta \varphi }}\left[{\displaystyle \frac{\mathrm{\Omega }}{2\eta }}{\displaystyle \frac{\delta }{\delta \varphi }}^2\varphi +{\displaystyle \frac{\delta V}{\delta \varphi }}\right],`$ (28)
leading to the eigenvalue functional equation
$`\left\{{\displaystyle \frac{\delta }{\delta \varphi }}\left[{\displaystyle \frac{\mathrm{\Omega }}{2\eta }}{\displaystyle \frac{\delta }{\delta \varphi }}^2\varphi +{\displaystyle \frac{\delta V}{\delta \varphi }}\right]+E_N\right\}P[\varphi ]=0,`$ (29)
$`\eta _tP[\varphi ]=E_NP[\varphi ].`$ (30)
Again the canonical distribution is the eigenfunctional $`P_0`$ corresponding to a zero eigenvalue, i.e.
$`E_0=0;P_{eq}[\varphi ]=𝒩^{}e^{\beta {\scriptscriptstyle d^Dx{\scriptscriptstyle \frac{1}{2}}(\varphi )^2}+V(\varphi )},`$ (31)
where as before the Einstein relation $`\mathrm{\Omega }=2\eta /\beta `$ must hold.
## II Approach to equilibrium
A complete dynamical solution of the ensemble (under average over the noise) of fields obeying the Langevin equation (3-4) can be obtained in principle by solving the corresponding Fokker-Planck equation. Unfortunately the task of solving for these non-linear functional equations in general is quite monumental. The main difficulty is connected with the non-local form of the equations: in x-space this arises from the Laplacian term, while in k space the difficulty is transfered to the non-linear terms.
Below I give closed form solutions to the equations in the fields and their conjugate momenta for the particular case of the harmonic potential. These are still interesting since they offer an alternative to the direct solutions of the Langevin equation, which, clearly, can also be found in the linear case. In the former picture, however, the averages over the stochastic fields are already taken into account.
It is simpler to build some intuition for the solution of the overdamped Fokker-Planck equation first to which I now turn.
### A The overdamped Fokker-Plank equation and its solutions
The advantage of the harmonic case is that in Fourier space different modes decouple in the Langevin equation:
$$\eta _t\varphi _k(t)=(k^2+m^2)\varphi _k(t)+\xi _k(t)$$
(32)
where the stochastic fields obey
$`\xi _k(t)=0,\xi _k(t)\xi _k(t^{})=\mathrm{\Omega }\delta (tt^{})`$ (33)
and
$`\varphi (x,t)={\displaystyle \frac{d^Dk}{(2\pi )^D}\varphi _k(t)e^{ik.x}}.`$ (34)
The fact that $`\varphi (x,t)`$ is real implies $`\varphi _k=\varphi _k^{}`$ and similarly for $`\xi _k`$. Because of these relations it is more convenient to work with $`\varphi _k^R=\mathrm{Re}(\varphi _k)`$ and $`\varphi _k^I=\mathrm{Im}(\varphi _k)`$, which obey the Langevin Eq. (32) with $`\xi _k^{R,I}`$:
$`\xi _k(t)^{R,I}=0,\xi _k^{R,I}(t)\xi _k^{R,I}(t^{})={\displaystyle \frac{\mathrm{\Omega }}{2}}\delta (tt^{}).`$ (35)
The Fokker-Planck equation for the modes now is
$`\eta _tP_k=_kP_k`$ (36)
$`_k={\displaystyle \frac{\delta }{\delta \varphi _k^{R,I}}}\left[{\displaystyle \frac{\mathrm{\Omega }}{4\eta }}{\displaystyle \frac{\delta }{\delta \varphi _k^{R,I}}}+\left(k^2+m^2\right)\varphi _k^{R,I}\right].`$ (37)
which holds, as indicated, both for the real and imaginary components of $`\varphi _k`$. The solution of the Fokker-Planck equation then follows as
$`P=\mathrm{\Pi }_{k=0}^{\mathrm{}}\left[P_k(\varphi _k^R)P_k(\varphi _k^I)\right]`$ (38)
The time independent solution is clearly
$`P_k(\varphi _k^R)P_k(\varphi _k^I)\mathrm{exp}\left[\beta (k^2+m^2)\left(\varphi _{k}^{R}{}_{}{}^{2}+\varphi _{k}^{I}{}_{}{}^{2}\right)\right]`$ (39)
which leads to
$`P_0=`$ $`𝒩\mathrm{exp}\left[\beta {\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{d^Dk}{(2\pi )^D}}(k^2+m^2)\left(\varphi _{k}^{R}{}_{}{}^{2}+\varphi _{k}^{I}{}_{}{}^{2}\right)\right]`$ (40)
$`=`$ $`𝒩\mathrm{exp}\left[\beta {\displaystyle _{\mathrm{}}^+\mathrm{}}{\displaystyle \frac{d^Dk}{(2\pi )^D}}\varphi _k{\displaystyle \frac{k^2+m^2}{2}}\varphi _k\right].`$ (41)
$`=`$ $`𝒩\mathrm{exp}\left[\beta {\displaystyle _{\mathrm{}}^{\mathrm{}}}d^Dx{\displaystyle \frac{1}{2}}\left(\varphi (x)\right)^2+{\displaystyle \frac{m^2}{2}}\varphi ^2(x)\right].`$ (42)
To compute the eigenfunctions corresponding to non-zero eigenvalues we need to solve the eigenvalue problem
$`\left[{\displaystyle \frac{1}{2\beta }}{\displaystyle \frac{\delta ^2}{\delta \varphi _{k}^{R,I}{}_{}{}^{2}}}(k^2+m^2)\varphi _k^{R,I}{\displaystyle \frac{\delta }{\delta \varphi _k^{R,I}}}+E_N\right]F_N=0`$ (43)
where we wrote $`P_k=P_0F_N`$. This equation has a familiar solution. To see this explicitly we perform a change of variables to obtain
$`\left[{\displaystyle \frac{\delta ^2}{\delta X}}2X{\displaystyle \frac{\delta }{\delta X}}+{\displaystyle \frac{2E_N}{k^2+m^2}}\right]F_N=0`$ (44)
where $`X=\sqrt{\beta (k^2+m^2)}\varphi _k^{R,I}`$. Eq. (44) has the solution
$`E_N=N(k^2+m^2)`$ (45)
$`F_N=H_N(X)=H_N(\sqrt{\beta (k^2+m^2)}\varphi _k^{R,I})`$ (46)
where $`H_N`$ is the Hermite polynomial of order $`N`$.
From the properties of the Hermite polynomials we know that these solutions are orthogonal under the measure $`e^{X^2}`$, which is just the k-part of the canonical distribution. The set of functions $`\{H_N(X)e^{X^2/2}\}`$ thus constitutes an orthogonal basis. It can additionally be shown trivially that this basis is complete.
The non-zero Eigenvalues $`E_N`$ set the time scales for the approach to thermal equilibrium, which we will discuss in more detail below.
### B The second order Fokker-Planck equation
For the second time-derivative harmonic evolution the Langevin equation becomes
$`\varphi _k=\pi _k`$ (47)
$`\pi _k=\eta \pi _k(k^2+m^2)\varphi _k+\xi _k,`$ (48)
which as in subsection II A can be written in terms of real and imaginary components. The Fokker-Planck Hamiltonian for the latter is
$`_k=`$ $`{\displaystyle \frac{\mathrm{\Omega }}{4}}{\displaystyle \frac{\delta ^2}{\delta \pi _{k}^{R,I}{}_{}{}^{2}}}+\pi _k^{R,I}{\displaystyle \frac{\delta }{\delta \varphi _k^{R,I}}}`$ (50)
$`{\displaystyle \frac{\delta }{\delta \pi _k^{R,I}}}\left[\eta \pi _k^{R,I}+(k^2+m^2)\varphi _k^{R,I}\right].`$
To find the excited states we proceed as before to factor out the canonical distribution $`P=F_NP_0`$ to get
$`\{{\displaystyle \frac{\mathrm{\Omega }}{4}}{\displaystyle \frac{\delta ^2}{\delta \pi _{k}^{R,I}{}_{}{}^{2}}}+[\eta \pi _k^{R,I}(k^2+m^2)\varphi _k^{R,I}]{\displaystyle \frac{\delta }{\delta \varphi _k^{R,I}}}`$ (51)
$`+{\displaystyle \frac{\delta }{\delta \varphi _k^{R,I}}}E_N\}F_N=0.`$ (52)
We can bring Eq. (52) to the form of Eq. (44) by making the change of variables $`X_\pm =a_\pm \pi _k^{R,I}+b_\pm \varphi _k^{R,I}`$, with
$`a_\pm ^2={\displaystyle \frac{\beta }{2}}\left[1\pm \sqrt{14{\displaystyle \frac{k^2+m^2}{\eta ^2}}}\right],`$ (53)
$`b_\pm ^2={\displaystyle \frac{2\beta }{\eta ^2}}(k^2+m^2)^2/\left[1\pm \sqrt{14{\displaystyle \frac{k^2+m^2}{\eta ^2}}}\right].`$ (54)
The solutions of Eq. (52) then are
$`E_N^\pm =N{\displaystyle \frac{a_\pm ^2\eta }{\beta }}=N{\displaystyle \frac{\eta }{2}}\left[1\pm \sqrt{14{\displaystyle \frac{k^2+m^2}{\eta ^2}}}\right]`$ (55)
$`F_N^\pm =H_N[X_\pm ].`$ (56)
These solutions can be brought to products of Hermite polynomials in $`\pi _k^{R,I}`$ and $`\varphi _k^{R,I}`$ alone through the use of the property:
$`H_N[x+y]=2^{N/2}{\displaystyle \underset{i=0}{\overset{N}{}}}\left(\stackrel{i}{N}\right)H_{Ni}[x\sqrt{2}]H_i[y\sqrt{2}].`$ (57)
This property allow us to see that the solutions
$`P_N=\left(F_N[X_+]e^{E_N^+t}+F_N[X_{}]e^{E_N^{}t}\right)P_0`$ (58)
form a complete orthogonal basis.
A few properties of these solutions are worth pointing out. The coefficients $`a_\pm ,b_\pm `$ and the eigenvalues $`E_N^\pm `$ can now be complex numbers. This happens for $`4(k^2+m^2)>\eta ^2`$. In this case $`X_+=X_{}^{}`$, $`E_N^+=(E_N^{})^{}`$ and the combination in (58) remains real as required.
The overdamped limit is recovered as $`\eta ^2>>k^2+m^2`$. Then
$`E_N^{}=N{\displaystyle \frac{k^2+m^2}{\eta }},`$ (59)
$`a_{}^2={\displaystyle \frac{\beta }{2}}{\displaystyle \frac{k^2+m^2}{\eta ^2}}<<b_{}^2={\displaystyle \frac{\beta }{2}}(k^2+m^2).`$ (60)
The eigenfunctionals associated with $`E_N^+N\eta [1+(k^2+m^2)/\eta ^2]`$ are, in contrast, damped away rapidly.
The characteristic thermalization times of the system are thus set by $`t_N=1/E_N`$. For the long wave length modes ($`k0`$) we have that the equilibration time $`t_{eq}\eta /(k^2+m^2)`$. In the converse limit, of short wavelengths, the decay of all states is controlled at leading order by $`t_{eq}2/\eta `$, which is scale independent. In this sense the evolution of the short wavelength modes is quantitatively different from the overdamped Langevin system.
## III A mean-field approximation to the non-linear Fokker-Planck dynamics
It would no doubt be desirable to be able to solve the Fokker-Planck equations for more general (interacting) potentials. At the next level of complexity it is possible to attempt a solution of the Fokker-Planck equation in a mean-field approximation, which we discuss in this section.
In order to illustrate the procedure let us consider an interaction potential of the form $`V_{int}=\frac{\lambda }{4}\varphi (x)^4`$. In momentum space this potential generates a term in the Fokker-Planck equation for $`\varphi _k`$ of the form
$`{\displaystyle \frac{\delta V_{int}}{\delta \varphi _k}}=\lambda {\displaystyle \frac{d^Dp}{(2\pi )^D}\frac{d^Dq}{(2\pi )^D}\varphi _p\varphi _q\varphi _{kpq}}.`$ (61)
A mean-field approximation to this term can be obtained by taking $`p=q`$ and the noise average such that
$`{\displaystyle \frac{\delta V_{MF}}{\delta \varphi _k}}=3\lambda \left[{\displaystyle \frac{d^Dp}{(2\pi )^D}\varphi _p\varphi _p}\right]\varphi _k3\lambda G(t)\varphi _k.`$ (62)
Under this approximation the evolution equation can be written in terms of a mean-field Fokker-Planck operator where $`V_{MF}`$ substitutes the potential $`V`$, with
$`V_{MF}=V_0+{\displaystyle \frac{3\lambda }{4}}G(t)^2,`$ (63)
where $`V_0`$ corresponds to the purely harmonic case, together with the self-consistency condition
$`G(t)={\displaystyle \frac{d^Dp}{(2\pi )^D}𝑑\varphi _p\varphi _p\varphi _pP[\pi _q,\varphi _p,t]}.`$ (64)
It is then clear that the under this Gaussian approximation the theory remains harmonic, but with a time-dependent (self-consistently determined) mass.
Unfortunately the time dependence of the mean field potential destroys the separability of the solution into a function of time, and a functional of the fields and their conjugate momenta. Closed form solutions are thus difficult to construct, which is of course also the case for the Langevin equation under a similar approximation. Numerical solutions of the this mean field Fokker Planck equation can nevertheless be easily obtained.
As a starting point let us again consider the overdamped case first. Mean field approximations correspond in general to Gaussian probability distributions. In the overdamped case, this can only lead to
$`P_{MF}[\varphi _k,t]=𝒩_k(t)\mathrm{exp}\left[\varphi _k{\displaystyle \frac{A_k^2(t)}{2}}\varphi _k\right],`$ (65)
where the normalization $`𝒩(t)`$ is a function of $`A_k^2(t)`$
$`𝒩_k(t)=A_k(t)/\sqrt{\pi },`$ (66)
and must therefore be time-dependent. This does not of course introduce any additional dynamical freedom.
The mean-field Fokker-Planck equation for $`P_{MF}[\varphi _k,t]`$ translates into an equation for the coefficients $`A_k(t)`$ of the form
$`\eta {\displaystyle \frac{dA_k}{dt}}=\left[k^2+m^2+3\lambda G(t)A_k^2T\right]A_k,`$ (67)
$`G(t)={\displaystyle \frac{d^Dp}{(2\pi )^D}\frac{1}{A_p^2}}.`$ (68)
Clearly solving the Fokker-Planck equation in this mean-field approximation requires the solution of a system of coupled ordinary differential equations, one for every mode $`k`$. This is similar in effort to the solution of the Langevin mean-field equation with the advantage that the stochastic average is already taken into account. This latter step corresponds to the functional integral necessary to compute $`G(t)`$ in the Fokker-Planck picture, which can be given in closed form because $`P_{MF}`$ is Gaussian.
The canonical (mean-field) distribution is of course a static solution of Eqs (67-68) where $`A_k`$ obeys
$`A_k^2={\displaystyle \frac{k^2+m^2}{T}}+3{\displaystyle \frac{\lambda }{T}}{\displaystyle \frac{d^Dp}{(2\pi )^D}\frac{1}{A_p^2}},`$ (69)
It is then clear that for vanishing interactions $`\lambda =0`$, the free Boltzmann propagator is the solution of Eq. (69). In general because the mean-field contribution is momentum independent Eq. (69) are equivalent to the ’gap’ equation
$`A_k^2`$ $`={\displaystyle \frac{k^2+m^2+\mathrm{\Delta }m^2}{T}},`$ (70)
$`\mathrm{\Delta }m^2`$ $`=3\lambda {\displaystyle \frac{d^Dp}{(2\pi )^D}\frac{T}{k^2+m^2+\mathrm{\Delta }m^2}}.`$ (71)
For the second order Fokker-Planck equation the most general Gaussian function of the field and momentum modes is of the form
$`P_{MF}[\pi _k,\varphi _k,t]=`$ (72)
$`𝒩_k(t)\mathrm{exp}\left\{\left[{\displaystyle \frac{A_k^2}{2}}\varphi _k\varphi _k+{\displaystyle \frac{B_k^2}{2}}\pi _k\pi _k+C_k^2\mathrm{Re}(\pi _k\varphi _k)\right]\right\}.`$ (73)
The normalization is now a function of all $`A_k^2,B_k^2,C_k^2`$
$`𝒩_k={\displaystyle \frac{1}{\pi }}\sqrt{A_k^2B_k^2C_k^4}.`$ (74)
The equations of motion for $`A_k^2,B_k^2`$ and $`C_k^2`$ then follow from the second order Fokker-Planck equation:
$`{\displaystyle \frac{dA_k}{dt}}={\displaystyle \frac{C_k}{A_k}}\left[k^2+m^2+3\lambda G(t)\eta TC_k^2\right]`$ (75)
$`{\displaystyle \frac{dB_k}{dt}}=\eta B_k\left[1TB_k^2\right]`$ (76)
$`{\displaystyle \frac{dC_k}{dt}}={\displaystyle \frac{\eta C_k}{2}}\left[12TB_k^2\right]`$ (77)
$`{\displaystyle \frac{1}{2C_k}}\left[A_k^2\left(k^2+m^2+3\lambda G(t)\right)B_k^2\right],`$ (78)
where
$`G(t)={\displaystyle \frac{d^Dp}{(2\pi )^D}\frac{B_p^2}{A_p^2B_p^2C_p^4}}.`$ (79)
This system has as static solutions the mean-field canonical distribution namely
$`A_k^2={\displaystyle \frac{k^2+m^2+\mathrm{\Delta }m^2(T)}{T}},`$ (80)
$`B_k^2=1/T=\beta `$ (81)
$`C_k^2=0,`$ (82)
where the equation for $`A_k`$ is simply (69). Other static solutions exist though. Additional stationary distributions distinct from thermal equilibrium also occur in the microcanonical evolution of the class of models considered here and seemingly result from the truncation of the system at any finite order in a hierarchy of correlators. The overdamped limit, $`k^2+m^2+\mathrm{\Delta }m^2(t)<<\eta `$ is obtained when $`C_k^2A_k^2/\eta <<B_k^2,A_k^2`$.
Again we see that the effort in solving the mean-field second order Fokker-Plank equation is comparable to its Langevin counterpart, but with the stochastic averages taken into account. A renormalization scheme for rendering these equations ultraviolet cutoff independent can be constructed, if desired, following standard techniques for mean-field evolutions .
## IV The uses of the closed form solutions of the Fokker-Planck equation
The solutions presented above permit some analysis of several important dynamical situations.
In many cases in field theory the harmonic solutions can reveal the essential physics of the (short time) evolution. This is true for example for the motion of soft fields in the vicinity of a second order phase transition, where only the lowest relevant operators are needed, and for some situations when (spinodal) instabilities can develop, such as in the early stages of reheating after cosmological inflation or in the aftermath of a violent pressure quench, in some condensed matter experiments.
In this section I analyze these two situations in the context of the solutions of section II.
### A Critical Dynamics of the fields and the theory of topological defect formation at second order transitions
In the vicinity of a symmetry breaking second order transition the physical mass squared of the fields will approach zero, with a universal critical exponent $`\nu `$,
$`m^2(T)m^2|{\displaystyle \frac{TT_c}{T_c}}|^\nu ,`$ (83)
where $`T_c`$ is the critical temperature. This can be seen already at first order in perturbation theory, albeit only with the mean-field value of $`\nu =1/2`$. In 3D we have
$`m^2(T)m^2+3\lambda {\displaystyle ^\mathrm{\Lambda }}{\displaystyle \frac{d^3k}{(2\pi )^3}}{\displaystyle \frac{T}{k^2}}=m^2+{\displaystyle \frac{3\lambda }{2\pi ^2}}\mathrm{\Lambda }T.`$ (84)
$`\mathrm{\Lambda }`$ is the ultraviolet cutoff which in many cases has physical meaning because generally scalar field theories are effective low energy models. An example is the scale of separation between the true quantum behavior of excitations at large $`k^2`$ and the effective classical dynamics for the long wave-length modes which we describe stochastically. In this case, at high temperature and small coupling $`\lambda `$, $`\mathrm{\Lambda }T`$.
The interesting fact about $`m^2(T)0`$ at criticality is that it supplies us with an arbitrarily long time scale separation between the thermalization time for the short and long wave-length modes. As we discussed in section II the long wave length modes $`k^20`$ thermalize in a typical time-scale
$`t_{eq}{\displaystyle \frac{\eta }{m^2(T)}}_{TT_c}+\mathrm{},`$ (85)
which is the expression of critical slowing down in our stochastic system. On the other hand the short-wave length modes thermalize instead on a characteristic time scale
$`t_{eq}2/\eta .`$ (86)
It is worth commenting of the form (85). It shows that even in the context of a second order in time Langevin evolution the long wavelength modes are effectively overdamped in the critical domain. This is the essence of the perhaps more familiar TDGL evolution, which is an expansion in the lowest number of relevant field operators that can only be justified rigorously in the critical domain of a second order transition. We see therefore that our Langevin description encapsulates the TDGL dynamics but in addition also applies more generally.
Now, particularly in $`D=3`$ (as compared to lower dimensions) the temperature corrections to the mass are dominated by the small wave-length modes. Imagine then that the bath temperature is changed to a new value $`T_f`$ in the vicinity of the phase transition. Over a time $`t2/\eta `$ the thermal mass will then adapt to its new small value $`m^2(T_f)`$. In contrast the long wave-length modes require a much longer time to rethermalize and stay away from thermal equilibrium for a time $`t\eta /m^2(T_f)>>\eta `$.
This imbalance is at the heart of our current understanding of the dynamics of second order phase transitions and constitutes in particular the essence of the Kibble-Zurek theory of topological defect formation. The familiar response time $`\tau (T)`$ is here given by Eq. (85). Incidentally the temperature dependence of (84) also implies that whatever the time behavior of the external bath temperature the dependence of $`m^2(T)`$ on time will be the same, i.e. for example for a linear external drive of the bath temperature one also obtains $`m^2(t)t`$, again on timescales $`t\stackrel{>}{}2/\eta `$.
The consequences of this imbalance are very important. By crossing a second order phase transition at some finite rate $`\tau _Q`$ imposed externally there will be a time in the vicinity $`T_c`$ when the long-wave length modes fall out of thermal equilibrium. Because the system always rethermalizes in its small scales first, soft long wavelength non-trivial field configurations can persist for a long time. In particular topological defects can be “formed” in this way.
### B Spinodal instabilities
Another interesting application of the harmonic solutions of section II is the onset of (spinodal) instabilities. This situation assumes an initial instability characterized eg. by a coherent (i.e. quasi-spatially homogeneous) field, close to the origin $`\varphi =0`$ of a double-well potential. In addition there may be fluctuations (thermal and/or quantum) about this mean field but these must be tuned to be small and are often taken to be Gaussian and white.
In these circumstances $`m^2`$ is initially negative and the long wave-length modes $`k^2<m^2`$ will develop instabilities characterized by exponential growth of their amplitude at early times. In the limit of vanishing stochastic fields this behavior can be captured by solving the linearized Langevin equation , whose solutions are of course simple exponentials of $`\pm i\sqrt{k^2m^2}t`$. Linear dependences of $`m^2`$ on time can also be captured in similar ways.
The advantage of using the set of solutions (55) in this context is that the noise effects are then better taken into account. For the linearized problem the set of eigenfunctions and eigenvalues given in (55) is exact. The behavior of these solutions in the overdamped case, where the eigenvalues are
$`E_N=N{\displaystyle \frac{\left(k^2m^2\right)}{\eta }},`$ (87)
is obvious: whenever $`k^2m^2<0`$ all excited states display exponential growth. It is worth noting that it is the highest lying states (i.e. those with largest $`N`$) that grow the fastest, leading to spectacularly out of equilibrium configurations, at least relative to the initial state. The short wavelength modes are of course still dissipated away on a time scale $`t2/\eta `$ so these excited states are highly biased towards the infrared as is well known from direct numerical studies.
For the second order Fokker-Planck equation the eigenvalues are instead (55). It is then clear that there will be growing modes as long as $`k^2<m^2`$. In these circumstances, for small but negative $`m^2/\eta `$, the eigenvalues coincide with (87). For the second time derivative evolution however there will be both exponentially growing and decaying modes.
The validity of the linear approximation breaks down when the 2-point function $`G(t)`$, which controls the leading corrections to the mass squared, becomes of order the negative bare mass squared. In the Fokker-Planck picture this is an average which correctly accounts for the noise effects on the short wavelength modes, which are quasi-Gaussian (perturbative) in any case. The value of $`G`$ in linearized evolution can clearly be computed in closed form since all integrations are over polynomials of the fields times a Gaussian.
Computations of eg. $`G(x,t)`$ at the point where the linearized approximation fails have been shown to be useful , eg. in the computation of topological defect formation through the insertion of $`G(x,t)`$ into the well-known Halperin formula for counting field zeros, provided there is little energy left in the system at this point to cause substantial subsequent transport towards short wavelengths, (see for counterexamples).
## V Discussion and Conclusions
In this paper I discussed general properties of the second order Langevin scalar field evolutions as well as of some of its related models. The associated Fokker-Planck equations were solved in closed form for the harmonic potential and a mean-field approximation was devised for both second order and overdamped cases. The analogous treatment for the stochastic nonlinear Schrödinger equation was given in Appendix A. The latter can be solved numerically with an effort that is comparable to the mean-field Langevin equation, but has the advantage of already including the effects of the averages over the stochastic fields.
Although limited to special cases closed form solutions of the Fokker-Planck equations, have been shown to be very useful in providing us with analytical qualitative (and sometimes quantitative) insights into the evolution of the fully non-linear field theoretical system, especially in determining time scales for thermalization of several parts of the system. Through this analysis I showed how fundamental ingredients of the dynamics of second order phase transitions are embodied in these effective models and how they generalize the well known TDGL evolution in the critical domain.
In particular in the vicinity of $`T_c`$, where $`m^2(T)0`$, one obtains true time scale separation, with the small wavelength field modes thermalizing in a time $`t_{eq}\eta ^{}1`$, while the long wave length modes take a time $`t_{eq}\eta /m^2(T)>>1/\eta `$. In these circumstances, over times $`1/\eta <<t<<\eta /m^2(T)`$, one can take the short wavelength modes to be thermalized at the temperature of the stochastic thermal bath, and include their effects on the long-wavelength modes through the temperature dependence of the parameters in their potential. This treatment allows us to discuss in qualitative correct terms the evolution of the long wavelength modes in the critical domain, the phenomenon of critical slowing down and the ingredients of the theory of topological defect formation in these models.
## VI Acknowledgments
It is a pleasure to thank N. Antunes, S. Habib, G. Lythe and W. H. Zurek for useful comments. This research was supported by the U.S. Department of Energy, under contract W-7405-ENG-36.
## Apendix A
In this appendix I repeat the analysis applied to the Langevin equations (1) to the case of a stochastic non-linear Schrödinger equation (sNLS), known in some instances also as the stochastic Gross-Pitaevsii equation. This model is important in the study of non-relativistic many body systems. Many examples exist that are thought to obey this effective dynamics, ranging from strong type II (or hard) superconductors and superfluid <sup>4</sup>He to atomic Bose-Einstein condensates and light propagation in non-linear media.
The sNLS is of the form
$`\left(i+\eta \right)_t\psi =^2\psi +{\displaystyle \frac{\delta V}{\delta \psi ^{}}}+\xi (x,t),`$ (88)
where $`\psi =\psi (x,t)`$ is a complex field in space and time. The potential $`V`$ is usually taken to be
$`V\left[|\psi |^2\right]=m^2|\psi (x,t)|^2+{\displaystyle \frac{\lambda }{2}}|\psi (x,t)|^4,`$ (89)
where $`m^2`$ can be positive or negative. This theory is essentially the non-relativistic version of the second order in time Langevin system for a complex scalar field. The theory is manifestly invariant under global phase $`U(1)`$ transformations. Its Hamiltonian is
$`H={\displaystyle d^Dx\left\{|\psi |^2+V\left[|\psi |^2\right]\right\}}.`$ (90)
Although the stochastic terms in Eq. (88) may be representative of additional intrinsic degrees of freedom, it may also be possible to actually build an experimental situation in which the many-body system is driven externally and has losses to the outside so as to realize (88). In this latter picture the requirements are that the driving field should be phase incoherent (at least over some short characteristic time and spatial scales) and that the losses would be proportional to the frequency of the excitations in the system. Perhaps the most challenging feature would be to find a trap or mirror whose transmission is linear in the range of frequencies expected in the system.
The Fokker-Planck equation corresponding to this evolution is
$`_tP=`$ $`\{{\displaystyle \frac{1}{i+\eta }}{\displaystyle \frac{\delta }{\delta \psi }}[2{\displaystyle \frac{\delta H}{\delta \psi ^{}}}+{\displaystyle \frac{\mathrm{\Omega }}{i+\eta }}{\displaystyle \frac{\delta }{\delta \psi ^{}}}]`$ (92)
$`+{\displaystyle \frac{1}{i+\eta }}{\displaystyle \frac{\delta }{\delta \psi ^{}}}[2{\displaystyle \frac{\delta H}{\delta \psi }}+{\displaystyle \frac{\mathrm{\Omega }}{i+\eta }}{\displaystyle \frac{\delta }{\delta \psi }}]\}P`$
where $`\mathrm{\Omega }=2\eta T`$ as before. The system thermalizes to its canonical distribution
$`P_0=𝒩{\displaystyle D\psi D\psi ^{}\mathrm{exp}\left[\beta H\right]}.`$ (93)
The set of eigenvalues and eigenvectors for the harmonic problem are
$`E_N=2N\eta {\displaystyle \frac{k^2+m^2}{1+\eta ^2}},`$ (94)
$`F_N=H_N\left[\sqrt{\beta (k^2+m^2)\psi _k\psi _k^{}}\right].`$ (95)
The factor of 2 in $`E_N`$ relative to (59) accounts for 2 real fields in the complex quantity $`\psi `$, instead of one. Thus, the equilibration time scale is
$`t_{eq}{\displaystyle \frac{1+\eta ^2}{2\eta \left(k^2+m^2\right)}},`$ (96)
which again displays critical slowing down for $`k^20`$ and $`m^2(T)0`$.
It is clear that the overdamped limit is recovered for large $`\eta `$, as expected. We therefore see that in the overdamped limit all three models considered in this paper lead to the same characteristic time scale for equilibration as could be expected on general grounds. The differences arise for the short wavelength modes in the system and/or for small $`\eta `$.
An experimental realization of (88) would allow for the driving of the physical system across its phase transition and more generally to canonical equilibrium at an arbitrary temperature, by changing the intensity of the phase incoherent driving field $`\xi `$. If $`\xi `$ had a phase coherent component, in addition to the incoherent piece, the drive would make the system behave like an XY magnet at finite temperature in the presence of a magnetic field in the XY plane. This would allow for experiments in non-linear optics (where this effect is known as phase pinning) or atomic Bose-Einstein condensates to explore regimes similar to those realized in eg. high-Tc superconductors.
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# Néel ordered versus quantum disordered behavior in doped spin-Peierls and Haldane gap systems
## 1 Introduction
Doping quasi one dimensional (1D) antiferromagnets with a spin gap has become experimentally possible since the discovery of several inorganic quasi 1D oxides. One of these compounds is CuGeO<sub>3</sub> having a spin-Peierls transition at $`T_{SP}14`$ K . Below $`T_{SP}`$, the spin-phonon coupling induces a dimerization of the lattice, and the opening of a gap in the spin excitation spectrum. The Haldane gap in spin-$`1`$ chains in another example of a spin gap state in low dimensional magnets . Two inorganic spin-$`1`$ Haldane gap antiferromagnets have been discovered in the recent years: (i) PbNi<sub>2</sub>V<sub>2</sub>O<sub>8</sub> having a spin gap $`28`$ K . (ii) Y<sub>2</sub>BaNiO<sub>5</sub> having a spin gap $`100`$ K . The spin-Peierls compound CuGeO<sub>3</sub> and the two Nickel oxides PbNi<sub>2</sub>V<sub>2</sub>O<sub>8</sub> and Y<sub>2</sub>BaNiO<sub>5</sub> can be doped in a very controlled fashion. Substituting the magnetic Cu sites (having $`S=1/2`$) of the spin-Peierls compound CuGeO<sub>3</sub> with a variety of ions (Ni – a spin-1 ion –, Co – a spin-$`3/2`$ ion –, Zn or Mg – non magnetic ions –), or substituting the Ge sites with Si leads to the formation of an antiferromagnetic phase (AF) at low temperature. Moreover, in CuGeO<sub>3</sub>, there is AF long range order even with an extremely weak concentration of Zn impurities . On the other hand, the two Nickel oxides PbNi<sub>2</sub>V<sub>2</sub>O<sub>8</sub> and Y<sub>2</sub>BaNiO<sub>5</sub> have been the subject of an important experimental interest recently. It has been shown that substituting the spin-$`1`$ Ni sites of the PbNi<sub>2</sub>V<sub>2</sub>O<sub>8</sub> compound with Mg – a spin-$`0`$ ion – leads to AF long range order. In the Y<sub>2</sub>BaNiO<sub>5</sub> compound, the Ni sites can be substituted with Zn or Mg – non magnetic ions –. In this case, no sign of AF long range order has been reported, even at extremely low temperature . Instead, it has been found that the susceptibility has a power-law temperature dependence . Experiments therefore show that doping quasi 1D antiferromagnets with a spin gap can lead to very different situations: either antiferromagnetism, or a power-law susceptibility without AF ordering. The purpose of the present article is to describe these experimental observations in a unified theoretical framework, and provide a detailed theoretical analysis of the different phases of the model.
On the theoretical side, a lot of efforts have been devoted to understand the behavior of random spin chains. The theoretical tool usually used to study these models is the cluster renormalization group (RG) , which is a perturbation theory in the inverse of the strength of the strongest exchange in the chain, and leads to a certain number of exact results at low temperature because the exchange distribution becomes extremely broad. A lot of different models have been solved using this approach. For instance: the Ising chain in a transverse magnetic field , the random spin-$`1/2`$ chain , the spin-$`1/2`$ chain with random ferromagnetic and antiferromagnetic bonds , the dimerized chain with random bonds , the disordered Haldane gap chain . These studies have revealed that 1D disordered magnets can be controlled by several types of Griffiths phases: (i) the random singlet phase with a diverging susceptibility and algebraic correlation; (ii) the “weakly disordered” phase with a diverging susceptibility and short range correlations.
An important question is to understand the relation between the available experiments on quasi one dimensional oxides and the available theories of disordered 1D magnets. This type of approach followed in the present article incorporates realistic constraints such as interchain couplings and a finite temperature.
The starting point of such our description has been already established in previous works . In the doped spin-Peierls systems, non magnetic impurities generate solitonic spin-$`1/2`$ degrees of freedom distributed at random with a concentration $`x`$. The solitons are confined close to the impurities because of interchain interactions , and interact with the Hamiltonian
$$=\underset{i,j}{}J_{i,j}𝐒_i.𝐒_j.$$
(1)
The exchange between two spin-$`1/2`$ moments at positions $`(x_i,y_i)`$ and $`(x_j,y_j)`$ is mediated by virtual excitations of the gaped medium and therefore decays exponentially with distance:
$$J_{i,j}=()^{x_ix_j+y_iy_j}\mathrm{\Delta }\mathrm{exp}\left(\sqrt{\left(\frac{x_ix_j}{\xi _x}\right)^2+\left(\frac{y_iy_j}{\xi _y}\right)^2}\right),$$
(2)
where $`\xi _x`$ and $`\xi _y`$ are the correlation lengths in the direction of the chains and perpendicular to the chains respectively . The form Eq. 2 of the exchange incorporates a correlation length in the transverse direction shorter than in the longitudinal direction ($`\xi _y=\xi _x/10`$, and $`\xi _x10`$ in CuGeO<sub>3</sub> ). The exchange Eq. 2 is staggered because the dimerized pattern propagates staggered antiferromagnetic correlations. The Hamiltonian Eqs. 12 is therefore strongly disordered but unfrustrated.
Now, the model relevant to describe doping in a Haldane gap system is almost identical. It is well known that an impurity in a Haldane gap chain generates two “edge” spin-$`1/2`$ moments: one at the right and one at the left of the impurity site . The spin-1 chain can be thought in terms of a Valence Bond Solid (VBS) . Introducing a paramagnetic site breaks two VBS bonds, therefore resulting in two “edge” spin-$`1/2`$ moments. At energies far below the Haldane gap, only these edge moments are the relevant degrees of freedom (see Fig. 1).
The two edge moments in the same unit (i.e. generated by the same impurity) interact with a ferromagnetic exchange $`J_2`$ originating from the coupling to neighboring chains, with therefore the same order of magnitude as the interchain interaction: $`J_2J_{}`$ . The edge moments belonging to different units are coupled by the staggered exchange Eq. 2.
We find that, depending on the doping concentration and interchain interactions, the model Eqs. 12 has two regimes: a Néel ordered region and a quantum disordered region. In the Néel ordered region, relevant for CuGeO<sub>3</sub> and PbNi<sub>2</sub>V<sub>2</sub>O<sub>8</sub>, quantum mechanics plays little role, and we are lead to replace the spin variables in Eq. 12 by classical Ising spins. We propose here that this type of disordered Ising model is equivalent to another type of disordered Ising model, and solve the hierarchical lattice version of the latter. Using this treatment, we can compute the renormalized exchange distribution. In spite of a strongly disordered initial Hamiltonian (see Eqs. 12), we find that at large scale, the problem behaves as if it were non disordered. This appears to be consistent with susceptibility experiments showing a well-defined transition even with a small doping concentration .
In the quantum disordered region of the phase diagram, the physics is dominated by the formation of randomly distributed singlets. We show that the susceptibility has a power-law behavior, which turns out to be in agreement with existing experiments on Y<sub>2</sub>BaNiO<sub>5</sub> . We also find the existence of two distinct “quantum disordered” phases. The high temperature “quantum disordered” phase appears to have been observed experimentally in Y<sub>2</sub>BaNiO<sub>5</sub> . There appears to be another low temperature “quantum disordered” phase in which the edge spins in the same unit (see Fig. 1) are frozen into spin-$`1`$ objects. Finally, we show that even in the quantum disordered regime of the model, part of the quantum disordered behavior is already contained in the classical disordered magnet.
The article is organized as follows: the phases of the model Eq. 12 are discussed in section 2. We show the existence of two phases: a Néel ordered region and a quantum disordered region. The nature of these two phases is next discussed in details in sections 3 and 4. Concluding remarks are given in section 5.
## 2 Phases of the model
In this section, we use a phenomenological approach to derive the phase diagram of the model. The calculation of the relevant energy scales in the problem is based on the analysis of a two-spin model. There is a first temperature scale (being a fraction of the spin gap $`\mathrm{\Delta }`$) below which magnetic correlations start to develop inside the chains. There is a second energy scale $`T_{\mathrm{typ}}`$, equal to the typical exchange, associated to singlet formation. There is a third energy scale $`T_{\mathrm{Stoner}}`$ associated to long range AF ordering. The behavior of the model depends strongly on whether $`T_{\mathrm{Stoner}}`$ is larger or smaller than $`T_{\mathrm{typ}}`$.
### 2.1 Onset of magnetic correlations
Magnetic correlations start to appear inside the chains when the temperature is a fraction of the spin gap $`\mathrm{\Delta }`$. To show this, let us consider a simple model in which two spins at a distance $`l`$ are coupled antiferromagnetically: $`=J(l)𝐒_1.𝐒_2`$. We use an exchange decaying exponentially with distance (see Eq. 2): $`J(l)=\mathrm{\Delta }\mathrm{exp}(l/\xi )`$, and a Poisson bond length distribution $`𝒫(l)=x\mathrm{exp}(xl)`$. Rigorously, the spacing $`l`$ is a discrete quantity, distributed according to a geometrical distribution. However, the physics will turn out to be controlled by the large-$`l`$ behavior and it is legitimate to consider $`l`$ as a continuous variable, and replace the geometrical distribution by the Poisson distribution.
The internal energy of the two-spin model reads
$$U(T)=\frac{3}{4}x\xi \mathrm{\Delta }^{x\xi }T^{x\xi +1}_0^{\beta \mathrm{\Delta }}u^{x\xi }\frac{\mathrm{exp}(3u/4)\mathrm{exp}(u/4)}{\mathrm{exp}(3u/4)+3\mathrm{exp}(u/4)}𝑑u.$$
(3)
This expression can be expanded in $`T`$:
$$U(T)\frac{3}{4}\frac{x\xi \mathrm{\Delta }}{1+x\xi }+\frac{9}{2}x\xi \left(\frac{T}{\mathrm{\Delta }}\right)^{x\xi }T+\mathrm{}$$
(4)
Magnetic correlations start to appear in the two-spin model in the low temperature regime in which $`U(T)`$ is linear in $`T`$. This regime appears when $`T`$ is a fraction of $`\mathrm{\Delta }`$ (see Fig. 2).
### 2.2 Singlet formation
To discuss in what temperature range is the physics controlled by the quantum mechanical ground state (being a singlet), we need to calculate the probability $`𝒫_s(T)`$ to find the two spins in a singlet state at a finite temperature $`T=1/\beta `$. We have
$`𝒫_s(T)`$ $`=`$ $`{\displaystyle 𝑑l𝒫(l)\frac{\mathrm{exp}[3\beta J(l)/4]}{\mathrm{exp}[3\beta J(l)/4]+3\mathrm{exp}[\beta J(l)/4]}}`$ (5)
$`=`$ $`13x\xi \left({\displaystyle \frac{T}{\mathrm{\Delta }}}\right)^{x\xi }{\displaystyle _0^{\beta \mathrm{\Delta }}}u^{1+x\xi }{\displaystyle \frac{\mathrm{exp}(u/4)}{\mathrm{exp}(3u/4)+3\mathrm{exp}(u/4)}}𝑑u,`$ (6)
where we used the dimensionless parameter $`u=\beta J`$. The integral in Eq. 6 is dominated by the small exchanges and we have $`𝒫_s(T)1\kappa (T/\mathrm{\Delta })^{x\xi }`$, with $`\kappa `$ a numerical factor. We are lead to conclude that the ground state occupancy is close to unity below the energy scale $`T_{\mathrm{typ}}\mathrm{\Delta }\mathrm{exp}(1/(x\xi ))`$. $`T_{\mathrm{typ}}`$ is nothing but the typical exchange, already identified in a previous work on the 1D model .
The calculation of the ground state occupancy can be made even simpler by noticing that only the disorder configurations in which the exchange is larger than $`T`$ are in a singlet configuration. This leads to
$$𝒫_s(T)_0^{\xi \mathrm{ln}(\beta \mathrm{\Delta })}x\mathrm{exp}(xl)𝑑l1\left(\frac{T}{\mathrm{\Delta }}\right)^{x\xi }.$$
It is remarkable that the energy scale $`T_{\mathrm{typ}}`$ arising from the two-spin model is identical to the one obtained previously from the exact solution of the 1D effective model of the spin-Peierls chain . This shows that a model with only two spins contains already the relevant physics.
### 2.3 Néel ordered versus quantum disordered behavior
At low temperature, correlations between the chains induce a long range ordering of the spin system. The simplest phenomenological description of long range ordering is provided by a Stoner model, already considered in Ref. . In CuGeO<sub>3</sub>, there is a succession of three regimes (see Fig. 3-(a)): (i) a paramagnet at high temperature, (ii) intrachain correlations develop when the temperature is a fraction of $`\mathrm{\Delta }`$ (iii) interchain correlations give rise to long range antiferromagnetism below $`T_{\mathrm{Stoner}}=J_{}x\xi `$.
In PbNi<sub>2</sub>V<sub>2</sub>O<sub>8</sub>, the relevant parameters are $`\mathrm{\Delta }30`$ K, $`x0.02`$, $`J_{}1.1`$ K. We approximate the correlation length in the Haldane gap phase to be $`\xi 6`$. The true correlation length is expected to be larger than this value because PbNi<sub>2</sub>V<sub>2</sub>O<sub>8</sub> is close to a transition to an Ising ordered antiferromagnet . We find $`T_{\mathrm{Stoner}}=0.13`$ K, $`T_{\mathrm{typ}}=6`$ mK, showing that the same succession of regimes occur in PbNi<sub>2</sub>V<sub>2</sub>O<sub>8</sub> and CuGeO<sub>3</sub> (see Fig. 3-(a)). This is compatible with existing experiments in PbNi<sub>2</sub>V<sub>2</sub>O<sub>8</sub> .
Therefore, in CuGeO<sub>3</sub> and PbNi<sub>2</sub>V<sub>2</sub>O<sub>8</sub>, one has $`T_{\mathrm{Stoner}}T_{\mathrm{typ}}`$. This implies that singlet formation plays little role in the physics of the antiferromagnetic transition. The quantum two-spin model can then be well mimicked by the classical two-spin model. For instance, the internal energy of the classical and quantum two-spin models have an identical temperature dependence (see Fig. 2). This indicates that one might expect to obtain a reasonable description of antiferromagnetism in CuGeO<sub>3</sub> and PbNi<sub>2</sub>V<sub>2</sub>O<sub>8</sub> on the basis of a classical model, which we analyze in details in section 3.
Now, the situation is different in Y<sub>2</sub>BaNiO<sub>5</sub>, where one has $`\mathrm{\Delta }100`$ K, $`\xi 6`$, $`J_{}0.3`$ K, and $`x0.04`$ . We find $`T_{\mathrm{typ}}=1.6`$ K, and $`T_{\mathrm{Stoner}}=0.07`$ K. What is new compared to CuGeO<sub>3</sub> and PbNi<sub>2</sub>V<sub>2</sub>O<sub>8</sub> is that a well defined quantum disordered regime is present below $`T_{\mathrm{typ}}`$. In between $`J_{}`$ and $`T_{\mathrm{typ}}`$, there is singlet formation in the chain direction, and below $`J_{}`$, the singlets develop in the transverse direction. Below the energy scale $`T_{\mathrm{typ}}`$, the staggered susceptibility is well described by the one of the random spin-$`1`$ chain
$$\chi (T)xT^{\alpha 1}/T_{\mathrm{typ}}^\alpha .$$
(7)
Using a cluster RG calculation, we will calculate the susceptibility in section 4.2 and show that it has indeed a power-law temperature dependence. The Stoner criterion leads to the ordering temperature
$$\frac{T_{\mathrm{Stoner}}^{}}{T_{\mathrm{typ}}}=\left(\frac{J_{}x\xi }{T_{\mathrm{typ}}}\right)^{1/(1\alpha )}.$$
This shows that the quantum disordered model transits to a reentrant antiferromagnetic ground state below $`T_{\mathrm{Stoner}}^{}`$. The different regimes of the model are shown on Fig. 3-(b).
The existence of two classes of models is best summarized by calculating the ratio
$$\frac{T_{\mathrm{Stoner}}}{T_{\mathrm{typ}}}=\frac{J_{}}{\mathrm{\Delta }}x\xi e^{1/(x\xi )},$$
which can be smaller or larger than unity, controlling whether the model has a Néel transition at $`T_{\mathrm{Stoner}}`$ or is in a quantum disordered regime. These two behaviors, as well as a comparison between the model and existing experiments on CuGeO<sub>3</sub>, PbNi<sub>2</sub>V<sub>2</sub>O<sub>8</sub> and Y<sub>2</sub>BaNiO<sub>5</sub>, have been reported on the phase diagram on Fig. 4.
## 3 Nature of the antiferromagnetic transition
### 3.1 Motivation of the hierarchical lattice study
In low doping experiments in CuGeO<sub>3</sub>, Manabe et al. have measured the doping dependence of the Néel temperature, and found that the experimental data in the range $`x>0.1\%`$ could be well fitted by the behavior $`T_NA\mathrm{exp}(B/x)`$ . This suggests that there is no critical concentration associated to antiferromagnetism. As already pointed out in Refs. , the model Eqs. 12 shows sings of compatibility with these experiments. The main unsolved question raised by the experiments by Manabe et al. is to determine whether the maximum in the susceptibility is really the signature of an antiferromagnetic phase transition with a diverging staggered correlation length. In fact, the susceptibility experiments give no information about the existence / absence of a diverging correlation length, and there are are no neutron experiments with a concentration of impurities of order $`0.1\%`$. On the theoretical side, we used previously several approaches to describe the nature of the antiferromagnetic phase of the model Eqs. 12: a Stoner model , a decimation method, a cluster RG , and a Bethe-Peierls solution of the classical model . It appears that different treatments of the model have lead to different answers. For instance, there is a well defined transition in the Stoner criterion and the Bethe-Peierls treatments . On the contrary, there is no transition in the decimation method where the Hamiltonian Eqs. 12 is mapped onto a percolation model . A possible approach to this problem would be to generalize the work in Ref. : instead of considering the model Eqs. 12 with infinite range exponential interactions, it is possible to approximate the problem by considering the Voronoi lattice model with exponential interactions, in which each lattice site has a finite number of neighbors. This model is well suited for carrying out numerical simulations, and avoids the difficulty that the initial model Eqs. 12 has infinite range interactions.
Here, we would like to follow a different route and replace the original Hamiltonian Eqs. 12 by a simplified one. This is done by noticing that the essential feature of the Hamiltonian Eqs. 12 is that the exchange between two spins at a random position is distributed according to
$$𝒫(J)=\frac{x\xi }{\mathrm{\Delta }}\left(\frac{J}{\mathrm{\Delta }}\right)^{x\xi 1},$$
(8)
where we used the spacing distribution $`P(l)=x\mathrm{exp}(xl)`$ and the exchange $`J(l)=\mathrm{\Delta }\mathrm{exp}(l/\xi )`$. We replace the original model Eqs. 12 by another model in which the spins are on the sites of a regular square lattice, and have a random nearest neighbor exchange in the distribution (8). Because the square lattice model with the exchange distribution Eq. 8 and the original model Eqs. 12 are controlled by the same type of disorder, it is natural to conjecture that the two models have an identical physics. One way to study the square lattice model with the exchange distribution Eq. 8 would be to perform large scale Monte Carlo simulations. There is however a more direct way to handle the model, which consists in replacing the square lattice by a recursive hierarchical lattice (see Fig. 5), where the Migdal Kadanoff RG equations can be obtained in an exact form. We will show that a non trivial physics is going on in the hierarchical lattice model, which is an indication there is also a non trivial physics in the square lattice model with the exchange distribution Eq. 8.
### 3.2 Renormalization group equations
Let us derive the RG equations of the Ising model with the exchange distribution Eq. 8 on a hierarchical lattice. The partition function associated to the exchange configuration on Fig. 6-(a) reads
$$Z(\mathrm{\Sigma },\mathrm{\Sigma }^{})=\underset{\sigma ,\sigma ^{}}{}\mathrm{exp}\beta \left[J_1\mathrm{\Sigma }\sigma +J_2\mathrm{\Sigma }\sigma ^{}+J_3\mathrm{\Sigma }^{}\sigma +J_4\mathrm{\Sigma }^{}\sigma ^{}\right],$$
(9)
and we impose that Eq. 9 be identical to the partition function associated to the exchange configuration on Fig. 6-(b), up to a proportionality factor:
$$Z(\mathrm{\Sigma },\mathrm{\Sigma }^{})=𝒩\mathrm{exp}(\beta \stackrel{~}{J}\mathrm{\Sigma }\mathrm{\Sigma }^{}).$$
Using the relation
$$\stackrel{~}{J}=\frac{1}{2\beta }\mathrm{ln}\left[\frac{Z(+,+)}{Z(+,)}\right],$$
we find $`\stackrel{~}{J}=\stackrel{~}{J}_{13}+\stackrel{~}{J}_{24}`$, where
$$\stackrel{~}{J}_{13}=\frac{1}{2\beta }\mathrm{ln}\left[\frac{\mathrm{cosh}(\beta (J_1+J_3))}{\mathrm{cosh}(\beta (J_1J_3))}\right]\text{ , }\stackrel{~}{J}_{24}=\frac{1}{2\beta }\mathrm{ln}\left[\frac{\mathrm{cosh}(\beta (J_2+J_4))}{\mathrm{cosh}(\beta (J_2J_4))}\right].$$
(10)
Now that we have determined the renormalization of the exchanges, we iterate the exchange distribution. Noting $`P(J_1)`$$`P(J_4)`$ the distribution of the exchanges $`J_1`$$`J_4`$, and $`P(\stackrel{~}{J})`$ the distribution of the renormalized exchange $`\stackrel{~}{J}`$, we have
$$\stackrel{~}{P}(\stackrel{~}{J})=𝑑J_1𝑑J_2𝑑J_3𝑑J_4P(J_1)P(J_2)P(J_3)P(J_4)\delta \left[\stackrel{~}{J}\stackrel{~}{J}_{13}\stackrel{~}{J}_{24}\right].$$
(11)
It will be useful to change variables to the bond lengths $`l=\xi \mathrm{ln}(\mathrm{\Delta }/J)`$ and iterate the bond length distribution $`p(l)`$ instead of the exchange distribution $`P(J)`$.
To calculate numerically the iteration of the exchange distribution Eq. 11, we use a discrete bond length $`l=1,\mathrm{},N`$, and we introduce an upper cut-off for the bond length. This is valid if we can check that the RG flow does not depend on $`N`$.
### 3.3 Analysis of the RG flow
The RG flow of the model is shown on Fig. 7. We have shown on this figure the evolution of the average bond length $`l`$ and the width of the bond length distribution $`\sqrt{(ll)^2}`$. At low temperature, the average bond length renormalizes to zero (ordered phase) while at high temperature it renormalizes to infinity (paramagnetic phase). Therefore, there is a well-defined transition in the model. We have checked that the transition temperature is independent on the cut-off $`N`$ used to iterate the bond length distribution. The existence of a thermodynamic transition could have been anticipated on the basis of the simplest possible approximation of the RG flow (see Appendix A). What is less obvious is that, after a transient in the first RG iterations, the width of the exchange distribution becomes much smaller than the average exchange: in spite of a broad initial exchange distribution in which $`l=\sqrt{(ll)^2}`$, the system renormalizes to an almost disorder-free exchange distribution in which $`\sqrt{(ll)^2}l`$. Therefore, at large scale, the spin system looks ordered while inhomogeneities are visible only at small scale. This type of behavior may explain why there is a pronounced maximum in the temperature dependence of the susceptibility even at very low doping .
## 4 Nature of the quantum disordered region
Now, we would like to investigate the behavior of the model in the quantum disordered region, and compare it to experimental data on the Haldane gap compound Y<sub>2</sub>BaNiO<sub>5</sub>. More specifically, we would like to determine whether the behavior of the model is compatible with the susceptibility experiments by Payen et al. , who reported that the susceptibility of Y<sub>2</sub>BaNiO<sub>5</sub> has a power-law temperature dependence. In our opinion, it is an important question to determine which ingredients should be incorporated in the theoretical model to describe the existing experiments. A first strategy, followed by Batista et al. is to look for the “most realistic possible” model. The first step in this approach is to consider that the relevant Hamiltonian for Y<sub>2</sub>BaNiO<sub>5</sub> takes the form
$$=\underset{i}{}\{J𝐒_i.𝐒_{i+1}+D(S_i^z)^2+E[(S_i^x)^2(S_i^y)^2]\}.$$
(12)
The anisotropy parameters in Eq. 12 have been determined by fitting inelastic neutron scattering experiments . Next, a density matrix renormalization group (DMRG) method has been used to treat the Hamiltonian Eq. 12 in the presence of magnetic impurities. The authors of Ref. are then able to reproduce specific heat experiments, and also arrive to an agreement with susceptibility experiments . The main objection that one might be tempted to formulate is that the Hamiltonian of the spin-$`1`$ chain relevant for Y<sub>2</sub>BaNiO<sub>5</sub> contains already three adjustable parameters, and that two more additional gyromagnetic factors have been added to describe the susceptibility experiments . Here, it is proposed that the power law temperature dependence of the susceptibility is a generic feature of the quantum disordered region of the phase diagram, that can be explained in a model with only a minimal number of ingredients.
### 4.1 Two-spin model
As we already explained in the Introduction, in the low energy model relevant to describe the doped Haldane gap compound, each paramagnetic impurity generates a unit of two “edge” spin-$`1/2`$ moments (see Fig. 1). The two spin-$`1/2`$ moments in the same unit are coupled by a ferromagnetic exchange $`J_2`$ the magnitude of which is of order of interchain interactions . As shown on Fig. 1, there is a staggered exchange $`()^l\mathrm{exp}(l/\xi )`$ coupling two edge moments at a distance $`l`$.
Let us first consider the simplest model in which the impurities are assumed to cut the chain into finite segments: $`J_2=0`$. This is a valid model if the temperature is larger than $`J_2`$, or equivalently, than the strength of interchain interactions. Consider two edge spins at distance $`l`$, coupled with a Heisenberg Hamiltonian $`=J(l)𝐒_1.𝐒_2h(𝐒_1^z+𝐒_2^z)`$, and the exchange $`J(l)=\mathrm{\Delta }()^l\mathrm{exp}(l/\xi )`$, with $`𝒫(l_+)=x\mathrm{exp}(xl_+)`$ the distribution of even length segments and $`𝒫(l_{})=x\mathrm{exp}(xl_{})`$ the distribution of odd length segments. It is easy to calculate the average magnetization in a magnetic field $`h`$:
$$M(h,T)=\frac{1}{2}x_0^\mathrm{\Delta }𝒫(J)\left[M_+(h,T,J)+M_{}(h,T,J)\right]𝑑J,$$
(13)
with $`𝒫(|J|)=(x\xi /\mathrm{\Delta })(|J|/\mathrm{\Delta })^{1+x\xi }`$ the exchange distribution, and
$$M_+(h,T,J)=\frac{e^{\beta h}e^{\beta h}}{e^{\beta h}+e^{\beta h}+1+e^{\beta J}}\text{ , }M_{}(h,T,J)=\frac{e^{\beta h}e^{\beta h}}{e^{\beta h}+e^{\beta h}+1+e^{\beta J}}$$
(14)
the magnetization of the spins coupled by an antiferromagnetic or ferromagnetic exchange $`J`$ at a finite temperature $`T=1/\beta `$. To calculate the magnetization Eqs. 1314, it is convenient to integrate by parts:
$$M(h,T)=xA+x\left(\frac{T}{\mathrm{\Delta }}\right)^{x\xi }_0^{\mathrm{\Delta }/T}u^{x\xi }f(u,h/T),$$
(15)
with
$`A`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left[{\displaystyle \frac{e^{\beta h}e^{\beta h}}{e^{\beta h}+e^{\beta h}+1+e^{\mathrm{\Delta }/T}}}+{\displaystyle \frac{e^{\beta h}e^{\beta h}}{e^{\beta h}+e^{\beta h}+1+e^{\mathrm{\Delta }/T}}}\right]`$ (16)
$`f(u,{\displaystyle \frac{h}{T}})`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \frac{d}{du}}\left[{\displaystyle \frac{e^{\beta h}e^{\beta h}}{e^{\beta h}+e^{\beta h}+1+e^u}}+{\displaystyle \frac{e^{\beta h}e^{\beta h}}{e^{\beta h}+e^{\beta h}+1+e^u}}\right].`$ (17)
The resulting susceptibility is shown on Fig. 8, where it is visible that the product $`T\chi (T)T^\alpha `$ has a power-law behavior at low temperature, like what has been observed in experiments on Y<sub>2</sub>BaNiO<sub>5</sub> . The magnetization in an applied magnetic field is shown on Fig. 9, and is close to the experimental observation . This shows that this model with only two spins contains much of the physics as far as the temperature is above the interchain coupling $`J_{}`$.
### 4.2 Cluster RG
#### 4.2.1 Quantum disordered phase I
Now, let us consider the cluster RG of the model with a finite $`J_2`$, and first consider the behavior of the model when the temperature is larger than $`J_2`$. I refer the reader to Refs. for an explanation of the method, and just present here the results. First, if the temperature is above $`J_2`$, we find that the exponent of the power-law Curie constant $`T\chi (T)T^\alpha `$ is identical to the one of the two-spin model (see Fig 8). This is not a surprise because above $`J_2`$ the edge moments in the same unit remain uncoupled and the cluster RG contains the same physics as the two-spin model. Moreover, the two-spin model is an exact treatment while the cluster RG is approximate. The agreement between the two methods shows the validity of the cluster RG method.
#### 4.2.2 Quantum disordered phase II
Now, I consider the cluster RG of a single chain at a temperature smaller than $`J_2J_{}`$. The temperature dependence of the Curie constant is shown on Fig. 10. It is visible that the Curie constant increases strongly with decreasing the temperature below $`J_2`$. This is because when $`TJ_2`$, the survival spin-$`1/2`$ moments in the same unit are frozen into spin-$`1`$ moments. This phenomenon is not quantum in nature because it occurs also in the classical disordered model, which is analyzed in details below.
Once the freezing into spin-1 units has been done at the energy scale $`J_2`$, the resulting effective model is again the one of spin-$`1`$ objects with random exchanges, being ferromagnetic or antiferromagnetic. This results in the low temperature quantum disordered phase II on Fig. 10.
In the presence of a finite interchain coupling, one still expects the appearance of two types of quantum disordered regions. It is also expected that the quantum disordered phase II on Fig. 10 is a 3D random singlet state, with singlet formation between spins belonging to different chains. The cross-over to a 3D regime is however not expected to change the shape of the temperature dependence of the susceptibility (see Fig. 10).
### 4.3 Classical disordered model
Now, let us determine to what extend the physics of the classical magnet resembles the physics of the quantum magnet. The Ising chain Hamiltonian reads
$$=\underset{i}{}J_{i,i+1}\sigma _i\sigma _{i+1},$$
with the Ising variables $`\sigma _i`$ corresponding to the consecutive edge moments, and $`J_{i,i+1}`$ determined according to the rules on Fig. 1. The even bonds correspond to a ferromagnetic exchange $`J_2`$ and the odd bonds correspond to an exchange $`J(l)=\mathrm{\Delta }()^l\mathrm{exp}[l/\xi ]`$, with the spacing $`l`$ drawn in the Poisson distribution $`P(l)=x\mathrm{exp}(xl)`$.
As shown in Appendix B, the Curie constant crosses over from $`2x`$ at high temperature to $`4x`$ at low temperature, when the temperature decreases below $`J_2`$. Therefore, in the classical model, the product $`T\chi (T)`$ increases monotonically with a decreasing temperature, while the opposite is observed experimentally in Y<sub>2</sub>BaNiO<sub>5</sub> . This is not unexpected because the classical model cannot describe the gapless Haldane phases. Nevertheless, the increase in the susceptibility below $`J_2`$ is very reminiscent of the behavior of the quantum model in the same temperature range. This is because the freezing of the edge moments of the same unit in a ferromagnetic alignment occurs both in the classical and quantum models.
The scaling function of the magnetization can be calculated easily by iterating numerically the magnetization distribution Eq. 20 and using the relation
$$𝒫(M,h)=\frac{𝒫(M,h=0)e^{\beta hM}}{_M^{}𝒫(M^{},h=0)e^{\beta hM^{}}}$$
to obtain the magnetization distribution in a finite magnetic field. The magnetization takes the form $`M(h,T)=T^\gamma G(x,h/T)`$ in a given temperature range where the susceptibility can be approximated by $`\chi (T)T^{1\gamma }`$. The scaling function is shown on Fig. 12, where it is visible that it has qualitatively the correct behavior in spite of the exponent $`\gamma `$ having the wrong sign compared to experiments. Therefore, the shape of the scaling function of the magnetization does not appear to be a crucial test to the model.
## 5 Conclusions
To conclude, the present work was intented to describe doping a spin-Peierls and a Haldane gap state in a unified framework. We have first shown how the relevant energy scales in the problem could be calculated from the analysis of a two-spin model, which allowed to discuss the phases of the model as a function of the doping concentration and interchain interactions. In the relevant temperature window, there are two distinct regions depending on the doping concentration and interchain interactions: (i) an antiferromagnetic region; (ii) a quantum disordered region. remarkably, this type of phase diagram compares well with the known behavior of the spin-Peierls compound CuGeO<sub>3</sub> and the two Nickel oxides PbNi<sub>2</sub>V<sub>2</sub>O<sub>8</sub> and Y<sub>2</sub>BaNiO<sub>5</sub>.
Next, we used our approach to investigate in more details the two possible phases of the model. We have shown that the physics in the antiferromagnetic region of the phase diagram is classical in nature and therefore we were lead to study the corresponding Ising model. We have replaced the original Hamiltonian Eqs. 12 by another Hamiltonian having the same features, and presented the solution of the latter Hamiltonian on a hierarchical lattice structure. Interestingly, we find that the renormalized problem is non disordered. This is in agreement with the presence of a well-defined cusp associated to antiferromagnetism in the susceptibility .
In the “quantum disordered” region of the phase diagram, the physics is strongly controlled by quantum fluctuations. Already in a model with two spins only does the susceptibility have a power-law temperature dependence, very similar to the experimental observation . We have suggested that there is no need to introduce many coupling constants to reach an agreement between the model and experiments. Going beyond the level of a two-spin model, we have found the existence of another quantum disordered phase at low temperature. We have analyzed the behavior of the classical model and shown it contains already a physics relevant to the quantum model.
## Appendix A Projection of the RG flow on a trial exchange distribution
We would like to present the simplest possible approximation of the RG equations obtained in section 3.2 in which we project the RG flow onto the single parameter distribution $`p_n(l)=\delta (lL_n)`$. The initial distribution is obtained via the relation $`p_0(l)𝑑l=x\mathrm{exp}(xl)𝑑l`$, with $`p_0(l)=\delta (lL_0)`$. This leads to
$$L_0=\xi \mathrm{ln}\left(\frac{1+x\xi }{x\xi }\right).$$
(18)
Next, we start from the distribution $`p_n(l)`$, make one RG transformation Eq. 11, and impose that the iterated distribution has the same first moment as $`p_{n+1}(L)`$, from what we can determine the parameter $`L_{n+1}`$:
$$L_{n+1}=\xi \mathrm{ln}\frac{\beta \mathrm{\Delta }}{\mathrm{ln}\mathrm{cosh}\left[2\beta \mathrm{\Delta }\mathrm{exp}(L_n/\xi )\right]}.$$
(19)
Since we want to discuss the stability of the paramagnetic phase, we consider Eq. 19 in the limit of a large bond length: $`L_{n+1}\xi \mathrm{ln}(2\beta \mathrm{\Delta })+2L_n`$. In this limit, we find $`L_n=\xi \mathrm{ln}(2\beta \mathrm{\Delta })+\left(L_0\xi \mathrm{ln}(2\beta \mathrm{\Delta })\right)2^n.`$ Using Eq. 18, we obtain a phase transition at the temperature
$$T_c=\frac{2\mathrm{\Delta }x\xi }{1+x\xi }.$$
The transition temperature is in agreement with what has been found in previous works with different methods (the Stoner criterion and the Bethe-Peierls method ). As we show in the body of the article, the renormalized exchange distribution can be well approximated by the distribution $`p_n(l)=\delta (lL_n)`$ in the sense that the problem renormalizes to a non disordered one.
## Appendix B Solution of the classical analog of the doped spin-$`1`$ chain
We consider a finite chain with $`N`$ edge moments in which the end spin at site $`N`$ is frozen in the direction $`+`$, and note $`𝒫_N^+(M)`$ the corresponding magnetization distribution. We note $`x_i=\mathrm{exp}(\beta J_i)/[\mathrm{exp}(\beta J_i)+\mathrm{exp}(\beta J_i)]`$ the probability to find the spins $`\sigma _i`$ and $`\sigma _{i+1}`$ in an antiparallel alignment. We have
$$𝒫_{N+1}^+(M)=(1x_N)𝒫_N^+(M1)+x_N𝒫_N^{}(M1).$$
(20)
Using the relation $`𝒫_N^+(M)=𝒫_N^{}(M)`$, we get
$`M_{N+1}^+`$ $`=`$ $`(12x_N)M_N^++1`$ (21)
$`M^2_{N+1}`$ $`=`$ $`M^2_{N+1}+2(12x_N)M_N^++1.`$ (22)
These relations can be solved analytically. For this purpose, let us separate the even bonds coupled by the ferromagnetic exchange $`J_2`$ and note $`x_F=e^{\beta J_2}/[e^{\beta J_2}+e^{\beta J_2}]`$. The odd bonds are ferromagnetic or antiferromagnetic and we are lead to define
$$y=\underset{l=1}{\overset{+\mathrm{}}{}}𝒫(l)\frac{e^{\beta J(l)}}{e^{\beta J(l)}+e^{\beta J(l)}}.$$
The magnetization of a chain with an even number of sites $`N=2p`$ is found to be
$$M_{2p}^+=\frac{2(1x_F)}{1X}\left[1X^p\right],$$
(23)
with $`X=(12x_F)(12y)`$. The correlation length is $`\xi _T=2/\mathrm{ln}X`$. With an odd number of sites $`N=2p+1`$, the magnetization is
$$M_{2p+1}^+=\frac{2(1y)}{1X}\frac{X+12y}{1X}X^p.$$
(24)
The expressions of the first moment Eqs. 2324 are next used to solve for the second moment:
$$M^2_{2p}2p\left\{1+\frac{2(12y)(1x_F)}{1X}+\frac{2(12x_F)(1y)}{1X}\right\},$$
to leading order in the chain length $`2p`$. At high temperature, one has $`y1/2`$, $`x_F1/2`$ and $`X0`$, and therefore a susceptibility scaling like $`\chi 2x/T`$. At low temperature, one has $`x_F0`$, $`y1/2`$ and $`X0`$, and a susceptibility scaling like $`\chi 4x/T`$. Therefore, the Curie constant crosses over from the high temperature value $`2x`$ to $`4x`$ at low temperature, because the spin-$`1/2`$ moments in the same unit are frozen ferromagnetically at a temperature $`TJ_2`$ (see Fig. 11).
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# Moduli spaces of maximally supersymmetric solutions on noncommutative tori and noncommutative orbifolds
## 1 Introduction
The idea that String theory leads to some sort of fuzzy or noncommutative microscopic structure of space-time has been around for quite a while. A new boost to this idea emerged after the paper . It was shown in that paper that noncommutative tori arise as particular compactifications of M(atrix) theory () that is conjectured to be a nonperturbative definition of String theory (see \- and references therein for the subsequent development). Later in a number of important results concerning relations between String theory and noncommutative geometry were obtained. In particular the conditions under which noncommutative geometry arises within perturbative open string theory were clarified.
On the mathematics side noncommutative tori are the best studied examples of noncommutative spaces (see for a good overview). An important notion in noncommutative geometry is that of Morita equivalence that gives some equivalence relation between algebras of functions on noncommutative spaces. A striking result first observed in and later on proved rigorously in is that noncommutative world volume field theories living on noncommutative tori are invariant under duality transformations generated by Morita equivalences of noncommutative tori. This duality is directly related to T-duality of perturbative string theory compactifications (, ).
It seems to be interesting and important to study compactifications on other noncommutative spaces and see how much of noncommutative geometry techniques that proved to be useful for tori can be extended (see , , , for some work done in that direction). In paper we studied M(atrix) theory compactifications on noncommutative toroidal orbifolds. That paper primarily concentrated on two-dimensional $`𝐙_2`$ orbifolds. In the present paper we continue investigation of M(atrix) theory compactification on noncommutative spaces. The main topic of this paper is the structure of projective modules over noncommutative orbifolds that admit a constant curvature Yang-Mills field and moduli spaces of all such fields. Such modules and constant curvature connections on them describe configurations preserving half of the unbroken supersymmetries. All our results concern classical aspects only. However counting of quantum states with specified brane charges can be made in terms of a supersymmetric sigma model on the appropriate classical moduli space provided we have a sufficient number of supersymmetries.
In the commutative case configurations with vanishing $`SU\left(N\right)`$ part of the curvature and their interpretations in terms of D-branes were considered in a number of papers (see and references therein). D0 branes on toroidal $`𝐙_2`$ orbifolds were studied in , . Whenever our results can be compared with the commutative results we observe a complete agreement.
The paper is organized as follows. Section 2 contains some generalities on noncommutative tori and matrix theory as well as an explanation of our general strategy regarding a moduli space problem. In essence the novelty of our approach is in the following. In the usual approach one first fixes a module, i.e. a representation of algebra of functions on a noncommutative space, and then considers all constant curvature connections modulo gauge transformations. Instead we fix the connection $`_i`$ and then look at all equivalence classes of the torus representations compatible with that connection. If $`[_i,_j]=f_{ij}\mathrm{𝟏}`$ where $`f_{ij}`$ is a non-degenerate matrix then the algebra generated by $`_i`$ is isomorphic to Heisenberg algebra. We can use then the well known results about complete reducibility of Heisenberg algebra representations. The case of generic constant curvature connection can be reduced to the nondegenerate case by use of Morita equivalence provided the torus dimension is even. The Morita equivalence technique by itself is not directly relevant to the present paper; its discussion is relegated to the appendix. We restrict ourselves to the even dimensional case in this paper. However after appropriate modification our method also works for odd dimensions.
In section 3 we apply our general strategy to noncommutative tori and prove that the moduli space is isomorphic to $`\left(\stackrel{~}{T}^d\right)^r/S_r`$ where $`\stackrel{~}{T}^d`$ is a commutative $`d`$-dimensional torus and $`r`$ is a greatest common divisor of topological integers (D-brane charges) characterizing the module. For $`d=2`$ this result was proved by A. Connes and M. Rieffel in .
In section 4 we introduce noncommutative toroidal orbifolds and outline how the approach used for tori can be extended to the study of moduli space of equivariant constant curvature connections.
In sections 5 and 6 we study in detail $`𝐙_2`$ and $`𝐙_4`$ orbifolds respectively. For this cases we work out a general construction of a module that admits a constant curvature equivariant connection. In each case a module is built out of standard blocks that can be interpreted to describe D0 particles and various membranes stuck at exceptional points of the orbifold. The corresponding moduli spaces are proved to be isomorphic to $`(\stackrel{~}{T}^d/𝐙_2)^m)/S_m`$ and $`\left(\stackrel{~}{T}^d/𝐙_4\right)^n/S_n`$ respectively. Here $`m`$ and $`n`$ are some integers depending on topological numbers of the module. In the two-dimensional case $`𝐙_2`$ and $`𝐙_4`$ orbifolds of noncommutative tori were studied by S. Walters in a number of papers (). It would be interesting to calculate the topological numbers introduced in those papers for the modules we consider.
In section 7 we add scalar fields to the discussion and describe how Coulomb branches vary over the moduli space of constant curvature connections.
Finally the appendix contains general discussion of Morita equivalence for toroidal orbifolds and details of its application to $`𝐙_2`$ and $`𝐙_4`$ cases.
## 2 1/2 BPS configurations on noncommutative tori
We start this section by discussing some general aspects of constant curvature connections on noncommutative tori.
We define an algebra $`A_\theta `$ of smooth functions on a $`d`$-dimensional noncommutative torus in the following way. Let $`D𝐑^d`$ be a lattice $`D𝐙^d`$ and $`\theta _{ij}`$ be an antisymmetric $`d\times d`$ matrix. The algebra $`A_\theta `$ is an associative algebra whose elements are formal series
$$\underset{𝐧D}{}C\left(𝐧\right)U_𝐧$$
where $`C\left(𝐧\right)`$ are complex numbers and $`U_𝐧`$ satisfy the relations
$$U_𝐧U_𝐦=U_{𝐧+𝐦}e^{\pi in_j\theta ^{jk}m_k}.$$
(1)
Here the coefficients $`C\left(𝐧\right)`$ are assumed to decrease faster than any power of $`\left|𝐧\right|`$ as $`𝐧\mathrm{}`$. We will denote generators corresponding to standard basis vectors $`\left(n^i\right)_j=\delta _j^i`$ by $`U_i`$. Any $`U_𝐧`$ can be expressed as a product of $`U_i`$’s times a numerical phase factor.
The notion of a vector bundle (or rather that of a space of sections of a vector bundle) in noncommutative geometry is replaced by a projective module over $`T_\theta `$. Unless specified we will assumed that we work with a right module (i.e. we have a right action of $`T_\theta `$). Let $`L`$ be a $`d`$-dimensional commutative Lie algebra acting on $`T_\theta `$ by means of derivations
$$\delta _jU_𝐧=2\pi in_jU_𝐧.$$
(2)
A connection on a projective module $`E`$ over $`T_\theta `$ is defined in terms of operators $`_i:EE`$ satisfying a Leibniz rule
$$_i\left(ea\right)=\left(_ie\right)a+e\left(\delta _ia\right)$$
for any $`eE`$ and any $`aT_\theta `$.
It was first shown in () how (super)Yang-Mills theory on a noncommutative torus arises as a particular compactification of M(atrix) theory (, ). The Minkowski action functional of M(atrix) theory compactified on a noncommutative torus $`T_\theta `$ can be written as
$`S=`$ $`{\displaystyle \frac{V}{4g_{YM}^2}}\mathrm{Tr}\left(\left(F_{ij}+\varphi _{ij}\mathrm{𝟏}\right)^2+[_i,X_I]^2+[X_I,X_J]^2\right)`$ (3)
$`+{\displaystyle \frac{iV}{2g_{YM}^2}}\mathrm{Tr}\left(\overline{\psi }\mathrm{\Gamma }^i[_i,\psi ]+\overline{\psi }\mathrm{\Gamma }^I[X_I,\psi ]\right)`$
where $`_i`$, $`i=1,\mathrm{},d`$ is a connection on a $`T_\theta `$-module $`E`$, $`\psi `$ is a ten-dimensional Majorana-Weyl spinor taking values in the algebra $`End_{T_\theta }E`$ of endomorphisms of $`E`$, $`\varphi _{ij}`$ is an antisymmetric tensor, $`X_IEnd_{T_\theta }E`$, $`I=d+1,\mathrm{}10`$ are scalar fields, $`\mathrm{\Gamma }_\mu `$ are ten-dimensional gamma-matrices.
The action (3) is invariant under the following supersymmetry transformations
$`\delta _i={\displaystyle \frac{i}{2}}\overline{ϵ}\mathrm{\Gamma }_i\psi `$
$`\delta X_I={\displaystyle \frac{i}{2}}\overline{ϵ}\mathrm{\Gamma }_I\psi `$
$`\delta \psi =\stackrel{~}{ϵ}{\displaystyle \frac{1}{4}}\left(F_{ij}\mathrm{\Gamma }^{ij}+2[_i,X_I]\mathrm{\Gamma }^{iI}+[X_I,X_J]\mathrm{\Gamma }^{IJ}\right)ϵ`$ (4)
where $`ϵ`$, $`\stackrel{~}{ϵ}`$ are 10D Majorana-Weyl spinors parameterizing the transformation.
Classical configurations preserving 1/2 of these supersymmetries satisfy the equations
$`[_j,_k]=2\pi if_{jk}\mathrm{𝟏},`$
$`[_j,X_I]=0,[X_I,X_J]=0`$ (5)
where $`f_{ij}`$ are constants and $`\mathrm{𝟏}`$ is a unit endomorphism. We call solutions to the first equation modulo gauge transformations a Higgs branch of the 1/2-BPS moduli space. The whole moduli space of solutions to (2) can be viewed as a fibration over the Higgs branch. We will first describe the Higgs branch which is the moduli space of constant curvature connections and then take into account the scalar fields.
Let us outline here the strategy we will take addressing the moduli space problem. The complete set of equations that gives a module over a noncommutative torus and a constant curvature connection on it reads
$`U_jU_k=e^{2\pi i\theta _{jk}}U_kU_j,`$
$`[_j,U_k]=2\pi i\delta _{jk}U_k,`$
$`[_j,_k]=2\pi if_{jk}`$ (6)
If $`U_i`$, $`_j`$ are operators in Hilbert space $`E`$ satisfying these equations and $`Z:EE`$ is a unitary linear (over $`𝐂`$) operator, then $`ZU_iZ^1`$, $`Z_iZ^1`$ also solve (2). The usual approach to equations (2) is to fix generators $`U_i`$, i.e. fix a module, and then look for $`_i`$ satisfying the last two equations in (2). The transformations $`Z:EE`$ preserving $`U_i`$’s are unitary endomorphisms. We will take a different approach in this paper and fix the representation of $`_i`$. The gauge transformations in this picture are given by unitary operators $`Z`$ that commute with all $`_i`$. The moduli space is then a space of solutions $`U_i`$ to the first two equations in (2) modulo gauge transformations. Note that in this approach inequivalent modules, that in our case are modules that have distinct Chern (and/or other topological numbers in the case of orbifolds), are treated simultaneously. Fixing topological numbers then corresponds to choosing a connected component of the total moduli space of solutions to (2).
The (super)Yang-Mills theories (3) on Morita equivalent noncommutative tori $`T_\theta `$, $`T_{\widehat{\theta }}`$ are physically equivalent (see ). To any module $`E`$ over $`T_\theta `$ there corresponds a module $`\widehat{E}`$ over $`T_{\widehat{\theta }}`$ and to any connection $``$ on $`E`$ there corresponds a connection $`\widehat{}`$ and vice versa. The correspondence of connections respects gauge equivalence relation and maps a constant curvature connection into a constant curvature one. This means that moduli spaces of constant curvature connections are isomorphic for Morita equivalent tori. A crucial fact that we are going to exploit analyzing representations of (2) is that by use of Morita equivalence one can reduce the problem to the situation when $`f_{ij}`$ is a nondegenerate matrix (see Appendix).
## 3 Constant curvature connections on noncommutative tori
Let $`A_\theta `$ be a $`d=2g`$-dimensional noncommutative torus. Let $`E`$ be a projective module over $`A_\theta `$ admitting a constant curvature connection $`_i`$, $`i=1,\mathrm{},d`$
$$[_j,_k]=2\pi if_{jk}\mathrm{𝟏}$$
(7)
where $`f_{jk}`$ is a constant nondegenerate antisymmetric matrix with real entries. (We assume that $`_i`$ are antihermitian operators.) Up to normalizations algebra (7) is the Heisenberg algebra specified by canonical commutation relations for $`g`$ degrees of freedom. According to the Stone-von Neumann theorem there is a unique irreducible representation $``$ of algebra (7). We assume that $`E`$ can be decomposed into a direct sum of a finite number of irreducible components: $`E^N𝐂^N`$.
We fix the representation $``$ as follows. First let us bring the matrix $`f_{ij}`$ to a canonical block-diagonal form
$$\left(f_{ij}\right)=\left(\begin{array}{cccc}f_1ϵ& 0& \mathrm{}& 0\\ 0& f_2ϵ& \mathrm{}& 0\\ 0& 0& \mathrm{}& 0\\ 0& 0& \mathrm{}& f_gϵ\end{array}\right)$$
(8)
where
$$ϵ=\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)$$
(9)
is a $`2\times 2`$ matrix and $`f_i`$ are positive numbers.
Then we can define a representation space as $`L_2\left(𝐑^g\right)`$ and the operators $`_i`$ as
$$_j=\sqrt{f_{\left(j+1\right)/2}}_j,j\text{odd},_j=i\sqrt{f_{j/2}}x_{j1},j\text{even}$$
(10)
where $`_j`$, $`x_k`$, $`j,k=1,\mathrm{},g`$ are derivative and multiplication by $`x^k`$ operators acting on smooth functions $`f\left(x\right)L_2\left(𝐑^g\right)`$. An arbitrary representation of the torus generators $`U_i`$, $`i=1,\mathrm{},d`$ has the form $`U_i=U_i^{st}u_i`$ where $`U_i^{st}`$ is some standard representation satisfying
$$[_j,U_k^{st}]=2\pi i\delta _{jk}U_k^{st}$$
and $`u_i`$ is an $`N\times N`$ unitary matrix. This form of representation of $`U_i`$ follows from the irreducibility of representation $``$. A straightforward calculation shows that one can take $`U_i^{st}`$ to be
$$U_k^{st}=e^{\left(f^1\right)^{kl}_l}.$$
(11)
These operators satisfy
$$U_j^{st}U_k^{st}=e^{2\pi i\left(f^1\right)^{jk}}U_k^{st}U_j^{st}.$$
(12)
Since $`U_i=U_i^{st}u_i`$ must give a representation of a noncommutative torus it follows from (12) that so must do the operators $`u_i`$. But the last ones are finite-dimensional matrices so they can only represent a noncommutative torus whose noncommutativity matrix has rational entries, i.e. $`u_i`$’s have to satisfy
$$u_iu_j=e^{2\pi in^{ij}/N}u_ju_i$$
(13)
where $`N`$ is a positive integer and $`n^{ij}`$ is an integer valued antisymmetric matrix. Putting the formulas (12) and (13) together one finds that $`U_i`$’s give a representation of a noncommutative torus $`T_\theta `$ with
$$\theta ^{ij}=\left(f^1\right)^{ij}+n^{ij}/N.$$
(14)
It follows from the results obtained by M. Rieffel () that for finite $`N`$ (i.e. when $`E`$ decomposes into a finite number of irreducible components) the module $`E`$ endowed with $`U_i=U_i^{st}u_i`$ as above is a finitely generated projective module over $`T_\theta `$ with $`\theta `$ given in (14). Conversely one can show that the finiteness of $`N`$ is required by the condition of $`E`$ to be finitely generated and projective. (See for a detailed discussion of modules admitting a constant curvature connection.)
Topological numbers.
Let us calculate here the topological numbers of the modules constructed above. We assume here that the matrix $`\theta _{ij}`$ given in (14) has irrational entries. Then a projective module $`E`$ is uniquely characterized by an integral element $`\mu \left(E\right)`$ of the even part of Grassmann algebra $`\mathrm{\Lambda }^{even}\left(𝐑^d\right)`$. In order to calculate $`\mu \left(E\right)`$ we can use Elliot’s formula
$$\mu \left(E\right)=exp\left(\frac{1}{2}\frac{}{\alpha ^i}\theta ^{ij}\frac{}{\alpha ^j}\right)ch\left(E\right)$$
(15)
where
$$ch\left(E\right)=Dexp\left(\frac{1}{2\pi i}\alpha ^if_{ij}\alpha ^j\right)$$
(16)
and $`D`$ is a nonnegative real number that plays a role of the dimension of a module in noncommutative geometry. From (15) applying a Fourier transform we obtain
$`\mu \left(E\right)`$ $`=`$ $`D\mathrm{Pfaff}\left(f\right){\displaystyle 𝑑\beta exp\left(\frac{1}{2}\beta _i\left(\left(f^1\right)^{ij}+\theta ^{ij}\right)\beta _j+\alpha ^i\beta _i\right)}=`$ (17)
$`=`$ $`D\mathrm{Pfaff}\left(f\right){\displaystyle 𝑑\beta exp\left(\frac{1}{2}\beta _in^{ij}\beta _j/N+\alpha ^i\beta _i\right)}.`$
At this point it is convenient to assume that the matrix $`n_{ij}`$ is brought to a canonical block-diagonal form similar to (8) with integers $`n_i`$, $`i=1,\mathrm{},g`$ on the diagonal by means of an $`SL(d,𝐙)`$ transformation (this is always possible, see ). Then we can explicitly do the integration in (17) and obtain
$$\mu \left(E\right)=C\left(\frac{n_1}{N}+\alpha ^1\alpha ^2\right)\left(\frac{n_2}{N}+\alpha ^3\alpha ^4\right)\mathrm{}\left(\frac{n_g}{N}+\alpha ^{2g1}\alpha ^{2g}\right)$$
(18)
where $`C=D\mathrm{Pfaff}\left(f\right)`$ is a constant that can be determined by the requirement that $`\mu \left(E\right)`$ is an integral element of Grassmann algebra $`\mathrm{\Lambda }`$ (i.e. each coefficient is an integer). By looking at the term with a maximal number of $`\alpha `$’s in (18) we immediately realize that $`C`$ must be an integer. We will prove below that $`C=N`$.
Let us introduce numbers
$$N_i=\frac{N}{g.c.d.(N,n_i)},\stackrel{~}{N}_i=\frac{n_i}{g.c.d.(N,n_i)}$$
so that for each $`i=1,\mathrm{},g`$ the pair $`N_i`$, $`\stackrel{~}{N}_i`$ is relatively prime. Then we can rewrite (18) as
$$\mu \left(E\right)=\frac{C}{N_1N_2\mathrm{}N_g}\underset{i=1}{\overset{g}{}}\left(\stackrel{~}{N}_i+N_i\alpha ^{2i1}\alpha ^{2i}\right)\frac{C}{N_1N_2\mathrm{}N_g}\mu _0\left(E\right).$$
(19)
For any integral element $`\nu `$ of Grassmann algebra $`\mathrm{\Lambda }`$ let us introduce a number $`g.c.d.\left(\nu \right)`$ which is defined to be the largest integer $`k`$ such that $`\nu =k\nu ^{}`$ where $`\nu ^{}`$ is also integral. It is a simple task to prove by induction in $`g`$ that $`g.c.d.\left(\mu _0\left(E\right)\right)=1`$. Hence, $`C`$ must be an integer divisible by the product $`N_1N_2\mathrm{}N_g`$. Moreover $`C=g.c.d.\left(\mu \left(E\right)\right)N_1N_2\mathrm{}N_g`$. It is known (for example see ) that the dimension of an irreducible representation of the algebra (13) is equal to the product $`N_1N_2\mathrm{}N_g`$. Thus, necessarily this product divides $`N`$, i.e. $`N=N_1N_2\mathrm{}N_gN_0`$ where $`N_0`$ is an integer equal to the number of irreducible components in the representation $`𝐂^N`$ of the algebra (13). Evidently $`N_0`$ divides $`g.c.d.\left(\mu \left(E\right)\right)`$. We are going to show below that $`g.c.d.\left(\mu \left(E\right)\right)`$ cannot be bigger than $`N_0`$. This will prove that $`C=N`$.
Let us look at some particular examples of formula (18). If the matrix $`n^{ij}`$ is nondegenerate then $`\mu \left(E\right)`$ is a quadratic exponent:
$$\mu \left(E\right)=pexp\left(\frac{1}{2}\alpha ^i\left(n^1\right)_{ij}\alpha ^jN\right),p=N_0\stackrel{~}{N}_1\mathrm{}\stackrel{~}{N}_g$$
(20)
where $`p=N\mathrm{Pfaff}\left(n/N\right)`$ is written in a form where it is manifestly an integer. If $`n^{ij}`$ is degenerate then $`\mu \left(E\right)`$ is a so called generalized quadratic exponent (see and , Appendix D). For example if $`n^{ij}=0`$ for all $`i`$ and $`j`$ then we obtain from (17)
$$\mu \left(E\right)=N\alpha ^1\alpha ^2\mathrm{}\alpha ^d.$$
(21)
Moduli space of constant curvature connections in terms of irreps of rational n.c. tori.
In the previous subsection we showed that modules endowed with a constant curvature connection correspond to representations of matrix algebra (13). The residual gauge transformations preserving (10) correspond to $`N\times N`$ unitary transformations acting on the $`𝐂^N`$ factor of $`E`$. Thus, we see that the moduli space of constant curvature connections on a module with fixed $`(N,n^{ij})`$ (or fixed $`\mu \left(E\right)`$, which is the same) can be described as a space of inequivalent representations of the matrix algebra (13). The center of algebra (13) is spanned by elements $`u_𝐤`$ with $`𝐤D𝐙^d`$ satisfying $`k_in^{ij}m_j/N𝐙`$ for any $`𝐦D`$. Such elements correspond to a sublattice of $`D`$ that we denote $`D^{}`$ (the notation reveals the fact that $`D^{}`$ can be viewed as a dual lattice to $`D`$). To describe the center in a more explicit way it is convenient to choose the basis we used above, in which the matrix $`n^{ij}`$ is brought to a block-diagonal canonical form. In this basis generators of the center can be chosen to be elements $`\left(u_i\right)^{M_i}`$ where we set $`M_i=N_{\left(i+1\right)/2}`$, $`i`$-odd and $`M_i=N_{i/2}`$, $`i`$-even. Thus, in an irreducible representation $`\left(u_i\right)^{M_i}=\lambda _i𝐂^\times `$ are constants of absolute value 1.
Using the substitution
$$u_ic_iu_i$$
(22)
where $`c_i`$ are constants, $`\left|c_i\right|=1`$, we obtain an irreducible representation with $`\lambda _i^{}=\lambda _i\mathrm{}`$ with values of center generators $`\lambda _i`$ into the one with $`\lambda _i^{}=\lambda _ic_i^{M_i}`$. By means of this substitution one can transform any irrep into the one with $`\lambda _i=1`$. The last one corresponds to a representation of the algebra specified by relations (13) along with the relations $`\left(u_i\right)^{M_i}=1`$. This algebra has unique irreducible representation of dimension $`N_1N_2\mathrm{}N_g`$ (for example see ). Therefore, the space of irreducible representations of algebra (13) is described by means of $`d`$ complex numbers $`\lambda _i`$ with absolute value 1, i.e. is isomorphic to a (commutative) torus $`\stackrel{~}{T}^d𝐑^d/D^{}`$. We denote the corresponding irreps by $`E_\mathrm{\Lambda }`$, $`\mathrm{\Lambda }=(\lambda _1,\mathrm{},\lambda _d)`$. In general for any noncommutative torus $`T_\theta `$ one can construct a group $`L_\theta `$ of automorphisms isomorphic to a commutative torus of the same dimension by means of (22). This torus acts naturally on the space of unitary representations of $`T_\theta `$. If $`\theta `$ is rational we obtain a transitive action of this automorphism group on the space of irreducible representations. In this case one can consider $`L_\theta `$ as a finite covering of $`\stackrel{~}{T}^d`$.
Let us assume now that the space $`𝐂^N`$ is decomposed into irreducible representations of algebra (13)
$$𝐂^N=E_{\mathrm{\Lambda }_1}\mathrm{}E_{\mathrm{\Lambda }_{N_0}}.$$
(23)
Note that in the picture we are working with, gauge transformations are given by unitary linear operators acting on $`E`$ that commute with all $`_i`$’s that is by unitary $`N\times N`$ matrices. The matrices representing central elements are diagonalized in the basis specified by decomposition (23). There are residual gauge transformations corresponding to permutations of diagonal entries. Thus, we see that in general the moduli space is isomorphic to $`\left(\stackrel{~}{T}^d\right)^{N_0}/S_{N_0}`$. As it was noted in the previous subsection $`N_0`$ divides $`g.c.d.\left(\mu \left(E\right)\right)`$. On the other hand as we know from , any module $`E`$ over a noncommutative torus $`T_\theta `$ admitting a constant curvature connection $`_i`$ can be represented as a direct sum of $`k`$ identical modules $`E=E^{}\mathrm{}E^{}`$ with $`k=g.c.d.\left(\mu \left(E\right)\right)`$. This implies that the moduli space of constant curvature connections necessarily contains a subset isomorphic to $`\left(\stackrel{~}{T}^d\right)^k/S_k`$. Thus, on dimensional grounds we conclude that $`k=g.c.d.\left(\mu \left(E\right)\right)=N_0`$.
## 4 General toroidal orbifolds
In this section we will consider how the considerations above can be generalized to the case of toroidal orbifolds. We will work out in detail the particular cases of $`𝐙_2`$ and $`𝐙_4`$ orbifolds in the subsequent sections. Let $`D𝐑^d`$ be a $`d`$-dimensional lattice embedded in $`𝐑^d`$ and let $`G`$ be a finite group acting on $`𝐑^d`$ by linear transformations mapping the lattice $`D`$ to itself. For an element $`gG`$ we will denote the corresponding representation matrix $`R_i^j\left(g\right)`$. One can write down constraints describing compactification of M(atrix) theory on the orbifold $`T^d/G`$, where $`T^d=𝐑^d/D`$:
$$X_j+\delta _{ij}2\pi \mathrm{𝟏}=U_i^1X_jU_i,$$
(24)
$$X_I=U_i^1X_IU_i\psi _\alpha =U_i^1\psi _\alpha U_i,$$
(25)
$$R_i^j\left(g\right)X_j=W^1\left(g\right)X_iW\left(g\right),$$
(26)
$$\mathrm{\Lambda }_{\alpha \beta }\left(g\right)\psi _\beta =W^1\left(g\right)\psi _\alpha W\left(g\right),X_I=W^1\left(g\right)X_IW\left(g\right).$$
(27)
Here $`i,j=1,\mathrm{},d`$ are indices for directions along the torus , $`I=d+1,\mathrm{},9`$ is an index corresponding to the transverse directions, $`\alpha `$ is a spinor index; $`\mathrm{\Lambda }_{\alpha \beta }\left(g\right)`$ is the matrix of spinor representation of $`G`$ obeying $`\mathrm{\Lambda }^{}\left(g\right)\mathrm{\Gamma }^i\mathrm{\Lambda }\left(g\right)=R_{ij}\left(g\right)\mathrm{\Gamma }_j`$; $`U_i`$, $`W\left(g\right)`$ \- unitary operators. One can check that the quantities $`U_iU_jU_i^1U_j^1`$ commute with all $`X_i`$, $`X_I`$, and $`\psi _\alpha `$. It is natural to set them to be proportional to the identity operator. This gives us defining relations of a noncommutative torus
$$U_jU_k=e^{2\pi i\theta _{jk}}U_kU_j.$$
It is convenient to work with linear generators $`U_𝐧`$ that can be expressed in terms of products of $`U_i`$. One can further check that expressions $`W\left(gh\right)W^1\left(g\right)W^1\left(h\right)`$ and $`W^1\left(g\right)U_𝐧W\left(g\right)U_{R^1\left(g\right)𝐧}^1`$ also commute with all fields $`X_i`$, $`X_I`$, $`\psi _\alpha `$. We assume that these expressions are proportional to the identity operator. This leads us to the following relations
$`W\left(g\right)W\left(h\right)=W\left(gh\right)e^{i\varphi (g,h)},`$
$`W^1\left(g\right)U_𝐧W\left(g\right)=U_{R^1\left(g\right)𝐧}e^{i\chi (𝐧,g)}`$ (28)
where $`\varphi (g,h)`$, $`\chi (𝐧,g)`$ are constants. The first equation means that operators $`W\left(g\right)`$ furnish a projective representation of $`G`$. It follows from these equations that the matrix $`\theta `$ is invariant under the group action $`R\left(g\right)`$. Below we will confine ourselves to the case of vanishing cocycles $`\varphi `$ and $`\chi `$. We refer the reader to papers , for a discussion of cases when cocycle $`\varphi `$ does not vanish. (Note that for cyclic groups both cocycles are always trivial. This means that they can be absorbed into redefinitions of generators.)
One can define an algebra of functions on a noncommutative orbifold as an algebra generated by the operators $`U_𝐧`$ and $`W\left(g\right)`$ satisfying (1) and
$$W^1\left(g\right)U_𝐧W\left(g\right)=U_{R^1\left(g\right)𝐧},$$
(29)
$$W\left(g\right)W\left(h\right)=W\left(gh\right).$$
(30)
These equations mean that the algebra at hand is a crossed product $`T_\theta _RG`$. Again we remark here that allowing central extensions gives a more general case of twisted crossed products. In this paper we will concentrate on the untwisted case.
The algebra $`T_\theta _RG`$ can be equipped by an involution $``$ by setting $`U_𝐧^{}=U_𝐧`$, $`W^{}\left(g\right)=W\left(g\right)`$. This makes it possible to embed these algebras into a general theory of $`C^{}`$ algebras. A projective module over an orbifold can be considered as a projective module $`E`$ over $`T_\theta `$ equipped with operators $`W\left(g\right)`$, $`gG`$ satisfying (29). The equations (24), (26) mean that $`X_i`$ specifies a $`G`$\- equivariant connection on $`E`$, i.e. $`_i`$ is a $`T_\theta `$ connection satisfying
$$R_i^j\left(g\right)_j=W^1\left(g\right)_iW\left(g\right).$$
(31)
The fields $`X_I`$ are endomorphisms of $`E`$, commuting both with $`U_𝐧`$ and $`W\left(g\right)`$ and the spinor fields $`\psi _\alpha `$ can be called equivariant spinors.
Let us comment here on the supersymmetry of these compactifications. The surviving supersymmetry transformations are transformations (2) corresponding to invariant spinors $`ϵ`$, $`\stackrel{~}{ϵ}`$, i.e. the ones satisfying $`\mathrm{\Lambda }\left(g\right)ϵ=ϵ`$. For $`d=4`$, $`6`$ this equation has a nontrivial solution provided the representation $`R\left(g\right)`$ lies within an $`SU\left(2\right)`$, $`SU\left(4\right)`$ subgroup respectively. The possible finite groups $`G`$ that can be embedded in this way are well known. Those include the examples of $`𝐙_2`$ and $`𝐙_4`$ four-dimensional orbifolds to be considered below. In those cases when the supersymmetry is not broken completely equivariant connections of constant curvature correspond to configurations preserving half of the unbroken supersymmetries.
Now we restrict ourselves to modules admitting a constant curvature equivariant connection (7). All the steps of the analysis made above for the case of tori leading to the decomposition (23) can be repeated in a straightforward way. We should add now to that analysis operators $`W\left(g\right)`$. Equation (31) implies that the curvature tensor $`f_{ij}`$ is invariant under the action of $`G`$. As above we fix a representation of the Heisenberg algebra (7) in the form (10). Then a connection $`_i^g=R_i^j\left(g\right)_j`$ gives a representation of the same Heisenberg algebra (7). By uniqueness of irreducible representation $``$ there exists a set of operators $`W^{st}\left(g\right):`$ satisfying (31). It follows from the definition (11) of $`U_j^{st}`$ that $`W^{st}\left(g\right)`$ and $`U_j^{st}`$ commute as in (29). This implies that a general set of operators $`W\left(g\right):EE`$ has a form
$$W\left(g\right)=W^{st}\left(g\right)w\left(g\right)$$
(32)
where $`w\left(g\right)`$ are $`N\times N`$ matrices satisfying
$$w^1\left(g\right)u_𝐧w\left(g\right)=u_{R^1\left(g\right)𝐧}.$$
(33)
Here $`u_𝐧`$ are (linear) generators of the rational torus (13). In general relation (31) only implies that the operators $`W^{st}\left(g\right)`$ form a projective representation of $`G`$. Then, the commutation relations for $`w\left(g\right)`$ are twisted by an opposite cocycle. The problem of describing moduli space of equivariant constant curvature connections now boils down to the study of irreducible representations of the matrix algebra generated by $`u_i`$, $`w\left(g\right)`$.
As it was explained above an irreducible representation of a rational torus is labeled by a point of a commutative torus $`\stackrel{~}{T}^d=𝐑^d/D^{}`$. The lattice $`D^{}`$ is a sublattice of $`D`$ which is preserved by $`G`$. Therefore $`G`$ acts on the torus $`\stackrel{~}{T}^d`$. We denote this action by $`R^{}\left(g\right)`$. It follows from (33) that an irreducible representation $`E_\mathrm{\Lambda }`$ from decomposition (23) is mapped by $`w\left(g\right)`$ into $`E_{R^{}\left(g\right)\mathrm{\Lambda }}`$. Therefore, the set of $`\mathrm{\Lambda }_i`$ in the decomposition (23) has to be invariant under the action $`R^{}\left(g\right)`$. The irreducible representations of the algebra generated by $`u_i`$, $`w\left(g\right)`$ can be grouped according to the orbits of $`G`$ action on $`\stackrel{~}{T}^d`$. For a generic point one has an orbit of length $`\left|G\right|`$. There are also exceptional points that include fixed points and points whose orbits length is a nontrivial divisor of $`\left|G\right|`$. If the exceptional points are isolated the corresponding representations can be interpreted in terms of branes that are stuck at the exceptional points. Below we will consider in detail the structure of the aforementioned representations for the orbifold groups $`𝐙_2`$ and $`𝐙_4`$.
## 5 Noncommutative $`𝐙_\mathrm{𝟐}`$ orbifolds
As in the case of tori we start by fixing the representation of connection (10). In addition to representation of the torus generators (11), (13) now we need to represent the $`𝐙_\mathrm{𝟐}`$ generator $`W`$. As in the case of $`U_i`$’s a generic representation of $`W`$ has the form $`W=W^{st}w`$ where $`W^{st}`$ is some standard representation satisfying
$$W^{st}_iW^{st}=_i,W^{st}U_i^{st}W^{st}=\left(U_i^{st}\right)^1,\left(W^{st}\right)^2=1$$
and $`w`$ is a $`N\times N`$ matrix that obeys the relations
$$wu_iw=u_i^1,w^2=1.$$
(34)
It is easy to check that one can take
$$W^{st}:f_i\left(x\right)f_i\left(x\right)$$
where $`f_i\left(x\right)EL_2\left(𝐑^d\right)𝐂^N`$ ($`i=1,\mathrm{},N`$). As before consider the decomposition of $`𝐂^N`$ into irreducible representations of the algebra (13) given by (23). The operator $`w`$ maps $`E_\mathrm{\Lambda }E_{\mathrm{\Lambda }^1}`$ where $`\mathrm{\Lambda }^1(\lambda _1^1,\mathrm{},\lambda _d^1)`$. Thus, the space $`𝐂^N`$ carries a representation of the algebra specified by (13) and (34) only if the set of $`\left\{\mathrm{\Lambda }_i,i=1,\mathrm{},k\right\}`$ in (23) is invariant under the inversion $`\mathrm{\Lambda }_i\mathrm{\Lambda }_i^1`$. We see that decomposition (23) can contain summands in the form of couples $`\stackrel{~}{E}_\mathrm{\Lambda }E_\mathrm{\Lambda }E_{\mathrm{\Lambda }^1}`$ and possible exceptional representations $`E_{\mathrm{\Lambda }_ϵ}`$ with $`\mathrm{\Lambda }_ϵ=(ϵ_1,\mathrm{},ϵ_d)`$, $`ϵ_i=\pm 1`$, i.e. we have
$$𝐂^N=\left(\stackrel{~}{E}_{\mathrm{\Lambda }_1}\mathrm{}\stackrel{~}{E}_{\mathrm{\Lambda }_r}\right)\underset{ϵ}{}E_{\mathrm{\Lambda }_ϵ}^{\tau \left(ϵ\right)}$$
(35)
where $`\tau \left(ϵ\right)`$ are nonnegative integers. This decomposition can be chosen in such a way that the matrix $`w`$ has a block diagonal form with $`r`$ blocks $`w_i:\stackrel{~}{E}_{\mathrm{\Lambda }_i}\stackrel{~}{E}_{\mathrm{\Lambda }_i}`$ of the form
$$w_i=\left(\begin{array}{cc}0& \stackrel{~}{w}_i\\ \stackrel{~}{w}_i^1& 0\end{array}\right)$$
where $`\stackrel{~}{w}_i:E_{\mathrm{\Lambda }_i}E_{\mathrm{\Lambda }_i^1}`$ and a number of blocks $`w_ϵ:E_{\mathrm{\Lambda }_ϵ}E_{\mathrm{\Lambda }_ϵ}`$.
Let us first look at the components specified by pairs $`(E_{\mathrm{\Lambda }_ϵ},w_ϵ)`$. It follows from the irreducibility of the representation $`E_{\mathrm{\Lambda }_ϵ}`$ and the fact that $`w_ϵ^2=1`$ that $`w_ϵ`$ is defined up to an overall sign. Let $`w_ϵ^{\left(0\right)}`$ be a standard choice of $`w_ϵ`$. Denote by $`F_ϵ^\pm `$ the representation of the algebra (13), (34) specified by $`(E_{\mathrm{\Lambda }_ϵ},\pm w_ϵ^{\left(0\right)})`$.
As for the blocks $`(\stackrel{~}{E}_{\mathrm{\Lambda }_i},w_i)`$ we first note that due to the irreducibility of $`E_{\mathrm{\Lambda }_i}`$ and $`E_{\mathrm{\Lambda }_i^1}`$ the operator $`\stackrel{~}{w}_i`$ is defined up to a constant factor. This constant factor can be gauged away by conjugating $`w_i`$ with a suitable rescaling transformation
$$\left(\begin{array}{cc}\mathrm{𝟏}\mu _1& 0\\ 0& \mathrm{𝟏}\mu _2\end{array}\right)$$
where $`\mu _1`$, $`\mu _2`$ are complex numbers of absolute value 1 (these rescalings are allowed gauge transformations). Thus, we get a single representation of the algebra specified by $`(u_i,w)`$ that we denote $`F_{\mathrm{\Lambda }_i}`$. This representation is irreducible except the cases when $`\mathrm{\Lambda }_i=\mathrm{\Lambda }_ϵ`$ for some $`ϵ`$. Then we have $`\stackrel{~}{w}_i=\pm w_ϵ^{\left(0\right)}`$ and one can readily check that $`F_{\mathrm{\Lambda }_ϵ}F_ϵ^+F_ϵ^{}`$. This permits one to decompose $`𝐂^N`$ into the components
$$𝐂^N=\left(F_{\mathrm{\Lambda }_1}\mathrm{}F_{\mathrm{\Lambda }_r}\right)\underset{ϵ}{}\left(F_ϵ^{\eta _ϵ}\right)^{\tau \left(ϵ\right)}$$
(36)
where $`\eta _ϵ=\pm `$, $`\tau \left(ϵ\right)`$ are nonnegative integers specifying the multiplicities with which the corresponding representations enter the decomposition. Note that the set of numbers $`\mathrm{\Lambda }_i`$, $`r`$, $`\tau \left(ϵ\right)`$ is uniquely determined by the given $`B_\theta =A_\theta 𝐙_2`$ module $`E`$ and an equivariant constant curvature connection $`_i`$ on it. The residual gauge transformations preserving the decomposition (36) can be represented as compositions of transpositions acting inside each $`F_{\mathrm{\Lambda }_i}`$ block and sending $`\mathrm{\Lambda }_i\mathrm{\Lambda }_i^1`$ and permutations of different $`F_{\mathrm{\Lambda }_i}`$ blocks. Thus, the moduli space of equivariant constant curvature connections is isomorphic to $`\left(T^d/𝐙_\mathrm{𝟐}\right)^r/S_r`$ where $`r`$ is some integer. Using the relation
$$g.c.d.\left(\mu \left(E\right)\right)=2r+\underset{ϵ}{}\tau \left(ϵ\right)$$
one can find an estimate from above on $`r`$: $`r\left[\frac{g.c.d.\left(\mu \right)}{2}\right]`$.
Let us comment here briefly on the invariance of the above results under Morita equivalence. It follows from the definition that the mapping of modules and connections induced by (complete) Morita equivalence relation preserves the gauge equivalence relation and maps constant curvature connections into constant curvature connections. Hence the moduli spaces over the modules related by Morita equivalence are isomorphic. In particular the dimension $`r`$ of the moduli space is preserved under Morita equivalences. As the number $`g.c.d.\left(\mu \left(E\right)\right)`$ is also preserved we see that $`_ϵ\tau \left(ϵ\right)=g.c.d.\left(\mu \left(E\right)\right)2r`$ is preserved under Morita equivalence. In fact the set of pairs $`(\tau \left(ϵ\right),\eta _ϵ)`$ can only get permuted under Morita equivalence transformations. One can show that for modules admitting a constant curvature connections the set of parameters $`\left\{(\tau \left(ϵ\right),\eta _ϵ)\right\}`$ is in a one-to-one correspondence with additional topological numbers characterizing the module (see ).
## 6 $`𝐙_4`$ orbifolds
In this section we will consider the case of toroidal $`𝐙_4`$ orbifolds. For definiteness let us concentrate on the four-dimensional case. It is not hard to extend the results we obtain to other even-dimensional $`𝐙_4`$ toroidal orbifolds. For $`d=4`$ without loss of generality one can define a $`𝐙_4`$ action on $`𝐑^4`$ by $`\rho :(x,y)(y,x)`$. Here we assume that $`𝐑^4`$ has a product structure $`𝐑^2\times 𝐑^2`$ and $`x`$, $`y`$ belong to the first and second $`𝐑^2`$ factor respectively. This action preserves the orthogonal lattice $`D𝐙^4𝐑^4`$ and thus descends to an action on the torus $`𝐑^4/D`$. We can consider a noncommutative four-torus $`T_\theta `$ constructed by means of lattice $`D`$ and an antisymmetric two-form $`\theta _{ij}`$ that is assumed to be decomposed into a $`2\times 2`$ block form relative to the $`𝐑^2\times 𝐑^2`$ product structure. One can easily see that such a form is invariant with respect to the above defined $`𝐙_4`$ action. Thus, we can consider a noncommutative toroidal orbifold $`T_\theta _\rho 𝐙_4`$.
Let $`E`$ be a projective module over this orbifold and let $`Y`$ denotes a representation of the $`𝐙_4`$ generator. Applying the general construction of section … we arrive at the decomposition (23) and a representation of a matrix algebra generated by $`u_i`$ and $`y`$, satisfying (13) and
$$y^4=1,y^1u_𝐧y=u_{\rho ^1𝐧}.$$
(37)
The total representation space $`𝐂^N`$ splits into irreducible representations of this matrix algebra. To describe those we first classify the orbits of the $`𝐙_4`$ action $`\rho ^{}`$ on the dual torus $`\stackrel{~}{T}^4𝐑^d/D^{}`$. For a generic point of $`\stackrel{~}{T}^4`$ the orbit consists of 4 distinct points. There are 16 exceptional points. Those include the four fixed points: $`(0,0;0,0)`$, $`(0,1/2;0,1/2)`$, $`(1/2,0;1/2;0)`$, and $`(1/2,1/2,1/2,1/2)`$, and 12 points whose orbit is of length 2. The last ones have coordinates $`0`$ or $`1/2`$ and they complete the $`𝐙_4`$ action fixed points written above to the whole set of 16 $`𝐙_2`$ action fixed points (the $`𝐙_2`$ action is specified by $`\rho ^2`$). Hence, we have 4 orbits of order 1 and 6 orbits of order 2. Denote the fixed points $`\mathrm{\Lambda }_ϵ`$, $`ϵ=1,\mathrm{},4`$ and the pairs of points whose orbits are of order 2 by $`(\mathrm{\Lambda }_\nu ^{},\mathrm{\Lambda }_\nu ^{\prime \prime })`$, $`\nu =1,\mathrm{},6`$.
Let us first consider the representations corresponding to generic points. Denote
$$\stackrel{~}{E}_\mathrm{\Lambda }=E_\mathrm{\Lambda }E_{\rho ^{}\mathrm{\Lambda }}E_{\left(\rho ^{}\right)^2\mathrm{\Lambda }}E_{\left(\rho ^{}\right)^3\mathrm{\Lambda }}.$$
Relative to this decomposition the generator $`y`$ can be written in a block form as
$$y=\left(\begin{array}{cccc}0& y_{12}& 0& 0\\ 0& 0& y_{23}& 0\\ 0& 0& 0& y_{34}\\ y_{41}& 0& 0& 0\end{array}\right)$$
where $`y_{ij}:E_{\left(\rho ^{}\right)^{i1}\mathrm{\Lambda }}E_{\left(\rho ^{}\right)^i\mathrm{\Lambda }}`$. It follows from irreducibility of representations $`E_\mathrm{\Lambda }`$ that the blocks $`y_{ij}`$ are determined uniquely up to constant factors. The last ones can be gauged away by gauge transformations that multiply each of the $`E_{\left(\rho ^{}\right)^n\mathrm{\Lambda }}`$ components of $`\stackrel{~}{E}_\mathrm{\Lambda }`$ by a constant. Thus, we obtain a representation $`F_\mathrm{\Lambda }`$ of the matrix algebra (13), (37) (we hope that using the same notation as the one used in the previous section when studying the $`𝐙_2`$ case will not cause any confusion).
Next let us look at representations labeled by pairs of exceptional points $`(\mathrm{\Lambda }_\nu ^{},\mathrm{\Lambda }_\nu ^{\prime \prime })`$. The representation of rational torus is $`\widehat{E}_{\mathrm{\Lambda }_\nu }=E_{\mathrm{\Lambda }_\nu ^{}}E_{\mathrm{\Lambda }_\nu ^{\prime \prime }}`$. A general form of the generator $`y`$ is
$$y=\left(\begin{array}{cc}0& \mu _1y_1\\ \mu _2y_2& 0\end{array}\right)$$
where $`y_1:E_{\mathrm{\Lambda }_\nu ^{\prime \prime }}E_{\mathrm{\Lambda }_\nu ^{}}`$, $`y_2:E_{\mathrm{\Lambda }_\nu ^{}}E_{\mathrm{\Lambda }_\nu ^{\prime \prime }}`$ are fixed and $`\mu _1`$, $`\mu _2`$ are constants satisfying $`\left(\mu _1\mu _2\right)^2=1`$. Using gauge transformations one can bring $`y`$ to one of the two forms specified by $`\mu _1=\mu _2=1`$ and $`\mu _1=\mu _2=1`$. Therefore, we have two inequivalent representations of the algebra (13), (37) denoted $`G_\nu ^\pm `$ with $`\pm `$ standing for the sign of the product $`\mu _1\mu _2`$.
Finally, consider the fixed points $`\mathrm{\Lambda }_ϵ`$. Since each $`E_{\mathrm{\Lambda }_ϵ}`$ is an irreducible representation of a rational torus the operator $`y`$ acting on it is defined uniquely up to a multiplication by a 4-th root of unity $`\xi _k=exp\left(2\pi ik/4\right)`$, $`k=0,1,2,3`$ . Thus, for every fixed point we have 4 different representations $`F_ϵ^k`$, $`k=0,\mathrm{},3`$.
The generic representation $`F_\mathrm{\Lambda }`$ is irreducible unless $`\mathrm{\Lambda }`$ hits one of the exceptional points. If it hits a fixed point the representation splits as $`F_{\mathrm{\Lambda }_ϵ}_{k=0}^3F_ϵ^k`$. In this case the representation of $`𝐙_4`$ is a tensor product of a regular representation with a representation that acts on a single copy of $`E_{\mathrm{\Lambda }_ϵ}`$. If $`\mathrm{\Lambda }`$ coincides with one of $`\mathrm{\Lambda }_\nu ^{}`$ or $`\mathrm{\Lambda }_\nu ^{\prime \prime }`$ the corresponding representation splits as $`F_{\mathrm{\Lambda }_\nu ^{}}F_{\mathrm{\Lambda }_\nu ^{\prime \prime }}G_\nu ^+G_\nu ^{}`$. Using this equivalences we can decompose a general representation as
$$𝐂^N=\left(F_{\mathrm{\Lambda }_1}\mathrm{}F_{\mathrm{\Lambda }_r}\right)\underset{\nu }{}\left(G_\nu ^{\zeta _\nu }\right)^{s\left(\nu \right)}\underset{ϵ}{}\left(\left(F_ϵ^{\eta _ϵ^1}\right)^{\tau ^1\left(ϵ\right)}\left(F_ϵ^{\eta _ϵ^2}\right)^{\tau ^2\left(ϵ\right)}\left(F_ϵ^{\eta _ϵ^3}\right)^{\tau ^3\left(ϵ\right)}\right)$$
(38)
where $`\zeta _\nu =\pm `$; $`\eta _ϵ^i`$, $`i=1,2,3`$ are a triple of distinct integers from 0 to 4; $`\tau ^i\left(ϵ\right)`$ and $`s\left(\nu \right)`$ are nonnegative integers standing for the multiplicities of modules. The numbers $`r`$, $`\tau ^i\left(ϵ\right)`$, $`s\left(\nu \right)`$, $`\zeta _\nu =\pm `$, $`\eta _ϵ^i`$ are uniquely determined by a given module. In is straightforward to generalize considerations of the previous section to show that the moduli space of equivariant constant curvature connections is isomorphic to $`\left(\stackrel{~}{T}^d/𝐙_4\right)^r/S_r`$ where $`r`$ is some integer such that $`r\left[\frac{g.c.d.\left(\mu \right)}{4}\right]`$. The last restriction follows from the relation
$$g.c.d.\left(\mu \right)=4r+2\left(\underset{\nu }{}s\left(\nu \right)\right)+\underset{ϵ}{}\underset{i=1}{\overset{3}{}}\tau ^i\left(ϵ\right).$$
## 7 Coulomb branches of the moduli space
Once we described moduli space of constant curvature connections it is not hard to add scalar fields to the discussion. Let us first consider the case of tori. The equations for scalars that we have are
$$[_i,X_I]=0,[X_I,X_J]=0,[U_i,X_I]=0.$$
(39)
Here $`I=1,\mathrm{},9d`$. For the fixed representation of $`_i`$ (10) the first equation in (39) implies that $`X_I`$ are constant $`N\times N`$ matrices. It follows then from the last two equations that the matrices $`X_I`$, $`u_i`$ can be simultaneously brought to the form when $`X_I`$’s are diagonal and $`u_i`$’s are block diagonal corresponding to the decomposition (23). Moreover the $`X_I`$’s are constant on each irreducible component $`E_{\mathrm{\Lambda }_i}`$. Quotienting over $`S_{N_0}`$ residual gauge transformations gives us Coulomb branches isomorphic to $`\left(𝐑^{9d}\right)^{N_0}/S_{N_0}`$ that matches with a dual (Morita equivalent) description of this system as a system of $`N_0=g.c.d.\left(\mu \left(E\right)\right)`$ D0-branes.
For the case of a $`𝐙_2`$ orbifold in addition to equations (39) we have the condition
$$[W,X_I]=0.$$
(40)
This implies that the matrices $`X_I`$ commute with the matrix $`w`$. Thus, the matrices $`u_i`$, $`w`$, and $`X_I`$’s can be simultaneously brought to the form that corresponds to the decomposition (35) and $`X_I`$ are in a block diagonal form with $`r`$ blocks $`x_I^{\left(i\right)}:F_{\mathrm{\Lambda }_i,\eta _i}F_{\mathrm{\Lambda }_i,\eta _i}`$, $`i=1,\mathrm{},r`$ and blocks of the form $`x_I^ϵ:F_ϵF_ϵ`$. The last ones are necessarily constants due to irreducibility of $`E_{\mathrm{\Lambda }_i}`$. If none of the points $`\mathrm{\Lambda }_i`$ coincides with one of the fixed points $`\mathrm{\Lambda }_ϵ`$ we obtain Coulomb branches of the form
$$\left(𝐑^{\left(9d\right)}\right)^r/S_r\times \underset{ϵ}{}\left(\left(𝐑^{\left(9d\right)}\right)^{\tau \left(ϵ\right)}/S_{\tau \left(ϵ\right)}\right).$$
If $`\mathrm{\Lambda }_{i_k}=\mathrm{\Lambda }_ϵ`$ for some subset of indices $`i_k`$, $`k=1,\mathrm{},p`$ and for some $`ϵ`$, each representation $`F_{\lambda _{i_k},\eta }`$ breaks into a direct sum of two representations $`F_ϵ`$ and $`F_ϵ^{}`$ and instead of the factor $`\left(𝐑^{\left(9d\right)}\right)^r/S_r\times \left(\left(𝐑^{\left(9d\right)}\right)^{\tau \left(ϵ\right)}/S_{\tau \left(ϵ\right)}\right)`$ we get a factor of
$$\left(𝐑^{\left(9d\right)}\right)^{rp}/S_{rp}\times \left(𝐑^{\left(9d\right)}\right)^p/S_p\times \left(𝐑^{\left(9d\right)}\right)^{\tau \left(ϵ\right)+p}/S_{\tau \left(ϵ\right)+p}.$$
(41)
This picture has an interpretation suggested in in terms of splitting of a D0 particle into a membrane-antimembrane pair that occurs once the D0 particle hits a fixed point. In terms of this interpretation (41) corresponds to $`p`$ membranes (or antimembranes depending on the value of $`\eta _ϵ`$) and $`p+\tau \left(ϵ\right)`$ antimembranes (membranes) sitting at the fixed point $`\mathrm{\Lambda }_ϵ`$.
Now let us look at the $`𝐙_4`$ case. Again we have a set of $`9d`$ $`N\times N`$ matrices $`X_I`$. These matrices commute between themselves and with the matrix $`y`$. This leads us to a block decomposition of each of the $`X_I`$ relative to the decomposition (38) in which $`X_I`$ is constant on every irreducible representation of the algebra generated by $`u_i`$, $`y`$. Thus, for a generic point in the moduli space of constant curvature connection the Coulomb branch is
$$\left(𝐑^{\left(9d\right)}\right)^r/S_r\times \underset{\nu }{}\left(𝐑^{\left(9d\right)}\right)^{s\left(\nu \right)}/S_{s\left(\nu \right)}\underset{ϵ}{}\underset{i=1}{\overset{3}{}}\left(𝐑^{\left(9d\right)}\right)^{\tau ^i\left(ϵ\right)}/S_{\tau ^i\left(ϵ\right)}.$$
If one of $`\mathrm{\Lambda }_i`$ in (38) coincides with one of exceptional points $`\mathrm{\Lambda }_\nu ^{}`$, $`\mathrm{\Lambda }_\nu ^{\prime \prime }`$ or $`\mathrm{\Lambda }_ϵ`$ representations $`F_{\mathrm{\Lambda }_i}`$ split accordingly as $`F_{\mathrm{\Lambda }_\nu ^{}}F_{\mathrm{\Lambda }_\nu ^{\prime \prime }}G_\nu ^+G_\nu ^{}`$, or $`F_{\mathrm{\Lambda }_ϵ}_{k=0}^3F_ϵ^k`$. In that case the block decomposition of $`X_I`$ is different and the answer is best described using brane terminology. In general the moduli space coincides with the one describing a system consisting of a number of free D0 particles, four different types of membranes that are stuck at fixed points and two types of couples of membranes sitting at points $`(\mathrm{\Lambda }_\nu ^{},\mathrm{\Lambda }_\nu ^{\prime \prime })`$. It seems to be natural to identify different types of membranes with a discrete $`B`$-field flux carried by them.
Acknowledgements
We are grateful to J. Blum, M. Rieffel and C. Schweigert for stimulating discussions.
## Appendix. Morita equivalence of toroidal orbifolds
In this appendix we will show how Morita equivalences known for noncommutative tori can be extended to noncommutative toroidal orbifolds. We will give a general construction for arbitrary orbifold group $`G`$ and then specialize to the $`𝐙_2`$ and $`𝐙_4`$ cases.
We start with a reminding of some basic definitions concerning Morita equivalence (see for details). Let $`A`$ and $`\widehat{A}`$ be two involutive associative algebras. A $`(A,\widehat{A})`$-bimodule $`P`$ is said to establish a Morita equivalence between $`A`$ and $`\widehat{A}`$. This means that the projective bimodule $`P`$ obeys the conditions
$$\overline{P}_AP\widehat{A},P_{\widehat{A}}\overline{P}A$$
(42)
where $`\overline{P}`$ is a $`(\widehat{A},A)`$-bimodule that is complex conjugated to $`P`$ that means that $`\overline{P}`$ consists of elements of $`P`$ and multiplications are defined as $`\widehat{a}\left(e\right):=\left(e\right)\widehat{a}^{}`$, $`\left(e\right)a:=a^{}\left(e\right)`$ where multiplications on the right hand sides are those defined for bimodule $`P`$. The algebras $`A`$ and $`\widehat{A}`$ are said to be Morita equivalent if such $`P`$ satisfying (42) exists The bimodule $`P`$ determines a one to one correspondence between $`A`$-modules and $`\widehat{A}`$-modules by the rule
$$\widehat{E}=E_AP,E=\widehat{E}_{\widehat{A}}\overline{P}.$$
For the case of noncommutative tori one can define a notion of complete Morita equivalence () that allows one to transport connections between modules $`E`$ and $`\widehat{E}`$. Let us remind here the basic definitions. Let $`\delta _j`$, $`j=1,\mathrm{},d`$ be a set of derivations of $`A_\theta `$ specified by their action on generators
$$\delta _jU_𝐧=2\pi in_jU_𝐧.$$
A connection on a projective module $`E`$ over $`A_\theta `$ can be defined as a set of operators $`_i:EE`$ satisfying a Leibniz rule
$$_i\left(ea\right)=\left(_ie\right)a+e\left(\delta _ia\right)$$
for any $`eE`$ and any $`aT_\theta `$. We say that $`(A_\theta ,A_{\widehat{\theta }})`$ Morita equivalence bimodule $`P`$ establishes a complete Morita equivalence if it is endowed with operators $`_i^P`$ that determine a constant curvature connection simultaneously with respect to $`A_\theta `$ and $`A_{\widehat{\theta }}`$, i.e. satisfy
$`_i^P\left(ae\right)=a_i^Pe+\left(\delta _ia\right)e,`$
$`_i^P\left(e\widehat{a}\right)=\left(_i^Pe\right)\widehat{a}+e\widehat{\delta }_i\widehat{a},`$
$`[_i^P,_j^P]=\sigma _{ij}\mathrm{𝟏}.`$ (43)
Here $`\delta _i`$ and $`\widehat{\delta }_j`$ are standard derivations on $`A_\theta `$ and $`A_{\widehat{\theta }}`$ respectively. Sometimes for brevity we will omit the word Morita in the term (complete) Morita equivalence bimodule. If $`P`$ is a complete $`(A_\theta ,A_{\widehat{\theta }})`$ equivalence bimodule then there exists a correspondence between connections on $`E`$ and connections on $`\widehat{E}`$. An operator $`_i1+1_i^P`$ on $`E_\mathrm{C}P`$ descends to a connection $`\widehat{}_\alpha `$ on $`\widehat{E}=E_{A_\theta }P`$. The curvatures of $`\widehat{}_i`$ and $`_i`$ are connected by the formula $`F_{ij}^\widehat{}=\widehat{F}_{ij}^{}+\mathrm{𝟏}\sigma _{ij}`$.
Given a $`(A_\theta ,A_{\widehat{\theta }})`$ equivalence bimodule $`P`$ there is a possibility of extension of the equivalence relation that it defines to modules over orbifolds $`A_\theta _RG`$, $`A_{\widehat{\theta }}_RG`$. We will first describe a general construction and then discuss for which bimodules $`P`$ it exists.
Assume that $`P`$ is equipped with a set of operators $`W^P\left(g\right)`$, $`gG`$ satisfying
$`W^P\left(g\right)W^P\left(h\right)=W\left(hg\right),`$
$`\left(W^P\left(g\right)\right)^1U_𝐧W^P\left(g\right)=U_{R𝐧},\left(W^P\left(g\right)\right)^1\widehat{U}_𝐧W^P\left(g\right)=\widehat{U}_{R𝐧}`$ (44)
where $`U_𝐧`$ and $`\widehat{U}_𝐧`$ stand for actions of $`A_\theta `$ and $`A_{\widehat{\theta }}`$ generators respectively. The first equation in (Appendix. Morita equivalence of toroidal orbifolds) means that operators $`W^P\left(g\right)`$ give a right action of the group $`G`$ on $`E`$ that is just a choice of conventions. If $`F`$ is a right module over $`A_\theta _RG`$ specified by a $`A_\theta `$-module $`E`$ and operators $`W\left(g\right)`$ acting on it one defines a right $`A_{\widehat{\theta }}_RG`$ module $`\widehat{F}`$ as a $`A_{\widehat{\theta }}`$-module $`\widehat{E}=E\times _{A_\theta }P`$ equipped with operators $`\widehat{W}\left(g\right):=W\left(g\right)W^P\left(g\right)`$. Analogously one defines a mapping in the opposite direction.
Given a pair $`(P,\left\{W^P\left(g\right);gG\right\})`$ as above one can construct a true $`(A_\theta _RG,A_{\widehat{\theta }}_RG)`$ Morita equivalence bimodule in the sense of general definition given at the beginning of this section. The construction goes as follows. To shorten notations denote $`B_\theta =A_\theta _RG`$, $`B_{\widehat{\theta }}=A_{\widehat{\theta }}_RG`$. Elements of the $`B_\theta ,B_{\widehat{\theta }})`$ bimodule that we denote $`Q`$ are pairs $`p_\mathrm{C}_gc\left(g\right)g`$ where $`pP`$ and $`_gc\left(g\right)g`$ are formal linear combinations of elements of the group $`G`$ with complex coefficients. Multiplication by complex numbers on $`Q`$ is defined in the obvious way. We further define left and right actions of generators of $`B_\theta `$ and $`B_{\widehat{\theta }}`$ as
$`U_𝐧\left(pg\right)=\left(U_{R^1\left(g\right)𝐧}p\right)g,W\left(h\right)\left(pg\right)=p\left(hg\right),`$
$`\left(pg\right)\widehat{U}_𝐧=\left(p\widehat{U}_𝐧\right)g,\left(pg\right)\widehat{W}\left(h\right)=\left(pW^P\left(h\right)\right)\left(gh\right).`$ (45)
One can check that $`Q`$ satisfies the defining properties of Morita equivalence bimodule, i.e.
$$\overline{Q}_{B_\theta }QB_{\widehat{\theta }},Q_{B_{\widehat{\theta }}}\overline{Q}B_\theta .$$
Namely, let fix us the isomorphisms $`I:\overline{P}_{A_\theta }PA_{\widehat{\theta }}`$, $`\overline{I}:P_{A_{\widehat{\theta }}}\overline{P}A_\theta `$, then one can define a mapping $`J:\overline{Q}_{B_\theta }QB_{\widehat{\theta }}`$ as
$$J\left(\left(\overline{p}\overline{g}\right)_{B_\theta }\left(pg\right)\right)=I\left(\overline{p}\left(pW^P\left(g^1\overline{g}\right)\right)\right)\left(\overline{g}^1g\right).$$
It is easy to check that $`J`$ is an isomorphism of the corresponding $`(B_{\widehat{\theta }},B_{\widehat{\theta }})`$ bimodules. An isomorphism $`\overline{J}:Q_{B_{\widehat{\theta }}}\overline{Q}B_\theta `$ can be constructed in a similar way.
We further need to construct a correspondence between equivariant connections on $`B_\theta `$ and $`B_{\widehat{\theta }}`$ modules $`F`$, $`\widehat{F}`$. We call an $`A_\theta `$-connection $`_i`$ defined on a $`B_\theta `$-module $`E`$ equivariant if it satisfies a constraint
$$W^1\left(g\right)_iW\left(g\right)=R_i^j\left(g\right)_j.$$
(46)
Now let us assume that we are given a triple $`(P,^P,W^P\left(g\right))`$ where $`P`$ is a $`(A_\theta ,A_{\widehat{\theta }})`$ Morita equivalence bimodule, operators $`W^P\left(g\right)`$ acting on $`P`$ satisfy (Appendix. Morita equivalence of toroidal orbifolds), $`_i^P`$ satisfies (Appendix. Morita equivalence of toroidal orbifolds) and an additional equivariance constraint (46). The couple $`(P,^P)`$ establishes a complete Morita equivalence of the algebras $`A_\theta `$. It is simple to check that due to condition (46) a mapping of $`A_\theta `$\- connections that is defined by this couple preserves the equivariance condition. We will say that a triple $`(P,^P,W^P\left(g\right))`$ specifies a $`(B_\theta ,B_{\widehat{\theta }})`$ complete equivalence. If a triple $`(P,_\alpha ^P,W^P\left(g\right))`$ specifies a $`(B_\theta ,B_\theta ^{})`$ complete equivalence and a triple $`(P^{},_\alpha ^P^{},W^P^{}\left(g\right))`$ specifies a $`(B_\theta ^{},B_{\theta ^{\prime \prime }})`$ complete equivalence, then the tensor product $`P^{\prime \prime }=P_{B_\theta ^{}}P^{}`$ along with the appropriate connection $`_\alpha ^{P^{\prime \prime }}`$ and involution $`W^{P^{\prime \prime }}`$ determines $`(B_\theta ,B_{\theta ^{\prime \prime }})`$ complete equivalence. This means that we can consider a groupoid of equivalences.
Let us discuss now the existence of triples $`(P,^P,W^P\left(g\right))`$. First we would like to note that not every (complete) Morita equivalence of tori can be lifted to orbifolds. There is an obvious constraint that the matrix $`\theta `$ should stay invariant under the orbifold group action. Morita equivalence of $`A_\theta `$ algebras is governed by the group $`SO(d,d|𝐙)`$. An element $`gSO(d,d|𝐙)`$ can be represented by a $`2d\times 2d`$ block matrix
$$g=\left(\begin{array}{cc}M& N\\ R& S\end{array}\right)$$
and the action on $`\theta `$ is given by a fractional transformation
$$g:\theta \left(M\theta +N\right)\left(R\theta +S\right)^1.$$
(47)
The orbifold group $`G`$ naturally acts on $`\theta `$ as
$$h:\theta R^t\left(h\right)\theta R\left(h\right),hG.$$
The condition that two actions commute specifies a certain subgroup of $`SO(d,d|𝐙)`$ and can be formalized as follows. We can embed both groups into $`O(d,d|𝐑)`$. For the first group the embedding is obvious, for the orbifold group we embed an element specified by the matrix $`R\left(g\right)`$ into
$$R\left(g\right)\left(\begin{array}{cc}R^t\left(g\right)& 0\\ 0& R^t\left(g^1\right)\end{array}\right).$$
Now the commutation condition has a precise meaning. We proceed to the construction of $`(P,^P,W^P\left(g\right))`$ triples. The group $`SO(d,d|𝐙)`$ is generated by the following transformations. The first type of transformations is a subgroup $`SL(d,𝐙)`$ of modular transformations. The second type consists of shifts $`\theta \theta +N`$ where $`N`$ is an arbitrary antisymmetric matrix with integer entries. To generate the whole $`SO(d,d|𝐙)`$ one has to add a “flip” transformation $`\sigma `$ that inverts a $`2\times 2`$ nondegenerate block in the matrix $`\theta `$. Namely, without loss of generality we can assume that $`\theta `$ has a block form
$$\theta =\left(\begin{array}{cc}\theta _{11}& \theta _{12}\\ \theta _{21}& \theta _{22}\end{array}\right).$$
where $`\theta _{11}`$ is a $`2p\times 2p`$ nondegenerate matrix. Then a flip $`\sigma _{2p}`$ sends $`\theta `$ into
$$\sigma _2\left(\theta \right)=\left(\begin{array}{cc}\theta _{11}^1& \theta _{11}^1\theta _{12}\\ \theta _{21}\theta _{11}^1& \theta _{22}\theta _{21}\theta _{11}^1\theta _{12}\end{array}\right).$$
(48)
One can check that modular transformations, shifts and the flip $`\sigma _2`$ generate the whole $`SO(d,d|𝐙)`$. For any such transformation that commutes with $`G`$ one can construct a triple $`(P,^P,W^P\left(g\right))`$. Let $`ASL(d,𝐙)`$ and $`A`$ commute with $`G`$. The corresponding $`(A_\theta ,A_{A^t\theta A}`$ equivalence bimodule consists of elements $`aA_\theta `$ and the actions of generators are defined as
$$U_𝐧^l\left(a\right)=U_𝐧a,\left(a\right)U_𝐧^r=aU_{A𝐧}.$$
(49)
where to avoid confusion we denoted by $`U_𝐧`$ elements of $`A_\theta `$, and by $`U_𝐧^l`$, $`U_𝐧^r`$ left and right actions of the corresponding tori. This bimodule can be completed to a triple by adding the following operators $`_i^P`$ and $`W^P\left(g\right)`$
$$W^P\left(g\right)\left(a\right)=\rho _g\left(a\right),_i^P\left(a\right)=\delta _i\left(a\right)$$
(50)
where $`\rho _g`$ stands for automorphisms of algebra $`A_\theta `$ induced by the action $`R\left(g\right)`$ on the lattice.
Let $`N=\left(N^{ij}\right)`$ be an antisymmetric $`G`$-invariant matrix with integer entries. The corresponding $`(A_\theta ,A_{\theta +N})`$ equivalence bimodule consists of elements $`aA_\theta `$ with tori actions defined as
$$U_𝐧^l\left(a\right)=U_𝐧a,\left(a\right)U_𝐧^r=aU_𝐧\left(1\right)^{_{i<j}n_iN^{ij}n_j}.$$
(51)
The operators $`_i^P`$, $`W^P\left(g\right)`$ are the same as in (49), (50).
Finally we need to define a triple corresponding to a G-invariant flip (48). Although it is not hard if only somewhat lengthy to describe triples corresponding to generic $`G`$ invariant flips of the type $`\sigma _2`$ we will confine ourselves to the “total” flip that inverts the whole matrix $`\theta `$ (provided the last one is invertible). The formulas for that case are most succinct and elegant. Besides this case alone will be sufficient for all our future needs. The operators $`_j^P`$ should satisfy
$$[_j^P,_k^P]=2\pi i\theta _{jk}\mathrm{𝟏}.$$
(52)
Since $`\theta `$ is nondegenerate the last relation defines a Heisenberg algebra. The $`(A_\theta ,A_{\theta ^1})`$ bimodule is modeled on $`L_2\left(𝐑^{d/2}\right)`$ space that is assumed to carry an irreducible representation of the Heisenberg algebra (52). The tori actions are defined to be given by the following operators on $`L_2\left(𝐑^{d/2}\right)`$
$$U_j^l=exp\left(_j^P\right),U_j^r=exp\left(\left(\theta ^1\right)^{jk}_k^P\right).$$
(53)
Now the matrix $`\theta `$ is assumed to be $`G`$-invariant. This assumption implies that the transformation $`_jR_j^k_k^P`$ is an isomorphism of the Heisenberg algebra (52). There exists then a set of unitary operators $`W^P\left(g\right)`$ satisfying (46). It follows from (46) and (53) that the second equation in (Appendix. Morita equivalence of toroidal orbifolds) holds. As for the first equation, in general it is satisfied up to a phase factor, i.e. $`W^P\left(g\right)`$ determine a projective representation of $`G`$. This possibility can be easily embedded into a general theory discussed above. However we will assume that the first equation in (Appendix. Morita equivalence of toroidal orbifolds) is satisfied precisely which is always the case for cyclic orbifolds. This finishes the construction of the corresponding triple.
Now we specialize to the cases of $`𝐙_2`$ and $`𝐙_4`$ orbifolds. We will show that by use of Morita equivalence any module admitting a constant curvature equivariant connection can be mapped to a module with a nondegenerate curvature tensor (the dimension $`d`$ is assumed to be even). First consider $`𝐙_2`$ orbifolds. In this case the mapping $`xx`$ always preserves the lattice $`D`$ and any antisymmetric two-form $`\theta `$ on it. Moreover any transformation from the group $`SO(d,d|𝐙)`$ commutes with the $`𝐙_2`$ transformation. In general under transformation (47) the curvature tensor $`f=\left(f_{ij}\right)`$ transforms as
$$\stackrel{~}{f}=AfA^t+RA^t$$
(54)
where $`A=R\theta +S`$. The matrix $`\stackrel{~}{f}`$ is nondegenerate if the matrix $`f+A^1R`$ is nondegenerate. One can check that there exist $`SO(d,d|𝐙)`$ transformations such that the matrix $`A^1R`$ is invertible and its matrix norm can be made arbitrarily large. It follows from this that $`f+A^1R`$ can be made nondegenerate.
In the $`𝐙_4`$ case we can choose a basis in which a $`𝐙_4`$ generator has a form
$$\left(\begin{array}{cc}0& 1_{2\times 2}\\ 1_{2\times 2}& 0\end{array}\right)$$
where for simplicity we set $`d=4`$. An arbitrary $`𝐙_4`$ invariant matrix then has a form
$$\left(\begin{array}{cc}A& B\\ B& A\end{array}\right)$$
where $`A`$, $`B`$ are arbitrary $`2\times 2`$ matrices. In particular we see that we have a big supply of invariant matrices with integral entries. This allows one to repeat the argument made above for the $`𝐙_2`$ case.
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# Interaction-induced Bose Metal in 2D
Two ground states are thought to exist for bosons at $`T=0`$: a superconductor with long-range phase coherence and an insulator in which the quantum mechanical phase is disordered. In this paper, we prove that the phase-disordered regime is actually not an insulator but rather a metal with a universal conductivity given by $`(2/\pi )4e^2/h`$. This result arises from a conspiracy: In the quantum-disordered regime, the population of bosons is exponentially suppressed; however, so is the scattering rate between bosons. But because the conductivity is a product of the density and the scattering time, the exponentials cancel, giving rise to a finite conductivity at $`T=0`$. Any finite amount of dissipation, however, reinstates the insulator. Consequently, the traditional arguments for the insulating phase of bosons based on quantum fluctuations of the phase must be reconsidered. Quantum fluctuations alone are insufficient to yield an insulating phase as illustrated clearly in Fig. (1). Generically, quantum fluctuations lead to a Bose metal phase in 2D.
To establish our result, we rely on the Landau-Ginzburg formalism. While this approach has had much success in elucidating the critical properties of thin films and Josephson junction arrays at the point of transition, the transport properties have proven to be more elusive. For example, simple physical considerations place the conductivity at the point of transition at a universal value of $`\sigma _Q=4e^2/h`$, whereas in experiments, $`\sigma _Q`$ ranges anywhere from $`2\sigma _Q`$ to $`\sigma _Q/3`$. Likewise, theoretical calculations based on the Landau-Ginsburg approach have yielded values ranging from $`\pi \sigma _Q/8`$ to $`1.037\sigma _Q`$. In addition, Damle and Sachdev have shown that the frequency and temperature tending to zero limits do not commute. The lack of commutativity of these two limits arises from the fundamental fact that close to the transition point, the conductivity is a universal function of $`\mathrm{}\omega /T`$. Consequently, $`\sigma (\omega 0,T=0)\sigma (\omega =0,T0)`$. The original theoretical work was all based on the former limit which physically represents the coherent regime.
In such calculations, the dc-conductivity acquires a singular Drude part on the insulating side which can be regularized by the phenomenological inclusion of dissipation. Experimentally, however, it is $`\sigma (\omega =0,T0)`$, or equivalently the hydrodynamic regime, that is relevant. In this regime, collisions between quasiparticles dominates the conductivity which can be obtained from the quantum kinetic equation. While such an approach has been used in the quantum critical regime, no corresponding study has been made in the quantum-disordered regime. It is precisely this limit that we study here using the quantum kinetic approach. We show explicitly that the quartic interaction between the bosons, whose dispersion relation is gapped in the phase-disordered regime, ultimately gives rise to a metallic rather than the anticipated insulating phase in the quantum-disordered regime. The Bose metal that we find is distinct from the Bose metal of Das and Doniach which is purported to obtain once the size of the Josephson superconducting grains exceeds a critical value. Further, our result establishes that the Bose metal is the generic ground state of bosons lacking phase coherence in the absence of dissipation.
Although our microscopic system is an array of Josephson junctions, we coarse-grain over the phase associated with each junction and hence use as our starting point the Landau-Ginzburg action,
$`F[\psi ]`$ $`=`$ $`{\displaystyle \underset{𝐤,\omega _n}{}}(k^2+\omega _n^2+m^2)|\psi (𝐤,\omega _n)|^2`$ (3)
$`+{\displaystyle \frac{U}{2N\beta }}{\displaystyle \underset{\omega _1,\mathrm{},\omega _4;𝐤_1,\mathrm{},𝐤_4}{}}\delta _{\omega _1+\mathrm{}\omega _4,0}\delta _{𝐤_1+\mathrm{}𝐤_4,0}`$
$`\psi _\nu (\omega _1,𝐤_1)\psi _\nu (\omega _2,𝐤_2)\psi _\mu (\omega _3,𝐤_3)\psi _\mu (\omega _4,𝐤_4)`$
where $`\psi (𝐫,\tau )`$ is the complex Bose order parameter whose expectation value is proportional to $`\mathrm{exp}(i\varphi )`$, where $`\varphi `$ is the phase of a particular junction. The summation in the action over discrete Matsubara frequencies, $`\omega _i=2\pi n_iT`$, and integration over continuous wavevectors, $`𝐤`$, is assumed. The parameter $`m^2`$ is the inverse square of the correlation length. In writing the action in this fashion, we have already included the one-loop renormalization arising from the quartic term. In the quantum-disordered regime, $`mT`$ and hence it is the quantum fluctuations that dominate the divergence of the correlation length.
Our goal is to calculate the conductivity in the quantum disordered regime in the collision-dominated limit. Following the quantum kinetic approach of Damle and Sachdev, we introduce the distribution function, $`f(𝐤,t)`$ for quasiparticles. We have suppressed the distinction between particles and holes as they have identical distribution functions. The quantum kinetic equation for the quasiparticle distribution function $`f(𝐤,t)`$ takes the following form
$`{\displaystyle \frac{}{t}}\delta f(𝐤,t)+e^{}𝐄(t){\displaystyle \frac{}{𝐤}}n(ϵ_𝐤)=I\{\delta f\}`$ (4)
when it is linearised in the correction $`\delta f(𝐤,t)`$ to the equilibrium Bose distribution $`n(ϵ_k)=(e^{ϵ_k/T}1)^1`$ that is induced by the electric field $`𝐄(t)`$. The gapped dispersion relation for the bosons is $`ϵ_𝐤=\sqrt{k^2+m^2}`$. The resultant linearised collision integral
$`I\{\delta f\}`$ $`=`$ $`{\displaystyle \frac{1}{N(2\pi )^2}}[{\displaystyle _{ϵ_{𝐪+𝐤}>ϵ_𝐤}}d^2q{\displaystyle \frac{1}{ϵ_{𝐪+𝐤}ϵ_𝐤}}\mathrm{Im}{\displaystyle \frac{1}{\mathrm{\Pi }(𝐪,ϵ_{𝐪+𝐤}ϵ_𝐤)}}`$ (13)
$`[\delta f(𝐤)(n(ϵ_{𝐪+𝐤}ϵ_𝐤)n(ϵ_𝐪))`$
$`\delta f(𝐪+𝐤)(1+n(ϵ_𝐤)+n(ϵ_{𝐪+𝐤}ϵ_𝐤))]`$
$`+{\displaystyle _{ϵ_{𝐪+𝐤}<ϵ_𝐤}}d^2q{\displaystyle \frac{1}{ϵ_{𝐪+𝐤}ϵ_𝐤}}\mathrm{Im}{\displaystyle \frac{1}{\mathrm{\Pi }(𝐪,ϵ_𝐤ϵ_{𝐪+𝐤})}}`$
$`\times [\delta f(𝐤)(1+n(ϵ_𝐪)+n(ϵ_𝐤ϵ_{𝐪+𝐤}))`$
$`+\delta f(𝐪+𝐤)(n(ϵ_𝐤)n(ϵ_𝐤ϵ_{𝐪+𝐤}))]`$
$`+{\displaystyle _q}d^2q{\displaystyle \frac{1}{ϵ_𝐪ϵ_𝐤}}\mathrm{Im}{\displaystyle \frac{1}{\mathrm{\Pi }(𝐪+𝐤,ϵ_𝐪+ϵ_𝐤)}}`$
$`[\delta f(𝐤)(n(ϵ_𝐪)n(ϵ_𝐤+ϵ_𝐪))`$
$`+\delta f(𝐪)(n(ϵ_𝐤)n(ϵ_𝐤+ϵ_𝐪))]]`$
is a sum of all incoming and outgoing quasiparticle scattering processes. A central quantity appearing in the collision integral is the polarization function
$`\mathrm{\Pi }(𝐪,i\mathrm{\Omega }_m)=T{\displaystyle \underset{n}{}}{\displaystyle \frac{d^2p}{(2\pi )^2}G_0(𝐩+𝐪,\omega _n+\mathrm{\Omega }_m)G_0(𝐩,\omega _n)}`$ (14)
where the bare field propagator $`G_0(𝐩,\omega _n)=(p^2+\omega _n^2+m^2)^1`$. As it is the imaginary part of $`\mathrm{\Pi }`$ that is required in the collision integral, we must perform an analytical continuation $`\mathrm{\Omega }_ni\mathrm{\Omega }_n\delta `$ with $`\delta `$ a positive infinitesimal. Unlike Damle and Sachdev, we have found it convenient to work directly in d=2 within a large $`N`$ expansion. Eq. (4) represents the linear integral equation that must be solved in the generic case to determine the role of the quartic term on the collision-induced conductivity.
To compute the conductivity, we work in the relaxation-time approximation in which $`I\{\delta f\}\delta f(𝐤)/\tau _𝐤`$. In this approximation, the conductivity emerges as a momentum integration of the form,
$`\sigma =2{\displaystyle \frac{(e^{})^2}{\mathrm{}}}{\displaystyle \frac{d^2k}{(2\pi )^2}\frac{k_x^2}{ϵ_𝐤^2}\tau _𝐤\left(\frac{n(ϵ_𝐤)}{ϵ_𝐤}\right)}.`$ (16)
We must extract then the relaxation time, $`\tau _𝐤`$ from the collision integral. In the quantum-disordered regime, the statistics of the quasiparticles becomes Boltzmannian because $`mT`$. This property results in a suppression of each subsequent $`1/N`$ correction by a factor $`e^{m/T}`$. Having set $`N=2`$, we obtain then from Eq. (13) that
$`{\displaystyle \frac{1}{\tau _𝐤}}`$ $`=`$ $`{\displaystyle \frac{1}{2(2\pi )^2}}[{\displaystyle }{\displaystyle \frac{d^2q}{ϵ_{𝐪+𝐤}ϵ_𝐤}}\left(\mathrm{Im}{\displaystyle \frac{1}{\mathrm{\Pi }(𝐪,ϵ_{𝐪+𝐤}ϵ_𝐤)}}\right)n(ϵ_{𝐪+𝐤}ϵ_𝐤)`$ (18)
$`+{\displaystyle }{\displaystyle \frac{d^2q}{ϵ_𝐪ϵ_𝐤}}\left(\mathrm{Im}{\displaystyle \frac{1}{\mathrm{\Pi }(𝐪+𝐤,ϵ_𝐪+ϵ_𝐤)}}\right)n(ϵ_𝐪)]`$
to leading accuracy. Here we took into account the fact that from Eq. (14), $`\mathrm{Im}\mathrm{\Pi }(𝐪,\mathrm{\Omega })`$ is an odd function of $`\mathrm{\Omega }`$ and used the identity $`1+n(ϵ_𝐤ϵ_{𝐪+𝐤})=n(ϵ_{𝐪+𝐤}ϵ_𝐤)`$.
The essence of our central result is that to leading order in $`T/m`$, the inverse relaxation time $`1/\tau _𝐤`$ is momentum-independent and given by
$`{\displaystyle \frac{1}{\tau }}=\pi Te^{m/T}.`$ (19)
Substitution of this expression into Eq. (16) leads to the mutual cancellation of the exponential factors yielding the remarkable result
$`\sigma (T0)={\displaystyle \frac{2}{\pi }}{\displaystyle \frac{4e^2}{h}}.`$ (20)
It is curious to note that a similar cancellation of exponential factors (from the mean free path and the density of states) arises in the context of the quasiparticle conductivity in a dirty d-wave superconductivity yielding the identical numerical prefactor $`2/\pi `$.
To establish this result, we obtain a workable expression for the inverse polarization function. In the leading approximation in $`T/m`$, the real part of the polarization function can be found from its value at $`T=0`$:
$`\mathrm{\Pi }(𝐪,\mathrm{\Omega })={\displaystyle \frac{1}{4\pi \sqrt{q^2\mathrm{\Omega }^2}}}\mathrm{arctan}{\displaystyle \frac{\sqrt{q^2\mathrm{\Omega }^2}}{2m}}`$ (21)
It is easy to show from the above expression that when $`\sqrt{\mathrm{\Omega }^2q^2}>2m`$, the leading term of the imaginary part is temperature independent and equal to
$`\mathrm{Im}\mathrm{\Pi }(𝐪,\mathrm{\Omega })={\displaystyle \frac{1}{8\sqrt{\mathrm{\Omega }^2q^2}}},`$ (22)
while in all other cases, such a contribution is absent and the leading temperature dependence comes from the first term in
$`\mathrm{Im}\mathrm{\Pi }(𝐪,\mathrm{\Omega })`$ $`=`$ $`{\displaystyle \frac{1}{16\pi }}{\displaystyle }{\displaystyle \frac{d^2p}{ϵ_𝐩ϵ_{𝐩+𝐪}}}\{|n(ϵ_{𝐩+𝐪})n(ϵ_𝐩)|`$ (25)
$`\delta (|ϵ_{𝐩+𝐪}ϵ_𝐩|\mathrm{\Omega })+(n(ϵ_{𝐩+𝐪})+n(ϵ_𝐩))`$
$`\delta (ϵ_{𝐩+𝐪}+ϵ_𝐩\mathrm{\Omega })\},`$
which is obtained by applying the Poisson summation formula to Eq. (14). The $`\delta `$ functions are eliminated upon an angular integration. As we are concerned with the quantum-disordered regime, we focus on the limits $`q,k\sqrt{mT}m`$. Because in the first term in Eq. (6), $`q>|ϵ_{𝐪+𝐤}ϵ_𝐤|`$, we must use Eq. (25) to determine the imaginary part of the inverse polarization function. Using the fact that in the region of momentum integration, $`|\mathrm{Re}\mathrm{\Pi }(𝐪,\mathrm{\Omega })||\mathrm{Im}\mathrm{\Pi }(𝐪,\mathrm{\Omega })|`$, we obtain that
$`\mathrm{Im}{\displaystyle \frac{1}{\mathrm{\Pi }(𝐪,ϵ_{𝐪+𝐤}ϵ_𝐤)}}`$ $``$ $`{\displaystyle \frac{8\pi m^2}{q}}\sqrt{{\displaystyle \frac{\pi T}{2m}}}e^{m/T}e^{(q+(\mathrm{𝐪𝐤})/q)^2/2mT}`$ (27)
$`\left(e^{(ϵ_{𝐪+𝐤}ϵ_𝐤)/T}1\right).`$
In the second term of the Eq. (18) we have that $`[(ϵ_𝐪+ϵ_𝐤)^2(𝐪+𝐤)^2]>2m`$. So, it would be sufficient to use Eq. (22) for the imaginary part, and the real part can be obtained from Eq. (21) yielding
$`\mathrm{Im}{\displaystyle \frac{1}{\mathrm{\Pi }(𝐪+𝐤,ϵ_𝐪+ϵ_𝐤)}}{\displaystyle \frac{16m}{\frac{4}{\pi ^2}\mathrm{ln}^2\frac{4m}{|𝐪𝐤|}+1}}.`$ (28)
This suggests that in the limit $`m/T\mathrm{}`$, the contribution to the relaxation time arising from the second term is logarithmically suppressed compared to the contribution from the first term. Substitution of Eq. (27) into Eq. (18) and performing the $`q`$integration we arrive at the advertised result, Eq. (19), to leading accuracy in $`T/m`$.
This result is truly remarkable, because, as one can see, the conductivity in the quantum disordered regime (in the leading approximation) depends neither on temperature nor on the distance from the transition point. This means that the dc conductivity in this model, $`\sigma =\sigma _Qg(m/T)`$, where $`g`$ is a universal function close to the transition point, displays a crossover upon lowering $`T`$ from the universal value in the quantum critical regime, $`\sigma _Q`$, to the smaller value, Eq. (20) in the quantum-disordered regime. Though we focused on the relaxation time approximation, it is possible to verify a posteriori that the incoming term in the collision integral provides a contribution subdominant in $`T/m`$. Indeed, seeking the solution of the kinetic equation by means of consecutive approximations $`\delta f(𝐤)=\delta f^{(0)}(𝐤)+\delta f^{(1)}(𝐤)+\mathrm{}`$, where
$`\delta f^{(0)}(𝐤)=e^{}\tau _𝐤{\displaystyle \frac{(\mathrm{𝐄𝐤})}{ϵ_𝐤}}\left({\displaystyle \frac{n}{ϵ_𝐤}}\right),`$ (29)
we obtain after substitution into the kinetic equation, that $`\delta f^{(1)}(𝐤)`$ is proportional to higher powers of $`T/m`$ than are the outgoing terms. Consequently, this subdominance justifies the relaxation time approximation in which only outgoing terms are retained.
The universal value for the conductivity, Eq. (20), was obtained for the region near the zero- temperature IST point, where $`m1`$. However, it is not difficult to see that the above results can be generalized for the case $`mO(1)`$, keeping $`T1`$. One needs only to make the substitution in the collision integral
$`{\displaystyle \frac{1}{N}}\mathrm{Im}{\displaystyle \frac{1}{\mathrm{\Pi }(𝐪,\mathrm{\Omega })}}\mathrm{Im}{\displaystyle \frac{U}{1+UN\mathrm{\Pi }(𝐪,\mathrm{\Omega })}}.`$ (30)
The subsequent steps are identical to those described above and yield the value of the conductivity
$`\sigma (T0)={\displaystyle \frac{2}{\pi }}{\displaystyle \frac{U^2}{(4\pi m+U)^2}}{\displaystyle \frac{4e^2}{h}}.`$ (31)
As expected, this value is not universal, and both $`m`$ and $`U`$ are the functions of parameters of the initial microscopic Hamiltonian. So, the leveled resistivity, generally speaking, depends on the distance from the transition point. In the limit $`m1`$, Eq. (31) reduces to the universal value (20). Note this result is specific to 2D. Because the conductivity is not universal for the 3D system, the 3D case deserves special attention.
In the approximation that the inverse quasiparticle scattering rate is small relative to its energy, the general form for the conductivity, Eq. (16), can be obtained from the standard Kubo formula. With this realization, the central result, Eq. (20), can be obtained from Eq. (8) of Ref. (18) by making the substitution $`\eta 1/\tau `$ ($`\kappa =1`$), with $`1/\tau `$ given by Eq. (19).
It is worth exploring how robust our bose metal is as it has not yet been observed generically in experiments. For s-wave pairing, the pair amplitude survives in the presence of weak disorder. Consequently, the Bose metal remains intact in this limit. As a result, we have constructed a concrete example of an interaction-induced metal that persists in the limit of weak disorder. This is important as metallic phases have been observed in thin films which should nominally exhibit only insulating or superconducting phases. For example, Jaeger, et. al. have observed a downturn followed by a leveling of the resistance in Ga and In thin films with a saturation value ranging between $`0.5k\mathrm{\Omega }50k\mathrm{\Omega }`$. Our value of $`(\pi /2)h/4e^210k\mathrm{\Omega }`$ is certainly within the experimental range over which the Bose metal phase has been observed. Hence the Bose metal is a serious candidate to explain these experimental observations. In addition, the new metallic phase in a dilute 2D electron gas can also be explained by the Bose metal phase if Cooper pairs form at the melting boundary of a 2D Wigner crystal as has been suggested previously.
Experimentally, measurements of the dc-conductivity are always made at finite frequency. As the frequency-dependent conductivity has a Lorentzian-type peak at $`\omega =0`$ with a width of order $`1/\tau `$, the theoretical constraint on the experimental observation of the Bose metal is that $`\omega 1/\tau `$. For a temperature of $`0.1K`$, the relaxation time, Eq. (19) is $`10^{10}\mathrm{exp}(m/T)\mathrm{sec}^1`$, where $`mT`$. Typical experimental frequencies under which the dc-conductivity is measured are on the order of several Hz. Hence, the Bose metal can be observed provided that $`T>0.05m`$. Typically, $`T>0.1m`$. Consequently, there does not appear to be any experimental constraint regarding the frequency that prohibits the observation of the Bose metal phase.
What about dissipation? We have assumed in thus far that the only source of quasiparticle relaxation is the quartic term in the Ginzburg-Landau action. Meanwhile, other dissipative mechanisms are present in real experimental systems exhibiting the IST. For example, considering the dissipation-tuned IST, Ohmic dissipation of the form $`\eta |\omega _n|`$ is always a relevant perturbation arising generically from coupling to a lattice or an Ohmic resistor. For d-wave pairing, static disorder leads to Ohmic dissipation. Provided that $`\eta <m`$, in the lowest approximation, the inverse relaxation time $`1/\tau `$ can be replaced by $`1/\tau _{eff}=1/\tau +\eta (k,T)`$, where in the case of Ohmic dissipation $`\eta (k,T)=\eta `$, and our formula for the conductivity reads in the quantum disordered regime
$`\sigma (T)=2{\displaystyle \frac{(e^{})^2}{h}}{\displaystyle \frac{T}{\eta e^{m/T}+\pi T}}.`$ (32)
We see, that already for very small $`\eta `$ the conductivity quickly starts to decrease exponentially with lowering temperature. This result is certainly intriguing and raises fundamental questions regarding the origin of the insulating state based on quantum fluctuations. We propose then that the more correct phase diagram for insulator-superconductor transitions is Fig. (1) in which it is clearly depicted that quantum fluctuations alone are insufficient to drive the insulating state. The finite temperature line starting at $`T_{\mathrm{KT}}`$ and terminating at $`g_c`$ describes the standard Kosterlitz-Thouless transition. In the absence of dissipation, the $`T_{\mathrm{KT}}`$ line separates a superconductor and a Bose metal, not a superconductor and an insulator as had traditionally been thought. Systematic studies are needed in the presence of weak dissipation are needed to see if an insulator replaces the Bose metal phase at low temperatures.
###### Acknowledgements.
We thank the NSF grant No. DMR98-96134 for funding this work and A. Castro-Neto, S. Sachdev, E. Fradkin, A. Kapitulnik, and A. Yazdani for crucial discussions.
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# Refined 𝑞-trinomial coefficients and character identities
## 1 Introduction
### 1.1 Rodney Baxter
Rodney Baxter is justly famous for his many beautiful discoveries in mathematics and physics. The 8-vertex model, Yang–Baxter equation, corner transfer matrix and hard-hexagon model are among his most envied mathematical trophies. This paper deals with a less-well-known discovery of Rodney Baxter (made together with George Andrews), that of the $`q`$-trinomial coefficients . My main aim will be to (for once) prove Baxter (and Andrews) wrong, and show that the statement
> “The literature is sparse on trinomial coefficients perhaps because they lack both depth and elegance.”,
made in the introduction of , is not at all justified.
To have any chance of succeeding, I have omitted all proofs in this paper (which are lacking elegance indeed!). These will be given in a forthcoming longer paper on the same topic.
### 1.2 $`q`$-Trinomial coefficients
In their joined work on a generalization of the hard-hexagon model, Andrews and Baxter were led to consider $`q`$-deformations of the numbers appearing in the following generalized Pascal triangle:
$$\begin{array}{ccccccccc}& & & & 1& & & & \\ & & & 1& 1& 1& & & \\ & & 1& 2& 3& 2& 1& & \\ & 1& 3& 6& 7& 6& 3& 1& \\ .& .& .& .& .& .& .& .& .\end{array}$$
The generating function for the numbers appearing in the $`(L+1)`$th row is $`(1+x+x^2)^L`$, so that an explicit expression for the trinomial coefficients can be found by double application of Newton’s binomial expansion. Explicitly,
$$(1+x+x^2)^L=\underset{a=L}{\overset{L}{}}\left(\genfrac{}{}{0pt}{}{L}{a}\right)_2x^{a+L},$$
with
$$\left(\genfrac{}{}{0pt}{}{L}{a}\right)_2=\underset{k0}{}\left(\genfrac{}{}{0pt}{}{L}{k}\right)\left(\genfrac{}{}{0pt}{}{Lk}{k+a}\right).$$
(The effective range of summation is from $`\mathrm{max}\{0,a\}`$ to $`\mathrm{min}\{L,(La)/2\}`$ so that one indeed finds a nonzero number for $`|a|L`$ only.)
Andrews and Baxter introduced several $`q`$-analogues of the trinomial coefficients. Here we shall restrict ourselves to the simplest two given by \[2, Eq. (2.7); $`B=A`$\]
$$\left[\genfrac{}{}{0.0pt}{}{L;q}{a}\right]_2=\left[\genfrac{}{}{0.0pt}{}{L}{a}\right]_2=\underset{k0}{}q^{k(k+a)}\left[\genfrac{}{}{0.0pt}{}{L}{k}\right]\left[\genfrac{}{}{0.0pt}{}{Lk}{k+a}\right]$$
and
$$T(L,a;q)=T(L,a)=q^{\frac{1}{2}(La)(L+a)}\left[\genfrac{}{}{0.0pt}{}{L;q^1}{a}\right]_2.$$
(1.1)
(This is $`T_0(L,a;q^{1/2})`$ of .) Here
$$\left[\genfrac{}{}{0.0pt}{}{n}{a}\right]=\{\begin{array}{cc}\frac{(q)_n}{(q)_a(q)_{na}}\hfill & \text{for }0an\hfill \\ 0\hfill & \text{otherwise,}\hfill \end{array}$$
is a $`q`$-binomial coefficient or Gaussian polynomial, with $`(a;q)_n=(a)_n=_{j=0}^{n1}(1aq^j)`$ for $`n1`$ and $`(a;q)_0=(a)_0=1`$. A convenient explicit expression for $`T(L,a)`$ is given by \[2, Eq. (2.60)\]
$$T(L,a)=\underset{\begin{array}{c}n=0\\ n+a+L\text{ even}\end{array}}{\overset{L|a|}{}}\frac{q^{\frac{1}{2}n^2}(q)_L}{(q)_{\frac{Lan}{2}}(q)_{\frac{L+an}{2}}(q)_n}.$$
Some useful properties of the $`q`$-trinomial coefficients are the symmetries $`\left[\genfrac{}{}{0.0pt}{}{L}{a}\right]_2=\left[\genfrac{}{}{0.0pt}{}{L}{a}\right]_2`$ and $`T(L,a)=T(L,a)`$, and the large $`L`$ limits
$$\underset{L\mathrm{}}{lim}\left[\genfrac{}{}{0.0pt}{}{L}{a}\right]_2=\frac{1}{(q)_{\mathrm{}}}$$
(1.2)
and
$$\underset{\begin{array}{c}L\mathrm{}\\ L+a+\sigma \text{ even}\end{array}}{lim}T(L,a)=\underset{\begin{array}{c}n=0\\ n+\sigma \text{ even}\end{array}}{\overset{\mathrm{}}{}}\frac{q^{\frac{1}{2}n^2}}{(q)_n}=c_\sigma (q).$$
(1.3)
Here $`c_0`$ and $`c_1`$ are (normalized) level-1 string functions of $`\text{A}_1^{(1)}`$, which admit the alternative representations
$`c_\sigma (q)`$ $`={\displaystyle \frac{(q^{1/2};q)_{\mathrm{}}+(1)^\sigma (q^{1/2};q)_{\mathrm{}}}{2(q;q)_{\mathrm{}}}}`$ (1.4)
$`={\displaystyle \frac{q^{\frac{1}{2}\sigma }}{(q;q)_{\mathrm{}}(q^{32\sigma },q^4,q^{5+2\sigma };q^8)_{\mathrm{}}(q^{2+4\sigma },q^{144\sigma };q^{16})_{\mathrm{}}}}`$
with the convention that $`(a_1,\mathrm{},a_k;q)_n=(a_1;q)_n\mathrm{}(a_k;q)_n`$.
## 2 A refinement of the $`q`$-trinomial coefficients
For integers $`L,M,a`$ and $`b`$ we define the polynomial
$$\begin{array}{c}𝒯(L,M,a,b;q)=𝒯(L,M,a,b)\hfill \\ \hfill =\underset{\begin{array}{c}n=0\\ n+a+L\text{ even}\end{array}}{\overset{\mathrm{min}\{L|a|,M\}}{}}q^{\frac{1}{2}n^2}\left[\genfrac{}{}{0.0pt}{}{M}{n}\right]\left[\genfrac{}{}{0.0pt}{}{M+b+(Lan)/2}{M+b}\right]\left[\genfrac{}{}{0.0pt}{}{Mb+(L+an)/2}{Mb}\right].\end{array}$$
Some trivial properties of $`𝒯`$ are
$$𝒯(L,M,a,b)=0\text{if }|a|>L\text{ or }|b|>M$$
(the if is not an iff), the symmetry
$$𝒯(L,M,a,b)=𝒯(L,M,a,b),$$
the duality
$$𝒯(L,M,a,b;1/q)=q^{abML}𝒯(L,M,a,b;q)$$
(2.1)
and the limit
$$\underset{M\mathrm{}}{lim}𝒯(L,M,a,b)=\frac{T(L,a)}{(q)_L}.$$
(2.2)
What is perhaps less evident is that $`𝒯`$ can be viewed as a refinement of both types of $`q`$-trinomial coefficients in the following sense:
$$\underset{i=|b|}{\overset{L|ab|}{}}q^{\frac{1}{2}(i^2b^2)}𝒯(Li,i,ab,b)=T(L,a)$$
(2.3)
and
$$\underset{i=|b|}{\overset{L|ab|}{}}q^{\frac{1}{2}(i^2b^2)}𝒯(i,Li,b,ab)=\left[\genfrac{}{}{0.0pt}{}{L}{a}\right]_2.$$
(2.4)
Here it is assumed that $`ab0`$ or $`ab0`$ in both formulas. We note that the second equation follows from the first by application of (1.1) and (2.1). Equation (2.3) results after taking $`M\mathrm{}`$ in Theorem 3.1 of the next section.
As an example of (2.3) let us calculate $`T(4,2)`$ in three different ways. When $`b=0`$ in (2.3) we get
$`T(4,2)`$ $`=𝒯(4,0,2,0)+q^{1/2}𝒯(3,1,2,0)+q^2𝒯(2,2,2,0)`$
$`=1+q(1+q+q^2)+q^2(1+q+2q^2+q^3+q^4),`$
when $`b=1`$,
$`T(4,2)`$ $`=𝒯(3,1,1,1)+q^{3/2}𝒯(2,2,1,1)+q^4𝒯(1,3,1,1)`$
$`=1+q+q^2+q^2(1+q)^2+q^4(1+q+q^2)`$
and, finally, when $`b=2`$,
$`T(4,2)`$ $`=𝒯(2,2,0,2)+q^{5/2}𝒯(1,3,0,2)+q^6𝒯(0,4,0,2)`$
$`=1+q+2q^2+q^3+q^4+q^3(1+q+q^2)+q^6.`$
Simplifying each of these three expressions correctly yields $`T(4,2)=1+q+2q^2+2q^3+2q^4+q^5+q^6`$.
To conclude this section we remark that (2.4) is a bounded analogue of the following summation \[3, Eq. (4.3); $`q^n\mathrm{}`$\], \[4, Eq. (2.10)\]
$$\underset{i=0}{\overset{\mathrm{}}{}}q^{\frac{1}{2}i^2}\frac{T(i,b)}{(q)_i}=\frac{q^{\frac{1}{2}b^2}}{(q)_{\mathrm{}}}$$
(2.5)
as can be seen by taking $`L`$ to infinity in (2.4) using (1.2) and (2.2). Hence (2.4) should be compared with \[13, Eq. (10)\]
$$\underset{i=0}{\overset{\mathrm{}}{}}q^{\frac{1}{2}i^2}\left[\genfrac{}{}{0.0pt}{}{L}{i}\right]T(i,b)=q^{\frac{1}{2}b^2}\left[\genfrac{}{}{0.0pt}{}{2L}{Lb}\right],$$
(2.6)
which also yields (2.5) in the large $`L`$ limit.
## 3 $`𝒯`$-invariance
The important question to be addressed is whether the refined $`q`$-trinomial $`𝒯`$ is at all relevant. The answer to this is a clear “yes”. Not only did we find that almost any result for $`q`$-trinomial coefficients has an analogue for the polynomials $`𝒯`$ ((2.6) appears to be an exception), but, thanks to the following theorem, $`𝒯`$ has perhaps even more depth than the $`q`$-trinomials.
###### Theorem 3.1.
For $`L,M,a,b`$ integers such that $`a,b0`$ or $`a,b0`$ there holds
$$\underset{i=|b|}{\overset{\mathrm{min}\{L|a|,M\}}{}}q^{\frac{1}{2}i^2}\left[\genfrac{}{}{0.0pt}{}{L+Mi}{L}\right]𝒯(Li,i,a,b)=q^{\frac{1}{2}b^2}𝒯(L,M,a+b,b).$$
(3.1)
This is a very powerful summation formula that allows one to iterate identities involving $`𝒯`$, thereby generating an infinite chain of $`𝒯`$-identities. By the limit (2.2) this then produces an infinite chain of $`q`$-trinomial identities, and hence (by (1.1)–(1.3)) of $`q`$-series identities.
## 4 Three exceptional examples
Taking simple $`𝒯`$-identities such as
$$\underset{j}{}q^{j(j+1)}\{𝒯(L,M,2j,j)𝒯(L,M,2j+2,j)\}=\delta _{L,0}\delta _{M,0}$$
or
$$\underset{j}{}(1)^jq^{j(j+1)/2}\{𝒯(L,M,j,j)𝒯(L,M,j+1,j)\}=\delta _{M,0}$$
as input, and using the $`𝒯`$-invariance of Theorem 3.1 to iterate these, we have proved large classes of identities for doubly bounded analogues of Virasoro characters. Many of the limiting character identities are known and many are new. In this paper we shall, however, not prove a single identity using (3.1). Instead, we shall only try to demonstrate the power of the theorem to generate new identities. As input we take three identities for which, at present, we have not a clue to a proof. However, accepting these initial conjectures, six beautiful series of identities follow, which would have been almost impossible to conceive without Theorem 3.1.
### 4.1 Some preliminaries
Let $`𝔤`$ be any of the simply laced Lie algebras whose Dynkin diagram is shown in Figure 1. Given the Dynkin diagram of $`𝔤`$ together with its labelling of vertices, we define a corresponding incidence matrix $`_𝔤`$ with entries
$$(_𝔤)_{i,j}=\{\begin{array}{cc}1\hfill & \text{if vertices }i\text{ and }j\text{ are connected by an edge}\hfill \\ 0\hfill & \text{otherwise,}\hfill \end{array}$$
where $`i,j=1,\mathrm{},r_𝔤`$ with $`r_𝔤`$ the rank of $`𝔤`$ (i.e., the number of vertices of the diagram). For any $`𝔤`$ we define an $`(m,n)`$-system as the set of $`r_𝔤`$ linear coupled equations
$$m+n=\frac{1}{2}(_𝔤m+Ne_i).$$
(4.1)
Here $`N`$ is a nonnegative integer, $`m,n,e_i`$ are vectors in $`_+^{r_𝔤}`$, with $`e_i`$ the unit vector ($`(e_i)_j=\delta _{i,j}`$) associated with the $`i`$th vertex of the Dynkin diagram of $`𝔤`$. Only those labels $`i`$ will occur that correspond to the marked vertices (drawn as open circles) in Figure 1. This fixes $`i`$ for all $`𝔤`$ other than $`\text{E}_7`$.
For given $`N`$ and $`i`$, $`m`$ determines $`n`$ and vice versa. If $`C_𝔤`$ is the Cartan matrix of $`𝔤`$, i.e., $`C_𝔤=2I_𝔤`$, we find explicitly that
$$n=\frac{1}{2}(Ne_iC_𝔤m)\text{and}m=C_𝔤^1(Ne_i2n).$$
Note though that not all $`m(n)`$ with integer entries will also yield an $`n(m)`$ with integer entries.
As an example let $`𝔤=\text{E}_7`$, $`N=6`$ and $`i=1`$. Then the only admissible solutions to (4.1) are
$`m`$ $`=5e_1+4e_2+3e_3+2e_4+e_7,`$ $`n`$ $`=e_5`$
$`m`$ $`=3e_1+4e_2+5e_3+6e_4+4e_5+2e_6+3e_7,`$ $`n`$ $`=2e_1`$
$`m`$ $`=5e_1+4e_2+5e_3+6e_4+4e_5+2e_6+3e_7,`$ $`n`$ $`=e_2`$ (4.2)
$`m`$ $`=7e_1+8e_2+9e_3+10e_4+6e_5+2e_6+5e_7,`$ $`n`$ $`=e_6`$
$`m`$ $`=9e_1+12e_2+15e_3+18e_4+12e_5+6e_6+9e_7,`$ $`n`$ $`=0`$
and
$`m`$ $`=0,`$ $`n`$ $`=3e_1`$
$`m`$ $`=2e_1,`$ $`n`$ $`=e_1+e_2`$
$`m`$ $`=4e_1+2e_2,`$ $`n`$ $`=e_3`$
$`m`$ $`=4e_1+4e_2+4e_3+4e_4+2e_5+2e_7,`$ $`n`$ $`=e_1+e_6`$ (4.3)
$`m`$ $`=6e_1+6e_2+6e_3+6e_4+4e_5+2e_6+2e_7,`$ $`n`$ $`=e_7`$
$`m`$ $`=6e_1+8e_2+10e_3+12e_4+8e_5+4e_6+6e_7,`$ $`n`$ $`=e_1,`$
where the first (second) set of solutions meets the criterion that $`n_1+n_3+n_7`$ is even (odd).
Finally we need polynomials associated with the algebras $`𝔤`$ depicted in the second column of Figure 1 as follows. Let $`p=3,5,1`$ for $`𝔤=\text{A}_5,\text{D}_6,\text{E}_7`$, respectively, so that $`p`$ corresponds to the marked vertex of $`𝔤`$. For $`M`$ a nonnegative integer, $`\sigma =0,1`$ and $`(m,n)`$-system
$$m+n=\frac{1}{2}(_𝔤m+2Me_p)$$
(4.4)
we define
$$F_{M;\sigma }^{\text{A}_5}(q)=\underset{\begin{array}{c}n_+^5\\ n_1+n_4n_2+n_5(mod3)\\ n_1+n_3+n_5+\sigma \text{ even}\end{array}}{}q^{nC_{\text{A}5}^1n}\left[\genfrac{}{}{0.0pt}{}{m+n}{n}\right]$$
and
$$F_{M;\sigma }^{\text{D}_6}(q)=\underset{\begin{array}{c}n_+^6\\ n_1+n_3+n_6\text{ even}\\ n_1+n_3+n_5+\sigma \text{ even}\end{array}}{}q^{nC_{\text{D}6}^1n}\left[\genfrac{}{}{0.0pt}{}{m+n}{n}\right]$$
and
$$F_{M;\sigma }^{\text{E}_7}(q)=\underset{\begin{array}{c}n_+^7\\ n_1+n_3+n_7+\sigma \text{ even}\end{array}}{}q^{nC_{\text{E}7}^1n}\left[\genfrac{}{}{0.0pt}{}{m+n}{n}\right].$$
Here we have used the abbreviations $`nC_𝔤^1n=_{i,j=1}^{r_𝔤}(C_𝔤^1)_{i,j}n_in_j`$ and $`\left[\genfrac{}{}{0.0pt}{}{m+n}{n}\right]=_{j=1}^{r_𝔤}\left[\genfrac{}{}{0.0pt}{}{m_j+n_j}{n_j}\right]`$ for $`m,n^{r_𝔤}`$. Similarly we will write $`(q)_n=_{j=1}^{r_𝔤}(q)_{n_j}.`$
As example consider again $`𝔤=\text{E}_7`$ and choose $`M=3`$. Then the only contributing terms to $`F_{3;0}^{\text{E}_7}(q)`$ and $`F_{3;1}^{\text{E}_7}(q)`$ correspond to the solutions of (4.1) listed in (4.1) and (4.1), respectively. Hence
$$F_{3;0}^{\text{E}_7}(q)=q^6+q^6\left[\genfrac{}{}{0.0pt}{}{5}{2}\right]+q^4\left[\genfrac{}{}{0.0pt}{}{5}{1}\right]+q^2\left[\genfrac{}{}{0.0pt}{}{3}{1}\right]+1$$
and
$$F_{3;1}^{\text{E}_7}(q)=q^{27/2}+q^{19/2}\left[\genfrac{}{}{0.0pt}{}{3}{1}\right]+q^{15/2}+q^{11/2}\left[\genfrac{}{}{0.0pt}{}{5}{1}\right]+q^{7/2}\left[\genfrac{}{}{0.0pt}{}{3}{1}\right]+q^{3/2}\left[\genfrac{}{}{0.0pt}{}{7}{1}\right].$$
### 4.2 An ($`\text{E}_7`$,$`\text{E}_8`$) series
Our first conjecture is the following polynomial identity involving $`\text{E}_7`$:
$$\begin{array}{c}\underset{j}{}\left\{q^{\frac{1}{2}j(15j+2)}𝒯(L,M,3j,5j)q^{\frac{1}{2}(3j+1)(5j+1)}𝒯(L,M,3j+1,5j+1)\right\}\hfill \\ \hfill =\underset{\begin{array}{c}n_+^7\\ n_1+n_3+n_7+L\text{ even}\end{array}}{}q^{nC_{\text{E}_7}^1n}\left[\genfrac{}{}{0.0pt}{}{\frac{1}{2}(L+M+m_1)}{2M}\right]\left[\genfrac{}{}{0.0pt}{}{m+n}{n}\right],\end{array}$$
(4.5)
with $`(m,n)`$-system (4.4) and $`𝔤=\text{E}_7`$ (so that $`p=1`$). The restriction on the sum in the right side guaranties that, given $`n_+^7`$, $`L+M+m_1`$ is even. Using Theorem 3.1 to iterate this conjecture we obtain an infinite series of polynomial identities. These identities are best expressed by turning $`\text{E}_7`$ into $`\text{E}_8`$ by the mechanism described in the caption of Figure 1. Specifically, for $`k1`$,
$$\begin{array}{c}\underset{j}{}\{q^{\frac{1}{2}j(5(5k+3)j+2)}𝒯(L,M,(5k+3)j,5j)\hfill \\ \hfill q^{\frac{1}{2}(5j+1)((5k+3)j+k+1)}𝒯(L,M,(5k+3)j+k+1,5j+1)\}\\ \hfill =\underset{r_+^{k1}}{}\left(\underset{a=0}{\overset{k2}{}}q^{\frac{1}{2}(r_ar_{a+1})^2}\left[\genfrac{}{}{0.0pt}{}{r_{a1}r_a+r_{a+1}}{r_a}\right]\right)\underset{n_+^8}{}q^{\frac{1}{4}mC_{\text{E}_8}m}\left[\genfrac{}{}{0.0pt}{}{r_{k2}\frac{1}{2}m_1}{r_{k1}}\right]\left[\genfrac{}{}{0.0pt}{}{m+n}{n}\right],\end{array}$$
(4.6)
with $`r_0=L`$, $`r_1=L+M`$ and $`(m,n)`$-system given by
$$m+n=\frac{1}{2}(_{\text{E}_8}m+r_{k1}e_1).$$
(4.7)
We now use this result to obtain $`q`$-series identities. First we let $`M`$ tend to infinity, which, by (2.2), turns (4.6) into an identity for $`q`$-trinomial coefficients. Then we either send $`L`$ to infinity using (1.3), or we first replace $`q1/q`$ and then send $`L`$ to infinity using (1.2). Omitting the actual calculations (which sometimes require variable changes to remove $`L`$-dependent terms in the exponent of $`q`$) we find two families of $`q`$-series identities, one of $`\text{E}_7`$ type and one of $`\text{E}_8`$ type. To present these identities in a neat form we recall the bosonic representation of the Virasoro characters
$$\chi _{r,s}^{(p,p^{})}(q)=\frac{q^{\frac{(p^{}rps)^21}{4pp^{}}}}{(q)_{\mathrm{}}}\underset{j=\mathrm{}}{\overset{\mathrm{}}{}}\left\{q^{j(pp^{}j+p^{}rps)}q^{(pj+r)(p^{}j+s)}\right\},$$
(4.8)
for coprime integers $`2p<p^{}`$ and $`r=1,\mathrm{},p1`$, $`s=1,\mathrm{},p^{}1`$. The somewhat unusual normalization of the Virasoro characters is chosen to simplify subsequent equations. Besides the Virasoro characters we also need the following (subset) of the branching functions corresponding to the coset $`(\text{A}_1^{(1)}\text{A}_1^{(1)},\text{A}_1^{(1)})`$ at levels $`2p/(p^{}p)2,2`$ and $`2p/(p^{}p)`$ :
$$\begin{array}{c}B_{r,s;\sigma }^{(p,p^{})}(q)=\frac{q^{\frac{(p^{}rps)^24}{8pp^{}}}}{(q)_{\mathrm{}}}\underset{j=\mathrm{}}{\overset{\mathrm{}}{}}\{q^{j(pp^{}j+p^{}rps)}c_{pj+\frac{1}{2}(rs)+\sigma }(q)\hfill \\ \hfill q^{(pj+r)(p^{}j+s)}c_{pj+\frac{1}{2}(r+s)+\sigma }(q)\},\end{array}$$
(4.9)
for integers $`p,p^{},r,s`$ in the same ranges as above, such that $`p^{}p`$ and $`rs`$ are even, and $`\mathrm{gcd}((p^{}p)/2,p^{})=1`$. The integer $`\sigma `$ takes the values $`0`$ or $`1`$ and the $`c_j`$ are the string functions of equations (1.3) and (1.4) with the obvious identification of $`c_{2j+\sigma }`$ with $`c_\sigma `$.
The various large $`L`$ and $`M`$ limits now lead to the following list of identities.
* Taking $`L+\sigma `$ even in (4.6), sending $`M`$ and $`L`$ to infinity using that $`c_\sigma (q)=\chi _{\sigma +1,1}^{(3,4)}(q)/(q)_{\mathrm{}}`$ yields for $`k1`$
$$\underset{n_1,\mathrm{},n_k0}{}\frac{q^{\frac{1}{2}(N_1^2+\mathrm{}+N_k^2)}F_{n_k;m_\sigma }^{\text{E}_7}(q)}{(q)_{n_1}\mathrm{}(q)_{n_{k1}}(q)_{2n_k}}=\{\begin{array}{cc}\chi _{\sigma +1,1}^{(3,4)}(q)\chi _{1,(k+1)/2}^{(5,(5k+3)/2)}(q)\hfill & k\text{ odd}\hfill \\ B_{1,k+1;\sigma }^{(5,5k+3)}(q)\hfill & k\text{ even,}\hfill \end{array}$$
where the following definitions have been employed:
$$N_a=n_a+\mathrm{}+n_k\text{and}m_\sigma \sigma +\underset{\begin{array}{c}a=1\\ a\text{ odd}\end{array}}{\overset{k}{}}n_a(mod2),$$
(4.10)
for $`m_\sigma \{0,1\}`$. There is a corresponding “$`k=0`$” identity obtained by taking the same limit as above, but now in the initial conjecture (4.5). Using the (from a $`q`$-series point of view nontrivial) relation $`B_{1,1;\sigma }^{(3,5)}=\chi _{2\sigma +1,1}^{(4,5)}`$, which follows from a symmetry of the A$`{}_{}{}^{(1)}{}_{1}{}^{}`$ branching functions, we find the well-known $`\text{E}_7`$ conjecture
$$\chi _{2\sigma +1,1}^{(4,5)}(q)=\underset{\begin{array}{c}n_+^7\\ n_1+n_3+n_7+\sigma \text{ even}\end{array}}{}\frac{q^{nC_{\text{E}_7}^1n}}{(q)_n}.$$
(4.11)
* If in (4.6) we send $`M`$ to infinity, replace $`q1/q`$ and then take the limit of large $`L`$ we obtain for $`k2`$,
$$\begin{array}{c}\underset{r_+^{k1}}{}\underset{m_+^8}{^{}}\frac{q^{\frac{1}{2}_{a=1}^{k1}(r_ar_{a1})^2}}{(q)_{r_1}}\left(\underset{a=2}{\overset{k1}{}}\left[\genfrac{}{}{0.0pt}{}{r_{a1}r_a+r_{a+1}}{r_a}\right]\right)q^{\frac{1}{4}mC_{\text{E}_8}m}\left[\genfrac{}{}{0.0pt}{}{m+n}{m}\right]\hfill \\ \hfill =\{\begin{array}{cc}\chi _{(k+1)/2,k}^{((5k+3)/2,5k2)}(q)\hfill & k\text{ odd}\hfill \\ \chi _{k/2,k+1}^{(5k/21,5k+3)}(q)\hfill & k\text{ even,}\hfill \end{array}\end{array}$$
(4.12)
with $`r_0=0`$, $`r_k=r_{k1}m_1/2`$ and $`(m,n)`$-system (4.7). The prime in the sum over $`m`$ denotes the restriction $`m_2m_4m_8r_{k1}(mod2)`$ with all other $`m_i`$ being even. When $`k=1`$ the resulting character formula takes a somewhat different form, and one obtains the $`\text{E}_8`$ identity
$$\chi _{1,1}^{(3,4)}(q)=\underset{n_+^8}{}\frac{q^{nC_{\text{E}_8}^1n}}{(q)_n}=\frac{1}{(q^3,q^4,q^5;q^8)_{\mathrm{}}(q^2,q^{14};q^{16})_{\mathrm{}}}.$$
(4.13)
### 4.3 A ($`\text{D}_6`$,$`\text{E}_7`$) series
In our second conjecture the role of $`\text{E}_7`$ is taken over by $`\text{D}_6`$,
$$\begin{array}{c}\underset{j}{}\left\{q^{\frac{1}{2}j(24j+2)}𝒯(L,M,4j,6j)q^{\frac{1}{2}(4j+1)(6j+1)}𝒯(L,M,4j+1,6j+1)\right\}\hfill \\ \hfill =\underset{\begin{array}{c}n_+^6\\ n_1+n_3+n_6\text{ even}\\ n_1+n_3+n_5+L\text{ even}\end{array}}{}q^{nC_{\text{D}_6}^1n}\left[\genfrac{}{}{0.0pt}{}{\frac{1}{2}(L+M+m_5)}{2M}\right]\left[\genfrac{}{}{0.0pt}{}{m+n}{n}\right],\end{array}$$
(4.14)
with $`(m,n)`$-system (4.4) where $`𝔤=\text{D}_6`$ (and $`p=5`$). The restriction that $`n_1+n_3+n_5+L`$ is even is necessary for $`L+M+m_5`$ to be even. The additional constraint on $`n_1+n_3+n_6`$ is to avoid an extra
$$q\underset{j}{}\left\{q^{\frac{1}{2}j(24j+14)}𝒯(L,M,4j+1,6j+2)q^{\frac{1}{2}(4j+3)(6j+1)}𝒯(L,M,4j+2,6j+3)\right\}$$
on the left-hand side. Of course, this means we have actually two conjectures, but the case when $`n_1+n_3+n_6`$ is odd will not be pursued here.
Using Theorem 3.1 to iterate the $`\text{D}_6`$ conjecture we obtain an infinite series of polynomial identities involving $`\text{E}_7`$. Specifically, for $`k1`$ there holds
$$\begin{array}{c}\underset{j}{}\{q^{\frac{1}{2}j(6(6k+4)j+2)}𝒯(L,M,(6k+4)j,6j)\hfill \\ \hfill q^{\frac{1}{2}(6j+1)((6k+4)j+k+1)}𝒯(L,M,(6k+4)j+k+1,6j+1)\}\\ \hfill =\underset{r_+^{k1}}{}\left(\underset{a=0}{\overset{k2}{}}q^{\frac{1}{2}(r_ar_{a+1})^2}\left[\genfrac{}{}{0.0pt}{}{r_{a1}r_a+r_{a+1}}{r_a}\right]\right)\\ \hfill \times \underset{\begin{array}{c}n_+^7\\ n_1+n_3+n_7\text{ even}\end{array}}{}q^{\frac{1}{4}mC_{\text{E}_7}m}\left[\genfrac{}{}{0.0pt}{}{r_{k2}\frac{1}{2}m_6}{r_{k1}}\right]\left[\genfrac{}{}{0.0pt}{}{m+n}{n}\right],\end{array}$$
(4.15)
where $`r_0=L`$, $`r_1=L+M`$ and
$$m+n=\frac{1}{2}(_{\text{E}_7}m+r_{k1}e_6).$$
(4.16)
As before we consider the various large $`L`$ and $`M`$ limits.
* Taking $`L+\sigma `$ even in (4.15) and sending $`M`$ and $`L`$ to infinity yields for $`k1`$ that
$$\underset{n_1,\mathrm{},n_k0}{}\frac{q^{\frac{1}{2}(N_1^2+\mathrm{}+N_k^2)}F_{n_k;m_\sigma }^{\text{D}_6}(q)}{(q)_{n_1}\mathrm{}(q)_{n_{k1}}(q)_{2n_k}}=\{\begin{array}{cc}\chi _{\sigma +1,1}^{(3,4)}(q)\chi _{1,(k+1)/2}^{(6,3k+2)}(q)\hfill & k\text{ odd}\hfill \\ B_{1,k+1;\sigma +k/2}^{(6,6k+4)}(q)\hfill & k\text{ even,}\hfill \end{array}$$
with the notation of equation (4.10) and with the identification of $`B_{r,s;\sigma +2j}^{(p,p^{})}`$ with $`B_{r,s;\sigma }^{(p,p^{})}`$.
The $`k=0`$ case, corresponding to the above limit taken in (4.14) yields
$$B_{1,1;\sigma }^{(4,6)}=\underset{\begin{array}{c}n_+^6\\ n_1+n_3+n_6\text{ even}\\ n_1+n_3+n_5\sigma \text{ even}\end{array}}{}\frac{q^{nC_{\text{D}_6}^1n}}{(q)_n}.$$
Although we were unable to prove this, it appears that the above branching function admits the following simplification
$$B_{1,1;\sigma }^{(4,6)}(q)=\{\begin{array}{cc}\frac{1}{(q)_{\mathrm{}}}\left(\underset{j=0}{\overset{\mathrm{}}{}}(q)^{j^2}+\underset{j=1}{\overset{\mathrm{}}{}}q^{6j^2}\right)\hfill & \sigma =0\hfill \\ \frac{q^{3/2}}{(q)_{\mathrm{}}}\underset{j=0}{\overset{\mathrm{}}{}}q^{6j(j+1)}=\frac{q^{3/2}(q^{24};q^{24})_{\mathrm{}}}{(q^{12};q^{24})_{\mathrm{}}(q;q)_{\mathrm{}}}\hfill & \sigma =1\text{.}\hfill \end{array}$$
For $`\sigma =1`$ this implies a new identity of the Rogers–Ramanujan type for the algebra $`\text{D}_6`$.
* If we send $`M`$ to infinity, replace $`q1/q`$ and then take the limit of large $`L`$ we obtain for all $`k2`$
$$\begin{array}{c}\underset{r_+^{k1}}{}\underset{m_+^7}{^{}}\frac{q^{\frac{1}{2}_{a=1}^{k1}(r_ar_{a1})^2}}{(q)_{r_1}}\left(\underset{a=2}{\overset{k1}{}}\left[\genfrac{}{}{0.0pt}{}{r_{a1}r_a+r_{a+1}}{r_a}\right]\right)q^{\frac{1}{4}mC_{\text{E}_7}m}\left[\genfrac{}{}{0.0pt}{}{m+n}{m}\right]\hfill \\ \hfill =\{\begin{array}{cc}\chi _{(k+1)/2,k}^{(3k+2,6k2)}(q)\hfill & k\text{ odd}\hfill \\ \chi _{k/2,k+1}^{(3k1,6k+4)}(q)\hfill & k\text{ even,}\hfill \end{array}\end{array}$$
(4.17)
with $`r_0=0`$, $`r_k=r_{k1}m_6/2`$ and $`(m,n)`$-system (4.16). The prime in the sum over $`m`$ denotes the restriction that $`m_1m_3m_5r_{k1}(mod2)`$ and that all other $`m_i`$ are even. The identity corresponding to $`k=1`$ is given by (4.11) with $`\sigma =0`$.
### 4.4 An ($`\text{A}_5`$,$`\text{E}_6`$) series
Our final conjecture is somewhat more involved than the previous two, as it is not possible to disentangle the two terms on the left-hand side below by an appropriate summation restriction on $`n`$,
$$\begin{array}{c}\underset{j}{}\left\{q^{\frac{1}{2}j(48j+2)}𝒯(L,M,6j,8j)q^{\frac{1}{2}(6j+1)(8j+1)}𝒯(L,M,6j+1,8j+1)\right\}\hfill \\ \hfill +q^3\underset{j}{}\left\{q^{\frac{1}{2}j(48j+34)}𝒯(L,M,6j+2,8j+3)q^{\frac{1}{2}(6j+1)(8j+7)}𝒯(L,M,6j+3,8j+4)\right\}\\ \hfill =\underset{\begin{array}{c}n_+^5\\ n_1+n_4n_2+n_5(mod3)\\ n_1+n_3+n_5L(mod2)\end{array}}{}q^{nC_{\text{A}_5}^1n}\left[\genfrac{}{}{0.0pt}{}{\frac{1}{2}(L+M+m_3)}{2M}\right]\left[\genfrac{}{}{0.0pt}{}{m+n}{n}\right],\end{array}$$
(4.18)
with $`(m,n)`$-system (4.4) where $`𝔤=\text{A}_5`$ (and $`p=3`$). Iterating this last conjecture using Theorem 3.1 one finds a series of $`\text{E}_6`$-type polynomial identities as follows ($`k1`$)
$$\begin{array}{c}\underset{j}{}\{q^{\frac{1}{2}j(8(8k+6)j+2)}𝒯(L,M,(8k+6)j,8j)\hfill \\ \hfill q^{\frac{1}{2}(8j+1)((8k+6)j+k+1)}𝒯(L,M,(8k+6)j+k+1,8j+1)\}\\ \hfill +q^{\frac{3}{2}(3k+2)}\underset{j}{}\{q^{\frac{1}{2}j(8(8k+6)j+48k+34)}𝒯(L,M,(8k+6)j+3k+2,8j+3)\\ \hfill q^{\frac{1}{2}(8j+7)((8k+6)j+k+1)}𝒯(L,M,(8k+6)j+4k+3,8j+4)\}\\ \hfill =\underset{r_+^{k1}}{}\left(\underset{a=0}{\overset{k2}{}}q^{\frac{1}{2}(r_ar_{a+1})^2}\left[\genfrac{}{}{0.0pt}{}{r_{a1}r_a+r_{a+1}}{r_a}\right]\right)\\ \hfill \times \underset{\begin{array}{c}n_+^6\\ n_1+n_4n_2+n_5(mod3)\end{array}}{}q^{\frac{1}{4}mC_{\text{E}_6}m}\left[\genfrac{}{}{0.0pt}{}{r_{k2}\frac{1}{2}m_6}{r_{k1}}\right]\left[\genfrac{}{}{0.0pt}{}{m+n}{n}\right],\end{array}$$
where $`r_0=L`$, $`r_1=L+M`$ and $`m+n=\frac{1}{2}(_{\text{E}_6}m+r_{k1}e_6)`$. Fortunately, the limiting character identities that follow from this monster are more manageable.
* Taking $`L+\sigma `$ even and letting $`M`$ and $`L`$ tend to infinity yields for positive $`k`$,
$$\begin{array}{c}\underset{n_1,\mathrm{},n_k0}{}\frac{q^{\frac{1}{2}(N_1^2+\mathrm{}+N_k^2)}F_{n_k;m_\sigma }^{\text{D}_5}(q)}{(q)_{n_1}\mathrm{}(q)_{n_{k1}}(q)_{2n_k}}\hfill \\ \hfill =\{\begin{array}{cc}\chi _{\sigma +1,1}^{(3,4)}(q)\chi _{1,(k+1)/2}^{(8,4k+3)}(q)+\chi _{2\sigma ,1}^{(3,4)}(q)\chi _{7,(k+1)/2}^{(8,4k+3)}(q)\hfill & k\text{ odd}\hfill \\ B_{1,k+1;\sigma +k/2}^{(8,8k+6)}(q)+B_{7,k+1;\sigma +k/2+1}^{(8,8k+6)}(q)\hfill & k\text{ even,}\hfill \end{array}\end{array}$$
with the notation of equation (4.10). The analogous limit taken in (4.18) leads to
$$B_{1,1;\sigma }^{(6,8)}(q)+B_{1,7;1\sigma }^{(6,8)}(q)=\underset{\begin{array}{c}n_+^5\\ n_1+n_4n_2+n_5(mod3)\\ n_1+n_3+n_5\sigma (mod2)\end{array}}{}\frac{q^{nC_{\text{A}_5}^1n}}{(q)_n}.$$
* If we send $`M`$ to infinity, replace $`q1/q`$ and then take the limit of large $`L`$ we obtain for $`k2`$,
$$\begin{array}{c}\underset{r_+^{k1}}{}\underset{m_+^6}{^{}}\frac{q^{\frac{1}{2}_{a=1}^{k1}(r_ar_{a1})^2}}{(q)_{r_1}}\left(\underset{a=2}{\overset{k1}{}}\left[\genfrac{}{}{0.0pt}{}{r_{a1}r_a+r_{a+1}}{r_a}\right]\right)q^{\frac{1}{4}mC_{\text{E}_6}m}\left[\genfrac{}{}{0.0pt}{}{m+n}{m}\right]\hfill \\ \hfill =\{\begin{array}{cc}\chi _{(k+1)/2,k}^{(4k+3,8k2)}(q)+\chi _{(k+1)/2,7k2}^{(4k+3,8k2)}(q)\hfill & k\text{ odd}\hfill \\ \chi _{k/2,k+1}^{(4k1,8k+6)}(q)+\chi _{7k/21,k+1}^{(4k1,8k+6)}(q)\hfill & k\text{ even,}\hfill \end{array}\end{array}$$
(4.19)
where $`m+n=\frac{1}{2}(_{\text{E}_6}m+r_{k1}e_6)`$ and $`r_0=0`$, $`r_k=r_{k1}m_6/2`$. The prime in the sum over $`m`$ denotes the restriction that $`m_1m_3m_5r_{k1}(mod2)`$ and that all other $`m_i`$ are even. Again $`k=1`$ is special, corresponding to the $`\text{E}_6`$ conjecture of ,
$$\chi _{1,1}^{(6,7)}(q)+\chi _{5,1}^{(6,7)}(q)=\underset{\begin{array}{c}n_+^6\\ n_1+n_4n_2+n_5(mod3)\end{array}}{}\frac{q^{nC_{\text{E}_6}^1n}}{(q)_n}.$$
(4.20)
## 5 Discussion
We hope that the examples presented in the previous section support our claim that $`q`$-trinomial coefficients, and their refinement introduced in this paper are mathematical objects of both depth and elegance.
It is quite intriguing to observe that the $`𝒯`$-invariance of Theorem 3.1 also appears to have physical significance. In a very recent paper by Dorey, Dunning and Tateo , new families of renormalization group flows between $`c<1`$ conformal field theories were proposed. Labelling such a theory by $`M(p,p^{})`$, in accordance with definition (4.8) of the Virasoro characters, Dorey et al. conjecture the following flows
$`M(p,p^{})+\varphi _{21}`$ $`M(p^{}p,p^{})`$ $`p<p^{}<2p`$ (5.1a)
$`M(p,p^{})+\varphi _{15}`$ $`M(p,4pp^{})`$ $`2p<p^{}<3p`$ (5.1b)
$`M(p,p^{})+\varphi _{15}`$ $`M(4pp^{},p)`$ $`3p<p^{}<4p`$ , (5.1c)
where $`\varphi _{rs}`$ is the perturbing operator of the $`M(p,p^{})`$ theory. Using these three flows we find the following chain ending in $`M(3,4)`$:
$$M(3,4)\stackrel{(\text{c})}{}M(4,13)\stackrel{(\text{a})}{}M(9,13)\stackrel{(\text{b})}{}M(9,23)\stackrel{(\text{a})}{}M(14,23)\stackrel{(\text{b})}{}\mathrm{},$$
where (a) denotes equation (5.1a) etc. But this flow diagram coincides with the chain of character identities given in (4.12) and (4.13)! In particular, $`M(3,4)`$ corresponds to the $`\text{E}_8`$ identity of (4.13), $`M(5n1,10n+3)`$ ($`n1`$) corresponds to (4.12) for $`k=2n`$, and $`M(5n+4,10n+3)`$ ($`n1`$) to (4.12) for $`k=2n+1`$.
In much the same way the flow diagram
$$M(4,5)\stackrel{(\text{c})}{}M(5,16)\stackrel{(\text{a})}{}M(11,16)\stackrel{(\text{b})}{}M(11,28)\stackrel{(\text{a})}{}M(17,28)\stackrel{(\text{b})}{}\mathrm{}$$
is in accordance with the chain of character identities given by (4.17) and (4.11), and
$$M(6,7)\stackrel{(\text{c})}{}M(7,22)\stackrel{(\text{a})}{}M(15,22)\stackrel{(\text{b})}{}M(15,38)\stackrel{(\text{a})}{}M(23,38)\stackrel{(\text{b})}{}\mathrm{}$$
is in one-to-one correspondence with (4.19) and (4.20).
To conclude we mention that more general applications of our refined $`q`$-trinomial coefficients will be presented in a future paper. Therein we will also discuss the $`𝒯`$-invariance in the broader context of the Bailey lemma , trinomial Bailey lemma and (generalized) Burge transform .
### Acknowledgements
This work was supported by a fellowship of the Royal Netherlands Academy of Arts and Sciences.
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# Analysis of Dislocation Mechanism for Melting of Elements: Pressure Dependence
## 1 Introduction
The idea that a proliferation of dislocations is associated with melting dates back to Mott . The very first theory of dislocation-mediated melting was a success, inasmuch as it predicted a first-order transition, as a consequence of incorporating the mutual screening of dislocations, in agreement with observations. Molecular dynamics and Monte Carlo calculations have more recently provided further evidence for the notion that dislocations drive the melting transition in three dimensions. There is also some experimental evidence that line defects are present in solids near melting .
In refs. we formulated a dislocation theory of melting in which dislocations near melt were modeled as non-interacting strings on a lattice. The possible configurations of a dislocation were taken to be closed random walks. Screening of long-range strain fields by other dislocations in a dense ensemble results in a $`\rho \mathrm{ln}\rho `$ dependence of the free energy on the dislocation density, $`\rho ,`$ and thus a first-order transition. We obtained the following relation between the melting temperature $`T_m,`$ the shear modulus, $`G,`$ the Wigner-Seitz volume, $`v_{WS},`$ the coordination number, $`z,`$ and the critical density of dislocations, $`\rho (T_m)`$ (in units where $`k_B=1):`$
$$T_m=\frac{\kappa \lambda Gv_{WS}}{8\pi \mathrm{ln}(z1)}\mathrm{ln}\left(\frac{\alpha ^2}{4b^2\rho (T_m)}\right).$$
(1)
Here $`b`$ is the length of the shortest perfect-dislocation Burgers vector, $`\kappa `$ is 1 for a screw dislocation and $`(1\nu )^13/2`$ for an edge dislocation $`(\nu `$ is the Poisson ratio), $`\lambda b^3/v_{WS}`$ and $`\alpha ,`$ which accounts for non-linear effects in the dislocation core, has a value of 2.9 . Experimental data on 51 elements show that
$$\frac{Gv_{WS}}{4\pi T_m\mathrm{ln}(z1)}=1.01\pm 0.17$$
(2)
at zero pressure . Eqs. (1) and (2) imply that the critical dislocation density at zero pressure is
$$\rho (T_m)=(0.61\pm 0.20)b^2.$$
(3)
This value is in good agreement with the critical density
$$\rho (T_m)=(0.66\pm 0.11)b^2,$$
(4)
obtained by applying our relation for the latent heat of fusion ,
$$L_m=\frac{1}{\lambda }b^2\rho (T_m)RT_m\mathrm{ln}(z1),$$
(5)
to data on latent heats for 75 elements. Hence, $`b^2\rho (T_m)`$ is approximately constant across the Periodic Table with the numerical value
$$b^2\rho (T_m)=0.64\pm 0.14,$$
(6)
which is the uncertainty-weighted average of Eqs. (3) and (4).
In this paper we investigate the validity of our melting relation, Eq. (1), up to pressures of order 100 GPa, by comparing to experimental melting curves, i.e., melting temperatures versus pressure, $`p.`$ This comparison requires $`v_{WS}(p),`$ or its equivalent, the pressure dependence of the compression, $`\eta V_0/V.`$ We obtain $`\eta (p)`$ from the bulk modulus, $`B(p),`$ which is extrapolated to high pressure using only its value and first pressure derivative at ambient conditions, viz., room temperature and zero pressure. The shear modulus $`G(p)`$ is similarly extrapolated to high pressure. Pressure derivatives of $`G`$ and $`B`$ are typically $`O(1),`$ so the extrapolation of the bulk modulus is expected to break down at pressures of order the ambient bulk modulus. The parameter $`\kappa `$ in Eq. (1), which depends on the Poisson ratio, varies by only a few percent between $`p=0`$ and 100 GPa (we discuss this in more detail in Section 2). Since the accuracy of our melting relation at zero pressure is 17%, we take $`\kappa `$ to be a constant. We also make the necessary but reasonable assumption that $`b^2\rho (T_m)`$ is also a pressure-independent constant. With this assumption, we find that our melting relation agrees well with experimental melting curves up to pressures $`B,`$ and, in fact, our extrapolation of $`T_m`$ is often in good agreement with data to pressures $`2B.`$ In addition to the good agreement with the existing melting curve data, we also predict the high-pressure melting curves of Ag, Au, Cr, Cu, Mo, Os, Ru, W, and several actinides.
## 2 Melting curve equation
We now consider the pressure dependences of the factors appearing in our melting relation, Eq. (1). The parameter $`\lambda `$ is constant by its definition, $`\alpha `$ is also assumed to be a constant, and $`\kappa `$ may be taken as constant provided that the Poisson ratio $`\nu `$ has a very weak pressure dependence. In fact, for an isotropic medium
$$\nu =\frac{1}{2}\frac{3B2G}{3B+G}.$$
(7)
Although both $`G`$ and $`B`$ vary with pressure, the ratio in Eq. (7) varies only weakly. Consider, for example, Cu, for which $`\nu 0.34`$ at $`p=0.`$ At $`p=100`$ GPa, we calculate the values of $`G`$ and $`B`$ with the help of Eqs. (13) and (14) below, with their pressure derivatives taken from ref. , and find $`\nu 0.38.`$ Therefore, in this case the corresponding values of $`1/\kappa 1\nu /2`$ are 0.83 and 0.81, respectively, so that the variation in the average value of $`1/\kappa `$ is $`2`$%. Thus, the pressure dependence of $`1/\kappa `$ can be safely neglected. (There exists an upper bound on the change in the value of $`1/\kappa `$ with pressure. In the ultra-high-pressure limit, $`p\eta ^{5/3},`$ in agreement with the theory of the free electron gas (Fermi gas). Therefore, $`BVdp/dV=\eta dp/d\eta \eta ^{5/3}G\eta ^{4/3},`$ and hence $`\nu 1/2,`$ in view of Eq. (7) (see also ). Thus, in contrast to the Poisson ratio which changes by $`50`$%, $`1/\kappa 1\nu /2`$ changes by $`10`$%: $`5/63/4.)`$
We assume further that the mean interdislocation spacing at the melting point, $`R1/\sqrt{\rho (T_m)},`$ scales with $`b,`$ independent of pressure, and hence $`b^2\rho (T_m)`$ is a pressure-independent constant (with a numerical value of $`0.64\pm 0.14,`$ in view of Eq. (6)). It then follows from Eq. (1) that, provided the coordination number does not change with pressure (i.e., the element either remains in the same crystalline phase, or changes phase without changing the coordination number, e.g., a face-centered cubic structure $``$ a hexagonal close-packed structure), the melting relation is given by
$$\frac{G(p,T_m(p))v_{WS}(p,T_m(p))}{T_m(p)}=\mathrm{const}.$$
(8)
The dependence of $`v_{WS}`$ on pressure and temperature is just the equation of state of the metal. Let us first focus on its temperature dependence. The fixed-pressure ratio of Wigner-Seitz volumes at $`T_m`$ and $`T=0`$ is equal to $`1+\beta T_m,`$ where $`\beta `$ is the volume expansivity. At $`p=0,`$ $`\beta `$ is typically of order $`10^5`$ K$`^1,`$ and melting temperatures are at most about 4000 K, so $`v_{WS}`$ changes by only a few percent between $`T=0`$ and $`T_m.`$ Assuming that $`\beta `$ does not increase appreciably with compression, we can use room-temperature values for $`v_{WS}.`$
In contrast to $`v_{WS},`$ the dependence of $`G`$ on $`T`$ is not necessarily weak. Its $`T`$-dependence involves two characteristic temperatures, namely the Debye temperature, $`T_D,`$ and the melting temperature. $`G`$ is always monotonically decreasing with $`T,`$ and is nonlinear for $`T\stackrel{<}{}T_D`$ and linear from $`T_D`$ to $`T_m.`$ However, there are no experimental data, no computer calculations, and no theoretical guidance that tells us how the temperature dependence of $`G`$ varies with pressure. In particular, how does the (negative) slope of the linear region vary with $`p`$? At this point we have no choice but to conjecture. We assume that $`G(p,T_m(p))/G(p,0)`$ is a slowly varying function of $`p,`$ so it can be considered constant up to moderate compressions, say, 20% to 30%. Thus, $`G(p,T_m)`$ is replaced by $`G(p,0)`$ in Eq. (8). In addition, data on the $`p=0`$ temperature dependence of shear moduli clearly show that $`G(p,300)G(p,0),`$ and therefore, we use the room temperature value of the shear modulus in our melting relation.
Subsequently, the explicit dependence of $`G`$ and $`v_{WS}`$ on $`T`$ will be dropped. It will be understood that $`G`$ and $`v_{WS}`$ are at room temperature. Our melting relation now reads
$$\frac{G(p)v_{WS}(p)}{T_m(p)}=\mathrm{const}.$$
(9)
Differentiating Eq. (9) with respect to $`p,`$ one finds
$$\frac{1}{T_m}\frac{dT_m}{dp}=\frac{1}{G}\frac{dG}{dp}\frac{1}{B},$$
(10)
where we have used the definition of the bulk modulus,
$$B(p)V\frac{dp}{dV}=v_{WS}\frac{dp}{dv_{WS}}.$$
(11)
Thus, upon integration, Eq. (10) gives
$$\frac{T_m(p)}{T_m(0)}=\frac{G(p)}{G(0)}\mathrm{exp}\left\{_0^p\frac{dp^{}}{B(p^{})}\right\}.$$
(12)
To proceed further, we have to specify $`G(p)`$ and $`B(p).`$
### 2.1 The shear modulus $`G`$ at finite pressure
For the shear modulus at all pressures, we use the relation
$$G=G_0+G_0^{}\frac{p}{\eta ^{1/3}},$$
(13)
where $`G_0^{}(dG/dp)_0.`$ The subscript 0 refers to ambient conditions: $`T300`$ K and $`p=0.`$
This equation satisfies the requirement that $`G\eta ^{4/3}`$ as $`\eta \mathrm{},`$ since $`p\eta ^{5/3}.`$ With the values of $`G_0^{}`$ for 32 elements tested in ref. Eq. (13) gives nearly the right value for the proportionality constant between $`G`$ and $`\eta ^{4/3}`$ at high compressions. Eq. (13) works well for a diverse selection of engineering metals covering many different crystal structures and nearly all groups of the Periodic Table .
### 2.2 Compression and the bulk modulus $`B`$ at finite pressure
Expanding the bulk modulus around $`p=0`$ we have
$$B(p)=B_0+B_0^{}p+\frac{1}{2}B_0^{\prime \prime }p^2+\mathrm{},$$
(14)
where $`B_0`$ and $`B_0^{}(dB/dp)_0,`$ $`B_0^{\prime \prime }(d^2B/dp^2)_0,\mathrm{}`$ can be extracted from equation of state data. Values of $`B_0^{\prime \prime }`$ are known for a few elements only (their determination is highly uncertain and involves an error of order 100% ), and besides, $`B_0^{\prime \prime }`$ first appears in the $`(p/B_0)^3`$ term in the power series expansion of $`\eta :`$
$$\eta =\mathrm{exp}\left\{_0^p\frac{dp^{}}{B(p^{})}\right\}=\left[\frac{2B_0+(B_0^{}+\sqrt{B_0^22B_0B_0^{\prime \prime }})p}{2B_0+(B_0^{}\sqrt{B_0^22B_0B_0^{\prime \prime }})p}\right]^{1/\sqrt{B_0^22B_0B_0^{\prime \prime }}}$$
$$=1+\left(\frac{p}{B_0}\right)\frac{B_0^{}1}{2}\left(\frac{p}{B_0}\right)^2+\frac{(B_0^{}1)(2B_0^{}1)B_0B_0^{\prime \prime }}{6}\left(\frac{p}{B_0}\right)^3+\mathrm{}.$$
(15)
Since only $`B_0`$ and $`B_0^{}`$ are generally known (for almost all the elements, see ref. ), we restrict ourselves instead to the first two terms in Eq. (14). Then the compression simplifies to
$$\eta =\left(1+\frac{B_0^{}}{B_0}p\right)^{1/B_0^{}}.$$
(16)
Eqs. (15) and (16) are two different approximations to the Murnaghan equation of state .
It then follows from Eqs. (12), (13) and (16) that the equation of the melting curve is
$$T_m(p)=T_m(0)\left(1+\frac{B_0^{}}{B_0}p\right)^{1/B_0^{}}\left[1+\frac{G_0^{}}{G_0}p\left(1+\frac{B_0^{}}{B_0}p\right)^{1/3B_0^{}}\right].$$
(17)
As discussed in Section 4, this equation is only valid for pressures $`p\stackrel{<}{}2B.`$
It follows from (17) that for $`pB_0`$
$$T_m(p)=T_m(0)\left[1+\left(\frac{B_0G_0^{}}{G_0}1\right)\left(\frac{p}{B_0}\right)\left(\frac{4}{3}\frac{B_0G_0^{}}{G_0}\frac{B_0^{}+1}{2}\right)\left(\frac{p}{B_0}\right)^2+\mathrm{}\right].$$
(18)
For the vast majority of the elements, $`B_0^{}>5/3`$ and $`B^{}`$ approaches 5/3 in the limit of large compressions. (In this limit $`p\eta ^{5/3},`$ and therefore $`BVdp/dV=\eta dp/d\eta =5p/3,`$ i.e., $`B^{}=5/3.)`$ In fact, the average value of $`B_0^{}`$ for the 65 elements analyzed in , except for Ce for which $`B_0^{}<0,`$ is $`4.30\pm 1.40.`$ Hence, if
$$\frac{G_0^{}}{G_0}>\frac{3}{8}\frac{B_0^{}+1}{B_0},$$
(19)
it follows from $`B_0^{}>5/3`$ that also $`G_0^{}/G_0>1/B_0,`$ i.e., Eq. (18) is of the form $`T_m(p)=T_m(0)(1+apbp^2+\mathrm{}),`$ $`a,b>0,`$ and describes melting curves for which melting temperatures increase with pressure . If, however,
$$\frac{G_0^{}}{G_0}<\frac{1}{B_0}$$
(20)
and $`B_0^{}>5/3,`$ then also $`G_0^{}/G_0<3/8(B_0^{}+1)/B_0,`$ i.e., Eq. (18) is of the form $`T_m(p)=T_m(0)(1ap+bp^2\mathrm{}),`$ $`a,b>0,`$ and describes melting curves for which melting temperatures initially decrease with pressure . For Si, for example, with the data from ref. we find $`G_0^{}/G_0<1/B_0`$ and $`B_0^{}=4.19,`$ in agreement with the negative initial slope of the experimental melting curve. Eqs. (19) and (20) plus $`B_0^{}>5/3`$ should be considered our criteria for the two types of melting curves discussed above.
## 3 Melting curves: comparison with data
In this section we compare our melting curve, Eq. (17), to some experimental melting curves, and predict a number of melting curves that can be compared with experiment in the not-so-distant future.
We have found 5 elements for which melting curves have been measured to higher pressures, $`pO(100`$ GPa): Al, Fe, Ni, Pb and U. We compare experimental data for these elements with our curves in Figs. 1-5. For Al, we also show the best fit to data in the form of the Simon equation, $`T_m(p)=T_m(0)(1+ap)^b`$ . For Fe, the experimental data are from ref. , and from ref. for Ni. For Pb, we combine the high-pressure data of ref. with the low-pressure data of ref. as corrected in ref. . For U, the high-pressure data of ref. are combined with the low-pressure data of ref. .
As claimed in ref. , the Simon equation may not be the best functional form for a fit to data. In fact, the initial slope provided by this equation for the Al melting curve is 80 K/GPa, in contrast to 59 and 65 K/GPa from the two previous low-pressure measurements . This accounts for the difference between the two curves in Fig. 1.
In Fig. 4, in addition to Fe, we also plot melting curves for Ru and Os, elements in the same column of the Periodic Table. Those curves should be considered predictions for these metals.
In Fig. 5, in addition to U for which there are high-pressure data, we also plot melting curves for the 5 actinides Am, Cm, Np, Pa and Th. For Np, we also show the low-pressure data of ref. . We do not show the low-pressure data of ref. for Am since they would overlay the low-pressure U data. We have checked that our melting curve is in agreement with the low-pressure Am data. For Cm, the values of $`B_0`$ and $`B_0^{}`$ are taken from ref. , and the value of $`G_0`$ is that estimated in ref. . For Pa, the values of $`B_0`$ and $`G_0`$ come from ref. . We estimate the values of $`G_0^{}`$ for Cm and Pa from Th and U, their neighbors in the same row in the Periodic Table. Our earlier $`G_0^{}`$ estimates for Am and Np lead to $`\gamma =1.05`$ and 1.09, respectively, in Eq. (24), which implies that such estimates are reliable. The values of $`B_0^{}`$ for Np and Pa are also estimated by interpolating between Am, Cm, Th and U. (We note that this estimation of $`B_0^{}`$ is justified by the pronounced periodic behavior of $`B_0^{}`$ in $`Z`$ .) The value of $`B_0^{}`$ for Am is taken from . We emphasize that the predicted melting curves assume constancy of coordination number along them. In the case of Am, e.g., there is still disagreement over the correct sequence of phases and their transition pressures , so this assumption may well be incorrect. For Th, however, it is claimed that there is a transition from a face-centered cubic structure to a body-centered tetragonal structure that changes coordination number . This transition occurs in the pressure range of $`70100`$ GPa , and thus our predictions for the Th melting curve up to 75 GPa should be quite reliable.
Although we can account for a decrease in melting temperature with pressure in our theoretical framework (Eqs. (17),(20)), we do not consider such cases here, among which there are Pu and Ce. It has been established that for Th, which is in the same column as Ce, $`\mathrm{}V>0,`$ and therefore, in view of Eq. (23), its melting temperature increases with pressure.
Fig. 1. Melting curve for Al. The dashed line is the Simon-fit to the data of ref. , which are not shown explicitly. The diamonds are the low-pressure data from ref. . The triangle is the shock-melting point at 125 GPa from ref. . The boxes are the points at 25, 69 and 137 GPa calculated in ref. from shock-melting data. They are assigned 20% error bars .
Fig. 2. Melting curve for Ni. The diamonds (with small error bars) are the data of ref. . The boxes are the points at 79 and 250 GPa calculated in ref. from shock-melting data. The corresponding error bars are not quoted in ref. .
Fig. 3. Melting curve for Pb. The diamonds are the low-pressure data of ref. corrected as in ref. . The triangles are the data from ref. , and the dashed line is a best fit to the data. The boxes are the points at 12, 34 and 68 GPa calculated in ref. from shock-melting data. The corresponding error bars are not quoted in ref. . The star is the point at 118 GPa calculated in ref. from shock-melting data.
Fig. 4. Melting curves for the elements of the iron group (Fe, Ru, Os). The diamonds are the data of ref. , and the dashed line is a best fit to the data. The boxes are the shock-melting points at 235 GPa and 300 GPa . The triangle is the shock-melting point at 240 GPa . The star is the shock-melting point at 243 GPa .
Fig. 5. Melting curves for the actinides Am, Cm, Np, Pa, Th and U. The data for U are from ref. , and for Np they come from ref. .
In Fig. 6 we plot the low-pressure data of ref. for W, and our melting curves for W, Mo and Cr. The initial slope of our melting curve of Mo, 26 K/GPa, is consistent with that predicted in ref. : $`(34\pm 6)`$ K/GPa. The same melting curve gives $`T_m9650`$ K at $`p=390`$ GPa, in good agreement with the shock-melting temperature $`10000`$ K at the same pressure, found in ref. .
Fig. 6. Melting curves for the elements of the chromium group (Cr, Mo, W). The data for W are from ref. .
In Fig. 7 we compare the low-pressure data of ref. for Cu, Ag and Au with our corresponding melting curves. Although the initial slopes of these curves are somewhat less than those of the data (the corresponding values of $`\gamma `$ in Fig. 1 are $`0.8),`$ they are in good agreement with the best extrapolation of data to higher pressures made in ref. , and with the calculation of ref. in the case of Cu.
Fig. 7. Melting curves for the elements of the copper group (Cu, Ag, Au). The diamonds, the stars, and the triangles are the low-pressure data of ref. in the inset, and the best extrapolations of these data to 20 and 30 GPa in the main plot for Cu, Ag and Au, respectively. The boxes are the points at 45 and 128 GPa calculated in ref. from shock-melting data for Cu. The corresponding error bars are not quoted . The gray diamond is the shock-melting point for Cu at 37 GPa .
In Fig. 8 we compare the low-pressure data on the noble gases to our corresponding melting curves. The unknown values of $`G_0^{}`$ for Ne, Ar, Kr and Xe are calculated with the help of Eq. (25) below using the measured values of $`B_0,`$ $`G_0`$ and $`\gamma _0`$ (the Grüneisen constant) . The values of $`B_0^{}`$ for Ne, Ar, Kr and Xe are taken from . In the case of Rn, for which $`B_0,`$ $`B_0^{},`$ $`G_0`$ and $`G_0^{}`$ have not been measured, we first calculate $`G_0`$ using the (approximate) relation $`GV/T_m=\mathrm{const}`$ for the noble-gas group, where $`V=v_{WS}N_A`$ is the molar volume. This relation follows from Eq. (1) provided that $`\kappa ,`$ $`\lambda ,`$ $`\alpha `$ and $`z`$ do not vary within this group. The value of the constant is determined by using the corresponding Ne, Ar, Kr and Xe data in this relation. We then calculate $`B_0`$ using Eq. (7) with the value of the Poisson ratio for Rn determined by extrapolating from the corresponding values for Ne, Ar, Kr and Xe. Finally, we determine both $`B_0^{}`$ and $`G_0^{}`$ by again extrapolating the Ne, Ar, Kr and Xe data.
Fig. 8. Melting curves for the noble gases Ne, Ar, Kr, Xe and Rn. The diamonds are the data of ref. . The stars are the data of ref. . The boxes and triangles come from ref. .
Finally, in Fig. 9 we show the experimental data and our theoretical melting curve for Mg.
Fig. 9. Melting curve for Mg. The low-pressure data are from ref. . The high-pressure data are the shock-melting points of ref. .
The 24 melting curves considered above constitute convincing evidence for the validity of our formula for melting temperature as a function of pressure, Eq. (17).
## 4 The range of validity of the new melting curve equation
In deriving our melting curve, Eq. (17), we have used both Eq. (13) for the pressure dependence of the shear modulus and the Murnaghan equation of state, Eq. (16). Since Eq. (13) has the correct zero-pressure limit (its Taylor series expansion in $`p`$ at $`p=0`$ is $`G=G_0+G_0^{}p(G_0^{}/3B_0)p^2+\mathrm{})`$ and is claimed to have the correct ultra-high-pressure limit , we assume that this equation is valid over the entire pressure range. In any event, we do not have data to either confirm or invalidate this assumption. It then follows that the range of validity of Eq. (17) depends crucially on the range of validity of the Murnaghan equation of state, Eq. (16).
The Murnaghan equation of state was examined in ref. , together with a number of different equations of state, by comparing with the theoretical results calculated by the augmented-plane-wave method and the quantum-mechanical model proposed by Kalitkin and Kuz’mina from low to ultra-high pressures. It was shown that the Murnaghan equation is in good agreement with the theoretical results up to $`V/V_00.7,`$ i.e., up to compressions $`1.41.5.`$ Since for the vast majority of the elements $`B_0^{}5`$ , we conclude, on the basis of Eq. (16), that the Murnaghan equation, and consequently, our equation for melting curve, Eq. (17), is valid up to pressures $`p2B_0.`$ The melting curves for Al and Pb in Figs. 1 and 3, respectively, and for Ne and Ar in Fig. 8 show that in some cases Eq. (17) is good to pressures even greater than $`2B_0.`$
For reliable predictions of melting curves to much higher pressures, $`p\stackrel{>}{}1`$ TPa, one has to use a better equation of state than Murnaghan’s. Hama and Suito claim that the Vinet equation of state is consistent with first-principles theoretical calculations to compressions $`\eta 5.`$ Reference also finds the Vinet equation of state to be most accurate among various suggested equations of state. In fact, we have calculated that the melting temperatures for Fe and Pb at pressures $`p50B_0`$ as given by Eq. (17) are about two times higher than those given by a relation that derives from Eqs. (12), (13), and the Vinet equation of state.
Another possible source of disagreement between the new melting curve, Eq. (17), and data may be inaccurate values of elastic constants and their pressure derivatives in some cases. The Murnaghan equation and its frequently used partner – the Birch equation – are derived from the second-order Taylor series expansion of the bulk modulus \[as in Eq. (14)\] or the elastic strain energy with respect to pressure or strain, respectively. Thus their validities are, in principle, restricted to a narrow range of compression. Extending this range would entail the inclusion of higher-order terms. This could explain why the values of $`B_0,`$ and especially those of $`B_0^{}`$ and $`B_0^{\prime \prime },`$ obtained from experiments which cover different ranges of compression by using a fitting method, are usually different. In many cases these differences between different experiments are small and can be safely neglected. In some cases, however, they are large, and so their use for predicting physical observables, such as melting temperature, is dubious. For example, in the case of Ni, we have used the value $`B_0^{}=6.20`$ given in . Reference , however, quotes $`B_0^{}30`$ (!). Similarly, for Mo we have used $`B_0^{}=4.4`$ of ref. , while ref. gives $`B_0^{}20.`$ (We note that the use of the values $`B_0^{}=6.20`$ for Ni and 4.4 for Mo is justified in view of the recent compilations of experimental data on $`B_0^{}`$ .) Although the numerical value of $`B_0^{}`$ does not matter at low $`p,`$ since it first appears in the $`(p/B_0)^2`$ term, in view of Eq. (15), it would strongly affect the predicted melting curve at pressures $`pO(B_0).`$
There are also inconsistencies in the values of $`G_0`$ quoted in the literature. For example, for Pb we use the value $`G_0=8.6`$ GPa from ref. , whereas ref. quotes $`G=5.5`$ GPa. (Our own calculation , based on the values of the elastic constants $`c_{11},`$ $`c_{12}`$ and $`c_{44},`$ shows that 8.6 GPa is preferred over 5.5 GPa.) Likewise the values of $`G_0`$ for K and Na from ref. are 0.9 and 1.98 GPa, whereas ref. quotes 1.3 and 3.5 GPa, respectively.
## 5 Relation of dislocation-based melting relation to the Lindemann criterion
The well-known Lindemann melting rule is based on the assumption that all elemental solids melt when the atomic vibrational amplitude is a fixed pressure-independent fraction of the interatomic distance. As shown by Lindemann , this implies the invariance of the Lindemann number
$$\theta _D\left(\frac{M}{T_m}\right)^{1/2}V^{1/3}=L$$
(21)
along the melting curve. Here $`\theta _D`$ is the density-dependent Debye temperature, $`V`$ is the molar volume, and $`M`$ is the molar mass. It is found that $`L150`$ .
There are compelling reasons to suppose that our dislocation-based melting relation is somehow equivalent to the Lindemann criterion. First of all, Eq. (21) gives melting curves that are typically very close to those predicted by our dislocation-based melting relation. (For example, the melting curve for Mg in Fig. 6.3 of ref. is very similar to our curve in Fig. 9.) Furthermore, the Lindemann number, which is proportional to the ratio of atomic vibrational amplitude to the lattice constant at melt, is analogous to $`b^2\rho (T_m),`$ since both are presumed constant along the melting curve. In the dislocation-based approach, melting is associated with a critical configuration of dislocations, and for any such configuration there is a corresponding mean displacement of atoms from their equilibrium positions. Hence $`L`$ and $`b^2\rho (T_m)`$ are clearly related, and therefore the left-hand sides of Eqs. (9) and (21) are related as well.
The mathematical equivalence of our melting relation and the Lindemann criterion would be established if it could be determined that the left-hand side of Eq. (21) is a fixed fraction of the left-hand side of Eq. (9). A search of the literature has turned up two results which show that $`L^2`$ is approximately proportional to $`Gv_{WS}/T_m.`$ For a Debye solid the relation is $`L^2=f(\nu (p,T))Gv_{WS}/T_m,`$ where $`f`$ is a complicated function of $`\nu .`$ Thus the two melting relations are not rigorously equivalent.
A second connection between the two melting formulas is provided by the following approximation for the Grüneisen constant ,
$$\gamma (p)=\frac{2}{3}\gamma _S(p)+\frac{1}{3}\gamma _L(p),$$
(22)
where
$$\gamma _S(p)=\frac{G^{}(p)}{2}\frac{B^T(p)}{G(p)}\frac{1}{6},$$
(23)
$$\gamma _L(p)=\frac{1}{2}\frac{B^T(p)}{B^S(p)+\frac{4}{3}G(p)}\frac{d(B^S(p)+\frac{4}{3}G(p))}{dp}\frac{1}{6},$$
(24)
are the contributions of the shear (transverse) and longitudinal acoustic modes. Here $`B^T`$ is the isothermal bulk modulus, which is equivalent to $`B`$ that we are using in this paper. Eqs. (22)-(24) follow from the two assumptions that (i) the only appreciable contribution to the heat capacity of a crystal arises from lattice vibrations, and (ii) averaging over all modes is equivalent to averaging only over the low-frequency acoustic modes. (I.e., the contribution of the optical modes is equal to that of the acoustic modes.) If in addition it is assumed that $`B^S(p),`$ the isentropic bulk modulus, is proportional to $`G(p),`$ then
$$\gamma _S(p)=\gamma _L(p)=\gamma (p)=\frac{G^{}(p)}{2}\frac{B(p)}{G(p)}\frac{1}{6}.$$
(25)
However, there is no basis for this assumption, i.e., $`\gamma _S(p)\gamma _L(p)`$ is to be expected. For example, in the ultra-high pressure limit, $`B^S(p)B^T(p)=5p/3`$ and $`G(p)p^{4/5},`$ quite different dependencies.
Integration of Eq. (25), using $`B(p)=dp/d\mathrm{ln}V(p)`$ and $`\gamma (p)=d\mathrm{ln}\theta _D(p)/d\mathrm{ln}V(p),`$ gives
$$\frac{\theta _D^2(p)V^{2/3}(p)/T_m(p)}{G(p)V(p)/T_m(p)}=\mathrm{const},$$
(26)
that is, $`L^2Gv_{WS}/T_m.`$ We emphasize that this proportionality is founded on a number of uncontrolled approximations.
Equation (25), which would ensure a rigorous mathematical equivalence of the defect and mechanical (Lindemann’s) approaches to melting, does not follow from first principles. This means that the defect and mechanical approaches to melting are basically different. Moreover, since the mechanical approach does not have a solid thermodynamic basis, it cannot, for example, predict the latent heat of fusion. In contrast, the defect approach predicts the latent heat of fusion, Eq. (9), which is in good agreement with data for three-quarters of the Periodic Table .
Finally, we wish to make the following comments on Eq. (25). We did not check extensively its validity at zero pressure, since that would go beyond the scope of this paper. We do, however, have some evidence that Eq. (25) is rather well satisfied: with the data from ref. , we calculate from the above relation $`\gamma _0=2.25`$ vs. measured 2.40 for Ag, 3.08 vs. 2.99 for Au, 1.66 vs. 1.78 for Fe, 1.28 vs. 1.29 for K, and 1.18 vs. 1.19 for Na. We have actually used Eq. (25) in Section 3 to calculate $`G_0^{}`$ for noble gases in order to get their melting curves via Eq. (17) and to compare with experiment. Good agreement between the calculated and experimental curves is another hint on the approximate validity of this formula. Also, Eq. (25) has the correct ultra-high-pressure limit in which $`\gamma 1/2`$ , since in this limit $`Gp^{4/5}`$ and $`B=5p/3.`$
## 6 Concluding remarks
We have extended the framework of melting as a string-mediated phase transition to non-zero pressure and derived a new equation for the melting curve, Eq. (17). As discussed above, with accurate experimental values of all the parameters involved, this equation reproduces the existing experimental melting data, and predicts unknown melting curves to pressures $`p\stackrel{<}{}2B_0.`$ For higher pressures, a better equation of state than Murnaghan’s should be used, e.g., the Vinet equation of state.
We have addressed the apparent equivalence of defect and mechanical approaches to melting curve, and demonstrated that both approaches are basically different. We have shown that their would-be rigorous mathematical equivalence must manifest itself in a new relation, Eq. (25), which we have not tested in detail.
To summarize, we have calculated melting curves for 24 elements: Al, Mg, Ni, Pb, the iron group (Fe, Ru, Os), the chromium group (Cr, Mo, W), the copper group (Cu, Ag, Au), the noble gases Ne, Ar, Kr, Xe and Rn, and the six actinides Am, Cm, Np, Pa, Th and U. These calculated melting curves are in good agreement with existing data.
## Acknowledgements
We wish to thank T. Goldman and A.Z. Patashinski for valuable discussions during the preparation of this work. One of us (L.B.) wishes to thank B.K. Godwal for useful correspondence.
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# Hypersequents and the Proof Theory of Intuitionistic Fuzzy Logic2000 Mathematics Subject Classification: Primary 03B50; Secondary 03B55, 03F05.Research supported by the Austrian Science Fund under grant P–12652 MAT
## 1 Introduction
Intuitionistic fuzzy logic IF was originally defined by Takeuti and Titani to be the logic of the complete Heyting algebra $`[0,1]`$. In standard many-valued terminology, IF is $`[0,1]`$-valued first-order Gödel logic, with truth functions as defined below. The finite-valued propositional versions of this logic were introduced by Gödel , and have spawned a sizeable area of logical research subsumed under the title “intermediate logics” (intermediate between classical and intuitionistic logic). The infinite-valued propositional Gödel logic was studied by Dummett , who showed that it is axiomatized by LC, i.e., intuitionistic propositional logic plus the linearity axiom $`(AB)(BA)`$.
Takeuti and Titani characterized IF by a calculus which extends the intuitionistic predicate calculus LJ by several axioms as well as the density rule
$$\begin{array}{cc}\mathrm{\Gamma }A(Cp)(pB)& tt^{}\\ \stackrel{\mathit{}}{\mathrm{\Gamma }A(CB)}\end{array}$$
This rule can be read as expressing the fact that the set of truth values is densely ordered. In this sense, the Takeuti-Titani axiomatization is the natural axiomatization of the $`[0,1]`$-valued Gödel logic. The valid formulas of IF are also characterized as those formulas valid in every first-order Gödel logic based on a linearly ordered set of truth-values (this is obvious for all logics based on truth value sets $`[0,1]`$, since a countermodel in such a truth-value set can be straightforwardly embedded in $`[0,1]`$. The general claim was established by Horn ). In this characterization, the density rule is not a natural assumption, since not every linearly ordered truth-value set is densely ordered. It follows from this characterization that the density rule is redundant for the axiomatization of IF, and completeness proofs without it have been given by Horn and Takano .<sup>1</sup><sup>1</sup>1Note that the corresponding axiom $`(p)((Ap)(pB))(AB)`$ is not redundant in quantified *propositional* $`[0,1]`$-valued Gödel logic. See . Takano posed the question of whether a syntactic elimination of the density rule is also possible.
More recently, another axiomatizable first-order extension of LC has been studied by Corsi and Avellone et al. . This extension is defined not via many-valued semantics but as the class of formulas valid in all linearly ordered intuitionistic Kripke models. It is different from IF; specifically, the formula $`()`$ below is not valid in it. IF can, however, also be characterized as the set of formulas valid in all linearly ordered Kripke models with constant domains (this was first observed by Gabbay \[7, §3\]).
The interest of IF lies in the fact that it combines properties of logics for approximate reasoning with properties of intuitionistic logic. On the one hand, IF is one of the basic $`t`$-norm logics (see Hájek ), on the other, it is an extension of intuitionistic logic which corresponds to concurrency (as has been argued by Avron ). We present here a calculus for IF which is adequate for further proof-theoretic study. The basic result in this regard is the cut-elimination theorem for this calculus, from which a midhypersequent-theorem can be derived. This theorem, in turn, corresponds to Herbrand’s Theorem in classical logic, and as such is a possible basis for automated theorem proving in IF.
The calculus also allows us to investigate the proof-theoretic effects of the Takeuti-Titani rule. We give a positive answer to Takano’s question, showing that the density rule can be eliminated from IF-proofs. A simple example illustrates the possible structural differences between proofs with and without the Takeuti-Titani rule.
## 2 Syntax and Semantics of Intuitionistic Fuzzy Logic
The language $`L`$ of IF is a usual first-order language with propositional variables and where free ($`a`$, $`b`$, …) and bound ($`x`$, $`y`$, …) variables are distinguished.
###### Definition 2.1
An IF-interpretation $`\mathrm{}=D,\text{s}`$ is given by the domain $`D`$ and the valuation function $`\text{s}`$. Let $`L^D`$ be $`L`$ extended by constants for each element of $`D`$. Then $`\text{s}`$ maps atomic formulas in $`\mathrm{Frm}(L^D)`$ into $`[0,1]`$, $`dD`$ to itself, $`n`$-ary function symbols to functions from $`D^n`$ to $`D`$, and free variables to elements of $`D`$.
The valuation function $`\text{s}`$ can be extended in the obvious way to a function on all terms. The valuation for formulas is defined as follows:
1. $`AP(t_1,\mathrm{},t_n)`$ is atomic: $`\mathrm{}(A)=\text{s}(P)(\text{s}(t_1),\mathrm{},\text{s}(t_n))`$.
2. $`A\neg B`$:
$$\mathrm{}(\neg B)=\{\begin{array}{cc}0\hfill & \text{if }\mathrm{}(B)0\hfill \\ 1\hfill & \text{otherwise.}\hfill \end{array}$$
3. $`ABC`$: $`\mathrm{}(BC)=\mathrm{min}(\mathrm{}(B),\mathrm{}(C))`$.
4. $`ABC`$: $`\mathrm{}(BC)=\mathrm{max}(\mathrm{}(A),\mathrm{}(B))`$.
5. $`ABC`$:
$$\mathrm{}(BC)=\{\begin{array}{cc}\mathrm{}(C)\hfill & \text{if }\mathrm{}(B)>\mathrm{}(C)\hfill \\ 1\hfill & \text{if }\mathrm{}(B)\mathrm{}(C)\text{.}\hfill \end{array}$$
The set $`\mathrm{Distr}_{\mathrm{}}(A(x))=\{\mathrm{}(A(d)):dD\}`$ is called the distribution of $`A(x)`$. The quantifiers are, as usual, defined by infimum and supremum of their distributions.
1. $`A(x)B(x)`$: $`\mathrm{}(A)=inf\mathrm{Distr}_{\mathrm{}}(B(x))`$.
2. $`A(x)B(x)`$: $`\mathrm{}(A)=sup\mathrm{Distr}_{\mathrm{}}(B(x))`$.
$`\mathrm{}`$ satisfies a formula $`A`$, $`\mathrm{}A`$, if $`\mathrm{}(A)=1`$. A formula $`A`$ is IF-valid if every IF-interpretation satisfies it.
Note that, as in intuitionistic logic, $`\neg A`$ may be defined as $`A`$, where $``$ is some formula that always takes the value 0.
## 3 Hypersequents and IF
Takeuti and Titani’s system IF is based on Gentzen’s sequent calculus LJ for intuitionistic logic with a number of extra axioms
$$\begin{array}{c}(AB)((AB)B)\\ (AB)B(BA)B\\ (AB)C(AC)(BC)\\ (A(BC))(AB)(AC)\\ (x)(A(x)B)(x)A(x)B\\ (x)A(x)C(x)(A(x)D)(DC)\end{array}$$
$`\begin{array}{c}\hfill (\text{Ax}1)\\ \hfill (\text{Ax}2)\\ \hfill (\text{Ax}3)\\ \hfill (\text{Ax}4)\\ \hfill ()\\ \hfill ()\end{array}`$
(where $`x`$ does not occur in $`B`$ or $`D`$) and the following additional inference rule:
$$\begin{array}{cc}\mathrm{\Gamma }A(Cp)(pB)& \mathrm{𝑡𝑡}^{}\\ \stackrel{\mathit{}}{\mathrm{\Gamma }A(CB)}\end{array}$$
where $`p`$ is a propositional eigenvariable (i.e., it does not occur in the lower sequent). It is known that the extra inference rule is redundant. In fact, the system H of Horn consisting of LJ plus the schemata
$$\begin{array}{c}(x)(A(x)B)(x)A(x)B\\ (AB)(BA)\end{array}$$
$`\begin{array}{c}\hfill ()\\ \hfill (D)\end{array}`$
is complete for IF (see also ). Neither of these systems, however, has decent proof-theoretic properties such as cut elimination, nor is a syntactic method for the elimination of the Takeuti-Titani rule ($`\mathrm{𝑡𝑡}^{}`$) known. Takano has posed the question of a syntactic elimination procedure of the Takeuti-Titani rule as an open problem.
We present a system which has the required properties, and which allows the syntactic elimination of the Takeuti-Titani rule. Our system is based on Avron’s cut-free axiomatization of LC using a hypersequent calculus.
###### Definition 3.1
A sequent is an expression of the form
$$\mathrm{\Gamma }\mathrm{\Delta }$$
where $`\mathrm{\Gamma }`$ and $`\mathrm{\Delta }`$ are finite multisets of formulas, and $`\mathrm{\Delta }`$ contains at most one formula. A hypersequent is a finite multiset of sequents, written as
$$\mathrm{\Gamma }_1\mathrm{\Delta }_1\mathrm{}\mathrm{\Gamma }_n\mathrm{\Delta }_n$$
The hypersequent calculus $`\mathrm{𝐇𝐈𝐅}`$ has the following axioms and rules:
Axioms: $`AA`$, for any formula $`A`$.
Internal structural rules:
$$\begin{array}{cc}G\mathrm{\Gamma }\mathrm{\Delta }& iw\\ \stackrel{\mathit{}}{GA,\mathrm{\Gamma }\mathrm{\Delta }}\end{array}\begin{array}{cc}G\mathrm{\Gamma }& iw\\ \stackrel{\mathit{}}{G\mathrm{\Gamma }A}\end{array}\begin{array}{cc}GA,A,\mathrm{\Gamma }\mathrm{\Delta }& ic\\ \stackrel{\mathit{}}{GA,\mathrm{\Gamma }\mathrm{\Delta }}\end{array}$$
External structural rules:
$$\begin{array}{cc}G& ew\\ \stackrel{\mathit{}}{G\mathrm{\Gamma }\mathrm{\Delta }}\end{array}\begin{array}{cc}G\mathrm{\Gamma }\mathrm{\Delta }\mathrm{\Gamma }\mathrm{\Delta }& ec\\ \stackrel{\mathit{}}{G\mathrm{\Gamma }\mathrm{\Delta }}\end{array}$$
Logical rules:
$$\begin{array}{cc}\begin{array}{cc}G\mathrm{\Gamma }A& \neg \\ \stackrel{\mathit{}}{G\neg A,\mathrm{\Gamma }}\end{array}& \begin{array}{cc}GA,\mathrm{\Gamma }& \neg \\ \stackrel{\mathit{}}{G\mathrm{\Gamma }\neg A}\end{array}\\ \begin{array}{cc}GA,\mathrm{\Gamma }\mathrm{\Delta }GB,\mathrm{\Gamma }\mathrm{\Delta }& \\ \stackrel{\mathit{}}{GAB,\mathrm{\Gamma }\mathrm{\Delta }}\end{array}& \begin{array}{cc}G\mathrm{\Gamma }AG\mathrm{\Gamma }B& \\ \stackrel{\mathit{}}{G\mathrm{\Gamma }AB}\end{array}\\ \begin{array}{cc}G\mathrm{\Gamma }A& _1\\ \stackrel{\mathit{}}{G\mathrm{\Gamma }AB}\end{array}& \begin{array}{cc}GA,\mathrm{\Gamma }\mathrm{\Delta }& _1\\ \stackrel{\mathit{}}{GAB,\mathrm{\Gamma }\mathrm{\Delta }}\end{array}\\ \begin{array}{cc}G\mathrm{\Gamma }B& _2\\ \stackrel{\mathit{}}{G\mathrm{\Gamma }AB}\end{array}& \begin{array}{cc}GB,\mathrm{\Gamma }\mathrm{\Delta }& _2\\ \stackrel{\mathit{}}{GAB,\mathrm{\Gamma }\mathrm{\Delta }}\end{array}\\ \begin{array}{cc}G\mathrm{\Gamma }_1AGB,\mathrm{\Gamma }_2\mathrm{\Delta }& \\ \stackrel{\mathit{}}{GAB,\mathrm{\Gamma }_1,\mathrm{\Gamma }_2\mathrm{\Delta }}\end{array}& \begin{array}{cc}GA,\mathrm{\Gamma }B& \\ \stackrel{\mathit{}}{G\mathrm{\Gamma }AB}\end{array}\\ \begin{array}{cc}GA(t),\mathrm{\Gamma }\mathrm{\Delta }& \\ \stackrel{\mathit{}}{G(x)A(x),\mathrm{\Gamma }\mathrm{\Delta }}\end{array}& \begin{array}{cc}G\mathrm{\Gamma }A(a)& \\ \stackrel{\mathit{}}{G\mathrm{\Gamma }(x)A(x)}\end{array}\\ \begin{array}{cc}GA(a),\mathrm{\Gamma }\mathrm{\Delta }& \\ \stackrel{\mathit{}}{G(x)A(x),\mathrm{\Gamma }\mathrm{\Delta }}\end{array}& \begin{array}{cc}G\mathrm{\Gamma }A(t)& \\ \stackrel{\mathit{}}{G\mathrm{\Gamma }(x)A(x)}\end{array}\end{array}$$
Cut:
$$\begin{array}{cc}G\mathrm{\Gamma }AGA,\mathrm{\Pi }\mathrm{\Lambda }& cut\\ \stackrel{\mathit{}}{G\mathrm{\Gamma },\mathrm{\Pi }\mathrm{\Lambda }}\end{array}$$
Communication:
$$\begin{array}{cc}G\mathrm{\Theta }_1,\mathrm{\Theta }_1^{}\mathrm{\Xi }_1G\mathrm{\Theta }_2,\mathrm{\Theta }_2^{}\mathrm{\Xi }_2& cm\\ \stackrel{\mathit{}}{G\mathrm{\Theta }_1,\mathrm{\Theta }_2^{}\mathrm{\Xi }_1\mathrm{\Theta }_1^{},\mathrm{\Theta }_2\mathrm{\Xi }_2}\end{array}$$
Density:
$$\begin{array}{cc}G\mathrm{\Phi }pp,\mathrm{\Psi }\mathrm{\Sigma }& \mathrm{𝑡𝑡}\\ \stackrel{\mathit{}}{G\mathrm{\Phi },\mathrm{\Psi }\mathrm{\Sigma }}\end{array}$$
The rules ($``$), $`()`$, and $`(\mathrm{𝑡𝑡})`$ are subject to eigenvariable conditions: the free variable $`a`$ and the propositional variable $`p`$, respectively, must not occur in the lower hypersequent. We denote the calculus obtained from $`\mathrm{𝐇𝐈𝐅}`$ by omitting the cut rule by $`\mathrm{𝐇𝐈𝐅}^{}`$, and that obtained by omitting (tt) by $`\mathrm{𝐇𝐈𝐅}^{}`$.
The semantics of IF can easily be extended to hypersequents by mapping a hypersequent $`H`$
$$\mathrm{\Gamma }_1\mathrm{\Delta }_1\mathrm{}\mathrm{\Gamma }_n\mathrm{\Delta }_n$$
to the formula $`H^{}`$
$$(\mathrm{\Gamma }_1\mathrm{\Delta }_1)\mathrm{}(\mathrm{\Gamma }_n\mathrm{\Delta }_n)$$
where $`\mathrm{\Gamma }_i`$ denotes the conjunction of the formulas in $`\mathrm{\Gamma }_i`$ or $``$ if $`\mathrm{\Gamma }_i`$ is empty, and $`\mathrm{\Delta }_i`$ the disjunction of the formulas in $`\mathrm{\Delta }_i`$ or $``$ if $`\mathrm{\Delta }_i`$ is empty. Deriving a formula $`A`$ in $`\mathrm{𝐇𝐈𝐅}`$ then is equivalent to deriving the sequent $`A`$: the translation of $`A`$, i.e., $`A`$ is equivalent to $`A`$.
###### Theorem 3.2 (Soundness)
Every hypersequent $`H`$ derivable in $`\mathrm{𝐇𝐈𝐅}`$ is IF-valid.
###### Proof
By induction on the length of the proof. It will suffice to show that the axioms are valid, and that the quantifier rules and (tt) preserve validity.
The soundness of the quantifier rules is established by observing that corresponding quantifier shifting rules are intuitionistically valid. For instance, since
$$\begin{array}{c}(x)(BA(x))(B(x)A(x))\\ (x)(BA(x))B(x)A(x)\end{array}$$
$`\begin{array}{c}\hfill ()\\ \hfill ()\end{array}`$
are intuitionistically valid, it is easily seen that $``$ is a sound rule. The only problematic rules are $`()`$ and $`()`$. Suppose $`G\mathrm{\Gamma }A(a)`$ is derivable in $`\mathrm{𝐇𝐈𝐅}`$. By induction hypothesis, $`G^{}(\mathrm{\Gamma }A(a))`$ is valid. Then certainly $`(x)(G^{}(\mathrm{\Gamma }A(x)))`$ is IF-valid. Since $`a`$ did not occur in $`G`$ or $`\mathrm{\Gamma }`$, we may now assume that $`x`$ does not either. Since the quantifier shift $`()`$, i.e.,
$$(x)(BA(x))(B(x)A(x)),$$
is valid in IF, we see that $`G^{}(x)(\mathrm{\Gamma }A(x))`$ is valid. The result follows since
$$(x)(BA(x))B(x)A(x)$$
is intuitionistically valid, and hence IF-valid.
The communication rule is sound as well. Suppose the interpretation $`\mathrm{}`$ satisfies the premises of (cm). The only case where the conclusion is not obviously also satisfied is if $`\mathrm{}(\mathrm{\Theta }_1^{})\mathrm{}(\mathrm{\Xi }_1)`$ and $`\mathrm{}(\mathrm{\Theta }_2^{})\mathrm{}(\mathrm{\Xi }_2)`$. If the left lower sequent is not satisfied, we have $`\mathrm{}(\mathrm{\Xi }_1)<\mathrm{}(\mathrm{\Theta }_2^{})`$, and hence $`\mathrm{}(\mathrm{\Theta }_1^{})\mathrm{}(\mathrm{\Xi }_2)`$, and thus the right lower sequent is satisfied. Similarly if the right lower sequent is not satisfied.
For (tt) we may argue as follows: Suppose that the hypersequent
$$H=G\mathrm{\Phi }pp,\mathrm{\Psi }\mathrm{\Sigma }$$
is IF-valid. Let $`\mathrm{}`$ be an interpretation, and let $`\mathrm{}_r`$ be just like $`\mathrm{}`$ except that $`\mathrm{}(p)=r`$. Since $`p`$ does not occur in the conclusion hypersequent
$$H^{}=G\mathrm{\Phi },\mathrm{\Psi }\mathrm{\Sigma }$$
we have $`\mathrm{}(H^{})=\mathrm{}_r(H^{})`$ and $`\mathrm{}(G)=\mathrm{}_r(G)`$. If $`\mathrm{}G`$ we are done. Otherwise, assume that $`\mathrm{}\overline{)}H^{}`$, i.e.,
$$r_1=\mathrm{min}\{\mathrm{}(\mathrm{\Phi }),\mathrm{}(\mathrm{\Psi })\}>\mathrm{}(\mathrm{\Sigma })=r_2$$
Let $`r=(r_1+r_2)/2`$. Now consider $`\mathrm{}_r`$: $`\mathrm{}_r\overline{)}G`$ by assumption; $`\mathrm{}_r\overline{)}\mathrm{\Phi }p`$, since $`\mathrm{}_r(\mathrm{\Phi })>r`$; and $`\mathrm{}_r\overline{)}p,\mathrm{\Psi }\mathrm{\Sigma }`$, since $`\mathrm{}_r(\mathrm{\Psi })>r>\mathrm{}_r(\mathrm{\Sigma })`$. Hence, $`\mathrm{}_r\overline{)}H`$, a contradiction. ∎
###### Theorem 3.3 (Completeness)
Every IF-valid hypersequent is derivable in $`\mathrm{𝐇𝐈𝐅}`$.
###### Proof
Observe that a hypersequent $`H`$ and its canonical translation $`H^{}`$ are interderivable using the cut rule and the following derivable hypersequents
$$\begin{array}{cc}ABAABB& AB,AB\\ ABA& AAB\end{array}$$
Thus it suffices to show that the characteristic axioms of IF are derivable; a simple induction on the length of proofs shows that proofs in intuitionistic predicate calculus together with the axioms (D) and $`()`$ can be simulated in $`\mathrm{𝐇𝐈𝐅}`$. The formula (D) is easily derivable using the communication rule.
$$\begin{array}{cc}\begin{array}{cc}\begin{array}{cc}\begin{array}{cc}\begin{array}{cc}\begin{array}{cc}AABB& \text{cm}\\ \stackrel{\mathit{}}{ABBA}\end{array}& \\ \stackrel{\mathit{}}{ABBA}\end{array}& \\ \stackrel{\mathit{}}{ABBA}\end{array}& \\ \stackrel{\mathit{}}{(AB)(BA)BA}\end{array}& \\ \stackrel{\mathit{}}{(AB)(BA)(AB)(BA)}\end{array}& ec\\ \stackrel{\mathit{}}{(AB)(BA)}\end{array}$$
The formula $`()`$ can be obtained thus:
$$\begin{array}{cc}\begin{array}{cc}\begin{array}{cc}\begin{array}{cc}\begin{array}{cc}\begin{array}{cc}\begin{array}{cc}A\left(a\right)A\left(a\right)BB& cm\\ \stackrel{\mathit{}}{BA\left(a\right)A\left(a\right)B}\end{array}\begin{array}{cc}BB& ew\\ \stackrel{\mathit{}}{BA\left(a\right)BB}\end{array}& \\ \stackrel{\mathit{}}{BA\left(a\right)BA\left(a\right)B}\end{array}\begin{array}{cc}A\left(a\right)A\left(a\right)& ew\\ \stackrel{\mathit{}}{A\left(a\right)A\left(a\right)BA\left(a\right)B}\end{array}& \\ \stackrel{\mathit{}}{BA\left(a\right)A\left(a\right)BA\left(a\right)B}\end{array}& \\ \stackrel{\mathit{}}{\left(x\right)\left(BA\left(x\right)\right)A\left(a\right)BA\left(a\right)B}\end{array}& \\ \stackrel{\mathit{}}{\left(x\right)\left(BA\left(x\right)\right)A\left(a\right)\left(x\right)\left(BA\left(x\right)\right)B}\end{array}& \\ \stackrel{\mathit{}}{\left(x\right)\left(BA\left(x\right)\right)\left(x\right)A\left(x\right)\left(x\right)\left(BA\left(x\right)\right)B}\end{array}& \\ \stackrel{\mathit{}}{\stackrel{\mathit{}}{\left(x\right)\left(BA\left(x\right)\right)B\left(x\right)A\left(x\right)}}\end{array}$$
The last line is obtained from the preceding by two ($``$) inferences, followed by an external contraction. We indicate this with the double inference line. ∎
Of course, the other axioms of Takeuti’s and Titani’s system are also derivable. We will leave the propositional axioms 1–4 as an exercise to the reader, and give the derivation on of ($``$) as another example:
$$\begin{array}{cc}\begin{array}{cc}\begin{array}{cc}\begin{array}{cc}\begin{array}{cc}\begin{array}{cc}\begin{array}{cc}A(a)A(a)DD& cm\\ \stackrel{\mathit{}}{A(a)DDA(a)}\end{array}& \\ \stackrel{\mathit{}}{A(a)DDA(a)}\end{array}& \\ \stackrel{\mathit{}}{(x)(A(x)D)DA(a)}\end{array}& \\ \stackrel{\mathit{}}{(x)(A(x)D)D(x)A(x)}\end{array}\begin{array}{cc}CC& ew\\ \stackrel{\mathit{}}{(x)(A(x)D)CC}\end{array}& \\ \stackrel{\mathit{}}{(x)(A(x)D)(x)A(x)C,DC}\end{array}& \\ \stackrel{\mathit{}}{(x)(A(x)D)(x)A(x)CDC}\end{array}& \\ \stackrel{\mathit{}}{\stackrel{\mathit{}}{(x)A(x)C(x)(A(x)D)(DC)}}\end{array}$$
## 4 Cut Elimination and Midhypersequent Theorem
###### Theorem 4.1 (Cut Elimination)
Any derivation of a hypersequent $`G`$ in $`\mathrm{𝐇𝐈𝐅}`$ can be transformed into a derivation of $`G`$ in $`\mathrm{𝐇𝐈𝐅}^{}`$.
This theorem is proved in the usual way by induction on the number of applications of the cut rule, using the following lemma.
###### Lemma 4.2
Suppose the hypersequents
$$H_1=G\mathrm{\Gamma }A\text{and}H_2=G\mathrm{\Pi }\mathrm{\Lambda }$$
are cut-free derivable. Then
$$H=G\mathrm{\Gamma },\mathrm{\Pi }^{}\mathrm{\Lambda }$$
where $`\mathrm{\Pi }^{}`$ is obtained from $`\mathrm{\Pi }`$ by removing all occurrences of $`A`$, is cut-free provable, and the number of applications of (ec) in the resulting proof is not more than the sum of applications of (ec) in $`\gamma `$ and $`\delta `$.
###### Proof
Let $`\gamma `$ and $`\delta `$ be the cut-free proofs of $`G`$ and $`H`$, respectively. We may assume, renaming variables if necessary, that the eigenvariables in $`\gamma `$ and $`\delta `$ are distinct. The proof follows Gentzen’s original Hauptsatz. Define the following measures on the pair $`\gamma ,\delta `$: the *rank* $`r=\mathrm{len}(\gamma )+\mathrm{len}(\delta )`$, the *degree* $`d=\mathrm{deg}(A)`$, and the *order* $`o`$ is the number of applications of the (ec) rule in $`\gamma `$, $`\delta `$. We proceed by induction on the lexicographical order of $`d,o,r`$.
If either $`H_1`$ or $`H_2`$ is an axiom, then $`H`$ can be derived from $`H_1`$ or $`H_2`$, respectively, using only weakenings. (This includes the case where $`r=2`$).
Otherwise, we distinguish cases according to the last inferences in $`\gamma `$ and $`\delta `$. The induction hypothesis is that the claim of the lemma is true whenever the degree is $`<d`$ or is $`=d`$ and either the order $`<o`$, or the order $`=o`$ and the rank $`<r`$.
(1) $`\gamma `$ or $`\delta `$ ends in an inference which acts on a sequent in $`G`$. We may invoke the induction hypothesis on the premises of $`H_1`$ or $`H_2`$, and $`H_2`$ or $`G_2`$, respectively.
(2) $`\gamma `$ or $`\delta `$ ends in ($`ec`$). For instance, $`\gamma `$ ends in
$$\begin{array}{cc}\begin{array}{cc}& \gamma ^{}\\ \mathrm{}\\ G\mathrm{\Gamma }A\mathrm{\Gamma }A\end{array}& ec\\ \stackrel{\mathit{}}{G\mathrm{\Gamma }A}\end{array}$$
Apply the induction hypothesis to $`\gamma ^{}`$ and $`\delta `$. The resulting proof $`\gamma ^{\prime \prime }`$ of
$$G\mathrm{\Gamma }A\mathrm{\Gamma },\mathrm{\Pi }^{}\mathrm{\Lambda }$$
has one less ($`ec`$) than $`\gamma `$ (although it may be much longer), and so the induction hypothesis applies again to $`\gamma ^{\prime \prime }`$ and $`\delta `$.
(3) $`\gamma `$ or $`\delta `$ end in another structural inference, (tt), or (cm): These cases are unproblematic applications of the induction hypothesis to the premises, followed by applications of structural inferences.
For example, assume $`\gamma `$ ends in (cm), i.e.,
$$\begin{array}{cc}\begin{array}{cc}& \gamma _1\\ \mathrm{}\\ G\mathrm{\Theta }_1,\mathrm{\Theta }_1^{}\mathrm{\Xi }_1\end{array}\begin{array}{cc}& \gamma _2\\ \mathrm{}\\ G\mathrm{\Theta }_2,\mathrm{\Theta }_2^{}A\end{array}& cm\\ \stackrel{\mathit{}}{G\mathrm{\Theta }_1,\mathrm{\Theta }_2^{}\mathrm{\Xi }_1\mathrm{\Theta }_1^{},\mathrm{\Theta }_2A}\end{array}$$
where $`\mathrm{\Gamma }=\mathrm{\Theta }_1^{},\mathrm{\Theta }_2`$. Apply the deduction hypothesis to the right premise and $`H_2`$ to obtain a cut-free proof of
$$G\mathrm{\Theta }_2,\mathrm{\Theta }_2^{},\mathrm{\Pi }^{}\mathrm{\Lambda }$$
Using applications of (ew) and (cm), we obtain the desired result.
The case of ($`tt`$) may be of special interest. Suppose $`\gamma `$ ends in(tt), with
$$\begin{array}{cc}G\mathrm{\Phi }pp,\mathrm{\Psi }A& tt\\ \stackrel{\mathit{}}{G\mathrm{\Phi },\mathrm{\Psi }A}\end{array}$$
Apply the induction hypothesis to the premises of $`H_1`$ and $`H_2`$, and apply (tt) to obtain the desired proof:
$$\begin{array}{cc}G\mathrm{\Phi }pp,\mathrm{\Psi },\mathrm{\Pi }^{}\mathrm{\Lambda }& tt\\ \stackrel{\mathit{}}{G\mathrm{\Phi },\mathrm{\Psi },\mathrm{\Pi }^{}\mathrm{\Lambda }}\end{array}$$
The case of $`\delta `$ ending in ($`tt`$) is handled similarly.
(4) $`\gamma `$ ends in a logical inference not involving the cut formula, or $`\delta `$ ends in a logical inference not involving the cut formula. These cases are easily handled by appeal to the induction hypothesis and application of appropriate logical and structural inferences. We outline the case where $`\gamma `$ ends in $`()`$:
$$\begin{array}{cc}\begin{array}{cc}& \gamma _1\\ \mathrm{}\\ GC,\mathrm{\Gamma }A\end{array}\begin{array}{cc}& \gamma _2\\ \mathrm{}\\ G\mathrm{\Gamma }B\end{array}& \\ \stackrel{\mathit{}}{GBC,\mathrm{\Gamma }A}\end{array}$$
We apply the induction hypothesis to the left premise and $`H_2`$, and apply ($``$):
$$\begin{array}{c}GC,\mathrm{\Gamma },\mathrm{\Pi }^{}\mathrm{\Lambda }G\mathrm{\Gamma }B\\ \stackrel{\mathit{}}{GBC,\mathrm{\Gamma },\mathrm{\Pi }^{}\mathrm{\Lambda }}\end{array}$$
(5) Both $`\gamma `$ and $`\delta `$ end in logical inferences acting on a cut formula. For instance, if $`A=BC`$ we have
$$\begin{array}{cc}\begin{array}{cc}& \gamma _1\\ \mathrm{}\\ GB,\mathrm{\Gamma }C\end{array}& \\ \stackrel{\mathit{}}{G\mathrm{\Gamma }BC}\end{array}\begin{array}{cc}\begin{array}{cc}& \delta _1\\ \mathrm{}\\ G\mathrm{\Pi }_1B\end{array}\begin{array}{cc}& \delta _2\\ \mathrm{}\\ GC,\mathrm{\Pi }_2\mathrm{\Lambda }\end{array}& \\ \stackrel{\mathit{}}{GBC,\mathrm{\Pi }_1,\mathrm{\Pi }_2\mathrm{\Lambda }}\end{array}$$
First we find proofs $`\delta _1^{}`$ and $`\delta _2^{}`$ of
$$G\mathrm{\Gamma },\mathrm{\Pi }_1^{}B\text{a}ndGC,\mathrm{\Gamma },\mathrm{\Pi }_2^{}\mathrm{\Lambda }$$
either by applying the induction hypothesis to $`\gamma `$ and $`\delta _1`$ or $`\delta _2`$ if $`\mathrm{\Pi }_1`$ or $`\mathrm{\Pi }_2`$, respectively, contain $`BC`$, or otherwise by adding (ic)-inferences to $`\delta _1`$ and $`\delta _2`$. Now apply the induction hypothesis based on the reduced degree of the cut formulas twice: first to $`\delta _1^{}`$ and $`\gamma _1`$ to obtain $`G\mathrm{\Gamma },\mathrm{\Gamma },\mathrm{\Pi }_1^{}C`$, and then to the resulting proof and $`\delta _2^{}`$ to obtain
$$G\mathrm{\Gamma },\mathrm{\Gamma },\mathrm{\Gamma },\mathrm{\Pi }_1^{},\mathrm{\Pi }_2^{}\mathrm{\Lambda }.$$
The desired result follows by several applications of (ic).
The other cases are similar and are left to the reader. ∎
Cut elimination is a basic prerequisite for proof theoretic and computational treatments of a logic. As an immediate consequence of cut elimination we have the subformula property: every IF-valid formula has a proof which only contains subformulas of the endformula (plus possibly propositional variables used in (tt)). Another important corollary is the midhypersequent theorem. It corresponds to Herbrand’s Theorem for classical logic and is thus the basis for any resolution-style automated proof method.
###### Theorem 4.3
Any hypersequent $`H`$ with only prefix formulas has a proof where no propositional inference follows a quantifier inference. Such a proof contains one or more hypersequents $`M`$, called midhypersequents, so that $`M`$ contains no quantifiers, all the inferences above $`M`$ are propositional or structural, and all the inferences below $`M`$ are either quantifier inferences of structural inferences.
###### Proof
This is proved exactly as for the classical and intuitionistic case (see Takeuti ). First, observe that all axioms are cut-free derivable from atomic axioms. The cut-elimination theorem thus provides us with a cut-free proof $`\pi `$ of $`H`$ from atomic axioms. Next, observe that the ($``$) rule can be simulated without using cuts by the rule
$$\begin{array}{cc}GA,\mathrm{\Gamma }\mathrm{\Delta }_1GB,\mathrm{\Gamma }\mathrm{\Delta }_2& ^{}\\ \stackrel{\mathit{}}{GAB,\mathrm{\Gamma }\mathrm{\Delta }_1AB,\mathrm{\Gamma }\mathrm{\Delta }_2}\end{array}$$
The rule can be derived as follows (we omit side sequents):
$$\begin{array}{cc}\begin{array}{cc}\begin{array}{cc}A,\mathrm{\Gamma }\mathrm{\Delta }_1B,\mathrm{\Gamma }\mathrm{\Delta }_2& \text{cm}\\ \stackrel{\mathit{}}{B,\mathrm{\Gamma }\mathrm{\Delta }_1A,\mathrm{\Gamma }\mathrm{\Delta }_2}\end{array}A,\mathrm{\Gamma }\mathrm{\Delta }_1& \\ \stackrel{\mathit{}}{AB,\mathrm{\Gamma }\mathrm{\Delta }_1A,\mathrm{\Gamma }\mathrm{\Delta }_2}\end{array}B,\mathrm{\Gamma }\mathrm{\Delta }_2& \\ \stackrel{\mathit{}}{AB,\mathrm{\Gamma }\mathrm{\Delta }_1AB,\mathrm{\Gamma }\mathrm{\Delta }_2}\end{array}$$
Of course, ($`^{}`$) together with (ec) simulates ($``$). We replace all applications of ($``$) by applications of ($`^{}`$) in our cut-free proof.
Define the order of a quantifier inference in $`\pi `$ to be the number of propositional inferences under it, and the order of $`\pi `$ as the sum of the orders of its quantifier inferences. The proof is by induction on the order of $`\pi `$. The only interesting case is of $`(^{})`$ occurring below a quantifier inference, since this case does not work for intuitionistic logic.
Suppose $`\pi `$ contains a ($``$) inference above a ($`^{}`$) inference, and so that all the inferences in between are structural. We have the following situation:
$$\begin{array}{cc}\begin{array}{c}\\ \mathrm{}\\ GA,\mathrm{\Gamma }\mathrm{\Delta }\end{array}\begin{array}{cc}\begin{array}{cc}\begin{array}{cc}& \delta ^{}\\ \mathrm{}\\ G^{}\mathrm{\Gamma }^{}A(a)\end{array}& \\ \stackrel{\mathit{}}{G^{}\mathrm{\Gamma }^{}(x)A(x)}\end{array}& \delta \\ \mathrm{}\\ GB,\mathrm{\Gamma }(x)A(x)\end{array}& ^{}\\ \stackrel{\mathit{}}{GAB,\mathrm{\Gamma }\mathrm{\Delta }AB,\mathrm{\Gamma }(x)A(x)}\end{array}$$
where $`\delta `$ contains only structural inferences. We reduce the order of $`\pi `$ by replacing this part of $`\pi `$ by:
$`\begin{array}{cc}\begin{array}{cc}\begin{array}{c}\\ \mathrm{}\\ GA,\mathrm{\Gamma }\mathrm{\Delta }\end{array}\begin{array}{cc}\begin{array}{cc}& \delta ^{}\\ \mathrm{}\\ G^{}\mathrm{\Gamma }^{}A(a)\end{array}& \delta \\ \mathrm{}\\ GB,\mathrm{\Gamma }A(a)\end{array}& ^{}\\ \stackrel{\mathit{}}{GAB,\mathrm{\Gamma }\mathrm{\Delta }AB,\mathrm{\Gamma }A(a)}\end{array}& \\ \stackrel{\mathit{}}{GAB,\mathrm{\Gamma }\mathrm{\Delta }AB,\mathrm{\Gamma }(x)A(x)}\end{array}`$
## 5 Elimination of the Takeuti-Titani Rule
The Takeuti-Titani rule is the least understood feature of the original Takeuti-Titani axiomatization of IF. We show below that the rule can be eliminated from proofs in $`\mathrm{𝐇𝐈𝐅}`$. This had been posed as a problem by Takano . The proof is by induction on the number of applications of (tt) and the length of the proof. The exact complexity of the elimination procedure is still to be investigated. The (tt) rule can have significant effects on proof structure. For instance, one of the calculi in Avron uses the split rule
$$\begin{array}{cc}G\mathrm{\Gamma },\mathrm{\Gamma }^{}\mathrm{\Delta }& \text{split}\\ \stackrel{\mathit{}}{G\mathrm{\Gamma }\mathrm{\Delta }\mathrm{\Gamma }^{}\mathrm{\Delta }}\end{array}$$
If this rule is added to $`\mathrm{𝐇𝐈𝐅}`$, it is possible to transform proofs so that each application of the communication rule has a premise which is a propositional axiom. This is not possible without (tt). The transformation works by replacing each occurrence of the communication rule by
$$\begin{array}{c}\begin{array}{cc}\begin{array}{cc}\begin{array}{cc}\begin{array}{cc}qq\begin{array}{cc}pp\begin{array}{cc}G_1\mathrm{\Gamma }_1,\mathrm{\Gamma }_1^{}A_1& \text{split}\\ \stackrel{\mathit{}}{G_1\mathrm{\Gamma }_1A_1\mathrm{\Gamma }_1^{}A_1}\end{array}& cm\\ \stackrel{\mathit{}}{G_1\mathrm{\Gamma }_1A_1\mathrm{\Gamma }_1^{}ppA_1}\end{array}& cm\\ \stackrel{\mathit{}}{G_1\mathrm{\Gamma }_1A_1\mathrm{\Gamma }_1^{}qpA_1qp}\end{array}\begin{array}{cc}\begin{array}{cc}G_2\mathrm{\Gamma }_2,\mathrm{\Gamma }_2^{}A_2& split\\ \stackrel{\mathit{}}{G_2\mathrm{\Gamma }_2A_2\mathrm{\Gamma }_2^{}A_2}\end{array}qq& cm\\ \stackrel{\mathit{}}{G_2\mathrm{\Gamma }_2qqA_2\mathrm{\Gamma }_2^{}A_2}\end{array}& cut\\ \stackrel{\mathit{}}{G_1G_2\mathrm{\Gamma }_1A_1\mathrm{\Gamma }_2qpA_1\mathrm{\Gamma }_2pqA_2\mathrm{\Gamma }_2^{}A_2}\end{array}& tt\\ \stackrel{\mathit{}}{G_1G_2\mathrm{\Gamma }_1A_1\mathrm{\Gamma }_2A_2pA_1\mathrm{\Gamma }_2p\mathrm{\Gamma }_2^{}A_2}\end{array}& tt\\ \stackrel{\mathit{}}{G_1G_2\mathrm{\Gamma }_1A_1\mathrm{\Gamma }_2A_2\mathrm{\Gamma }_2A_1\mathrm{\Gamma }_2^{}A_2}\end{array}\\ \stackrel{\mathit{}}{\stackrel{\mathit{}}{G_1G_2\mathrm{\Gamma }_1,\mathrm{\Gamma }_2^{}A_1\mathrm{\Gamma }_1^{},\mathrm{\Gamma }_2A_2}}\end{array}$$
###### Proposition 5.1
Let $`\delta `$ be a $`\mathrm{𝐇𝐈𝐅}^{}`$-derivation of hypersequent $`H`$ with length $`k`$, where $`H`$ is of the form
$$G\mathrm{\Gamma }_1,\mathrm{\Pi }_1\mathrm{\Delta }_1,\mathrm{\Pi }_1^{}\mathrm{}\mathrm{\Gamma }_n,\mathrm{\Pi }_n\mathrm{\Delta }_n,\mathrm{\Pi }_n^{}$$
and $`\mathrm{\Pi }_i\{p\}`$, $`\mathrm{\Pi }_i^{}=\mathrm{}`$, and $`p`$ does not occur in $`G`$, $`\mathrm{\Gamma }_i`$ or $`\mathrm{\Delta }_i`$ ($`\mathrm{\Pi }_i^{}=\{p\}`$, $`\mathrm{\Pi }_i=\mathrm{}`$, and $`p`$ does not occur in $`G`$, $`\mathrm{\Gamma }_i`$ or $`\mathrm{\Delta }_i`$).
Then the hypersequent $`G\mathrm{\Gamma }_{i_1}\mathrm{\Delta }_{i_1}\mathrm{}\mathrm{\Gamma }_{i_m}\mathrm{\Delta }_{i_m}`$ is derivable in length $`k`$.
###### Proof
Easy induction on $`k`$. Every occurrence of $`p`$ must arise from a weakening, simply delete all these weakenings.
###### Theorem 5.2
Applications of (tt) can be eliminated from $`\mathrm{𝐇𝐈𝐅}`$-derivations.
This follows from the following lemma by induction on the number of applications of (tt) in a given $`\mathrm{𝐇𝐈𝐅}^{}`$-derivation.
###### Lemma 5.3
If $`\delta `$ is an $`\mathrm{𝐇𝐈𝐅}^{}`$-derivation of
$$H=G\mathrm{\Phi }_1\mathrm{\Pi }_1\mathrm{}\mathrm{\Phi }_n\mathrm{\Pi }_n\mathrm{\Pi }_1^{},\mathrm{\Psi }_1\mathrm{\Sigma }_1\mathrm{}\mathrm{\Pi }_m^{},\mathrm{\Psi }_m\mathrm{\Sigma }_m,$$
where $`p`$ does not occur in $`G`$, $`\mathrm{\Phi }_i`$, $`\mathrm{\Psi }_i`$ or $`\mathrm{\Sigma }_i`$, and $`\mathrm{\Pi }_i\mathrm{\Pi }_i^{}\{p\}`$, then there is a $`\mathrm{𝐇𝐈𝐅}^{}`$-derivation of
$$H^{}=G\mathrm{\Phi }_1,\mathrm{},\mathrm{\Phi }_n,\mathrm{\Psi }_1\mathrm{\Sigma }_1\mathrm{}\mathrm{\Phi }_1,\mathrm{},\mathrm{\Phi }_n,\mathrm{\Psi }_m\mathrm{\Sigma }_m.$$
###### Proof
By induction on the length of $`\delta `$. We distinguish cases according to the last inference $`I`$ in $`\delta `$. For simplicity, we will write $`p`$ in what follows below instead of $`\mathrm{\Pi }_i`$ or $`\mathrm{\Pi }_i^{}`$ with the understanding that it denotes an arbitrary multiset of $`p`$’s.
(1) The conclusion of of $`I`$ is so that $`p`$ only occurs on the right side of sequents, or only on the left side. Then Prop. 5.1 applies, and the desired hypersequent can be derived without (tt).
(2) $`I`$ applies to sequents in $`G`$. Then the induction hypothesis can be applied to the premise(s) of $`I`$ and appropriate inferences added below.
(3) $`I`$ is structural inference other than (cut) and (cm), or a logical inference with only one premise, or a logical inference which applies to a $`\mathrm{\Sigma }_i`$. These cases are likewise handled in an obvious manner and are unproblematic. One instructive example might be the case of ($``$). Here the premises would be of the form, say,
$$\begin{array}{c}G\mathrm{\Phi }_1p\mathrm{\Phi }_2p\mathrm{}\mathrm{\Phi }_npp,\mathrm{\Psi }_1\mathrm{\Sigma }_1\mathrm{}p,\mathrm{\Psi }_m\mathrm{\Sigma }_mp,\mathrm{\Gamma }_1A\hfill \\ G\mathrm{\Phi }_1p\mathrm{\Phi }_2p\mathrm{}\mathrm{\Phi }_npp,\mathrm{\Psi }_1\mathrm{\Sigma }_1\mathrm{}p,\mathrm{\Psi }_m\mathrm{\Sigma }_mB,\mathrm{\Gamma }_2p\hfill \end{array}$$
Let $`\mathrm{\Phi }=\mathrm{\Phi }_1,\mathrm{},\mathrm{\Phi }_n`$. The induction hypothesis provides us with
$$\begin{array}{c}G\mathrm{\Phi },\mathrm{\Psi }_1\mathrm{\Sigma }_1\mathrm{}\mathrm{\Phi },\mathrm{\Psi }_m\mathrm{\Sigma }_m\mathrm{\Phi },\mathrm{\Gamma }_1A\hfill \\ GB,\mathrm{\Gamma }_2,\mathrm{\Phi },\mathrm{\Psi }_1\mathrm{\Sigma }_1\mathrm{}B,\mathrm{\Gamma }_2,\mathrm{\Phi },\mathrm{\Psi }_m\mathrm{\Sigma }_m\hfill \end{array}$$
We obtain the desired hypersequent by applying ($``$) successively $`m`$ times, together with some contractions.
(4) $`I`$ is a cut. There are several cases to consider, most of which are routine. The only tricky case is when the cut formula is $`p`$ and $`p`$ occurs both on the left and the right side of sequents in both premises of the cut. For simplicity, let us consider the cut rule in its multiplicate formulation
$$\begin{array}{c}G\mathrm{\Phi }_1p\mathrm{}\mathrm{\Phi }_npp,\mathrm{\Psi }_1\mathrm{\Sigma }_1\mathrm{}p,\mathrm{\Psi }_m\mathrm{\Sigma }_m\mathrm{\Gamma }p\hfill \\ G\mathrm{\Phi }_1p\mathrm{}\mathrm{\Phi }_npp,\mathrm{\Psi }_1\mathrm{\Sigma }_1\mathrm{}p,\mathrm{\Psi }_m\mathrm{\Sigma }_mp,\mathrm{\Pi }\mathrm{\Lambda }\hfill \end{array}$$
We want to find a derivation of
$$G\mathrm{\Phi },\mathrm{\Psi }_1\mathrm{\Sigma }_1\mathrm{\Phi },\mathrm{\Psi }_m\mathrm{\Sigma }_m\mathrm{\Gamma },\mathrm{\Pi }\mathrm{\Lambda }$$
where $`\mathrm{\Phi }=\mathrm{\Phi }_1,\mathrm{},\mathrm{\Phi }_n`$. The induction hypothesis applied to the premises of the cut gives us
$$\begin{array}{c}G\mathrm{\Gamma },\mathrm{\Phi },\mathrm{\Psi }_1\mathrm{\Sigma }_1\mathrm{}\mathrm{\Gamma },\mathrm{\Phi },\mathrm{\Psi }_m\mathrm{\Sigma }_m\hfill \\ G\mathrm{\Phi },\mathrm{\Psi }_1\mathrm{\Sigma }_1\mathrm{}\mathrm{\Phi },\mathrm{\Psi }_m\mathrm{\Sigma }_m\mathrm{\Phi },\mathrm{\Pi }\mathrm{\Lambda }\hfill \end{array}$$
We obtain the desired hypersequent by $`m`$ successive applications of (cm).
(5) $`I`$ is ($``$), or ($``$) applying to $`\mathrm{\Phi }_i`$ or $`\mathrm{\Psi }_i`$. Consider the case of ($``$), the others are treated similarly. The premises of $`I`$ are, for example,
$$\begin{array}{c}GA,\mathrm{\Phi }_1p\mathrm{\Phi }_2p\mathrm{}\mathrm{\Phi }_npp,\mathrm{\Psi }_1\mathrm{\Sigma }_1\mathrm{}p,\mathrm{\Psi }_m\mathrm{\Sigma }_m\hfill \\ GB,\mathrm{\Phi }_1p\mathrm{\Phi }_2p\mathrm{}\mathrm{\Phi }_npp,\mathrm{\Psi }_1\mathrm{\Sigma }_1\mathrm{}p,\mathrm{\Psi }_m\mathrm{\Sigma }_m\hfill \end{array}$$
By induction hypothesis, we obtain
$$\begin{array}{c}GA,\mathrm{\Phi }_1,\mathrm{},\mathrm{\Phi }_n,\mathrm{\Psi }_1\mathrm{\Sigma }_1\mathrm{}A,\mathrm{\Phi }_1,\mathrm{},\mathrm{\Phi }_n,\mathrm{\Psi }_m\mathrm{\Sigma }_m\hfill \\ GB,\mathrm{\Phi }_1,\mathrm{},\mathrm{\Phi }_n,\mathrm{\Psi }_1\mathrm{\Sigma }_1\mathrm{}B,\mathrm{\Phi }_1,\mathrm{},\mathrm{\Phi }_n,\mathrm{\Psi }_m\mathrm{\Sigma }_m\hfill \end{array}$$
It is not straightforwardly possible to derive the desired hypersequent from these. If $`\mathrm{\Psi }_i=\{P_{i1},\mathrm{},P_{ik_i}\}`$, let $`Q_i=P_{i1}\mathrm{}P_{ik_i}\mathrm{\Sigma }_i`$. Then we do easily obtain, however, the following by repeated application of ($``$), ($``$) and ($`ec`$):
$$\begin{array}{c}GA,\mathrm{\Phi }_1,\mathrm{},\mathrm{\Phi }_nQ_1\mathrm{}Q_m\hfill \\ GB,\mathrm{\Phi }_1,\mathrm{},\mathrm{\Phi }_nQ_1\mathrm{}Q_m\hfill \end{array}$$
Now a single application of ($``$), plus (ec) gives us
$$K=G\underset{\mathrm{\Gamma }}{\underset{}{AB,\mathrm{\Phi }_1,\mathrm{},\mathrm{\Phi }_n}}Q_1\mathrm{}Q_m$$
Then we derive, using $`m1`$ cuts:
$$\begin{array}{c}\begin{array}{c}\begin{array}{c}K\begin{array}{cc}& \delta _1\\ \mathrm{}\\ Q_1QQ_1Q_1QQ\end{array}\\ \stackrel{\mathit{}}{\stackrel{\mathit{}}{\mathrm{\Gamma }Q_1\mathrm{\Gamma }\underset{Q}{\underset{}{Q_2\mathrm{}Q_m}}}}\end{array}\\ \mathrm{}\\ \mathrm{\Gamma }Q_1\mathrm{}\mathrm{\Gamma }Q_{m1}Q_m\end{array}\begin{array}{cc}& \delta _{m1}\\ \mathrm{}\\ Q_{m1}Q_mQ_{m1}Q_{m1}Q_mQ_m\end{array}\\ \stackrel{\mathit{}}{\stackrel{\mathit{}}{\mathrm{\Gamma }Q_1\mathrm{}\mathrm{\Gamma }Q_m}}\end{array}$$
where $`\delta _i`$ is the derivation
$$\begin{array}{cc}Q_iQ_i\begin{array}{cc}\begin{array}{cc}QQQ_iQ_i& cm\\ \stackrel{\mathit{}}{QQ_iQ_iQ}\end{array}QQ& \\ \stackrel{\mathit{}}{QQ_iQ_iQQ}\end{array}& \\ \stackrel{\mathit{}}{Q_i\underset{Q}{\underset{}{Q_{i+1}\mathrm{}Q_m}}Q_iQ_i\mathrm{}Q_m\underset{Q}{\underset{}{Q_{i+1}\mathrm{}Q_m}}}\end{array}$$
The desired hypersequent is obtained by $`m`$ cuts with
$$Q_i,P_{i1},\mathrm{},P_{ik_i}\mathrm{\Sigma }_i$$
(6) $`I`$ is a communication rule. This is the most involved case, as several subcases have to be distinguished according to which of the two communicated sequents contains $`p`$. Neither of these cases are problematic. We present two examples:
(a) One of the communicated sequents contains $`p`$ on the right. Then the premises of $`I`$ are
$$\begin{array}{c}G\mathrm{\Phi }_1p\mathrm{}\mathrm{\Phi }_npp,\mathrm{\Psi }_1\mathrm{\Sigma }_1\mathrm{}p,\mathrm{\Psi }_m\mathrm{\Sigma }_m\mathrm{\Theta }_1,\mathrm{\Theta }_1^{}p\hfill \\ G\mathrm{\Phi }_1p\mathrm{}\mathrm{\Phi }_npp,\mathrm{\Psi }_1\mathrm{\Sigma }_1\mathrm{}p,\mathrm{\Psi }_m\mathrm{\Sigma }_m\mathrm{\Theta }_2,\mathrm{\Theta }_2^{}\mathrm{\Xi }_2\hfill \end{array}$$
where. The induction hypothesis applies to these two hypersequents. If we write $`\mathrm{\Phi }=\mathrm{\Phi }_1,\mathrm{},\mathrm{\Phi }_n`$, we have
$$\begin{array}{c}G\overline{)\mathrm{\Theta }_1,\mathrm{\Theta }_1^{},\mathrm{\Phi },\mathrm{\Psi }_1\mathrm{\Sigma }_1}\mathrm{}\mathrm{\Theta }_1,\mathrm{\Theta }_1^{},\mathrm{\Phi },\mathrm{\Psi }_m\mathrm{\Sigma }_m\hfill \\ G\overline{)\mathrm{\Theta }_2,\mathrm{\Theta }_2^{}\mathrm{\Xi }}\mathrm{\Phi },\mathrm{\Psi }_1\mathrm{\Sigma }_1\mathrm{}\mathrm{\Phi },\mathrm{\Psi }_m\mathrm{\Sigma }_m\hfill \end{array}$$
We obtain the desired result by applying $`m`$ instances of (cm), internal weakenings and external contractions as necessary, to obtain, in sequence
$$\begin{array}{c}G\mathrm{\Theta }_1,\mathrm{\Theta }_2^{},\mathrm{\Phi },\mathrm{\Psi }_1\mathrm{\Sigma }_1\mathrm{}\mathrm{\Theta }_1,\mathrm{\Theta }_1^{},\mathrm{\Phi },\mathrm{\Psi }_m\mathrm{\Sigma }_m\mathrm{\Theta }_1^{},\mathrm{\Theta }_2\mathrm{\Xi }\hfill \\ \mathrm{}\\ G\mathrm{\Theta }_1,\mathrm{\Theta }_2^{},\mathrm{\Phi },\mathrm{\Psi }_1\mathrm{\Sigma }_1\mathrm{}\mathrm{\Theta }_1,\mathrm{\Theta }_2^{},\mathrm{\Phi },\mathrm{\Psi }_m\mathrm{\Sigma }_m\mathrm{\Theta }_1^{},\mathrm{\Theta }_2\mathrm{\Xi }\hfill \end{array}$$
The sequents participating in the application of (cm) are marked by boxes. The original end hypersequent follows from the last one by internal weakenings.
(b) The communicated sequents both contain $`p`$, once on the right, once on the left. The premises of $`I`$ are
$$\begin{array}{c}G\mathrm{\Phi }_1p\mathrm{}\mathrm{\Phi }_npp,\mathrm{\Psi }_1\mathrm{\Sigma }_1\mathrm{}p,\mathrm{\Psi }_m\mathrm{\Sigma }_m\mathrm{\Theta }_1,\mathrm{\Theta }_1^{}p\hfill \\ G\mathrm{\Phi }_1p\mathrm{}\mathrm{\Phi }_npp,\mathrm{\Psi }_1\mathrm{\Sigma }_1\mathrm{}p,\mathrm{\Psi }_m\mathrm{\Sigma }_mp,\mathrm{\Theta }_2,\mathrm{\Theta }_2^{}\mathrm{\Xi }\hfill \end{array}$$
We have proofs of
$$\begin{array}{c}G\mathrm{\Theta }_1,\mathrm{\Theta }_1^{},\mathrm{\Phi },\mathrm{\Psi }_1\mathrm{\Sigma }_1\mathrm{}\mathrm{\Theta }_1,\mathrm{\Theta }_1^{},\mathrm{\Phi },\mathrm{\Psi }_m\mathrm{\Sigma }_m\hfill \\ G\mathrm{\Phi },\mathrm{\Psi }_1\mathrm{\Sigma }_1\mathrm{}\mathrm{\Phi },\mathrm{\Psi }_m\mathrm{\Sigma }_m\mathrm{\Theta }_2,\mathrm{\Theta }_2^{},\mathrm{\Phi }\mathrm{\Xi }\hfill \end{array}$$
Again, a sequence of $`m`$ applications of (cm), together with internal weakenings and external contractions produces the desired end sequent. ∎
Note that in case (5), several new cuts are introduced. As a consequence, the elimination procedure does not directly work for cut-free proofs. If a proof with neither cut nor communication is required, the elimination procedure has to be combined with the cut-elimination procedure of Thm. 4.1. The additional cuts can be avoided by replacing ($``$) and ($``$) by the following generalized rules:
$$\begin{array}{c}\begin{array}{cc}GA,\mathrm{\Gamma }_1\mathrm{\Delta }_1\mathrm{}A,\mathrm{\Gamma }_n\mathrm{\Delta }_nGB,\mathrm{\Gamma }_1\mathrm{\Delta }_1\mathrm{}B,\mathrm{\Gamma }_n\mathrm{\Delta }_n& ^{}\\ \stackrel{\mathit{}}{GAB,\mathrm{\Gamma }_1\mathrm{\Delta }_1\mathrm{}AB,\mathrm{\Gamma }_n\mathrm{\Delta }_n}\end{array}\\ \begin{array}{cc}GA(a),\mathrm{\Gamma }_1\mathrm{\Delta }_1\mathrm{}A(a),\mathrm{\Gamma }_n\mathrm{\Delta }_n& ^{}\\ \stackrel{\mathit{}}{G(x)A(x),\mathrm{\Gamma }_1\mathrm{\Delta }_1\mathrm{}(x)A(x),\mathrm{\Gamma }_n\mathrm{\Delta }_n}\end{array}\end{array}$$
These rules, however, cannot be simulated by the ordinary rules without using cut (the simulation with cut is given in case (5)). By changing case (5) accordingly, the elimination procedure will transform a cut-free $`\mathrm{𝐇𝐈𝐅}`$-derivation into a cut-free one without (tt), but with ($`^{}`$) and ($`^{}`$).
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# Annihilating Cold Dark Matter
\[
## Abstract
Structure formation with cold dark matter (CDM) predicts halos with a central density cusp, which are observationally disfavored. If CDM particles have an annihilation cross section $`\sigma v10^{29}(m/\mathrm{GeV})\mathrm{cm}^2`$, then annihilations will soften the cusps. We discuss plausible scenarios for avoiding the early Universe annihilation catastrophe that could result from such a large cross section. The predicted scaling of core density with halo mass depends upon the velocity dependence of $`\sigma v`$, and s-wave annihilation leads to a core density nearly independent of halo mass, which seems consistent with observations.
\] Introduction. The idea that the large-scale structure developed by gravitational instability from initially small-amplitude, adiabatic and nearly scale-invariant fluctuations is compatible with a number of observables across a wide range of length scales (e.g., from the cosmic microwave background anisotropy to the Lyman-$`\alpha `$ forest). Essential to this compatibility is the existence of cold dark matter: matter which is non-baryonic, has only very weak interactions with photons and baryons, and (prior to gravitational collapse) is cold.
The greatest challenge to this otherwise successful scenario comes from the apparent discrepancy between predicted dark-matter halo-density profiles and those inferred from observations. Simulations with non-interacting cold dark matter lead to halo density profiles that are singular at the center , whereas observations indicate uniform density cores. In this Letter we explore the possibility that the dark matter today has a large cross-section for annihilation which results in preferential destruction in high-density regions, softening halo cores .
Detecting and determining the properties of the dark matter is a major goal of observational cosmology. If annihilations are indeed altering the properties of dark matter halos, then we have a new means of studying the dark matter. The interactions of the CDM particles determine both the magnitude and velocity dependence of the annihilation cross section. For example, for s-wave annihilation, $`\sigma _A|v|`$ is independent of velocity and for p-wave annihilation, $`\sigma _A|v|`$ is proportional to $`v^2`$. These two different dependences result in different scaling relations between core density and halo velocity dispersion, which can be tested by current observations.
As we show below, current data for high velocity dispersion systems such as clusters of galaxies to low velocity dispersion systems such as galactic satellites are consistent with the same core density of about $`1\mathrm{GeV}/\mathrm{cm}^3(=0.026M_{}/\mathrm{pc}^3`$) . This scale invariance can be explained by s-wave annihilation with a cross-section $`\sigma v10^{29}(m/\mathrm{GeV})\mathrm{cm}^2`$, although future improvements in both the data and predictions will be necessary before such a statement can be made with confidence. As we shall discuss, the cosmological and astrophysical constraints on annihilating CDM point to a candidate beyond those currently favored (e.g., axion, neutralino).
Halos of Annihilating Dark Matter. Numerical simulations of structure formation in the CDM scenario show that the dark matter halos which form with a wide range of masses are all well-fit with the so-called NFW form for the density profile. This form has $`\rho r^3`$ at large $`r`$ and a cuspy inner region with $`\rho r^\alpha `$ with $`\alpha =1`$. More recent higher resolution simulations predict cusps that are even stronger, with $`\alpha 1.5`$. Nevertheless, in most of what follows, we use the NFW theory for simplicity.
To be precise, the NFW profile is
$$\rho (r=xr_s)=\rho _sx^1(1+x)^2,$$
(1)
where the value of $`\rho _s`$ is determined by the mean density of the Universe at the time the halo collapsed. In CDM theory, small objects collapse first, followed later by larger ones. Thus, there is an inverse relationship between $`\rho _s`$ and halo size. In Fig. 1 we show this scaling relation with halo size represented by velocity dispersion for the halo, estimated as $`\sigma _{\mathrm{vir}}=\sqrt{GM_{\mathrm{vir}}/2r_{\mathrm{vir}}}`$. (The virial radius, $`r_{\mathrm{vir}}`$, is defined such that the mean density inside the $`r_{\mathrm{vir}}`$ sphere is 200 times the present mean density of the Universe, and $`M_{\mathrm{vir}}`$ is the mass contained within $`r_{\mathrm{vir}}`$ .)
Annihilations will alter the halo profiles near the core where the density of the dark matter particles is the highest. The annihilation rate (per particle) $`\mathrm{\Gamma }=n\sigma |v|`$ depends on the velocity dispersion. We parameterize the velocity dependence as $`\mathrm{\Gamma }=(\rho /m)\sigma _Av^n`$ ($`n=0`$ for s-wave; $`n=2`$ for p-wave), where $`v`$ is the velocity dispersion and $`m`$ is the CDM particle mass.
Fig. 1 also shows in a qualitative way how annihilations affect the core structure of different objects. The annihilation lines drawn show whether or not annihilations are important at a NFW halo radius of $`0.1r_s`$ in different kinds of objects. (Note, since halo densities diverge, for any object annihilations become significant deep enough into the core). The annihilation lines are normalized to soften the cores of low surface-brightness (LSB) spiral galaxies. Because of how annihilations scale with velocity, for n = 0 clusters remain unaffected at $`r\begin{array}{c}>\hfill \\ \hfill \end{array}0.1r_s`$, while the cores of LSBs and smaller objects are dramatically softened. For n = 2, the opposite is true, which contradicts observations that indicate the NFW profile works well for clusters. We expect any $`n>1`$ to be inconsistent with observations. The case of $`n=0.5`$ is interesting since the annihilation line runs parallel to the structure line (for $`\sigma _{\mathrm{vir}}\begin{array}{c}<\hfill \\ \hfill \end{array}`$ 100 km/s), implying that all the systems will be smoothed off at the same value of $`r/r_s`$.
Model Building Constraints. For the annihilations to be effective in galaxy cores today, the annihilation rate must satisfy the (approximate) constraint:
$$\mathrm{\Gamma }\left(\rho /\rho _{\mathrm{LSB}}\right)\left(v/v_{\mathrm{LSB}}\right)^nH_0,$$
(2)
where the subscript LSB denotes the appropriate values for a typical LSB and $`H_0=100h\mathrm{km}\mathrm{s}^1`$ is the present expansion rate of the Universe. Outside collapsed objects today, the density of CDM is much lower and annihilations will be unimportant for $`n0`$. The early Universe is another matter as densities were much higher, $`\rho T^3`$, where $`T`$ is the cosmic background radiation temperature.
The figure of merit for the effectiveness of annihilations in the early Universe is measured by annihilation rate divided by the expansion rate: when $`\mathrm{\Gamma }/H>1`$ annihilations are effective (and vice versa). Assuming that the velocity dispersion of the CDM particles can be characterized by the background radiation temperature and normalizing the cross section to the desired value today, the temperature dependence of $`\mathrm{\Gamma }/H`$ is
$$\frac{\mathrm{\Gamma }}{H}10^9\left(\frac{T}{\mathrm{GeV}}\right)\left(\frac{T}{10^3m}\right)^n\sqrt{\frac{T}{T+T_{\mathrm{eq}}}},$$
(3)
where $`T_{\mathrm{eq}}1\mathrm{eV}`$ is the temperature at matter – radiation equality. There are three important things to note: (1) the large coefficient in front of this expression – annihilations in the early Universe are a significant consideration; (2) for $`n=1`$, the effectiveness of annihilations is epoch independent and disastrous; and (3) for $`n>1`$ annihilations were more important in the past.
Observational data suggest that if halos are made of annihilating CDM particles, their annihilation cross section is characterized by $`n\begin{array}{c}<\hfill \\ \hfill \end{array}1`$. Thus we will focus on $`n>1`$, where the danger of annihilations is in the past: $`\mathrm{\Gamma }/H>1`$ for
$$T>T_A10^{3(3+n)/(1+n)}\mathrm{GeV}(m/\mathrm{GeV})^{n/(1+n)},$$
(4)
or $`1\mathrm{eV}`$ for $`n=0`$. To ensure that early annihilations do not reduce CDM particles to negligible numbers, they must be protected against annihilation in the early Universe. We suggest two mechanisms; doubtless, there are other possibilities.
First, CDM particles could be produced late ($`T<T_A`$) by the decays of another massive particle. Note that this requires a long lifetime, $`\tau >t(T_A)10^5`$yrs, and the mass difference between the two particles should be small enough to ensure that the relativistic decay products do not make the Universe radiation dominated.
The second way of avoiding the early-Universe annihilation catastrophe is to make the mass of the annihilation product be dynamical. For example, a phase transition that takes place at $`T<T_A`$ could change annihilation from being kinematically impossible to possible if the mass of the annihilation product dropped below threshold after the phase transition (or if the mass of the CDM particle rose above threshold). A variation on this theme is coupling the annihilation produced particle to a scalar field, $`\varphi `$, with $`\varphi 0`$. As $`\varphi `$ decreases, either quickly to zero as a result of a symmetry-restoring phase transition, or slowly as $`\varphi `$ rolls to the minimum of its potential, the product particle’s mass may drop below threshold, opening up the new annihilation channel, at $`T<T_A`$.
Finally, the CDM annihilation products must not include photons because their $`\gamma `$-ray flux would far exceed observational limits. For example for 1 GeV CDM particles, the flux would be around $`10^5\mathrm{cm}^2\mathrm{sr}^1\mathrm{s}^1`$, some ten orders of magnitude above the observed diffuse $`\gamma `$-ray flux at 1 GeV.
Observational Constraints. We henceforth restrict ourselves to $`n=0`$. The contribution of annihilations to the evolution of the density profile is given by
$$d\left[\rho (r)/\rho _A\right]/dt=\left[\rho (r)/\rho _A\right]^2t_0^1,$$
(5)
where $`\rho _Am/(\sigma _At_0)`$ and $`t_0`$ is the age of the Universe today. Assuming the initial profile to be NFW, the resulting density profile is
$$\rho (r)=\rho _s\left[x(1+x)^2+\rho _s/\rho _{\mathrm{core}}\right]^1,$$
(6)
where $`xr/r_s`$, and a core of constant density, $`\rho _{\mathrm{core}}=\rho _A`$, is clearly evident. However, the mass loss due to annihilations results in adiabatic expansion of the core, such that the quantity $`M(r)r`$ is left invariant . This expansion results in a lower core density and one can estimate that the ratio $`\rho _{\mathrm{core}}/\rho _A`$ ranges from about 0.1 (dwarf galaxies) to about 0.3 (clusters). We have verified this by more detailed numerical work which allows us to determine $`\rho _{\mathrm{core}}/\rho _A`$ as a function of halo mass.
We now turn to the observable constraints on annihilating CDM. A robust prediction of the s-wave annihilation scenario is that the cores are more evident in smaller mass halos, as can be seen in Fig. 2. So we first turn to the galactic satellites in the Milky Way group , of which there are 11 known.
For a $`10^8M_{}`$ galactic satellite, the core radius produced by annihilations is about 1 kpc, which is about the same as the cut-off radius induced by tidal forces. Most of the galactic satellites have large velocity dispersions ($``$ 10 km/s) for their stellar content, which suggests that they are CDM dominated . If so, their internal velocity dispersions indicate that $`\rho _{\mathrm{core}}=𝒪(1\mathrm{GeV}/\mathrm{cm}^3)`$ .
We also looked at dwarf spiral galaxies and LSBs. One must use these with caution since van den Bosch et al. have recently claimed that most of the H $`\mathrm{I}`$ rotation curve data do not have sufficient spatial resolution to put meaningful constraints on the halo cusps. They do identify three nearby galaxies which have sufficient spatial resolution – NGC 247, DDO 154 and NGC 3109 . van den Bosch et al. find that $`0.55<\alpha <1.26`$ for the LSB (NGC 247), and $`\alpha <0.5`$ for the two dwarfs, at the 99.73% confidence level which at face value, argues for soft cores in low-mass systems. The annihilation scenario naturally explains this since the cores are more evident in low-mass systems (see Fig. 2). However, it should be noted that the error bars on the rotation velocity data are probably not a complete description of the total uncertainty and that a critical reevaluation might lead to a less stringent bound on $`\alpha `$, thus alleviating the discrepancy between the observed dwarf rotation curves and CDM predictions.
To estimate the cross-section required to achieve consistency with observations, we fit to the two dwarf galaxies identified above with the halo profile in Eq. 6, a thin stellar disk and the observed gas. We have included the effect of finite resolution. We find that $`\rho _A0.2M_{}/\mathrm{pc}^3`$ results in a good fit to both (see Fig. 3). In both cases, the outer parts of the halo (determined by $`\rho _s`$ and $`r_s`$) are consistent with NFW theory.
To see if structures on the largest scales are consistent with the annihilation scenario, we turn to strong gravitational lensing of background galaxies by clusters. Tyson et al. model the mass distribution in the cluster CL 0024, which produces multiple distorted images of a background galaxy, and find evidence for a compact soft core (of about $`35h^1\mathrm{kpc}`$) in the projected density. A similar value for the core radius was inferred earlier by Smail et al. for CL 0024 and other clusters. X-ray studies (Bohringer et al.) of CL 0024 are also in agreement with the above results. However, we urge caution in interpreting these results since the evidence for soft cores in clusters is largely based on just one cluster (CL 0024).
A value of $`\rho _A=0.2M_{}/\mathrm{pc}^3`$ would produce a core density of about $`0.06M_{}/\mathrm{pc}^3`$ in a cluster-sized object. We find that this core density is consistent with the surface density reconstruction of CL 0024 by Tyson et al.. The implied CDM mass within the arc radius (of $`107h^1\mathrm{kpc}`$) is in agreement with the quoted value of about $`1.66\times 10^{14}h^1M_{}`$ for the total mass within the arc radius .
Discussion. The s-wave annihilation scenario with a cross-section of $`\sigma |v|=10^{29}(m/\mathrm{GeV})\mathrm{cm}^2`$ produces a core density of about $`1\mathrm{GeV}/\mathrm{cm}^3`$ over widely different scales. Intriguingly, this seems to be consistent with observations.
Apart from cuspy cores , simulations of non-interacting CDM also predict a much larger number of sub-halos for a galactic size halo than the observed number of galactic satellites . Certainly, the s-wave annihilation scenario has a dramatic effect on the smallest halos, and this could contribute to their destruction. However, further study is required to test this hypothesis.
Another particle-physics solution in which CDM particles have a large cross section for self interaction ($`\sigma 10^{24}\mathrm{cm}^2`$ has been discussed . This possibility is being tested by numerical simulations . However, there are indications that self interactions lead to halos that are inconsistent with observations . The jury is still out.
The requirements on a model for annihilating CDM are stringent, but by no means impossible . They point to a particle beyond those currently being considered, and therefore, to new physics. While it is possible that the solution to the CDM cusp problem will involve the interpretation of the observations or less exotic astrophysics, it is appealing to think that the properties of halo cores may teach us about the fundamental properties of the CDM particle.
###### Acknowledgements.
We gratefully acknowledge support by the DOE, NASA and NSF at Chicago, and by the DOE and NASA grant NAG 5-7092 at Fermilab. We thank M. Valluri for useful discussions and pointers. We also thank S. Carroll, D. Eisenstein, S. Hannestad, J. Mohr, D. Spergel and L. Widrow.
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# An Exposition on Inflationary Cosmology
## I Introduction
The Standard Model of Cosmology has successfully predicted the nucleosynthesis of the light elements, the temperature and blackbody spectrum of the cosmic background radiation, and the observed redshift of light from galaxies which suggests an expanding universe. However, this model can not account for a number of initial value problems, such as the flatness and monopole problems. Inflationary cosmology resolves these concerns, while preserving the successes of the Big-Bang model. Inflation was originally introduced for this reason and its motivation relied on predictions from particle theory. In more recent times, inflation has been abstracted to a much more general theory. It continues to resolve the initial value problems, but also offers an explanation of the observed large-scale structure of the universe.
In this paper, the fundamentals of modern cosmology for an isotropic and homogeneous space-time, which is naturally motivated by observation, will be reviewed. The Friedmann equations are derived and the consequences for the dynamics of the universe are discussed. A brief introduction to the thermal properties of the universe is presented as motivation for a discussion of the horizon problem. Moreover, other issues suggesting a more general theory are presented and inflation is introduced as a resolution to this conundrum.
Inflation is shown to actually exist as a scenario, rather than a specific model. In the most general case one speaks of the inflaton field and its corresponding energy density. Models of inflation differ in their predictions and the corresponding evolution of an associated inflaton field can be explored in a cosmological context. The equations of motion are cast in a form that makes observational consequences manifest. The slow-roll approximation (SRA) is discussed as a more tractable and plausible evolution for the inflaton field and the slow-roll parameters are defined. Using the SRA, inflation predicts a near-Gaussian adiabatic perturbation spectrum resulting from quantum fluctuations in the inflaton field and the DeSitter space-time metric. These result in a predicted power spectrum of gravity waves and temperature anisotropies in the cosmic background, both of which will be detectable in future experiments.
Inflation is shown to be a rigorous theory that makes concise predictions in regards to a needed inflaton potential at the immediate Post-Planck or perhaps even the Planck epoch ($`10^{43}`$s). This offers the exciting possibility that inflation can be used to predict new particle physics or serve as a constraint for phenomenology from theories such as Superstring theory.
## II Standard Cosmology
### II.1 The Cosmological Principle
The Cosmological Principle (CP) is the rudimentary foundation of most standard cosmological models. The CP can be summarized by two principles of spatial invariance. The first invariance is isomorphism under translation and is referred to as homogeneity. An example of homogeneity can be seen in a carton of homogeneous milk. The milk or liquid, looks the same no matter where one is located within it. In the realm of cosmology, this corresponds to galaxies being uniformly distributed throughout the universe. This uniformity would be independent of the location one chooses to make the observations. Thus, a translation from one galaxy to another would leave the galactic distribution invariant (invariance under translation).
The next element of the CP is perhaps more difficult to be realized physically. This invariance is isomorphism under rotation and is referred to as isotropy. A simple way of visualizing isotropy is to say that direction, such as North or South, can not be distinguished. For example, if one were constrained to live on the surface of a uniform sphere, there would be no geometrical method to distinguish a direction in space. Although, as soon as features are introduced on the sphere (such as land masses or cracks in the surface of the sphere), the symmetry is lost and direction can be established. This fact gives a clue that isotropy, as you might have guessed, is closely related to homogeneity.
The concepts of homogeneity and isotropy may appear contradictory to local observation. The Earth and the solar system are not homogeneous nor isotropic. Matter clumps together to form objects like galaxies, stars, and planets with voids of near-vacuum in between. However, when one views the universe on a large scale, galaxies appear ‘smeared out’ and the CP holds.
Experimental proof of isotropy and homogeneity has been approached using a number of methods. One of the most convincing observations is that of the Cosmic Background Radiation (CBR). In the standard Big Bang model, the universe began at a singularity of infinite density and infinite temperature. As the universe expanded it began to cool allowing nucleons to combine and then atoms to form. About 300,000 years after the Big Bang, radiation decoupled from matter, allowing it to ‘escape’ at the speed of light. This radiation continues to cool to the present day and is observed as the CBR. As we will see, observations of the CBR gives a picture of the mass distribution at around 300,000 years. The temperature of the CBR, first predicted theoretically in the 1960’s by Alpher and Herman at $`5K`$ (alpher, ), and Gamow at a higher $`50K`$ (Gamow, ), was not taken seriously. A later prediction by Dicke, et. al. (dicke, ) yielded $`3K`$, but as Dicke and colleagues set out to measure this remnant radiation, they found someone had already made this measurement. Dicke remarked, “Well boys, we’ve been scooped” (partridge, ).
The first successful measurement of the CBR was made in 1964 by Penzias and Wilson, two scientists working on a satellite development project for Bell Labs (partridge, ). Their measurements revealed that the CBR was characteristic of a black-body with a corresponding temperature of around $`3K`$ as illustrated in Figure (1) (partridge, ). The measured wavelengths were on the order of $`7.35`$ cm, corresponding to the microwave range of the electromagnetic spectrum. The CBR in this range is referred to as the Cosmic Microwave Background (CMB)<sup>1</sup><sup>1</sup>1The significance in making this distinction will manifest itself later, but it is worth noting that other backgrounds are measurable and offer further evidence of the CP..
Another important observation of Penzias and Wilson is the fact that the CMB is uniform (homogeneous) in all directions (isotropic). Thus, the CMB offers an experimental proof of the isotropy and homogeneity of the universe.
Because of its importance, further measurements of the CBR have been carried out. One such project named COBE, for Cosmic Background Explorer, in $`1989`$, measured the CBR to have a temperature of $`2.73`$ K and a distribution that is isotropic to one part in $`10^5`$ (turner, ). COBE also has the distinction of being the first satellite dedicated solely to cosmology. Future measurements will be made by dedicated satellites like COBE, but these satellites will have much higher angular resolution. They are planned to be launched around the beginning of the century.
Balloon born experiments have been able to measure the background spectrum with greater resolution than COBE and the preliminary results seem to favor the type of spectrum predicted by the inflationary scenario, to be discussed later (Melchiorri:1999br, ),(Lange:2000iq, ). Several satellite projects are planned, MAP, for Microwave Anisotropy Probe<sup>2</sup><sup>2</sup>2For more info see: http://map.gsfc.nasa.gov/ will be launched at the end of this year by NASA and another named the PLANCK Explorer is planned for launched by the European Space Agency<sup>3</sup><sup>3</sup>3http://astro.estec.esa.nl/SA-general/Projects/Planck/ around the year 2006. The accurate measurement of the CBR offers an observational test of cosmological models, as well as, the CP.
In addition to these benefits of CBR observations, the CBR can also be used to setup a Cosmic Rest Frame (CRF). This concept is reminiscent to the ideas of Ernst Mach. One chooses a reference frame to coincide with the Hubble expansion, i.e., with the motion of the average distribution of matter in the universe. It is convenient to define our coordinates in this frame to save confusion in measurements such as the expansion of spacetime and the Hubble Constant; however, these coordinates are in no way ‘absolute’ coordinates. Using the CBR to define the CRF and taking galaxies as the test particles of the model serves to greatly simplify the dynamics in an expanding universe. The CRF is used to ease calculations and make the interpretation of the dynamics of an expanding universe more tractable.
The current and proposed measurements of the CBR offer a convincing test of the homogeneity of space. Measurements of the temperature of the CBR are uniform to one part in $`10^5`$. This suggests the universe is homogeneous and isotropic to a high degree of accuracy. However, since this measurement is taken from our (the Earth’s) vantage point, one can not assume the same conclusion from another vantage point. This can be remedied by considering how the CBR is related to the distribution of matter at the time the photons of the CBR decoupled. This offers a ‘snap shot’ of the inhomogeneities in the density of the universe. If these regions contained more inhomogeneity, galaxies would not be visible today. This idea will be discussed in more detail later; as an alternative one can introduce the Copernican Principle (CP).
The CP states that no observers occupy a special place in the universe. This appears to be a favorable prediction, based on the evidence above, as well as lessons coming from the past. For example, the correct model of the solar system was not realized until humans realized they were not the center of the solar system. This may be a bit humbling to the human ego, but the Copernican Principle, along with homogeneity and isotropy, serve to greatly simplify the number of possible cosmological models for the universe. Later, it will be seen that homogeneity follows naturally from inflation. If the universe went through a brief period of rapid expansion, the fact that galaxies exist at all will be a necessary and sufficient condition for a homogeneous universe.
There is also the proposal for cosmic ‘no-hair’ theorems. These theorems are similar to the ‘no-hair’ proposal of black holes, which predict that any object that contains an event horizon will yield a Schwartzschild spherically symmetric solution at the singularity. The Big-Bang singularity is no exception, and the event horizon is the Hubble distance to be explored in sections to come. For now, experiment suggests that it is safe to assume the Copernican Principle is valid.
Below is a brief descriptive summary of observational methods for testing the CP:
* Particle Backgrounds – These observations represent the strongest argument for isotropy and homogeneity. As the universe evolved it cooled allowing various particle species to become ‘frozen out’, meaning that the particles were freed from interactions. Photons, for example, became frozen out at the time of decoupling and are visible today as the CBR. These backgrounds serve as an important experimental test for predictions by various cosmological models.
* The Observed Hubble Law – This law states that the farther away a galaxy is, the faster it will be observed to recede<sup>4</sup><sup>4</sup>4One must be careful here, as we will see the spacetime between the galaxy and us is actually what is expanding, the galaxy itself is not really receding.. This phenomena is observed through a redshift of the light coming from the galaxy and will be described in a later section. The observed redshift, first witnessed by Edwin Hubble was the first indication that the universe obeys the CP.
* Source Number Counts – Of all methods this is the most uncertain at this time. This method requires collecting light from galaxies and inferring whether ‘clustering’ occurs. One debate over the accuracy of such methods is based on the idea that most matter in the universe might be of a non-luminous type, the so-called Dark Matter. Another problem is that current technology does not allow observations at distances far enough to get a good sample of the population. However, this technique shows promise for the future, and the SLOAN<sup>5</sup><sup>5</sup>5http://www.sdss.org/ Digital Sky Survey is a current project that will map in detail one-quarter of the entire sky, determining the positions and absolute brightness of more than $`100`$ million celestial objects. It will also measure the distances to more than a million galaxies and quasars.
* Inflation – Although it is premature at this point to discuss observational consequences of inflation, it will be shown that inflation predicts small perturbations in the universe that result in the large-scale structure observed today. It will be shown that if these perturbations were too large then the structure we observe today would not be possible. Thus, if inflation can be proved through observation, it would imply the universe must have been very homogeneous at the time of decoupling.
The established concepts of the CP aid in simplification of cosmological models, but a further simplification can be made by invoking the Perfect Cosmological Principle. This principle differs from the previous one in that it assumes temporal homogeneity and isotropy. This would imply a static universe, for if the universe were expanding or contracting it would not look the same now, as it did in the past. However, one exception that will prove important later is the case of a (anti or quasi) DeSitter Space. By the observations of Edwin Hubble and the theoretical work by Lamaître<sup>6</sup><sup>6</sup>6Lamaître will not be mentioned further but it is worth noting that his work and persistence, backed by the experimental efforts of Hubble, were instrumental in convincing Einstein that the universe was indeed expanding. After this persuasion, Einstein was quoted as saying this was the biggest mistake of his career (hawley, ). it was shown that the expansion of the universe is an accurate assumption. CP models further suggest that a static universe would be as stable as a pencil standing on its end. Thus, the Perfect Cosmological Principle does not appear to be an acceptable assumption within the standard model (hawking, )<sup>7</sup><sup>7</sup>7This is not totally correct. In some space-times, such as anti-DeSitter space, there exists temporal homogeneity. For a rigorous treatment of such space-times consult (hawking, )..
The last element to be discussed concerning the CP is the Weyl Postulate. This postulate formally states that, “the world lines of galaxies designated as ‘test particles’ form a 3-bundle of nonintersecting geodesics orthogonal to a series of spacelike hypersurfaces” (narlikar, ). In other words, the geodesics on which galaxies travel do not intersect. This adds another symmetry to the picture of the expanding universe allowing simplification of the spacetime metric and the Einstein equations.
### II.2 The Expanding Universe
In the mid-twenties, Edwin Hubble was observing a group of objects known as spiral nebulae<sup>8</sup><sup>8</sup>8It would later be found that most of these nebula were in fact galaxies (hawley, ).. These nebulae contain a very important class of stars known as Cepheid Variables. Because the Cepheids have a characteristic variation in brightness (bergstrom, ), Hubble could recognize these stars at great distances and then compare their observed luminosity to their known luminosity. This allowed him to compute the distance to the stars, since luminosity is inversely proportional to the square of the distance (bergstrom, ). The intrinsic, or absolute, luminosity is calculated from simple models that have been commensurate with observations of near Cepheids.
When Hubble compared the distance of the Cepheids to their velocities (computed by the redshift of their spectrum) he found a simple linear relationship,
$$\stackrel{}{V}_H=H\stackrel{}{r},$$
(1)
where $`\stackrel{}{V}_H`$ is the velocity of the galaxy, $`H`$ is the so-called Hubble Constant, and $`\stackrel{}{r}`$ is the displacement of the galaxy from the Earth. It will be shown later that the Hubble constant is not actually a constant, but can be a function of time depending on the chosen model. The standard notation is to adopt $`H_0`$ as the ‘current’ observed Hubble parameter, whereas $`H=H(t)`$ is referred to as the Hubble constant. The current accepted value of the Hubble parameter is,
$$H_0=100h_0\text{km}\text{s}^1\text{Mpc}^1\text{where}\mathrm{\hspace{0.33em}\hspace{0.33em}0.5}<h_0<0.8.$$
(2)
The unit of length, Mpc, stands for Megaparsec<sup>9</sup><sup>9</sup>9$`1Mpc=10^6parsecs3lightyears3\times 10^{16}meters.`$
A parsec is the distance to an object that has an angular parallax of $`1^{}`$ and a baseline of 1 A.U. For more on Observational Astronomy see (filippenko, ). .
Hubble’s interpretation of his data was crucial in helping determine the correct model for the universe. Hubble had found that the galaxies, on average, were receding away from us at a velocity proportional to their distance from us (1). This suggests a homogeneous, isotropic, and expanding universe. By this finding, the choices of cosmological models became greatly restricted.
Perhaps it is worth mentioning that the above analysis by Hubble is not quite as easily done as one might think. One factor that must be considered in the calculation of the Hubble velocity field (1) is the concept of peculiar velocity. This is the name given to the motion of a galaxy, relative to the CRF, due to its rotation and motion as influenced by the gravitational pull of nearby clusters. This speed, $`v_p\pm 500\text{km}\text{s}^1`$, can be neglected at far distances where the Hubble speed, $`V_H500\text{km}\text{s}^1`$. Thus, when Hubble conducted his survey most of the nebulae were too near to rule out an effect by the peculiar velocity. As a result, Hubble found $`H_0500\text{km}\text{s}^1\text{Mpc}`$, much greater than the value obtained today from surveys of type Ia supernovae<sup>10</sup><sup>10</sup>10Supernova Ia, like Cepheid Variables, have a known ‘signature’ and can therefore be used as ‘Standard Candles’, but unlike the Cepheids, supernovae are much more luminous and can therefore be seen at much greater distances (richtler, )..
#### II.2.1 The Hubble Law and Particle Kinematics
The Hubble law (1) is a direct result of the CP. Consider the expansion of the universe, which must occur in a homogeneous and isotropic manner according to the CP. The expansion can be visualized with the analogy of a balloon with a grid painted on it. Of course this should not be taken literally, since the spatial extent of the universe is three dimensional. Think of the grid as a network of meter sticks and clocks at rest with respect to the Hubble expansion, which corresponds to the Cosmic Rest Frame (CRF) mentioned earlier. Due to the expansion, two particles<sup>11</sup><sup>11</sup>11Remember that when one speaks of a cosmological model, the test particles are galaxies. initially separated by a distance $`l_0`$, will be separated by a distance $`l(t)=a(t)l_0`$ at some later time $`t`$, see Figure (2). Because of the CP, the function $`a(t)`$, known as the scale factor, can only be a function of time. From this relation, the speed of the observers relative to each other is,
$$v(t)=\frac{dl}{dt}=\dot{a}l_0=\left(\frac{\dot{a}}{a}\right)l(t)=H(t)l(t),$$
where $`\dot{a}`$ is the time derivative of the scale factor.
From this derivation of the Hubble law, it becomes manifest that the Hubble Constant can depend on time. In this new way of defining $`H(t)=\dot{a}(t)/a(t)`$, $`H(t)`$ measures the rate of change of the scale factor, $`a(t)`$, and offers a way to link observations (like Hubble’s) with a proposed model using the scale factor. For Hubble’s observations, the distance $`l(t)`$ was small and $`H(t)`$ could be estimated by a linear relation yielding equation (1).
To understand how particles ‘come to rest’ in the CRF, consider a particle starting out with a peculiar velocity $`v_pc`$. The particle passes a CRF observer ($`O_1`$) at time $`t`$ and travels a distance $`dl=v_pdt`$. At this time the particle passes another CRF observer ($`O_2`$), who has a velocity $`dv=Hdl=Hv_pdt`$ relative to $`O_1`$. $`O_2`$ measures the particle’s peculiar velocity as, $`v_p(t+dt)=v_p(t)dv`$. This shows that the peculiar velocity satisfies the equation of motion,
$$\frac{dv_p}{dt}=\frac{dv}{dt}=Hv_p=\left(\frac{\dot{a}}{a}\right)v_p.$$
(3)
Solving this differential equation yields,
$$v_p\frac{1}{a}.$$
This indicates that the peculiar velocity decreases as the scale factor increases. Indicating that as the universe expands, particles with peculiar velocities tend to go to zero meaning they ‘settle’ into the CRF.
#### II.2.2 The Robertson Walker Metric
The only metric compatible with Hubble’s findings and the Cosmological Principle is the Robertson Walker Metric (RWM) with the corresponding line element,
$$ds^2=c^2dt^2a^2(t)\left[\frac{dr^2}{1kr^2}+r^2d\theta ^2+r^2\mathrm{sin}^2\theta d\varphi ^2\right].$$
(4)
For a brief explanation consider the following:
* For the metric to be homogeneous, isotropic, and obey the Weyl postulate, the metric must be the same in all directions and locations,
$$g_{\mu \nu }=g_{\mu \mu }.$$
* For a uniform expansion we must have a scale factor $`a(t)`$ that is a function of time only.
* Allowance for any type of geometry (curvature) must be made. This is represented by the constant $`k`$, where $`k=0,k=1,`$ and $`k=1`$ corresponds to flat, spherical, and hyperbolic geometries, respectively.
There are a few subtleties that must be discussed. First, the $`r`$ that appears in the line element (4) is not the radius of the universe. The $`r`$ is a dimensionless, comoving coordinate that ranges from zero to one for $`k=1`$. The measurable, physical distance is given by the RWM above. Choosing a frame common to two distinct points, one obtains,
$$ds^2=c^2dt^2a^2(t)(1kr^2)^1dr^2,$$
for their separation. Where $`d\theta `$ and $`d\varphi `$ are zero, because one has freedom to arrange the axis and $`ds^2`$ represents their separation in spacetime. Thus, their spatial separation is found by considering spacelike hypersurfaces, that is $`dt^2=0`$. Thus, their separation is
$$d_p=a(t)_0^r(1kr^2)^{1/2}𝑑r.$$
Evidently for a $`k=0`$ flat universe, the distance is simply,
$$d_p=a(t)r.$$
(5)
Thus, $`a(t)`$ has units of length and depends on the geometry of the spacetime.
The next issue is that of curvature. The curvature of the universe is determined by the amount of energy and matter that is present. The space is one of constant curvature determined by the value of $`k`$. Because any arbitrary scaling of the line element (4) will not affect the sign of $`k`$, we have the following convention<sup>12</sup><sup>12</sup>12When the metric is invariant under multiplication by a scale factor, the metric is said to be conformally invariant.:
* k=1 represents positive, spherical geometry
* k=0 represents flat Minkowski space
* k=-1 represents negative, hyperbolic geometry
#### II.2.3 The Cosmological redshift
One observable prediction of an expanding universe is that of redshifting. When a light wave is traveling from a distant galaxy, to our own, it must travel through the intervening spacetime. This results in a stretching of the wavelength of light, since the spacetime is expanding. This longer wavelength results in the light being shifted to a ‘redder’ part of the spectrum. Of course light with wavelengths differing from visible light will not be visible to the human eye, but they will still be shifted to longer wavelengths.
To quantify this analysis, consider a light ray which must travel along a null geodesic ($`ds^2=0`$) in the comoving frame with constant $`\theta `$ and $`\varphi `$. Using (4),
$$0=ds^2=c^2dt^2a^2(t)\left[\frac{dr^2}{1kr^2}\right],$$
so,
$$cdt=\frac{a(t)dr}{\sqrt{1kr^2}}.$$
Integrating yields,
$$_{t_e}^{t_0}\frac{cdt}{a(t)}=_0^{r_e}\frac{dr}{\sqrt{1kr^2}}f(r_e),$$
(6)
where $`t_e`$ is the time the light pulse was emitted, $`t_0`$ was the time the light pulse was received, and $`r_e`$ is the distance to the galaxy. Thus, if one knows $`a(t)`$ and $`k`$, one can find the relation between the distance and the time. However, consider emitting successive wave crests in such a brief time that $`a(t)`$ is not given a chance to increase by a significant amount; i.e., the waves are sent out at times $`t_e`$ and $`t_e+\mathrm{\Delta }t_e`$ and received at times $`t_0`$ and $`t_0+\mathrm{\Delta }t_0`$, respectively. Then (6) becomes,
$$_{t_e+\mathrm{\Delta }t_0}^{t_0+\mathrm{\Delta }t_e}\frac{cdt}{a(t)}=_0^{r_e}\frac{dr}{\sqrt{1kr^2}}$$
Subtracting (6) from this equation and using the fact $`a(t)`$ doesn’t change, one can use the fundamental theorem of calculus to obtain,
$$\frac{c\mathrm{\Delta }t_0}{a(t_0)}\frac{c\mathrm{\Delta }t_e}{a(t_e)}=0,$$
or
$$\frac{c\mathrm{\Delta }t_0}{c\mathrm{\Delta }t_e}=\frac{a(t_0)}{a(t_e)}.$$
$`c\mathrm{\Delta }t`$ is just the wavelength, $`\lambda `$. Thus, it follows that the red shift, $`z`$, can be defined by
$$z=\frac{a(t_0)}{a(t_e)}1=\frac{\lambda _0}{\lambda _e}1=\frac{\mathrm{\Delta }\lambda }{\lambda _e}.$$
(7)
Here $`a(t_0)`$ is the scale factor of the universe as measured by a comoving observer when the light is received, $`a(t_e)`$ is the scale factor when the light was emitted in the comoving frame, $`\lambda _0`$ is the wavelength observed and $`\lambda _e`$ is the wavelength when emitted. It is clear that $`z`$ will be positive, since $`a(t_0)>a(t_e)`$, that is the universe is getting larger.
In addition to this cosmological redshift, which is due to the expanding universe, there can also be gravitational redshifts and Doppler redshifts. At great distances the former two can be neglected, but in local cases all three must be considered.
It must also be stressed that the Special Relativity (SR) formula for redshift can not be used. This is because SR only holds for ‘local’ physics. Attempting to use this across large distances can result in a contradiction. For example, the expansion rate of the universe can actually exceed the speed of light at great distances. This is not a violation of SR, because a ‘chain’ of comoving particles (galaxies) can be put together, spaced so the laws of SR are not violated. By summing together the measurements of each set of galaxies, one finds the expansion rate to exceed that of light, although locally SR holds locally (bergstrom, ). Another explanation is that in a universe described by SR, no matter or energy exists and the metric never changes. On the contrary, in an expanding spacetime none of these requirements are true. Although, SR continues to hold locally, since a ‘small enough’ region can always be chosen where the metric is approximately flat.<sup>13</sup><sup>13</sup>13One must be careful by what is meant by ‘small enough’. This technical point need not concern us with the present discussion, see (will, ).
### II.3 The Friedmann Models
To describe the expansion of the universe one must use the RWM along with the Einstein equations,
$$R_{\mu \nu }\frac{1}{2}g_{\mu \nu }R+\frac{\mathrm{\Lambda }g_{\mu \nu }}{c^2}=\frac{8\pi G}{c^4}T_{\mu \nu },$$
(8)
to determine the equations of motion. For reference, a summary of the metric coefficients, the Christoffel Symbols, and the Ricci Tensor components are presented in (collins, , Chapter 15). Note that in this book the scale factor $`a(t)`$ is written $`R(t)`$.
Before proceeding any further, an appropriate stress-energy tensor must be provided. This is the difficult part of the process. The composition of the known universe is a very controversial topic. The standard procedure is to consider simplified distributions of mass and energy to get an approximate model for how the universe evolves.
At this point, units are chosen such that the speed of light, $`c`$ is set equal to unity. This gives the simplification that the energy density, $`ϵ`$ is equal to the mass density, $`\rho `$ using, $`ϵ=\rho c^2=\rho .`$ This also allows mass and energy to be considered together, which is in the spirit of the stress-energy tensor. The mass/energy density will be referred to as the energy density for the remainder of this paper. The stress-energy tensor may be given as:
$$T_{\mu \nu }=(p+\rho )u_\mu u_\nu pg_{\mu \nu },$$
(9)
where $`p`$ is the pressure, $`\rho `$ is the density, and $`u_\mu `$ is the four-velocity.
At the earliest epoch of the universe, the contribution of photons to the energy density would have been appreciable. However, as the universe cooled below a critical temperature, allowing the photons to decouple from baryonic matter, the photon contribution became negligible. Thus, it is easier to consider different energy distributions for different epochs in the universe. The massive contribution to the energy density is usually referred to as the Baryonic contribution, since baryons (protons, neutrons, etc.) are significantly more massive than leptons (electrons, positrons, etc.) and leptons can therefore be disregarded as a major contributing factor to the total energy density. There is also the contribution of vacuum energy, which enters the Einstein equations through the cosmological constant, $`\mathrm{\Lambda }`$.
For each type of contribution, there is a corresponding density, $`\rho `$. The total density can be expressed as the sum of the different contributions as
$$\rho =\rho _M+\rho _R+\rho _\mathrm{\Lambda }.$$
(10)
Furthermore, assuming that one is dealing with a homogeneous and isotropic fluid, the density can be related to the pressure by a simple equation of state (see Table 1),
$$p=\alpha \rho .$$
(11)
Another useful relation involves the Conservation of Energy (1st Law of Thermodynamics). Assuming the ideal fluid expands adiabatically, one finds (finkelstein, ),
$$dE=pdV,$$
which may be rewritten as,
$$d(\rho a^3)=pd(a^3).$$
(12)
Relating (11) and (12) gives,
$$d(\rho a^3)=\alpha \rho d(a^3).$$
Using the product rule,
$$\rho d(a^3)+a^3d\rho =\alpha \rho d(a^3).$$
Which can be integrated,
$$\rho ^1𝑑\rho =(1+\alpha )a^3d(a^3).$$
$$\rho a^{3(1+\alpha )}.$$
(13)
In the Radiation epoch, where the energy density due to photons was appreciable (from about $`t`$=0 to approximately 300,000 years after the Big-Bang (dekel, )), the density due to massive particles can be neglected. The pressure is found to be equal to a third of the density, and we have a value of one-third for $`\alpha `$, so $`\rho _Ra^4`$ (bergstrom, ).
Following this epoch, the Matter Dominated epoch can be modeled after a ‘dust’ that uniformly fills space. Because the temperature of the universe had fallen to around 3000 K, most of the particles had non-relativistic velocities ($`vc`$). This corresponds to a negligible pressure and $`\alpha `$ is therefore zero, $`\rho _\mathrm{\Lambda }\text{constant}`$ (bergstrom, ).
The last case to consider is that of the vacuum energy. If the cosmological constant is indeed nonzero, this form of energy density will dominate. For this relation, the pressure is commensurate with that of a negative density. This would imply a value of -1 for $`\alpha `$, so $`\rho _Ma^3`$ (bergstrom, ). These results are summarized in Table 1.
Given an expression for the energy-momentum tensor, one can now proceed to find the equations of motion. The metric coefficients follow from the Robertson Walker line element, which is given by equation (4). Using these coefficients one can obtain the expression for the left side of the Einstein equations (8). Thus, from the Einstein equations one derives the Friedmann equations in their most general form:
$$\frac{\ddot{a}}{a}=\frac{4\pi G}{3}\left(\rho +3p\right)+\frac{\mathrm{\Lambda }}{3},$$
(14)
$$\left(\frac{\dot{a}}{a}\right)^2=\frac{8\pi G\rho }{3}\frac{k}{a^2}+\frac{\mathrm{\Lambda }}{3}.$$
(15)
Apparently, if the universe is in a vacuum dominated state $`p=\rho `$, (14) indicates the universe will be accelerating. This important conclusion will be the most general requirement for an inflationary model.
Now is the time to introduce a bit of machinery to make our calculations more tractable. Recall that the Hubble constant, $`H(t)`$, is defined as
$$H(t)=\frac{\dot{a}(t)}{a(t)}.$$
(16)
Next, one defines the Deceleration parameter (named for historical reasons) as,
$$q(t)=\frac{\ddot{a}(t)}{a(t)H^2(t)}.$$
(17)
To realize how this term arises, consider the Taylor expansion of the scale factor, about the present time, $`t_0`$,
$$a(t)=a_0+\dot{a}_0(tt_0)+\frac{1}{2}\ddot{a}_0(tt_0)^2+\mathrm{},$$
where the sub-zeros indicate the terms are evaluated at the present. Using equations (16) and (17), this becomes
$$a(t)=a_0\left[1+H_0(tt_0)\frac{1}{2}q_0H_0^2(tt_0)^2+\mathrm{}\right].$$
(18)
Remembering that the crux for obtaining Hubble’s Law (1) was measuring the luminosity distance, it is of interest to consider this calculation quantitatively. The flux $`F`$ (energy per time per area received by the detector) is defined in terms of the known luminosity $`L`$ (energy per time emitted in the star’s rest frame) and the luminosity distance $`d_L`$.
$$F=\frac{L}{4\pi d_L^2}.$$
(19)
The luminosity distance must take into account the expanding universe and can be written in terms of the redshift, $`z`$ as (turner, ),
$$d_L^2=a_0^2r^2(1+z)^2,$$
(20)
where $`a_0`$ is the present scale factor and $`r`$ is the comoving coordinate that parameterizes the space. Hubble used the measured flux and the known luminosity to find the distance to the objects he measured. The distance can then be compared with the known redshift of the object using (20) and the velocity can be approximated. However, $`r`$ in (20) is not a observable and it is of interest to examine the great amount of estimation that must be used to derive the desired result analytically.
Dividing (18) by $`a_0`$ and making use of (7) yields,
$$\frac{a_0}{a}=z=H_0(t_0t)+\left(1+\frac{q_0}{2}\right)H_0^2(t_0t)^2+\mathrm{},$$
(21)
which can be solved for $`(t_0t)`$,
$$(t_0t)=H_0^1\left[z\left(1+\frac{q_0}{2}\right)z^2+\mathrm{}\right].$$
(22)
One can also expand (6) in a power series,
$$f(r)=sinn^1(r),$$
(23)
where $`r_e`$ has been replaced by $`r`$ (for simplicity) and $`sinn^1(r)`$ is defined as $`sin^1(r)`$ for $`k=1`$, $`sinh^1(r)`$ for $`k=1`$ and $`r`$ for $`k=0`$. So to lowest order, (6) can be estimated as $`r`$, and the l.h.s. of (6) can be estimated as,
$$_{t_e}^{t_0}\frac{cdt}{a(t)}\frac{c(tt_0)}{a_0}.$$
(24)
Using the approximation from (23) and the above result we have,
$$\frac{c(tt_0)}{a_0}r.$$
Substitution of $`tt_0`$ from (22) and keeping only lowest order terms yields,
$$a_0rc(tt_0)\frac{cz}{H_0}.$$
At small redshift, $`z1`$ one finds $`zv/c`$. Thus, making this final approximation one obtains,
$$czva_0H_0rH_0d_p,$$
(25)
$$vH_0d_p,$$
(26)
where $`d_p`$ is the physical distance. Thus, we have obtained Hubble’s law (1) as an approximation. This derivation reflects the reason that the law only holds locally. The number of approximations that were needed to proceed was appreciable. Furthermore, one finds that this law deviates significantly at large $`z`$ as one would expect.
For a matter dominated model, one finds the exact Hubble relation to be given by (turner, ),
$$H_0d_l=q_0^2\left[zq_0+(q_01)\left(\sqrt{2q_0z+1}1\right)\right],$$
(27)
which depends on the deceleration parameter, $`q_0`$, which in turn relies on the curvature and the total mass density of the universe.
### II.4 Matter Dominated Models
The present epoch is best described by a matter dominated universe, so it is perhaps best to explore this model first. Again, matter domination corresponds to a non-relativistic, homogeneous, isotropic ‘dust’ filled universe with zero pressure. By setting $`p=0`$ in (15) and incorporating the $`\mathrm{\Lambda }`$ term into a total density, $`\rho _T=\rho +\frac{\mathrm{\Lambda }}{8\pi G}`$, the Friedmann equations for a matter dominated universe emerge,
$$\left(\frac{\dot{a}}{a}\right)^2+\frac{k}{a^2}=\frac{8\pi G\rho _T}{3}.$$
(28)
$$2\frac{\ddot{a}}{a}+\left(\frac{\dot{a}}{a}\right)^2+\frac{k}{a^2}=0.$$
(29)
Again, the value of $`k`$ describes the geometry of the space. Equation (29) is actually obtained by combining (14) and (15) and is often called the acceleration equation.
The idea of the total density, $`\rho _T`$, might be a bit confusing since it has been stated that the model is matter dominated. Although this is true, there can still be a small contribution in the form of radiation and other forms of energy, such as dark matter. The point is that any of these should be much less than $`\rho _M`$ for the model to be accurate. It will also be seen that $`\rho _M`$ can also be broken into different contributions as tacitly stated in the previous remark about dark matter. For the remainder of this section we take $`\rho `$ to mean $`\rho _T`$ to keep the notation as simple as possible.
#### II.4.1 The Einstein-DeSitter Model
The Einstein-DeSitter model is a matter dominated Friedmann model with zero curvature ($`k=0`$). This model corresponds to a Minkowski universe (zero curvature), in which the universe will continue to expand forever with just the right amount of energy to escape to infinity. It is analogous to launching a rocket. If the rocket is given insufficient energy, it will be pulled back by the Earth. However, if its energy exceeds a certain critical velocity (escape velocity), it will continue into space with ever increasing speed. If it has exactly the escape velocity, it will proceed to escape the Earth with a velocity going to zero as the rocket approaches spatial infinity. The Einstein-DeSitter model corresponds to the universe having exactly the right escape velocity provided by the Big-Bang to escape the pull of gravity due to the matter in the universe.
By substituting $`k=0`$ into (29) and (28), the Friedmann equations become
$$0=2\frac{\ddot{a}}{a}+\left(\frac{\dot{a}}{a}\right)^2=2q(t)H^2(t)+H^2(t),$$
(30)
$$\left(\frac{\dot{a}}{a}\right)^2=H(t)^2=\frac{8\pi G\rho }{3}.$$
(31)
By solving (31) for $`\rho `$, a critical density can be found for a flat universe. The critical density is the amount of matter required for the universe to be exactly flat ($`k=0`$) and is a function of time. The critical density at the present is defined as,
$$\rho _c=\frac{3H_0^2}{8\pi G}.$$
(32)
If the density of the universe exceeds the critical density, the universe is open. Conversely, if the density is below $`\rho _c`$ the universe is open. For the observed Hubble parameter as defined in (2), the critical density today corresponds to a value,
$$\rho _c2h_0^2\times 10^{23}\frac{\text{g}}{\text{m}^3}.$$
(33)
This is equivalent to roughly 10 hydrogen atoms per cubic meter. Although, this is incredibly small compared to Earthly standards, it must be remembered that most of space is empty and the concern is the total energy density.
Notice that the critical density depends on the Hubble constant. This means that the density required for a flat universe will change with time, in general, as the universe expands. For the universe to be ‘fine-tuned’ to this precision is highly improbable; yet, most observations suggest this type of geometry. This paradoxical issue is referred to as the Flatness problem and will lead to one of the claimed triumphs of Inflation theory. Because it is believed that the universe is so close to being flat, it is useful to define the density parameter, $`\mathrm{\Omega }`$. $`\mathrm{\Omega }_0`$ is the ratio of the density observed today, $`\rho _0`$, to the critical density, $`\rho _c`$. In general, $`\mathrm{\Omega }`$ is the ratio of the density to the critical value.
The quantity $`\mathrm{\Omega }`$ together with Equation (30), which implies $`q_0=\frac{1}{2}`$, can be used to discriminate between the possible geometries for the matter dominated universe (see Table 2).
From the previous result for a matter dominated energy density, we found $`\rho a^3`$. From this relation, conservation of energy follows,
$$\frac{d}{dt}(\rho a^3)=0\rho a^3=\text{constant}.$$
This can be used to obtain a useful relation for $`\rho `$,
$$\rho =\rho _0\frac{a_0^3}{a^3}.$$
Returning to the Friedmann equation (31) and substituting the above expression for $`\rho `$ one finds,
$$\left(\frac{\dot{a}}{a}\right)^2=\frac{8\pi G\rho }{3}=\frac{8\pi G\rho _0}{3}\frac{a_0^3}{a^3}.$$
Combining terms in $`a`$,
$$a\dot{a}^2=\frac{8\pi G\rho _0a_0^3}{3},$$
now integrating,
$$_0^a\sqrt{a}𝑑a=_0^t\sqrt{\frac{8\pi G\rho _0a_0^3}{3}}𝑑t,$$
$$a(t)=a_0\left(\frac{8\pi G\rho _0}{3}\right)^{\frac{1}{3}}t^{\frac{2}{3}},$$
$$at^{2/3}.$$
(34)
So for the Einstein-DeSitter Model, the scale factor evolves as $`t^{2/3}`$.
#### II.4.2 The Closed Model
The Closed Model is characterized by a positive curvature, $`k=1`$. Thus, the spatial structure is that of the 3-sphere, similar to the surface of a sphere, but in 3 dimensions instead of 2. This model corresponds to a universe that begins at a ‘Big-Bang’ and continues to expand until gravity finally halts the expansion. The universe will then collapse into a ‘Big-Crunch’, which will resemble the reverse process of the ‘Big-Bang’. The ability of the matter (or energy) in the universe to halt the expansion obviously depends on the density. If the matter-energy density is too low, the universe will have enough momentum from the ‘bang’ to escape the pull of gravity. In the Closed Model the density of the universe is great enough to halt the expansion and start a contraction. This corresponds to a value of $`\mathrm{\Omega }>1`$, which is evident from the use of the Friedmann equations with $`k=1`$. Plugging this $`k`$ value into the Friedmann equations (29),(28) and using $`\ddot{a}=qH^2a`$ one gets,
$$2\frac{\ddot{a}}{a}+\frac{\dot{a}^2}{a^2}+\frac{1}{a^2}=2\left(qH^2\right)+H^2+\frac{1}{a^2}=0.$$
This can be expressed as
$$\frac{1}{a^2}=H^2\left[2q1\right].$$
(35)
Equation (28) takes the form,
$$\frac{\dot{a}^2}{a^2}+\frac{1}{a^2}=\frac{8\pi G\rho }{3}.$$
(36)
Combining (35) and (36) gives,
$$H^2+\left[H^2\left(2q1\right)\right]=\frac{8\pi G\rho }{3},$$
or
$$2qH^2=\frac{8\pi G\rho }{3}.$$
Thus,
$$\rho =\frac{3H^2q}{4\pi G}.$$
(37)
Comparing (37) with the critical density (32) and the value of $`q>1/2`$ in Table 2, it is evident that $`\rho >\rho _c`$ for the universe to be closed. In terms of $`\mathrm{\Omega }`$, this gives
$$\mathrm{\Omega }=\frac{\rho }{\rho _c}>1.$$
(38)
The advantage of equation (37) above, is that the density is expressed all in quantities that can be measured. In that, if 2 of the 3 quantities are known the third may be found.
#### II.4.3 The Open Model
The so-called Open Model<sup>14</sup><sup>14</sup>14Although it is a standard practice to refer to this case $`(k=1)`$ as the ‘Open’ Model, it should be noted that the model can actually correspond to a closed universe. This is the result of a non-trivial topology, which results in geometry that can be hyperbolic; but, the topology can cause it to be contained in a finite space (cornish, ; luminet, ). is the case where $`k=1`$ and the geometry is said to be hyperbolic. Taking $`k=1`$ in the Friedmann equations (29),(28),
$$2\frac{\ddot{a}}{a}+\frac{\dot{a}^2}{a^2}\frac{1}{a^2}=0.$$
(39)
$$\frac{\dot{a}^2}{a^2}\frac{1}{a^2}=\frac{8\pi G\rho }{3}$$
(40)
Solving these equations (39),(40) yields the same value of the density (37) as the closed model, but in this case $`q<1/2`$ and $`\mathrm{\Omega }<1`$, as previously discussed.
### II.5 Summary
* All cosmological models are characterized by ‘test particles’, which are galaxies that are distributed in a homogeneous and isotropic manner in accordance with the CP.
* Open Models are characterized by $`\mathrm{\Omega }<1`$, negative curvature ($`k=1`$), hyperbolic geometry, a deceleration parameter $`q<1/2`$, and infinite spatial extent (ignoring topology).
* Closed Models are characterized by $`\mathrm{\Omega }>1`$, positive curvature ($`k=1`$), spherical geometry, a deceleration parameter $`q>1/2`$ and finite spatial extent.
* Flat Models are characterized by $`\mathrm{\Omega }=1`$, flat geometry with no curvature ($`k=0`$), infinite spatial extent, a deceleration parameter $`q=1/2`$ and with an age corresponding to Age=$`H_0^1`$, since the Hubble Constant is in-fact constant.
## III A Brief History of the Universe
One of the successes of the hot Big Bang model is its prediction of the light elements. These predictions are verified by observations of the structure and composition of the oldest stars, quasars, and other quasi-stellar remnants (e.g., QSOs) (9907128, ; 9904407, ; 9712031, ; 9904223, ). The process by which the elements form is referred to as nucleosynthesis.
The Hot Big-Bang model predicts a universe that will go through several stages of thermal evolution. As the universe expands adiabatically, the temperature cools, scaling as
$$Ta^{3(\gamma 1)}.$$
Here $`\gamma `$ is the ratio of specific heats and is equal to $`4/3`$ for a radiation dominated universe (peacock, ). This relation is manifest, since the temperature is equivalent to the energy density divided by the volume (in natural units $`\mathrm{}=c=k=1`$). Moreover, the radiation energy density is redshifted by an additional factor of $`1/a(t)`$ since the Hubble expansion stretches the wavelength, which is inversely proportional to the energy:
$$\rho _Ra^4,$$
and
$$\rho _R=a^3E.$$
Setting the Boltzman constant to unity,
$$E=Ta^1.$$
(41)
This can also be understood using the DeBroglie wavelength of the photon (for radiation) (peebles, ). The wavelength is inversely proportional to the energy in natural units. This raises the issue of a possible factor of redshifting for the DeBroglie wavelength of a massive particle. For a particle the simple relation $`E=pc`$ does not hold; thus, the velocity of the particle can decrease to preserve its wavelength. This redshifting is analogous to that of equation (3) in the first section and gives an alternative explanation of particles in motion settling into the cosmic rest frame. This also explains why one might expect to find primarily non-relativistic (cold) matter, which just means particles traveling at speeds much less than $`c`$.
From relation (41) for the temperature, one has a quantitative way to find critical temperature scales in the evolution of the universe. For a given value of the curvature the relation between the scale factor and time can yield an expression between temperature and time. For example, in a matter dominated, flat universe,
$$Tt^{2/3},$$
which was derived earlier.
In thermal physics one is usually interested in thermal equilibrium. This consideration is accounted for by the condition,
$$\frac{\mathrm{\Gamma }}{H}1.$$
This relations shows that the reaction rate, $`\mathrm{\Gamma }`$, must be much greater than the rate at which the universe expands for thermal equilibrium to be reached. $`\mathrm{\Gamma }`$ is related to the cross-section of the given particle interaction by
$$\mathrm{\Gamma }=n\sigma |v|,$$
where $`\sigma `$ is the cross-section, $`|v|`$ is the relative speed, and $`n`$ is the number density of the species<sup>15</sup><sup>15</sup>15For the reader interested in learning more about particle interactions see, (bergstrom, , Chapter 7),(peacock, , Chapter 9).
From equation (41), one can see that the universe began as a point of infinite temperature and zero size. This is a singular point for the history of the universe and the standard Big Bang model (General Relativity) breaks down at this singularity<sup>16</sup><sup>16</sup>16Supersting Theory offers solutions to the problem of a singularity by setting an ultimate smallness, the Planck length $`(10^{33}cm)`$. This is outside the scope of this paper, but the reader is referred to (greene, ) for an excellent popular account of strings and cosmology and in (gasperini, ) there are a number of papers with rigorous treatments of string cosmology.. However, after the Planck time $`(10^{43})`$s one can follow the evolution using the concepts of thermal physics and particle theory<sup>17</sup><sup>17</sup>17For an introductory survey of particle theory see (moyer, ; tipler, ) and for particle physics in cosmology rolfs is an excellent book. The remainder of this paper will assume a basic knowledge of particle theory and thermal physics..
In the earliest times following the Planck epoch, all matter existed as free quarks and leptons. The existence of free quarks (known as asymptotic freedom) is made possible by the high energy (temperature) during the early moments of the Big Bang. The universe cools, as indicated by (41), and the quarks begin to combine under the action of the strong force to form nucleons. This phase of formation is referred to as Baryogenesis, because the baryons (e.g., protons and neutrons) are created for the first time. The expansion continues to allow leptons, such as electrons, to interact with nucleons to form atoms. At this point, referred to as recombination, the photons in the universe are free to travel with virtually no interactions. For example, hydrogen is the most abundant element to form in nucleosynthesis and at the time of decoupling the temperature of the photons has dropped to around $`3000K`$. This corresponds to less than $`13.6\text{eV}`$, the energy needed to ionize the atoms. Therefore, there are no longer free electrons to interact with the photons and in fact the energy of a photon (around $`.26\text{eV}`$, at this temperature) is so low that it can not interact with the atoms. In this way the photons have effectively decoupled from matter and travel through the universe as the cosmic background discussed in Part II. A brief summary of the most significant events are encapsulated below,
* $`10^4`$ seconds: Baryogenesis occurs, quarks condense under strong interaction to form nucleons (e.g., Protons and Neutrons)
* 1 second: Nucleosynthesis occurs, universe cools enough (photon energies $`\mathrm{\hspace{0.33em}1}\text{MeV}`$) for light nuclei to form (e.g., deuterons, alpha particles).
* $`10^4`$ years: Radiation density becomes equal to matter density, since the radiation density has extra factor of $`a^1`$ due to red-shifting. Matter density is the dominate energy density after this epoch.
* $`10^5`$ years: Recombination occurs and electrons are combined with nucleons to form atoms. This time also coincides with the decoupling of photons from matter, giving rise to a surface of last scattering of the cosmic background radiation.
* $`10^{10}`$ years: The present.
## IV Problems with the Standard Cosmology
The hot Big Bang model has been very successful in predicting much of the phenomena observed in the universe today. The model successfully accounts for nucleosynthesis and the relative abundance of the light elements, (e.g., Hydrogen $``$ 75%, Helium $`25\%`$, Lithium (trace), Berylium (trace)). The prediction of the Cosmic Background Radiation and the fact that the universe is expanding (i.e., The Hubble Law), both represent successful predictions of the Big Bang theory. However, this model suggests questions which it can not answer, which brings about its own demise. These anomalies are discussed in the following sections.
### IV.1 The Horizon Problem
Why is the universe so homogeneous and isotropic on large scales? Radiation on opposite sides of the observable universe today appear uniform in temperature. Yet, there was not enough time in the past for the photons to communicate their temperature to the opposing sides of the visible universe (i.e., establish thermal equilibrium). Consider the comoving radius of the causally connected parts of the universe at the time of recombination compared to the comoving radius at the present, found from Equation (6) (remember $`c=\mathrm{}=1`$).
$$_0^{t_{rec}}\frac{dt}{a(t)}_{t_{rec}}^{t_0}\frac{dt}{a(t)}.$$
(42)
This means a much larger portion of the universe is visible today, than was visible at recombination when the CBR was ‘released’. So the paradox is how the CBR became homogeneous to 1 part in $`10^5`$ as we discussed in Part II. There was no time for thermal equilibrium to be reached. In fact, any region separated by more than 2 degrees in the sky today would have been causally disconnected at the time of decoupling (liddleeprint, ).
This argument can be made a bit more quantitative by consideration of the entropy, $`S`$, which indicates the number of states within the model. This can be used as a measure of the size of the particle horizon (turner, ).
$$S_{Horizon}^{RD}=s\frac{4}{3}\pi t^30.05g_{}^{1/2}(m_{pl}/T)^3,$$
(43)
$$S_{Horizon}^{MD}=s\frac{4}{3}\pi t^33\times 10^{87}(\mathrm{\Omega }_0h^2)^{3/2}(1+z)^{3/2},$$
(44)
where $`m_{pl}`$ is the Planck mass, $`s`$ is the entropy density, $`g_{}`$ is the particle degeneracy, and $`z`$ is the redshift. These equations for the entropy of the horizon in a radiation dominated (43) and matter dominated universe (44), are presented only to motivate the following estimates. For an explanation please consult (turner, ).
At the time of recombination ($`z1100`$), when the universe was matter dominated, equation (44) gives a value of about $`10^{83}`$ states. Compared with a value today of $`10^{88}`$ states, this is different by a factor of $`10^5`$. Thus, there are approximately $`10^5`$ causally disconnected regions to be accounted for in the observable universe today. The hot Big Bang offers no resolution for this paradox, especially since it is assumed to be an adiabatic (constant entropy) expansion.
### IV.2 The Problem of Large-Scale Structure
In contrast to the horizon problem, the fact that the Big Bang predicts no inhomogeneity is a problem as well. How are galactic structures to form in a perfectly homogeneous universe? The fact that galaxies have been shown to cluster locally with great voids on the order of 100 Mpc, is proof of the inhomogeneity of the universe. Moreover, the $`10^5`$ anisotropies (temperature differences) on angular scales of 10 degrees as measured by the COBE satellite, form a blueprint of the seeds of formation at the time of decoupling. However, there is no mechanism within the Big Bang theory to account for these ‘seeds’, or perturbations, that result in the large-scale structure. Not only does the Big Bang predict homogeneous structure, but it also had to ‘explode’ in just the right way to avoid collapse. This is often called the fine-tuning problem. Cosmologists would like to have a theory that does not require specific parameters to be put in the theory ad hoc. The density, the expansion rate, and the like, prove to be other unfavorable aspects of the hot Big Bang.
### IV.3 The Flatness Problem
The flatness problem is another example of a fine-tuning problem. The contribution to the critical density by the baryon density, based on calculations from nucleosynthesis and the observed abundance of light elements, are in good agreement with observations and give $`\mathrm{\Omega }_B<0.1`$. The radiation density is negligible and it is believed that non-baryonic dark matter, or quintessence/dark energy (non zero cosmological constant), will contribute the remainder of the critical density, yielding $`\mathrm{\Omega }=1`$. Although, an $`\mathrm{\Omega }`$ anywhere within the range of 1 causes a problem.
The Friedmann equation (28) can be used to take into account how $`\mathrm{\Omega }`$ changes with time. Noting that $`H=\dot{a}/a`$ and $`\mathrm{\Omega }=\rho /\rho _c`$, one can divide (28) by $`H^2`$ to obtain,
$$\mathrm{\Omega }(t)1=\frac{k}{a^2H(t)^2}.$$
Using the relationships between the scale factor and time,
$$\text{Matter Domination:}a_Mt^{2/3},$$
$$\text{Radiation Domination:}a_Rt^{1/2},$$
and using the definition of $`H`$ yields,
$$\text{Matter Domination:}\mathrm{\Omega }(t)1t^{2/3},$$
$$\text{Radiation Domination:}\mathrm{\Omega }(t)1t,$$
From these relations one can see that $`\mathrm{\Omega }`$ must be very fine-tuned at early times. For example, requiring $`\mathrm{\Omega }`$ to be one today, corresponds to a value of $`\mathrm{\Omega }(1)110^{16}`$ at the time of decoupling and a value of $`\mathrm{\Omega }(10^{43})110^{60}`$ at the Planck epoch. This value seems unnecessarily contrived and indicates that we live at a very special time in the universe. That is to say, when the universe happens to be flat. An alternative is that the universe has been, is, and always will be flat. However, this is a very special case and it would be nice to have a mechanism that explains why the universe is flat. The Big Bang offers no such explanation.
### IV.4 The Monopole Problem
At early times in the expansion ($`z>1000`$), the physics of the universe is described by particle theory. Many of these theories predict the creation of topological defects. These defects arise when phase transitions occur in particle models. Since the temperature of the universe cools as the expansion proceeds, these phase transitions are natural consequences of symmetry breakings that occur in particle models. Several types of defects are described briefly below (peacock, , Chapter 10),
* Domain Walls – Space divides into connected regions; one region with one phase and the other region exhibiting the other phase. The regions are separated by walls of discontinuity described by a certain energy per unit area.
* Strings – These are linear defects, characterized by some mass per unit length. They can be visualized at the present time as large strings stretched in space that possibly cause galaxies to form into groups. They serve as an alternative to inflation, for explaining the large-scale structure of the universe. However, at the moment they are not favored due to lack of observations of the gravitational-lensing effect they should exhibit<sup>18</sup><sup>18</sup>18Also note that these are not visible objects, they are distortions in the space-time fabric..
* Monopoles – These are point defects, where the field points radially away from the defect, which has a characteristic mass. These defects have a magnetic field configuration at infinity that makes them analogous to that of the magnetic monopole, hypothesized by Maxwell and others.
* Textures These objects are hard to visualize and are not expected to form in most theories. One can consider them as a kind of combination of all the other defects.
Out of all these defects, monopoles are the most prevalent in particle theories. It becomes a problem in the hot Big Bang model, when one calculates the number of monopoles produced in events, such as the electroweak symmetry breaking. One finds they would be the dominate matter in the universe. This is contrary to the fact that no monopole has ever been observed, directly or indirectly, by humans. These monopoles would effect the curvature of the universe and in turn the Hubble parameter, galaxy formation, etc. Therefore, unwanted relics, such as monopoles, remain an anomalous component of the hot Big Bang theory.
## V The Inflationary Paradigm
In past years, inflation has become more of a scenario than model. A plethora of models have been suggested, all of which share the common feature that the universe goes through a brief period of rapid expansion. This rapid expansion is manifested in the evolution of the scale factor, $`a(t)`$. In the case of inflation, $`a(t)t^n`$, where $`n>1`$ and the universe expands faster than light. This does not violate relativity, since the spacetime is the thing expanding (i.e., no information is being transferred). Since $`n`$ can take on any value greater than one, this is already an example of the flexibility of the theory.
In Section II it was shown by Equation (14) that if an equation of state $`p=\rho `$ is achieved and one has a positive cosmological constant, then the universe will accelerate. Incorporating the cosmological constant, $`\mathrm{\Lambda }`$, into an energy density, $`\rho _\mathrm{\Lambda }`$, and assuming it is the dominate one can use (15) and (16) to obtain,
$$H^2=\left(\frac{\dot{a}}{a}\right)^2=\frac{8\pi G\rho _\mathrm{\Lambda }}{3}\frac{k}{a^2},$$
(45)
One can choose to ignore the curvature term, since one anticipates a large increase in the scale factor. That is, the presence of the scale factor in the denominator of the $`k/a^2`$ term in the equation above will leave this term negligible. This is often referred to as the redshifting of the curvature, since the effect of the curvature can be ignored if the scale factor becomes large enough during a period of constant energy density $`\rho _\mathrm{\Lambda }`$. This is actually a glimpse of how the flatness problem will be resolved. So, ignoring the curvature term, we have
$$\left(\frac{\dot{a}}{a}\right)^2=\frac{8\pi G\rho _\mathrm{\Lambda }}{3}.$$
(46)
This is a differential equation with the solution,
$$a(t)e^{Ht},$$
(47)
where $`H=\left(\frac{8}{3}\pi G\rho _\mathrm{\Lambda }\right)^{1/2}`$ and since $`\rho _\mathrm{\Lambda }`$ is a constant, so is $`H`$. This model is referred to as the DeSitter model<sup>19</sup><sup>19</sup>19Not to be confused with the Einstein-DeSitter model.. By introducing a negative pressure, the flatness problem is solved. The crux of this argument is that $`\rho _\mathrm{\Lambda }`$ is a constant<sup>20</sup><sup>20</sup>20Actually, $`\mathrm{\Lambda }`$ does not have to be a constant; in-fact, it can be a function of time. Such vacuum energies are referred to as Quintessence, or Dark energy, and are the subject of much research. Unfortunately, time will not permit a discussion ((0001051, ),(9908518, ),(9912046, )).. This comes from the fact that $`\rho _\mathrm{\Lambda }`$ is an intrinsic property of the spacetime manifold. As the manifold is stretched, this vacuum energy does not change. Another way this can be explained is by that the Einstein equations are arbitrary up to a constant term $`\mathrm{\Lambda }`$. The disadvantage of this explanation is that it does not manifest the connection between cosmology and particle theory (more on this later).
Since $`\rho _\mathrm{\Lambda }`$ is taken to be the dominate form of energy, the other contributions to the density in the Friedmann equation (28) are also redshifted away, since $`\rho _Ma^3`$, $`\rho _Ra^4`$. This leads to the conclusion that no matter what the initial distribution of $`\rho _T=\rho _M+\rho _R+\rho _\mathrm{\Lambda }`$, the vacuum energy will eventually dominate. Thus, the assumption of $`\rho _\mathrm{\Lambda }`$ domination can actually be relaxed.
So given a constant vacuum term, the DeSitter scenario ‘drives’ the universe to a flat geometry, thus approaching $`\rho =\rho _c`$, where $`\rho _c`$ is the critical density, (i.e., $`\rho _c=3H^2/8\pi G`$). This evolution, if allowed to continue, will produce an empty universe with practically no radiation or matter. The fact that we live in a universe that is full of matter and radiation is why the original proposal, by DeSitter, was rejected and forgotten.
The revision of this idea was suggested by Guth in the early 1980’s (guth, ),(guth2, ). The crux to the modern inflationary scenario, in contrast to the DeSitter model, is to limit the amount of time that this rapid expansion (inflation) occurs. Guth explained the physical mechanism for such an inflationary period as corresponding to a phase transition in the early universe. By limiting the time of the quasi-exponential expansion, Guth was able to produce a universe more like our own. Unlike DeSitter’s model, which was based on a pure solution to Einstein’s equations, Guth’s idea was based on ideas from particle physics. Guth was studying a class of grand unified theories (GUTs) and the predictions they make about particle production in the universe. This suggested how cosmology could be united with particle physics in a phenomenological manner, which has become one of the most appreciated beauties in modern physics today.
### V.1 Particle Physics and Cosmology
To better understand the motivation behind inflation, it is important to outline a few aspects of particle physics. Often inflation is introduced in an abstract and unaesthetic manner. One speaks of an inflaton field, an arbitrary scalar field, for which there is no physical motivation. This is often the case because this type of introduction requires limited knowledge of the relevant topics. This includes, but is not limited to, the relativistic Schrödinger (Klein-Gordon) equation, the Dirac equation, scalar fields, symmetries, and group theory.
Since this paper is intended for undergraduates, a brief summary is presented on how one can pursue this knowledge in a qualitative and brief manner. A brief overview of the concepts in particle theory will be provided as needed. Thus, the reader is presented with a dilemma. One may choose to pause at this point and do a brief survey of the suggested texts or one may continue and plan to fill in the details at a later time. Both options have their advantages and disadvantages. I chose the former.
From the author’s experience, a student should read through all of the references to get an intuitive picture of the theory and then go back and comb through the details and ‘hairy’ calculations.
Three possible routes to obtaining the knowledge needed to continue are,
* Thorough Route (The one the author took)
(weinberg3, , Chapter 15-16)–Introduction to cosmology with general relativity
(tipler, , Chapter 13-14),(griffiths2, )–Elementary introduction to particle theory
(strange, , Chapter 1-6)–Introduction to relativistic quantum mechanics
(kaku, )–Introduction to quantum field theory
(peacock, ), (bergstrom, )–Bring the picture together
* Fast Route
(peacock, ),(bergstrom, )–Bergström extracts the particle physics to the appendix, so as not to interfere with the focus. Both of these books are excellent and I also recommend, (collins, ).
* Very fast route (bergstrom, , Appendix B and C)
#### V.1.1 A Brief Summary of the Modern Particle Physics
There are four fundamental forces in the realm of physics today; gravitation, electromagnetism, the weak force, and the strong force. For most of the twentieth century, physicists have worked vigorously to combine or unify these forces into one, in much the same way Maxwell combined the seemingly disparate forces of electricity and magnetism. Great progress has been made to unify three of the four forces, excluding the realm of gravitation. The first breakthrough came with the unification of electromagnetism and the weak force into the electroweak force<sup>21</sup><sup>21</sup>21As an aside to the interested reader, the electroweak force is not really a unified force, in the strict sense of the world, because the theory contains two couplings. See (kaku, ) for more. . This work was done primarily by Glashow, Salam, and Weinberg (salam, ),(weinberg2, ) in the late sixties. Although their theory was not realized until 1971, when the work of ’tHooft showed their theory and all other Yang-Mills theories could be renormalized (kaku, , Chapter 1). Later work was done to unify the strong and electroweak under the symmetry<sup>22</sup><sup>22</sup>22See (ryder, ) for a description of symmetries and how they relate to particle physics.,
$$SU(3)SU(2)U(1).$$
This model is referred to as the Standard Model and has made a number of predictions, which have been verified by experiment. However, there are many aspects of the model that suggest it is incomplete. The model produces accurate predictions for such phenomena as particle scattering and absorption spectra. Although, the model requires the input of some 19 parameters. These parameters consist of such properties as particle masses and charge. But one would hope for a model that could explain most, if not all of these parameters. This can be accomplished by taking the symmetry group of the standard model and embedding it in a higher group with one coupling. This coupling, once the symmetry is broken, would result in the parameters of the standard model. Theories of this type are often referred to as grand unified theories (GUTs). Many such models have been proposed along with some very different approaches. Some current efforts go by the interesting names; Superstring theory, Supersymmetry, Technicolor, SU(5), etc.
Of all the proposed theories the most promising at the current moment is Superstring theory. In addition to unifying the three forces, this theory can also include the fourth force, gravity. These theories (there’s more than one) can be summarized quite simply. In the standard model, and in all undergraduate physics courses, particles are considered points. If you have ever given any thought to this, it mostly likely has troubled you. You are not alone and the creators of string theory had this very idea as their motivation. String theory assumes that particles are not points, instead they are tiny vibrating strings. The modes of vibration of the string give rise to the particle masses, charges, etc. This simple picture, along with the idea of supersymmetry, produces a model that presents the standard model as a low energy approximation.
Supersymmetric theories differ from the standard model, by the existence of a supersymmetric partner for each particle in the standard model. For example, for each half-integer spin lepton there corresponds an integer spin slepton (thus, it is a boson). These supersymmetric partners are not observed today, because they are extremely unstable at low temperatures. However, some versions of the theory suggest a conservation of supersymmetric number. If this is the case, then all of the supersymmetric particles would be expected to decay into a lowest energy mode referred to as the neutralino. As a result, this particle is one of the leading candidates for cold dark matter (bergstrom, , Chapter 6).
The link with cosmology is further exhibited because the hot Big Bang model predicts that at some time in the past, the temperature was high enough for GUTs to be tested. Because it is impossible to recreate these temperatures today, the universe offers the only experimental apparatus to examine the physics of these unified theories<sup>23</sup><sup>23</sup>23This statement is not truly accurate. Particle theories, such as GUTs, will be further verified with the detection of the symmetry breaking, or Higgs particle. This particle should be detectable around 1Tev, which is currently possible.. As the universe expands, and thus cools ($`Ta^1`$), the supersymmetry is broken and the particles manifest themselves as the different particles that we observe today.
Superstring theorists have attempted to unify these supersymmetric models with gravity into a so-called Theory Of Everything (TOE). Some theories have relaxed the supersymmetric requirement and still produce TOEs by the addition of higher dimensions. Some proposed TOEs worth mentioning are: Superstrings, M-Theory, Supergravity (SUGRA), and Twistor Gravity. The details of these theories need not concern the reader at this point<sup>24</sup><sup>24</sup>24The reader is again referred to the electronic preprints at Los Alamos for the latest information on Superstring theory and the like: http://xxx.lanl.gov..
The common aspect of all of these theories is that they are usually associated with some sort of symmetry breaking mechanism, which in turn gives rise to a phase transition. In the late seventies, cosmologists explored the possibility that these effects may not be negligible (Lindecor1, ),(Lindecor2, ). In the case of an SU(5) GUT, the model predicted a world dominated by massive magnetic monopoles. In the early 1980’s Guth explored the possibilities of eliminating these relics and the associated cosmological consequences, which in turn leads to the concept of inflation.
One may argue that SU(5) is not known to be the correct theory. This is true. However, most physicists believe that any correct unified theory will exhibit symmetry breaking. Moreover, the electroweak theory has been verified experimentally and exhibits a symmetry breaking that could have given rise to inflation.
### V.2 Inflation As a Solution to the Initial Value Problems
It was discussed in the first part of this section that inflation solves the flatness problem because the universe expands at such a great rate that the curvature term is ‘redshifted’ away. Another way of stating this result is to define inflation as any period in the evolution of the universe in which the scale factor ($`a(t)`$) undergoes a period of acceleration; i.e., $`\ddot{a}(t)>0`$. This condition can be used to provide a further insight into what inflation means. Consider the quantity $`(Ha)^2`$. Knowing
$$H=\dot{a}/a,$$
it follows that
$$1/(Ha)^2=\dot{a}^2.$$
Now consider the time derivative of this quantity.
$$\frac{d}{dt}\left(\dot{a}^2\right)=2\dot{a}^3\ddot{a}<0,$$
given the conditions $`\dot{a}>0`$ and $`\ddot{a}>0`$. This implies,
$$\frac{d}{dt}\left(\frac{1}{H^2a^2}\right)<0.$$
(48)
Referring back to equation (45), and dividing through by $`H^2`$, one again gets the equation for the evolution of the density parameter, $`\mathrm{\Omega }=\rho /\rho _c`$,
$$\mathrm{\Omega }(t)1=\frac{k}{a^2H(t)^2}.$$
(49)
Comparing Equations (48) and (49) expresses the fact that the curvature decreases during inflation. More explicitly, as $`a`$ and $`H`$ increase by tremendous amounts during inflation, the right had side of (49) approaches zero since the denominator becomes large. Thus, $`\mathrm{\Omega }`$ is driven towards one and the universe is made flat by inflation. As the scale factor evolves under the condition $`\ddot{a}(t)>0`$ the density ($`\rho `$) approaches the critical density ($`\rho _c`$).
But (48) can also be written as, $`d(1/Ha)/dt<0`$ since $`H`$ and $`a`$ are both taken as positive quanitites. Recall that $`1/H`$ gives the particle horizon of a flat universe, so one can use Equation (5),
$$H^1=d_p=a(t)r1/Ha=r,$$
where $`r`$ is the comoving radial coordinate. Using $`d(1/Ha)/dt<0`$ gives the relation,
$$\frac{d}{dt}(r)<0.$$
What does this mean? This implies that during a period of inflation the comoving frame (parameterized by $`r,\theta ,\text{and}\varphi `$), SHRINKS! Remember that the comoving coordinates represent the system of coordinates that are at rest with respect to the expansion. In other words, instead of viewing the spacetime as expanding it is equally valid to view the particle horizon as shrinking. To visualize this, it is perhaps useful to again consider the idea of an expanding balloon (see Figures 3 and 4). Normally, in this example, one views two points on the surface of the balloon as getting farther apart because the balloon is expanding. However, if one chooses a frame in which the surface is not expanding this would mean that the metric, or way of measuring, would shrink. Thus, the distance between the points would get larger, since the comoving coordinates got smaller. Each frame of reference has its advantages. For the remainder of this paper I will choose the frame where the universe is seen to expand. This has the advantage that the Hubble length remains ‘almost’ constant during inflation, which eases the discussion in the analysis to follow.
Notice it is now justified to use the flat universe approximation, since inflation forces $`\mathrm{\Omega }=1`$ by the fact that $`1/a^2H^2`$ increases so rapidly compared to $`k`$ in Equation (49). Also note that $`\mathrm{\Omega }`$ doesn’t have to be entirely matter dominated. For example, $`\mathrm{\Omega }=\mathrm{\Omega }_M+\mathrm{\Omega }_\mathrm{\Lambda }=.3+.7=1`$ is an acceptable configuration in the inflation scenario.
So, the picture during inflation is that the spacetime background expands at an accelerating pace. This resolves the horizon problem, since causal regions in the early universe are stretched to regions much larger than the Hubble distance. This is because during inflation the scale factor evolves at super-luminal speeds, whereas the particle horizon (Hubble distance) is approximately constant. The particle horizon does expand at the speed of light (by definition), but this pales in comparison to the evolution of the scale factor. Remember the Hubble distance is the farthest distance light could have traveled from a source to reach an observer. Once inflation ends, the scale factor returns to its sub-luminal evolution leaving the particle horizon to “catch up”. This situation is illustrated in Figure 5. So as we look out at the sky today we are still seeing the regions of uniformity that were stretched outside the particle horizon during inflation.
A more quantitative argument is given by considering the physical distance light can travel during inflation compared to after.
$$a(t_{rec})_{t_{inf}}^{t_{rec}}\frac{dt}{a(t)}a(t_0)_{t_{rec}}^{t_0}\frac{dt}{a(t)},$$
(50)
where $`t_{inf}`$ marks the beginning of inflation, $`t_{rec}`$ is the time of recombination, and $`t_0`$ is today.
Equation (50) can be understood by making the following estimates. In the first integral, the scale factor during inflation is given by, $`a(t)e^{Ht}`$. Whereas, in the second integral one can assume the scale factor is primarily matter dominated $`a(t)t^{2/3}.`$ Furthermore, the integral on the right can be simplified by taking $`t_{rec}=0`$. Of course this only increases the integral. Lastly, $`t_{inf}`$ can be set equal to zero and then $`t_{rec}=\mathrm{\Delta }t`$ is the time inflation lasts. Thus,
$$e^{H\mathrm{\Delta }t}_0^{\mathrm{\Delta }t}\frac{dt}{e^{Ht}}t_0^{2/3}_0^{t_0}\frac{dt}{t^{2/3}}.$$
Evaluating the integrals and a bit of algebra gives,
$$H^1\left(e^{H\mathrm{\Delta }t}1\right)3t_0=2H^1,$$
(51)
where the last step uses $`H=\dot{a}/a=2/3t`$. So, we can see from (51) that as long as inflation lasts long enough ($`\mathrm{\Delta }t`$) then the horizon problem is solved.
With the discussion presented thus far, the monopole problem is solved trivially. The number of predicted monopoles per particle horizon at the onset of inflation is on the order of one (300years, ). As discussed previously, this would result in a density today that would force $`\mathrm{\Omega }1`$, which is not observed. As stated previously, the comoving (causal) horizon shrinks during inflation. Thus, if the universe starts with one monopole, it may contain that one monopole after inflation, but no more. However, this is highly unlikely if the universe inflates by an appreciative amount. Furthermore, inflation redshifts all energy densities. So, as long as the temperature does not go near the critical temperature after inflation, no additional monopoles may form.
This holds true for the other topological defects and unwanted relics associated with spontaneous symmetry breaking (SSB) in unified theories. This leads one to ask, why would the temperature increase after inflation? The mechanism by which this reheating of the universe takes place is related to the mechanisms that bring about the demise of the inflationary period. These mechanisms are understood through the dynamics of scalar fields, to be discussed in the next section.
One question has been left unresolved with reference to the problems of initial values in the hot Big Bang model. This is the problem of the origin of structure in the universe. It was pointed out that the DeSitter universe is left empty and cold with no stars or galaxies. The flatness and monopole problem were resolved by a redshifting of the various energy densities. But, if no energy is present, how can particle creation take place? This peculiar feature of inflation will be discussed in the next section, but here I would like to present a qualitative description.
At the end of inflation, all energy densities have become negligible except the vacuum density (or cosmological constant if you prefer). Where did the energy go? It went into the gravitational ‘potential’ of the universe, so energy is still conserved (guthp1, ),(guthp2, ),<sup>25</sup><sup>25</sup>25Actually, energy need not be conserved if we live in an open universe. However, this need not concern us here.. Thus, at the end of inflation there is a universe filled with vacuum energy, which takes the form of a scalar field. This scalar field is coupled to gauge fields, such as the photon. As the scalar field releases its energy to the coupled field, the universe goes through a reheating phase where particles are created as in the hot Big Bang model. The energy for this particle creation is provided by the ‘latent’ heat locked in the scalar field. More will be said on this later, but the important point is that the hot Big Bang model picks up where inflation leaves off. Thus, one may be inclined to say, inflation is a slight modification to the hot Big Bang model. One author refers to inflation as, “a bolt on accessory” (liddleeprint, ).
This all sounds very appealing, however reheating is a fragile topic for inflation and results in a number of different models. This derives from the fact that if the temperature is too high at the time of reheating, the unwanted particle relics could be re-introduced into the model! As a result, many different reheating scenarios have been proposed, along with many different models for the onset of inflation.
One surprise from inflation makes all of this worry worth it. Along with offering a solution to the various initial value problems of the hot Big Bang, inflation offers a mechanism to seed the large-scale structure of the universe. Depending on the model chosen, (e.g., reheating temperature, onset conditions, etc.) one gets predictions for the large-scale structure of the universe and the anisotropies in the cosmic background. As will be seen in the next section, this again demonstrates how the very small (quantum mechanics) can impact the very large (universal structure). In some models of inflation, a small fluctuation in the quantum foam of the Planck epoch ($`t<10^{43}`$) can give rise to the formation of galaxies, solar systems, and eventually human life! We are the result of pure chance! This is getting a little ahead of the game, so let us consider a quantitative and mechanical explanation of inflation.
### V.3 Inflation and Scalar Fields
As stated above, inflation is capable of solving many of the initial value, or ‘fine-tuning’, problems of the hot Big Bang model. This is assuming that there is some mechanism to bring about the negative pressure state needed for quasi-exponential growth of the scale factor. In the early 1980’s, Alan Guth (guth, ) was studying properties of grand unified theories or GUTs. It was found in the late 70’s that these theories predict a large number of topological defects (Lindecor1, ),(Lindecor2, ). Guth was specifically addressing the issue of monopole creation in the SU(5) GUT. It was found that the theory predicts a large number of these monopoles, and that they should ‘over-close’ the universe (Lindecor1, ),(Lindecor2, ). This means that the monopole contribution to $`\mathrm{\Omega }`$ is greater than the observed upper-bound on the density parameter, $`\mathrm{\Omega }>4`$, which comes from observation (turner, ). To remedy this, Guth suggested that the symmetry breaking associated with scalar fields in the particle theory cause the universe to enter a period of rapid expansion. This expansion ‘dilutes’ the density of the monopoles created, as stated above.
The first step in understanding the dynamics of scalar fields is to undertake the study of field theory. In field theory, one considers a Lagrangian density, as opposed to the usual Lagrangian from classical mechanics. This is because the scalar field is taken to be a continuous field, whereas the Lagrangian in mechanics is usually based on discrete particle systems. The Lagrangian ($`L`$) is related to the Lagrangian density ($``$) by,
$$L=d^3x.$$
(52)
Usually the scalar field is represented by a continuous function, $`\varphi (x,t)`$, which can be real or complex. Given a potential density of the field, $`V(\varphi )`$, $``$ takes the form,
$$=\frac{1}{2}_\mu \varphi ^\mu \varphi V(\varphi ).$$
(53)
The resulting Euler-Lagrange equations result from varying the action with respect to spacetime (ryder, ),
$$\frac{}{\varphi }\frac{d}{dx^\mu }\left(\frac{}{(_\mu \varphi )}\right)=0,$$
(54)
where $`x^\mu =(t,x^i)`$ as usual, and $`\mathrm{}=c=1`$. Also note, $`g_{\mu \nu }=\text{diag}(1,1,1,1)`$, the factor of $`\sqrt{g}`$ that usually appears in the action and other equations involving tensor densities will be $`\sqrt{(1)}=1`$ (Minkowski space). The resulting equation is
$$\ddot{\varphi }+3H\dot{\varphi }=V^{}(\varphi ).$$
(55)
The prime represents differentiation with respect to $`\varphi `$ and the term containing the Hubble constant serves as a kind of friction term resulting from the expansion. The field is taken to be homogeneous, which eliminates any gradient contributions. This homogeneity is a safe assumption, since physical gradients are related to comoving gradients by the scale factor,
$$_{physical}=a^1(t)_{comoving}.$$
(56)
Thus, the inhomogeneities in the field are redshifted away during inflation since the scale factor increases by a large amount.
One can also define the stress-energy tensor by use of Noether’s theorem (ryder, ),
$$T^{\mu \nu }=^\mu \varphi ^\nu \varphi g^{\mu \nu }.$$
(57)
This is useful, because it can be compared to $`T^{\mu \nu }`$ for a perfect fluid, namely,
$$T^{\mu \nu }=\text{diag}(\rho ,p,p,p).$$
Using (53) in (57) yields,
$$\rho =T^{00}=\frac{1}{2}\dot{\varphi }^2+V(\varphi )+\frac{1}{2}(\varphi )^2.$$
The calculation for the pressure is a bit more subtle,
$$p=(T^{11}+T^{22}+T^{33})/3.$$
Consider the first component of pressure, again making use of (53) and (57),
$$T^{11}=^1\varphi ^1\varphi g^{11}\left[\frac{1}{2}_\beta \varphi ^\beta \varphi V(\varphi )\right].$$
Since $`g^{11}=1`$ and one can use the metric to raise and lower indices,
$$T^{11}=(^1\varphi )^2V(\varphi )+\left[\frac{1}{2}g^{\beta \gamma }_\beta \varphi _\gamma \varphi \right].$$
Since the metric is diagonal this yields,
$$T^{11}=(^1\varphi )^2V(\varphi )+\frac{1}{2}\left[g^{00}_0\varphi _0\varphi +g^{11}_1\varphi _1\varphi +g^{22}_2\varphi _2\varphi +g^{33}_3\varphi _3\varphi \right],$$
$$=(^1\varphi )^2V(\varphi )+\frac{1}{2}\dot{\varphi }^2\frac{1}{2}(\varphi )^2.$$
Similarly, the $`T^{22}`$ and $`T^{33}`$ components may be found. So for the total pressure one finds,
$$p=\frac{1}{3}\underset{i}{}T^{ii},$$
or,
$$p=\frac{1}{2}\dot{\varphi }^2V(\varphi )\frac{1}{6}(\varphi )^2.$$
(58)
From the $`T^{00}`$ component we already found,
$$\rho =\frac{1}{2}\dot{\varphi }^2+V(\varphi )+\frac{1}{2}(\varphi )^2.$$
(59)
Equations (58),(59) for the pressure and the energy density, show that the equation, $`p=\rho `$ is not quite satisfied. A first resolution to this problem is to again assume that the scalar field ($`\varphi `$) is spatially homogeneous, allowing one to eliminate the gradient terms ($`\varphi `$). This assumption is only made at this point to simplify the analysis. As we have seen if one keeps the terms, it can be shown that the gradients are redshifted away by the expansion (56). Ignoring gradients, the equations become
$$p=\frac{1}{2}\dot{\varphi }^2V(\varphi ),$$
(60)
$$\rho =\frac{1}{2}\dot{\varphi }^2+V(\varphi ).$$
(61)
The first term $`\frac{1}{2}\dot{\varphi }^2`$ can be thought of as the kinetic energy and the second as the potential, or configuration energy. It is now possible to explicitly see where Equation (55) came from if we assume the field can be described as a ‘perfect’ fluid. This assumption allows us to use a continuity equation,
$$\dot{\rho }+3H(\rho +p)=0.$$
(62)
By plugging in the energy density of the field (61) and making use of the Friedmann equation (to get $`H`$) one obtains Equation (55) in perhaps a more enlightening way.
From the pressure and energy density derived above, we see that the requirement that $`p=\rho `$ can be approximately met, if one requires $`\dot{\varphi }V(\varphi )`$. This leads to what is called the slow-roll approximation (SRA), which provides a natural condition for inflation to occur<sup>26</sup><sup>26</sup>26In much of the literature on inflation, the slow-roll approximation is presented as a necessary and sufficient condition for inflation. However, in many new models of inflation this is not necessary. For a treatment of these models, see (guthp1, ),(guthp2, ),(lindep1, , and references therein).. To assure the constraint on $`\dot{\varphi }`$, one must also require that $`\ddot{\varphi }`$ be negligible. Given these requirements, we will to define the slow-roll parameters and introduce the Planck mass<sup>27</sup><sup>27</sup>27The Planck mass is easier to work with opposed to Newton’s constant $`G`$, since most of the interesting energy scales are on the order of GeV ($`1eV=1.6\times 10^{19}\text{Joules}`$). In these units, the Planck mass is $`10^{19}`$ GeV. (liddle3, ),
$$ϵ(\varphi )=\frac{M_p^2}{16\pi }\left(\frac{V^{}}{V}\right)^2,$$
(63)
$$\eta (\varphi )=\frac{M_p^2}{8\pi }\left(\frac{V^{\prime \prime }}{V}\right).$$
(64)
At this point, it is useful to distinguish $`\varphi `$ as the inflaton field. Inflaton is the name given to $`\varphi `$, since its origin does not have to originate with a specified particle theory. Although the original hope was that $`\varphi `$ would help determine the correct particle physics models, current model building does not necessarily require specific particle phenomenology. This is actually an advantage for inflation, it retains its power to solve the initial value problems, yet it could arise from any arbitrary source (i.e., any arbitrary inflaton).
However, observation of the large-scale structure of the universe and anisotropies in the cosmic background should be able to constrain the inflaton parameters to a particular region. This can then be used by particle theorists, as a motivation for some required scalar field. Observational aspects of inflation will be considered in the next section, but this property of the inflaton field manifests itself as one of the greatest contributions of cosmology commensurate with particle theory. Some examples of potentials that have been proposed for the inflaton are presented below (veneziano, ),(liddleeprint, ).
$`V(\varphi )=\lambda (\varphi ^2M^2)^2`$ Higgs potential (65)
$`V(\varphi )={\displaystyle \frac{1}{2}}m^2\varphi ^2`$ Massive scalar field (66)
$`V(\varphi )=\lambda \varphi ^4`$ Self-interacting scalar field (67)
$`V(\varphi )=2H_i^2\left(3{\displaystyle \frac{1}{s}}\right)e^{\varphi /\sqrt{s}}`$ Dilaton scalar field (string theory) (68)
In Guth’s original inflation scenario (guth, ), the inflaton field ($`\varphi `$) sits at a local minimum and is trapped in a false vacuum state (see Figure 7). The vacuum state of a field or particle is the lowest energy state available to the system. Some examples are the ground state of the hydrogen atom (-13.6 eV) and the ground state of the harmonic oscillator ($`\frac{1}{2}\mathrm{}\omega `$). The concept of ‘false’ vacuum comes from examination of Figure 7. If $`\varphi `$ is ‘placed’ in the potential well on the left, the lowest energy available is that of the false vacuum. The only way $`\varphi `$ can get out of this local minimum is by quantum tunneling, after some characteristic time. As tunneling takes place the universe inflates. Inflation halts when $`\varphi `$ reaches the false vacuum and bubbles of the false vacuum coalesce releasing the ‘latent’ heat that was stored in the field. This is much like the way bubble nucleation occurs when opening a bottle of compressed liquid (like soda). Energy escapes from the soda in the form of carbon dioxide and the liquid enters a lower more favorable energy state.
Tunneling that leads to bubble nucleation is a first order phase transition. This is very similar to processes that take place in the study of condensed matter physics, fluid dynamics, and ferromagnetism (see for example (maris, ) and (kittel, )). The bubbles experience a state of negative pressure. Once created, they continue to expand at an exponential rate. Each expanding bubble corresponds to an expanding domain. However, when one carefully investigates this situation, one finds that the bubbles can collide as they reach the false vacuum. Furthermore, the size of these bubbles expands at too great of a rate and the corresponding universe is left void of structure. One finds that too much inflation occurs and the visible universe is left empty. This is referred to as the ‘graceful exit’ problem. Again one is presented with an empty universe, which was the same reason that the DeSitter universe idea was abandoned.
Guth and others further tried to remedy these problems by fine-tuning the bubble formation. The problem with this is two fold. One, cosmologists and particle theorists don’t like fine-tuning. The idea is to form a model that gives our universe as a usual result that follows from natural consequences. By natural one means that the scales of the model are related to the fundamental constants of nature; e.g., quantum gravity should occur at the Planck scale, since this scale is the only one natural in units (c,$`\mathrm{}`$,G). Secondly, if the model is fine-tuned to agree with the observations of the anisotropies in the cosmic background, the bubbles would collide far too often. This results in the appearance of topological defects, like the monopoles. However, this was the whole reason inflation was invoked in the first place.
In 1982, a solution to the graceful exit problem was proposed by Linde (linde6, ) and independently by Steinhardt and Albrecht (steinhardt7, ). This New Inflation model solves the graceful exit problem by assuming the inflaton field evolves very slowly from its initial state, while undergoing a phase transition of second order. Figure 8 illustrates this by again considering the evolution of $`\varphi `$. If $`\varphi `$ ‘rolls’ down the potential at a slow rate, one obtains the amount of inflation needed to solve the initial value problems. After the universe cools to a critical temperature, $`T_c`$, $`\varphi `$ can proceed to its ‘true’ vacuum state energy. The transition of the potential is a second order phase transition, so this model does not require tunneling (300years, ). This type of transition is similar to the transitions that occur in ferromagnetic systems (kittel, ).
The majority of current models rely on another concept coined by Linde as Chaotic Inflation (lindecor3, ). This model differs from Old and New Inflation in that no phase transitions occur. In this scenario the inflaton is displaced from its true vacuum state by some arbitrary mechanism, perhaps quantum or thermal fluctuations. Given this initial state, the inflaton slowly rolls down the potential returning to the true vacuum (see Figure 9). This model has the advantage that no fine-tuning of critical temperature is required. This model presents a scenario, which can be fulfilled by a number of different models.
After the displacement of the inflaton, the universe undergoes inflation as the inflaton rolls back down the potential. Once the inflaton returns to its vacuum (true) state, the universe is reheated by the inflaton coupling to other matter fields. After reheating, the evolution of the universe proceeds in agreement with the Standard Big Bang model.
Although the inflaton could in principle be displaced by a very large amount, all the inflationist need be concerned with is the last moments of the evolution. This is when the perturbations in the scalar field are created that eventually lead to large-scale structure and anisotropies in the cosmic background. As long as the inflaton is displaced by a minimal amount (minimal to be defined in a moment) the initial value problems will be solved. When considering quantum fluctuations resulting in the displacement of $`\varphi `$, minimal displacement is easily achieved.
Successful evolution is only possible if $`V(\varphi )`$ is very flat and has minimal curvature. In terms of (63) and (64) this suggests that inflation will occur as long as the SRA requirements hold.
$$ϵ1,|\eta |1.$$
(69)
This method is successfully used in a number of inflationary models that make predictions in accordance with observation. It must be stated again that Chaotic Inflation results in a very general theory. The inflaton field originally proposed by Guth’s model was that of a grand unified theory, but within the Chaotic Inflationary scenario any inflaton field can be used that satisfies the SRA. With these general requirements, potentials used in supergravity, superstrings, and supersymmetry theories can be used to motivate inflation.
Using the energy density obtained in equation (61) one can restate the Friedmann equation (45) in terms of the scalar field.
$$H^2=\frac{8\pi G}{3}\left[V(\varphi )+\frac{1}{2}\dot{\varphi }^2\right],$$
(70)
Also, the equations of motion derived in (55) are restated here for convenience.
$$\ddot{\varphi }+3H\dot{\varphi }=V^{}(\varphi ).$$
(71)
Using the SRA, one can simplify these equations by eliminating the $`\dot{\varphi }^2`$ and $`\ddot{\varphi }`$ terms. This leaves the more tractable equations shown below, which remain valid until $`\varphi `$ approaches the true vacuum; i.e., $`ϵ1`$.
$$H^2=\frac{8\pi V(\varphi )}{3M_p^2},$$
(72)
$$3H\dot{\varphi }=V^{}(\varphi ),$$
(73)
where $`M_p`$ is the Planck mass and has been substituted for $`G`$,
$$M_p1/\sqrt{G},$$
remembering that $`\mathrm{}=c=1`$.
One can use equations (72) and (73) to manifest the connection between the slow-roll condition (69) and the generic definition of inflation, that is $`\ddot{a}>0`$. First note that
$$H=\frac{\dot{a}}{a}\dot{H}=\frac{\ddot{a}}{a}\left(\frac{\dot{a}}{a}\right)^2.$$
For inflation to take place means $`\ddot{a}(t)>0`$ and $`a(t)`$ is always positive thus,
$$\frac{\ddot{a}}{a}>0\dot{H}+H^2>0,$$
$$\frac{\dot{H}}{H^2}<1.$$
(74)
Using (72) and differentiating with respect to time,
$$2H\dot{H}=\frac{8\pi }{3M_p^2}\frac{d\varphi }{dt}\frac{d}{d\varphi }\left(V(\varphi )\right),$$
$$\dot{H}=\frac{8\pi \dot{\varphi }V^{}}{6M_p^2H}.$$
Plugging this result into (74) gives,
$$\frac{\dot{H}}{H^2}=\frac{\dot{\varphi }V^{}}{2HV}<1.$$
Solving (73) for $`\dot{\varphi }`$ and plugging the result into the last equation, we obtain
$$\dot{\varphi }=\frac{V^{}}{3H},$$
$$\frac{\dot{H}}{H^2}=\frac{(V^{})^2}{6H^2V}<1.$$
Lastly, substituting $`H^2`$ from (72) one obtains,
$$\frac{\dot{H}}{H^2}=\frac{M_p^2}{16\pi }\left(\frac{V^{}}{V}\right)^2<1.$$
(75)
But this is just the slow-roll condition $`ϵ1`$. So again one is reminded that inflation will take place until $`ϵ1`$, which has been shown to be equivalent to $`\ddot{a}>0`$.
#### V.3.1 Modeling the Inflaton Field
The equations of motion derived above (72),(73) describe the evolution of an arbitrary potential $`V(\varphi )`$ subject only to the constraint that $`V(\varphi )`$ conform to the slow roll conditions away from its minimum. As mentioned previously, the conditions for inflation are arbitrary and inflation will occur as long as $`\ddot{a}>0`$. There are three cases that are of particular interest (peacock, ).
* Polynomial Inflation. $`V(\varphi )\varphi ^n`$. This gives a scale factor which behaves quasi-exponentially.
$$a(t)\mathrm{exp}\left(\varphi ^{n/2}kt\right)\text{where}k=\sqrt{\frac{8\pi }{3M_p^2}}.$$
For particle theorist $`n=2,4`$ are most favorable, since they describe renormalizable particle theories (ryder, ).
* Power-law Inflation. $`V(\varphi )\mathrm{exp}\left(\sqrt{\frac{16\pi }{n}}\frac{\varphi }{M_p}\right)`$. This gives $`a(t)t^n`$. The only requirement being that $`n>1`$.
* Intermediate Inflation. $`V(\varphi )\varphi ^{}`$, where $`=4(n^11)`$. This yields $`a(t)e^{(t/t_0)^n}.`$
As an example consider a simple polynomial model with $`n=2`$.
$$V(\varphi )=\frac{1}{2}m^2\varphi ^2,$$
$$V^{}(\varphi )=m^2\varphi ,$$
$$V^{\prime \prime }(\varphi )=m^2.$$
The slow roll condition $`ϵ<1`$ implies
$$ϵ=\frac{M_p^2}{16\pi }\left(\frac{V^{}}{V}\right)^2=\frac{M_p^2}{16\pi }\left(\frac{m^2\varphi }{\frac{1}{2}m^2\varphi ^2}\right)^2<1,$$
$$\varphi >M_p.$$
In other words, the inflaton field must be larger than the Planck energy. This is actually the value one would expect for the cosmological constant, since the only natural scale at high energy is the Planck scale. However, inflation (as it has been presented) relies on the evolution of a classical field. If the value of the field is in the quantum regime, where the Planck scale lives, inflation can not be treated classically. Fortunately, there are many resolutions to this problem. First, what really matters in the field equations is the potential energy density of the field; namely,
$$V(\varphi )=\frac{1}{2}m^2\varphi ^2.$$
Remember that $`V`$ appeared in the Lagrangian density, thus $`V`$ is actually a density. From this equation one can see that the magnitude of the potential, and, therefore, the energy scale of the theory, depend on ‘$`m`$’ as well as $`\varphi `$. ‘$`m`$’ represents the mass of the inflaton and one resolution to the high value of $`\varphi `$ is to introduce a small mass for the inflaton; i.e., make ‘$`m`$’ small.
Another resolution to this scale problem is to limit considerations to the final part of the evolution of the inflaton. As stated before, the inflaton evolves until $`ϵ=1`$ when the inflaton then reheats the universe. The most important consequences of inflation are its resolution of the initial value problems and its predictions about large-scale structure and the cosmic background anisotropies. As we will see most of these phenomena only require analysis of the last moments of ‘e-foldings’ of inflation. E-folding is a way of measuring how the scale factor increases. Since,
$$a=a_0\mathrm{exp}(H\mathrm{\Delta }t),$$
one e-folding is defined as the amount of time for $`a`$ to grow by a factor of $`e`$:
$$\mathrm{\Delta }t=H^1.$$
It will be shown that only 60 e-foldings are needed to resolve the initial value problems and the scales that are important in determining structure in the universe only depend on modes that are present during these last 60 e-foldings.
However, it must be pointed out that many people find problems with these conditions on $`\varphi `$. The idea of fine-tuning the mass of the inflaton ‘$`m`$’ is certainly unappealing. One of the appealing aspects of inflation was its resolution of the initial value or fine-tuning problems of the Big Bang model. But, now we are again confronted with an initial value problem. This problem can be resolved by addressing inflation in the context of a quantum gravity theory such as superstring theory. I will not pursue such issues in this paper, although the reader is again referred to the eprint archive for recent efforts<sup>28</sup><sup>28</sup>28http://xxx.lanl.gov.
#### V.3.2 The Amount of Inflation
One can find the amount of inflation by considering the change of the scale factor. Considering the example of quasi-exponential expansion, meaning that the Hubble constant need not be constant. Then,
$`a(t)=a_0\mathrm{exp}(Ht)`$ (76)
$`\mathrm{ln}\left({\displaystyle \frac{a(t)}{a_0}}\right)=Ht,`$
$`𝒩\mathrm{ln}\left({\displaystyle \frac{a(t_{final})}{a(t_{initial})}}\right)={\displaystyle _{t_i}^{t_f}}H𝑑t`$ Number of e-foldings
Using the slow-roll equations, the number of e-foldings can be expressed in terms of the inflaton potential. Dividing (73) by (72) yields,
$`{\displaystyle \frac{3\dot{\varphi }}{H}}={\displaystyle \frac{3H\dot{\varphi }}{H^2}}={\displaystyle \frac{V^{}}{8\pi V}}\left(3M_p^2\right),`$
$`{\displaystyle \frac{\dot{\varphi }}{H}}={\displaystyle \frac{M_p^2}{8\pi }}\left({\displaystyle \frac{V^{}}{V}}\right)`$
Using this result, with the formula for $`𝒩`$ (76), one gets,
$$𝒩=_{t_i}^{t_f}H𝑑t=_{t_i}^{t_f}H\frac{dt}{d\varphi }𝑑\varphi =_{t_i}^{t_f}\frac{H}{\dot{\varphi }}𝑑\varphi =\frac{8\pi }{M_p^2}_{\varphi _i}^{\varphi _f}\frac{V}{V^{}}𝑑\varphi .$$
$$𝒩\frac{8\pi }{M_p^2}_{\varphi _f}^{\varphi _i}\frac{V}{V^{}}𝑑\varphi .$$
(78)
Here the fact that the SRA has been used is expressed using ‘$``$’ in (78).
For $`𝒩>60`$, which is needed to solve the initial value problems (peacock, ), we again find $`\varphi M_p`$. This can be seen from (78), where $`V^{}V/\varphi `$ using the SRA. This means that if one chooses a potential of $`\frac{1}{2}m^2\varphi ^2`$, one must choose the coupling, $`m^2`$ to be small. Given a self interacting potential term, $`\lambda \varphi ^4`$, the coupling must be extremely weak, $`\lambda 1`$. This coupling agrees nicely with theories of supergravity and certain string theories, although other potentials are ruled out because of their couplings, such as the weak coupling. This leaves the question. Can inflation be considered without a theory of quantum gravity? As mentioned before, the inflationist is often not concerned with these initial stages of inflation. The common standpoint is that chaotic inflation can present an evolution and then one studies the predictions of this evolution. As aesthetically displeasing as this may be, it allows cosmology to progress further without a quantum theory of gravity. Ultimately this issue will be addressed within the context of a theory of quantum gravity to create a complete picture of the creation of the universe. However, it has been argued that a complete understanding of the universe may be avoided in a scenario known as eternal inflation (Lindecor4, ),(guthp2, , and references within ),(guthp1, , and references within). Time does not permit to discuss these models in detail, however for a popular account (Lindecor5, ) is an excellent starting point.
Given the slow-roll conditions and the number of required e-foldings ($`𝒩`$), one can test a model inflaton to see if it is compatible with an inflationary scenario. With this generic framework that has been set forth, one can construct particle theories and then test their validity within the context of inflation theory. However, the slow-roll approximation and initial value problems ($`𝒩>60`$) are not the only constraints on the inflaton and therefore particle theory. The inflaton is further restricted by the predicted large-scale structure of the universe, along with the mechanisms involved with reheating of the universe at the end of inflation. The large-scale structure is determined by density perturbations resulting from quantum fluctuations in the evolution of the inflaton field. This analysis can actually be done without the advent of quantum gravity; however, it is outside the scope of this paper. Instead, the author hopes to manifest the stringency of these parameters on the inflaton field by addressing the observational consequences and predictions that inflation offers. In the next section these observational tests will be explored.
## VI Observational Tests of Inflation
It was shown in the last section that if the number of e-foldings exceeds $`60`$, then inflation can solve the horizon, flatness, and relic (monopole) problems. One generally favors this model over a Big-Bang model, because of its naturalness. That is to say, inflation offers a generic scenario for solving the initial value problems. However, this arbitrariness can also be viewed as a problem for inflation. For instance, throughout this paper it has been assumed that inflation necessarily leads to a flat universe ($`k=0`$). However, Hawking, Turok, Linde, and others have shown that inflation can result in a non-flat universe (open1, ),(open2, ). Models can be created that produce unwanted or wanted relics and contain inhomogeneities. Furthermore, we have seen inflation requires an inflaton field to drive the inflation. Where does it come from and what is its natural value? Originally Guth had the inflaton as corresponding to a GUT transition; however, today the preferred energy range is on the order of the Planck scale. Thus, to fully understand inflation one needs a full quantum theory of gravity. For these reasons one may ask; Is inflation a particular type of cosmological model, or is inflation an arbitrary constituent of any successful theory of the cosmos?
Inflation’s strength today can be seen from its predictions of large-scale structure. Different models predict different structure and this can be used to narrow the number of possible models. One can further constrain the inflationary models by cosmological parameters. The cosmic background observations, galaxy surveys, lensing experiments, and standard candles can be used pin-down the cosmological parameters. In this way, observational parameters (e.g., $`H`$, $`\mathrm{\Omega }_M`$,$`\mathrm{\Omega }_\mathrm{\Lambda }`$) can be given viable ranges and inflationary parameters can be determined based on these preferred ranges. With the cosmological parameters determined, inflation parameters depend only on the height and shape of the inflaton potential. The inflaton potential correspondingly yields predictions about the large-scale structure of the universe and the anisotropies in the cosmic background radiation.
The study of large-scale structure has been pursued for many years (zeldovich, ). The problem was that there was an appealing mechanism which could produce the types of perturbations needed to produce the observed large-scale structure. These perturbation types are manifested by the anisotropies in the cosmic background. The anisotropies result from acoustic oscillations in the baryon-photon fluid just before recombination. Therefore, the anisotropy spectrum offers a ‘snap-shot’ of the seeded inhomogeneities that eventually resulted in galactic structure. These inhomogeneities were first discovered when COBE mapped the cosmic background in the early 1990’s (smoot, ). In this intimate way, the cosmic background and galaxy surveys predict a scheme by which structure was formed. The type of perturbation that is needed only results from models which predict Gaussian, adiabatic, nearly scale-invariant perturbations (carroll42, ). The only known models that fall into this category are the inflationary models (guth, ),(linde6, ),(steinhardt7, ).
### VI.1 Perturbations and Large-Scale Structure
During the inflation epoch, perturbations (small fluctuations) of two types are generated: scalar (density) perturbations, and tensor (metric) perturbations. The scalar perturbations come from quantum fluctuations of the inflaton field before and during its evolution. Tensor perturbations arise from quantum fluctuations in the space-time metric within the quasi-DeSitter spacetime. Fluctuations of this type are a natural consequence of a DeSitter spacetime.
In DeSitter space there exists an event horizon. Consider the distance light can travel in comoving coordinates, which is given by (6) and $`a(t)\mathrm{exp}(Ht)`$.
$$r=_0^{\mathrm{}}\frac{cdt}{a(t)}=_0^{\mathrm{}}\mathrm{exp}(Ht)c𝑑t=c/H.$$
(79)
The presence of this event horizon suggests the presence of thermal fluctuations in the fields, similar to those present in a black hole. This can be understood by appealing to the uncertainty principle. The event horizon causes the ground state modes of any fields present to be restricted in spatial extent. The uncertainty principle then requires, $`\mathrm{}p\mathrm{}H/c`$, since the characteristic size is just $`c/H`$. This uncertainty in momentum gives rise to energy fluctuations and the corresponding Hawking temperature is given by (peacock, ),
$$kT_{DeSitter}=\frac{H}{2\pi },$$
(80)
where $`k`$ is the Boltzman constant relating the energy to the temperature.
This result provides a motivation for the existence of fluctuations in the metric and scalar field. As these perturbations are created during inflation they are inflated outside of the causal horizon (particle horizon). As mentioned in the previous section, the causal horizon is nearly stationary during inflation. Once the perturbation has been inflated outside the horizon, its ends are no longer in causal contact. In this way, the perturbations become ‘frozen-in’ as classical perturbations.
Ignoring the nonlinear, or super-horizon, effects of these perturbations may trouble the reader. However, as we shall see, ignoring these effects appears to be in agreement with the cosmic background data (hu2, ), (hu3, ), (hu4, ). As stated, one reason for choosing to ignore super-horizon evolution is the causal separation of the ends of the perturbation. However, this argument is far from rigorous and the study of nonlinear perturbations takes much care. Perturbation evolution relies on the extrinsic (super-horizon) properties of the spacetime manifold and is sensitive to the gauge of general relativity. The perturbation evolution can usually be ignored in regions where the pressure becomes negligible, which happens to be on the order of the horizon (peacock, ). This generally motivates one to ignore the super-horizon evolution, however for a complete treatment, see (branden, ),(ref31, ),(turner, , Chapter 8 and 9).
After inflation the expansion continues at sub-luminal speeds and these perturbations enter back inside the causal horizon. Thus, the most important perturbations for creating structure come from the ones that were exited near the end of the inflationary period (i.e., approximately the last 60 e-foldings). After the reheating process occurs, the inhomogeneities passing back inside the horizon cause fluctuations that seed the large-scale structure.
Pure exponential inflation, which corresponds to a DeSitter spacetime, has an interesting property. The spacetime is invariant under time translation. That is to say, there is no natural origin of time under true exponential expansion (peacock, , Chapter 11). The only fundamental size in the theory is that of the Hubble horizon ($`c/H`$). Thus, one expects that the amplitude of a ‘standing wave’ perturbation will be related to the horizon size, $`c/H`$, which is not changing. Therefore, we see why inflation predicts a scale-invariant spectrum for the perturbations.
This analysis can be illustrated through musical analogy. The fundamental mode of the perturbations are determined by the Hubble distance ($`c/H`$), much like the fundamental mode of a flute is determined by its length. Because the Hubble length (horizon) is nearly stationary during inflation, this means inflation predicts a scale-invariant, or Harrison-Zeldovich spectrum. Furthering this analogy, the ‘overtones‘ of the universe correspond to the inflaton potential that determines its behavior, much like overtones can be used to distinguish one instrument from another. However, as we have seen that inflation need not be exponential. The small deviations from DeSitter spacetime result in small deviations from a scale-invariant spectrum. These deviations can be used to successfully predict the correct potential for the inflaton.
The generated perturbations can be characterized by a power spectrum, $`\delta _H^2(k)`$. The $`H`$ indicates that the perturbation amplitude is taken to correspond to its value when it crossed the causal horizon. Quantitatively, this corresponds to $`k=aH`$, where $`k`$ is the wave number of the perturbation. Quantum field theory can be used to calculate an expression for $`\delta _H`$ similar to Equation (80) above (peacock, , Chapter 8).
$$\delta _H=\frac{H^2}{2\pi \dot{\varphi }},$$
(81)
where $`\varphi `$ is the inflaton. This formula manifests the connection between the inflaton potential and the perturbations generated during inflation.
One is generally interested in the scale dependence of the spectral index of these perturbations, since this dependence changes for different inflation models (lyth21, ).
$$n(k)12\frac{d\delta _H}{d\mathrm{ln}k}.$$
(82)
For an absolute scale-invariant spectrum, it follows that $`\delta `$ is independent of $`k`$ and the above relation gives $`n=1`$ as one would expect. For a spectrum that is nearly scale-invariant the amount $`n`$ differs from one is referred to as the tilt of the spectrum. Although it is not at all obvious, (81) and (82) can be used along with the slow roll parameters (64),(63) to express the tilt as,
$$1n=6ϵ2\eta .$$
(83)
The details of the calculation need not concern us here, for a derivation of this result see (dekel, ) or (peacock, ). This relation is only presented to demonstrate that the tilt, which is a discriminating factor between models, can be written in terms of the SRA parameters $`ϵ`$ and $`\eta `$. Therefore, if an experimental consequence of the tilt is observable, one can find the appropriate values for the SRA and reconstruct the inflaton potential (lyth34, ).
### VI.2 The Cosmic Background Anisotropies
The discovery of the anisotropies in the cosmic background by COBE created a new opportunity for verification of cosmological parameters and theories of large-scale formation. As discussed at the beginning of this paper, the CMB offers a ‘snap-shot’ of the universe at the time of recombination ($`z1000`$). The anisotropies that were present in the baryon-photon plasma at this time are manifested today by the temperature fluctuations in the spectrum. These fluctuations are representative of a nearly Gaussian, scale-invariant spectrum. As discussed previously, this is a unique prediction of inflation theories.
Although the quantitative details can become formible, the qualitative description of these temperature fluctuations is quite simple. During inflation, the perturbations formed must be of nearly the same amplitude and randomly distributed, as discussed above. After reheating takes place, these classical perturbations re-enter the horizon causing density fluctuations in the baryon-photon plasma. In over-dense regions, potential wells form that trap the plasma and cause it to heat up. At the same time, photon pressure induces a kind of restoring force to oppose the gravitational potential. In this way, a harmonic oscillator motion is set up in the plasma. These oscillations continue with no friction (viscosity) from the fluid. This is why they are referred to as adiabatic fluctuations. However, if this were not the case and the friction is deemed important, one obtains an isocurvature spectrum (huhuhu, ). These turn out to be indicative of cosmic strings, which are ruled-out by observation as a method of primordial structure formation. However, models containing cosmic strings that are produced during reheating following inflation may still play a major role in cosmological models (brandenberger, ).
At the time of recombination, when the photons were able to escape the fluid, they had to overcome the gravitational potentials. The picture is that the photons in these potential wells were hotter than the average, but this temperature difference was partially cancelled by the gravitational redshift resulting from the photons ‘climbing’ out of the potential well. This phenomena is know as the Sachs-Wolf effect. The result is that the photons that were in the wells have a slight temperature increase from those that were not. This variation is predicted by theory to be on the order of $`10^5`$ (peacock, ). These oscillations propagate through the fluid at the speed of sound. Thus, there is a acoustic horizon that is generated within the surface of last scattering and if present today would have an angular size of about one degree on the sky.
One concern with the simplicity of this analysis is what effects, such as reionization in the surface of last scattering, must be considered? It is important to consider the mean free path of the photon as it travels within the fluid before escaping. This could affect the energy and therefore temperature of the spectrum. However, it turns out that this effect only appears on small angular scales within the spectrum and can be ignored (hu3, ). Also, any effects from CMB scattering off interstellar gas only appear in the spectrum at very small angles. These observations are of course useful, but offer little insight into examination of the early universe and formation of large-scale structure.
Given the predicted anisotropies from the Sachs-Wolf effect, the next step is to examine the cosmic background spectrum through observation. The anisotropies in the temperature of the cosmic background spectrum can be expanded in spherical harmonics (carroll42, ),
$$\frac{\mathrm{\Delta }T}{T}=\underset{lm}{}a_{lm}Y_{lm}(\theta ,\varphi ),$$
(84)
The multipole coefficients are given by,
$$a_{lm}=𝑑AY_{lm}^{}\frac{\mathrm{\Delta }T}{T}.$$
(85)
The amount of anisotropy at multipole moment $`l`$ is expressed by the power spectrum,
$$C_l=|a_{lm}|^2.$$
(86)
The $`C_l`$’s measure the temperature anisotropy of two regions separated by angle $`\theta `$. This angle is related to the $`l`$’s by: $`\theta 180^^o/l.`$ Thus, $`l`$ allows one to express the temperature variations of regions separated by an angle $`180^^o/l`$. The $`l=0`$ represents the monopole contribution to the anisotropy, which is of course zero (this is comparing a point separated from itself by $`360^^o`$. The next moment, $`l=1`$, is the dipole moment, which compares regions separated by $`180^^o`$. This anisotropy originates from our peculiar velocity, the motion of the Earth relative to the cosmic background. This moment is usually taken out of the spectrum to leave the ‘true’ anisotropy. The $`l=2`$ moment is the quadrupole contribution, which marks the first non-trivial anisotropy for understanding structure formation.
When the COBE data is plotted with multipole $`l`$ versus temperature variation, a peak is found to occur around $`l=220`$ or $`1^^o`$ (Figure (10)). This peak corresponds to the angular size on the sky of the acoustic horizon discussed before and has been called the Doppler peak. Thus, there is a maximum temperature variation at precisely the angle predicted by the Sachs-Wolf effect. Since a mechanism of this type can only be explained by an inflationary model, one is presented with a strong argument for inflation (huhuhu, ).
However, the peak is actually sensitive to the cosmological parameters, such as the Hubble constant and the curvature of the universe. If the universe is non-flat then the null geodesics are found to converge (diverge) in the case of spherical (hyperbolic) geometry. The angle subtended is given by $`\theta _H\mathrm{\Omega }^{1/2}\mathrm{\hspace{0.25em}1}^^o`$. This means for the peak at $`l=220`$ we live in a universe which is flat. Although, this seems to represent a bit of circular logic. This difficulty can be remedied by calling upon other observational tests to constrain the parameters. These includ galaxy surveys, lensing experiments, or standard candle observations. When all of these methods are combined, strong constraints can be put on parameters and the best model can be determined.
The cosmic background is apparently richer in structure than was first realized. As we have discussed, inflation predicts fluctuations in both the scalar and tensor fields. This gives rise to slight differences in the anisotropies at small angles (large $`l`$). Also, there are multiple peaks in the spectrum following the peak at $`l=220`$ and the height and shape of the spectrum are related to the specifics of inflation models, such as the tilt.
To examine one aspect of the complexity involved, consider that the scalar and tensor perturbations can be fixed by their contribution to the quadrupole moment $`C_{l=2}`$ of the CMB (kamionkowski, ).
$$𝒮=6C_{l=2}^{scalar}=33.2[V^3/(V^{})^2],$$
(87)
$$𝒯=6C_{l=2}^{tensor}=9.2V,$$
(88)
where $`V=V(\varphi )`$ is the inflaton potential, $`𝒮`$ is the scalar contribution, and $`𝒯`$ is the tensor contribution. When the slow-roll approximation is considered and the determination of cosmological parameters is found by methods other than CMB analysis; e.g., for large-scale structure probing, one finds that the ratio $`𝒯/𝒮`$ is less than order unity. This restricts $`V5\times 10^{12}`$. Reformulating equation (82) in terms of the inflaton potential, one finds in Planckian units the relations for the scalar and tensor spectral indices, respectively.
$$1n_s=\frac{1}{8\pi }\left(\frac{V^{}}{V}\right)^2\frac{1}{4\pi },$$
(89)
$$n_t=\frac{1}{8\pi }\left(\frac{V^{}}{V}\right)^2.$$
(90)
Although most models of inflation predict a near scale-invariant, or Harrison-Zel’dovich spectrum, the small deviations from differing potentials $`V(\varphi )`$ give a way of testing inflation models. With future experiments such as MAP<sup>29</sup><sup>29</sup>29http://map.gsfc.nasa.gov and PLANCK<sup>30</sup><sup>30</sup>30http://astro.estec.esa.nl/SA-general/Projects/Planck, these spectral indices will be found with great precision and the shape and height of the spectrum can be used to manifest the correct inflaton potential. In this way, inflation will be used to predict new particle physics, instead of the original scenario which was vice versa.
However, the ultimate test of inflation is the precise determination of the tensor perturbations, that is the $`n_t`$’s. This can’t be deduced from the CMB spectrum because both the scalar and tensor perturbations contribute to the temperature anisotropy. However, if a method could be devised to separate out the tensor perturbations and this spectrum were detected, it would be concrete evidence of an inflationary period. This is because inflation is the only way metric perturbations can survive to the present. This is due to the structure of a DeSitter spacetime.
The tensor spectrum can be separated by creating a polarization map of the CMB. The separation is then possible because the tensor perturbations have an intrinsic axial component, and the angular dependence can be determined from the metric. Whereas, the scalar perturbations have no dependence on direction. Therefore, the polarization vector can be constructed out of two parts, a curl and a gradient.
$$\stackrel{}{P}(\theta ,\varphi )=\stackrel{}{}A+\stackrel{}{}\times \stackrel{}{B},$$
(91)
where $`\widehat{n}`$ gives the direction. Thus, to obtain the tensor terms one can take the divergence of this vector. Before proceeding further, it may be of interest to the reader why the perturbations only contain tensor and scalar contributions and vector type perturbations are absent. This is because massless vector fields are conformally invariant (turnertests, ). This invariance can be broken by introducing a mass term or by explicitly breaking the coupling. Although for a massive vector field the perturbations generated are far too small and die off far too quickly during the inflationary expansion (liddletests, ). Thus, a standard prediction of inflation is that there are no vector perturbations. But, isn’t the polarization a vector? No. The polarization is actually a $`2\times 2`$ trace free symmetric tensor. This tensor is written in terms of the Stokes parameters $`Q(\widehat{n})`$ and $`U(\widehat{n})`$, which give us the polarization in each direction. For an explanation of these parameters see, (jackson, , section 7.2).
$$𝒫_{ab}(\widehat{n})=\frac{1}{2}\left(\begin{array}{cc}Q(\widehat{n})& U(\widehat{n})\mathrm{sin}\theta \\ U(\widehat{n})\mathrm{sin}\theta & Q(\widehat{n})\mathrm{sin}^2\theta \end{array}\right).$$
(92)
The polarization tensor can be expanded in tensor spherical harmonics (ref14, ),
$$\frac{𝒫_{ab}(\widehat{n})}{T_0}=\underset{lm}{}\left[a_{(lm)}^GY_{(lm)ab}^G(\widehat{n})+a_{(lm)}^CY_{(lm)ab}^C(\widehat{n})\right]$$
(93)
The $`Y_{(lm)ab}^G`$ and $`Y_{(lm)ab}^C`$ represent a basis for the gradient (scalar) and curl (tensor) perturbation terms in this polarization mapping. The $`a_{(lm)}^G`$’s and $`a_{(lm)}^C`$’s are again just found by exploiting the orthogonality. One can use this spectrum to construct necessary requirements for the potentials of the inflaton field and this analysis can therefore be used to distinguish the various inflation models.
### VI.3 Summary
Scalar perturbations are the most easily detected form of perturbation. The scalar nature of these fluctuations arise from the fact that the perturbations are of mass fields in the primordial era, that is, before the time of decoupling. Different models of inflation and the corresponding reheating mechanisms differ in their predictions of mass variation. Regardless, these variations of the early universe eventually give rise to the structural formation of galaxy clusters, which can help distinguish the various theories of inflation and structure formation. Tensor perturbations are also detectable and these perturbations result from fluctuations in the space-time metric in the primordial era. Again, these perturbations are very small and would be very hard to detect. However, certain models of inflation predict wavelengths that could be detected by laser interferometry gravitational wave detectors, such as LIGO<sup>31</sup><sup>31</sup>31http://www.ligo.caltech.edu/. If these waves were detected it would help eliminate many inflation models and help narrow the region of viable theories. Furthermore, inflation is the only theory that can currently account for a gravitational wave spectrum. The detection of the spectrum would be a great success for the inflation theory.
Both of these types of perturbations contribute to the $`10^5`$ temperature fluctuation in the cosmic background. The biggest challenge for experimentalists is to separate the scalar and tensor contributions to the temperature fluctuations. In practice this is very difficult, if not impossible, and it becomes more practical to consider the polarization of the CBR.
For the inflationist, the goal of CRB measurements is to distinguish between the various models of inflation. A good way to begin, is to express many of the relations obtained thus far, in terms of $`ϵ`$ and $`\eta `$. The number of e-foldings ($`𝒩`$) can be expressed in terms of $`ϵ`$ using (63) and (78),
$$𝒩=\frac{2\sqrt{\pi }}{M_p}_{\varphi _i}^{\varphi _f}\frac{d\varphi }{\sqrt{ϵ}}.$$
(94)
Another useful relation, which may be found in the literature (ref31, ), gives a measure of when a given perturbations of wave length $`k`$ passes through the horizon and is therefore ‘frozen out’. This can be expressed as the number of e-foldings $`𝒩(k)`$ from the end of inflation.
$$𝒩(k)=62\mathrm{ln}\frac{k}{a_0H_0}\mathrm{ln}\frac{10^{16}\text{GeV}}{V_k^{1/4}}+\mathrm{ln}\frac{V_k^{1/4}}{V_e^{1/4}}\frac{1}{3}\mathrm{ln}\frac{V_e^{1/4}}{\rho _{RH}^{1/4}}.$$
(95)
$`V_k`$ is the potential when the mode $`k`$ leaves the horizon, $`V_e`$ is the potential at the end of inflation, and $`\rho _{RH}`$ is the energy density after reheating. This expression may appear formidable, however it can be used to begin understanding density fluctuations. For example, the modes $`k`$ entering the horizon today, left the horizon at $`𝒩(k)=5070`$. The uncertainty in this range manifests the lack of knowledge of the inflaton potential. Thus, once again different inflaton models make different predictions.
With the rapid advances in observational cosmology, cosmologists are able to use the abundance of data that is being obtained by the Hubble Space Telescope, balloon experiments, satellites (such as Chandra), etc. to narrow the parameters of the universe. Then with these values and the relations that have been presented in this section, one can use inflation to predict new physics for the pre-inflation or Planckian epoch. Ultimately this physics will need a quantum theory of gravity or supersting theory, but determination of the inflaton potential and the resulting large-scale structure will set stringent limits in which to test the predictions of these new theories. In this way, inflation offers the link between the innerspace of the quantum realm and the outerspace of the large-scale structure of the universe. The marvelous universe in which we live, the beauty that surrounds us, and even ourselves, will be the result of a quantum fluctuation or perhaps a chaotic mishap.
## Appendix A Acronyms
To hopefully ease the burden of dealing with the ubiquitous acronyms of cosmology, I have compiled a list of the most common to be encountered in this paper.
* CBR – Cosmic Background Radiation
* CMBR/CMB – Cosmic Microwave Background Radiation
* COBE – Cosmic Background Explorer
* CP – Cosmological Principle
* CRF – Cosmic Rest Frame
* eV,MeV – electron Volt, Mega-electron Volt
* MAP – Microwave Anisotropy Probe
* QSO’s – Quasi Stellar Remnants
* RWM – Robertson Walker Metric
* SRA – Slow Roll Approximation
* SR – Special Relativity
* VEV – Vacuum Expectation Value
## Appendix B Figures and Tables
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# Domain Wall Fermions with Exact Chiral Symmetry
## I Introduction
Recently a great deal of theoretical progress has been made in the construction of lattice regularizations of fermions with good chiral properties . For use in practical numerical simulations, though, approximations to these formulations are necessary. In the formulation using domain wall fermions (DWF) , the extent of an auxiliary fifth dimension has to be kept finite in numerical simulations while chiral symmetry holds strictly only in the limit of infinite fifth dimension. The violations of chiral symmetry are expected to be suppressed exponentially in the extent of the fifth dimension , but in practice the coefficient in the exponent can be quite small and the suppression correspondingly slow.
In the case of overlap fermions there is no such problem in principle. However, there is a problem of practicality: how to deal efficiently with $`ϵ(H)=H/\sqrt{H^2}`$, where $`H`$ is some auxiliary Hermitian lattice Dirac operator for large negative mass, but free of doublers. Most commonly, the Hermitian Wilson-Dirac operator $`H_w`$ is used. So far, the best methods use rational polynomial approximations of $`ϵ(H)`$ rewritten as a sum over poles
$`ϵ(H)={\displaystyle \frac{HP(H^2)}{Q(H^2)}}=H\left(c_0+{\displaystyle \underset{k=1}{\overset{n}{}}}{\displaystyle \frac{c_k}{H^2+b_k}}\right).`$ (1)
Variants of this are Neuberger’s polar decomposition and the optimal rational polynomial approximation of Ref. . In the application of $`ϵ(H)`$ on a vector $`\psi `$ the shifted inversions $`(H^2+b_k)^1\psi `$ are done simultaneously with a multi-shift CG inverter . This is referred to as the “inner CG”, since typically another, “outer”, CG method is used to compute propagators or eigenvalues of the overlap Dirac operator containing $`ϵ(H)`$.
In these methods, the main problem is the difficulty of getting an accurate approximation to $`ϵ(x)`$ for small $`x`$<sup>*</sup><sup>*</sup>*Using the scaling invariance $`ϵ(sH)=ϵ(H)`$ for any positive scale factor $`s`$ the upper range of the spectrum of $`H`$ can always be scaled to be in the range of good accuracy of the approximation to $`ϵ(x)`$.. As emphasized in Ref. accuracy for the lower range of the spectrum of $`H`$ can always be enforced by calculating a sufficient number of eigenvectors of $`H`$ with small eigenvalues and computing $`ϵ(H)`$ in the space spanned by these eigenvectors exactly. The approximation is then only used after projecting out these eigenvectors. This is referred to as “projection”.
Small eigenvalues of $`H_w`$ occur quite frequently on (quenched) lattices used in current simulations. Indeed, a non-zero density of eigenvalues near zero has been found . Projection is thus essential in efficient implementations of the overlap Dirac operator. It is known that the zero eigenvalues of $`H_w`$ correspond to unit eigenvalues of the transfer matrix $`T`$ that describes propagation along the fifth direction for domain wall fermions . The left and right handed physical fermions in the domain wall formulation occur along the two boundaries in the fifth direction. A unit eigenvalue of the transfer matrix then allows for unsuppressed interaction between left and right handed fermions and thus to a breaking of chiral symmetry. Therefore the same modes that make the approximation of $`ϵ(H)`$ difficult to achieve in overlap fermions, and need to be projected, are responsible for the chiral symmetry breaking for domain wall fermions at finite $`L_s`$.
If we could project out the modes with near unit eigenvalue of the transfer matrix and take their contribution with effectively infinite $`L_s`$ we could obtain a domain wall fermion action with no chiral symmetry breaking even at finite $`L_s`$. This would make domain wall fermions with finite $`L_s`$ equivalent to overlap fermions. It would then become a matter of computational efficiency as to which form is to be preferred for numerical simulations. Based on ideas of Boriçi we show in section II how projection for domain wall fermions can be achieved. In section III we illustrate how projection works, first on smooth instanton configurations and then on a configuration from a quenched simulation. We discuss the degree of chiral symmetry conservation that can be achieved and compare the performance (cost) of domain wall fermion implementations with and without projection to overlap fermions. We conclude the paper with a brief summary and some discussions in section IV.
## II Projected Domain Wall Fermion Actions
Kaplan’s method realizes a single massless fermion field through an infinite tower (or infinite number of flavors) of massive fermion fields with a particular flavor structure. The flavor index can be interpreted as an extra (here fifth) dimension, and the flavor structure as a defect along the fifth dimension to which the massless fermions are bound. Lattice calculations necessarily require a finite fifth dimensional extent and lattice spacing, as well as regularizations for the derivatives in the four dimensional space and in the fifth direction. These regularizations are not unique. In particular, higher order terms in the fifth dimensional lattice spacing can be added to the action, as long as they do not change the “defect structure”. Such changes can affect the effective four dimensional theory by altering the discretization effects in the four dimensional lattice spacing, but not the continuum limit.
The finite fifth dimensional extent induces some chiral symmetry violation. This can be eliminated by “projecting” an appropriate number of violation inducing states. We show how the projection method can be implemented for two different five dimensional domain wall actions. The Wilson fermion action is used for the four dimensional part of the action, but it should be understood that other actions realizing a single massive fermion field could be considered, and much of the following derivations do not depend on the specific form chosen.
### A Form of the actions
To fix our notation/conventions, we write the usual, 4-d, Wilson-Dirac operator as
$`D_w(M)`$ $`=`$ $`(4+M)\delta _{x,y}{\displaystyle \frac{1}{2}}{\displaystyle \underset{\mu =1}{\overset{4}{}}}\left[(1\gamma _\mu )U_\mu (x)\delta _{x+\mu ,y}+(1+\gamma _\mu )U_\mu ^{}(y)\delta _{x,y+\mu }\right]`$ (2)
$`=`$ $`\left(\begin{array}{cc}B+M& C\\ C^{}& B+M\end{array}\right).`$ (3)
Here, with $`\sigma _\mu =(\sigma _k,i\mathrm{𝟏})`$,
$`C`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{\mu =1}{\overset{4}{}}}\sigma _\mu \left[U_\mu (x)\delta _{x+\mu ,y}U_\mu ^{}(y)\delta _{x,y+\mu }\right],`$ (4)
$`B`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{\mu =1}{\overset{4}{}}}\left[2\delta _{x,y}U_\mu (x)\delta _{x+\mu ,y}U_\mu ^{}(y)\delta _{x,y+\mu }\right].`$ (5)
We are using a chiral basis: $`\gamma _\mu =\left(\begin{array}{cc}0& \sigma _\mu \\ \sigma _\mu ^{}& 0\end{array}\right)`$, $`\gamma _5=\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right)`$. We will also use the hermitian Wilson-Dirac operator
$`H_w(M)=\gamma _5D_w(M)=\left(\begin{array}{cc}B+M& C\\ C^{}& BM\end{array}\right).`$ (6)
Usually, we will use a large negative mass here, but often omit the argument from $`D_w`$ and $`H_w`$.
In this notation, the usual domain wall fermion action of Shamir reads, with an explicit fifth dimensional lattice spacing $`a_5`$ for the hopping term — the 4-d lattice spacing $`a`$ is kept fixed at $`a=1`$ throughout —,
$`S_{DW}=\overline{\mathrm{\Psi }}D_{DW}^{(5)}\mathrm{\Psi }={\displaystyle \underset{i=1}{\overset{L_s}{}}}\overline{\mathrm{\Psi }}_i\left\{\left[a_5D_w(M)+1\right]\mathrm{\Psi }_iP_{}\mathrm{\Psi }_{i+1}P_+\mathrm{\Psi }_{i1}\right\}`$ (7)
where the extent of the fifth dimension, $`L_s`$ has to be taken as even. The gauge fields in $`D_w(M)`$ are independent of the fifth coordinate, and the fermion fields satisfy the boundary conditions in the fifth direction:
$`P_{}\mathrm{\Psi }_{L_s+1}=mP_{}\mathrm{\Psi }_1,P_+\mathrm{\Psi }_0=mP_+\mathrm{\Psi }_{L_s}.`$ (8)
$`P_\pm `$ are the chiral projectors, $`P_\pm =\frac{1}{2}(1\pm \gamma _5)`$. $`0m<1`$ is proportional to the quark mass (for small $`m`$). The 4-d physical fermion degrees of freedom are identified with the fields at the boundaries as
$`q^R=P_+\mathrm{\Psi }_{L_s}=\mathrm{\Psi }_{L_s}^R,`$ $`q^L=P_{}\psi _1=\mathrm{\Psi }_1^L,`$ (9)
$`\overline{q}^R=\overline{\mathrm{\Psi }}_{L_s}P_{}=\overline{\mathrm{\Psi }}_{L_s}^R,`$ $`\overline{q}^L=\overline{\psi }_1P_+=\overline{\mathrm{\Psi }}_1^L.`$ (10)
As will be shown below, projection of low-lying eigenvalues of $`\mathrm{log}T`$, with $`T`$ the transfer matrix along the fifth direction, can be achieved by introducing an additional term in (7) so that the complete domain wall fermion action reads
$`S_{DWP}=\overline{\mathrm{\Psi }}D_{DWP}^{(5)}\mathrm{\Psi }=`$ $``$ $`{\displaystyle \underset{i=1}{\overset{L_s}{}}}\overline{\mathrm{\Psi }}_i\left\{\left[a_5D_w(M)+1\right]\mathrm{\Psi }_iP_{}\mathrm{\Psi }_{i+1}P_+\mathrm{\Psi }_{i1}\right\}`$ (11)
$`+`$ $`\overline{\mathrm{\Psi }}_1\widehat{A}(m)[P_{}\mathrm{\Psi }_1+P_+\mathrm{\Psi }_{L_s}].`$ (12)
Boriçi introduced an interesting variation of the domain wall fermion action, which differs from the usual one by terms of order $`𝒪(a_5)`$. He called it (5-d) truncated overlap action, since its effective 4-d action for the light fermions is just the polar decomposition approximation of order $`L_s/2`$ to the overlap Dirac operator introduced by Neuberger . In our notation, and introducing again an additional term to be used for projecting low-lying eigenvalues, Boriçi’s variant reads
$`S_{DW^{}}=\overline{\mathrm{\Psi }}D_{DW^{}}^{(5)}\mathrm{\Psi }={\displaystyle \underset{i=1}{\overset{L_s}{}}}\overline{\mathrm{\Psi }}_i`$ $`\left\{\right[a_5D_w(M)+1]\mathrm{\Psi }_i+[a_5D_w(M)1]P_{}\mathrm{\Psi }_{i+1}+`$ (14)
$`[a_5D_w(M)1\left]P_+\mathrm{\Psi }_{i1}\right\}+\overline{\mathrm{\Psi }}_1\widehat{A}(m)[P_{}\mathrm{\Psi }_1+P_+\mathrm{\Psi }_{L_s}].`$
The boundary conditions in the fifth direction remain as in (8). The kernel of the 5-d operator is, including the boundary conditions,
$`D_{DW^{}}^{(5)}=\left(\begin{array}{cccccccc}D_+\widehat{A}P_{}& D_{}P_{}& 0& 0& \mathrm{}& 0& 0& mD_{}P_+\widehat{A}P_+\\ D_{}P_+& D_+& D_{}P_{}& 0& \mathrm{}& 0& 0& 0\\ 0& D_{}P_+& D_+& D_{}P_{}& \mathrm{}& 0& 0& 0\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ 0& 0& 0& 0& \mathrm{}& D_{}P_+& D_+& D_{}P_{}\\ mD_{}P_{}& 0& 0& 0& \mathrm{}& 0& D_{}P_+& D_+\end{array}\right)`$ (15)
where
$`D_\pm =a_5D_w(M)\pm 1.`$ (16)
The kernel of the standard domain wall operator, but including the term to be used for projection, $`D_{DWP}^{(5)}`$, is simply obtained by the replacement $`D_{}1`$.
We now integrate out the heavy fermion degrees of freedom to arrive at a 4-d Dirac operator describing the light fermions. Following Ref. we define $`𝒫`$ by
$`𝒫_{jk}=\{\begin{array}{cc}P_{}\delta _{j,k}+P_+\delta _{j+1,k}\hfill & \text{for }j=1,\mathrm{},L_s1\hfill \\ P_{}\delta _{L_s,k}+P_+\delta _{1,k}\hfill & \text{for }j=L_s.\hfill \end{array}`$ (17)
This has an inverse $`𝒫^1=𝒫^{}`$, given by
$`𝒫_{jk}^1=\{\begin{array}{cc}P_{}\delta _{j,k}+P_+\delta _{j1,k}\hfill & \text{for }j=2,\mathrm{},L_s\hfill \\ P_{}\delta _{1,k}+P_+\delta _{L_s,k}\hfill & \text{for }j=1.\hfill \end{array}`$ (18)
Next, we introduce $`\chi _i`$’s through $`\mathrm{\Psi }_i=(𝒫\chi )_i`$ and define 4-d operators $`Q_\pm `$ as
$`Q_\pm =\{\begin{array}{cc}a_5H_wP_\pm \pm 1\hfill & \text{for the standard domain wall action}\hfill \\ a_5H_w\pm 1\hfill & \text{for Boriçi’s domain wall action}.\hfill \end{array}`$ (19)
Then, using $`\gamma _5P_+=P_+`$ and $`\gamma _5P_{}=P_{}`$ as well as the boundary conditions on the fermion fields, we can rewrite both domain wall actions as
$`S^{(5)}=`$ $``$ $`\{\overline{\mathrm{\Psi }}_1\gamma _5[Q_{}(P_{}mP_+)\chi _1\gamma _5\widehat{A}\chi _1+Q_+\chi _2]`$ (20)
$`+`$ $`{\displaystyle \underset{i=2}{\overset{L_s1}{}}}\overline{\mathrm{\Psi }}_i\gamma _5[Q_{}\chi _i+Q_+\chi _{i+1}]+\overline{\mathrm{\Psi }}_{L_s}\gamma _5[Q_{}\chi _{L_s}+Q_+(P_+mP_{})\chi _1]\}.`$ (21)
Next, we introduce $`\overline{\mathrm{\Psi }}_i=\overline{\chi }_iQ_{}^1\gamma _5`$ and
$`T^1=Q_{}^1Q_+.`$ (22)
This change of variables would give rise to a Jacobian in a dynamical simulation. However, it would be exactly canceled by the Jacobian of this transformation for the pseudo-fermion fields which are needed to cancel the bulk contribution from the 5-d fermions. The action for the pseudo-fermions is obtained by the replacement $`m1`$ from the fermion action. The Jacobians cancel since integration of fermions (Grassman fields) acts like differentiation. We note that for the standard domain wall action $`Q_\pm `$ do not commute and the ordering in (22) is important. From (6) we find
$`T^1=\left(\begin{array}{cc}1& a_5C\frac{1}{\stackrel{~}{B}}\\ 0& \frac{1}{\stackrel{~}{B}}\end{array}\right)\left(\begin{array}{cc}\stackrel{~}{B}& 0\\ a_5C^{}& 1\end{array}\right)=\left(\begin{array}{cc}\stackrel{~}{B}+a_5^2C\frac{1}{\stackrel{~}{B}}C^{}& a_5C\frac{1}{\stackrel{~}{B}}\\ a_5\frac{1}{\stackrel{~}{B}}C^{}& \frac{1}{\stackrel{~}{B}}\end{array}\right),`$ (23)
where $`\stackrel{~}{B}=1+a_5(BM)`$, and thus
$`T=\left(\begin{array}{cc}\frac{1}{\stackrel{~}{B}}& 0\\ a_5C^{}\frac{1}{\stackrel{~}{B}}& 1\end{array}\right)\left(\begin{array}{cc}1& a_5C\\ 0& \stackrel{~}{B}\end{array}\right)=\left(\begin{array}{cc}\frac{1}{\stackrel{~}{B}}& a_5\frac{1}{\stackrel{~}{B}}C\\ a_5C^{}\frac{1}{\stackrel{~}{B}}& a_5^2C^{}\frac{1}{\stackrel{~}{B}}C+\stackrel{~}{B}\end{array}\right).`$ (24)
$`T`$ is the usual domain wall fermion transfer matrix in our conventions.
For later use it will be convenient to introduce a 4-d Hamiltonian $`H_T`$ such that
$`T^1={\displaystyle \frac{1+a_5H_T}{1a_5H_T}}.`$ (25)
For Boriçi’s domain wall action we have simply $`H_T=H_w`$, while for the standard domain wall action one finds
$`H_T={\displaystyle \frac{1}{2+a_5H_w\gamma _5}}H_w=H_w{\displaystyle \frac{1}{2+a_5\gamma _5H_w}}.`$ (26)
In terms of the new fields $`\overline{\chi }`$ and $`\chi `$ the 5-d actions become
$`S^{(5)}=\overline{\chi }D_\chi ^{(5)}\chi =`$ $``$ $`\{\overline{\chi }_1[(P_{}mP_+)\chi _1Q_{}^1\gamma _5\widehat{A}\chi _1T^1\chi _2]`$ (27)
$`+`$ $`{\displaystyle \underset{i=2}{\overset{L_s1}{}}}\overline{\chi }_i[\chi _iT^1\chi _{i+1}]+\overline{\chi }_{L_s}[\chi _{L_s}T^1(P_+mP_{})\chi _1]\}.`$ (28)
Now, we integrate out, in succession, $`\chi _{L_s},\overline{\chi }_{L_s}`$, $`\chi _{L_s1},\overline{\chi }_{L_s1}`$, $`\mathrm{}`$, $`\chi _2,\overline{\chi }_2`$. For this we use, at the first step,
$`\overline{\chi }_{L_s}\chi _{L_s}\overline{\chi }_{L_s1}T^1\chi _{L_s}\overline{\chi }_{L_s}T^1(P_+mP_{})\chi _1=`$ (29)
$`\left[\overline{\chi }_{L_s}\overline{\chi }_{L_s1}T^1\right]\left[\chi _{L_s}T^1(P_+mP_{})\chi _1\right]\overline{\chi }_{L_s1}T^2(P_+mP_{})\chi _1`$ (30)
and at the $`(L_si)`$-th step
$`\overline{\chi }_{i+1}\chi _{i+1}\overline{\chi }_iT^1\chi _{i+1}\overline{\chi }_{i+1}T^{L_s+i}(P_+mP_{})\chi _1=`$ (31)
$`\left[\overline{\chi }_{i+1}\overline{\chi }_iT^1\right]\left[\chi _{i+1}T^{L_s+i}(P_+mP_{})\chi _1\right]\overline{\chi }_iT^{L_s+i1}(P_+mP_{})\chi _1.`$ (32)
With a change of variables $`\chi _{i+1}^{}=\chi _{i+1}T^{L_s+i}(P_+mP_{})\chi _1`$ and $`\overline{\chi }_{i+1}^{}=\overline{\chi }_{i+1}\overline{\chi }_iT^1`$ the integration over $`\chi _{i+1}^{},\overline{\chi }_{i+1}^{}`$ is trivial, giving a factor 1. At the end, we arrive at the 4-d action for $`\chi _1,\overline{\chi }_1`$
$`S^{(4)}=\overline{\chi }_1\left[(P_{}mP_+)T^{L_s}(P_+mP_{})Q_{}^1\gamma _5\widehat{A}\right]\chi _1=\overline{\chi }_1D^{(4)}(m)\chi _1.`$ (33)
The kernel is
$`D^{(4)}(m)`$ $`=`$ $`(P_{}mP_+)T^{L_s}(P_+mP_{})Q_{}^1\gamma _5\widehat{A}`$ (34)
$`=`$ $`\left[{\displaystyle \frac{1+m}{2}}\left(T^{L_s}+1\right)\gamma _5+{\displaystyle \frac{1m}{2}}\left(T^{L_s}1\right)+Q_{}^1\gamma _5\widehat{A}\right]`$ (35)
$`=`$ $`\left[\left(T^{L_s}+1\right)\gamma _5\right]\times \left[{\displaystyle \frac{1+m}{2}}+{\displaystyle \frac{1m}{2}}\gamma _5{\displaystyle \frac{T^{L_s}1}{T^{L_s}+1}}+\gamma _5\left[Q_{}(T^{L_s}+1)\right]^1\gamma _5\widehat{A}\right].`$ (36)
Finally, integrating out $`\chi _1,\overline{\chi }_1`$ and dividing by the pseudo-fermion determinant, obtained by the substitution $`m1`$ and requiring $`\widehat{A}(m=1)=0`$, we find
$`{\displaystyle \frac{detD^{(5)}(m)}{detD^{(5)}(1)}}={\displaystyle \frac{detD^{(4)}(m)}{detD^{(4)}(1)}}=det\left\{\left[D^{(4)}(1)\right]^1D^{(4)}(m)\right\}.`$ (37)
Now, from eq. (35), $`D^{(4)}(1)=\left(T^{L_s}+1\right)\gamma _5`$ and thus we get
$`\left[D^{(4)}(1)\right]^1D^{(4)}(m)={\displaystyle \frac{1}{2}}\left[1+m+(1m)\gamma _5{\displaystyle \frac{T^{L_s}1}{T^{L_s}+1}}+2\gamma _5\left[Q_{}(T^{L_s}+1)\right]^1\gamma _5\widehat{A}(m)\right].`$ (38)
Using (25) we note that
$`{\displaystyle \frac{T^{L_s}1}{T^{L_s}+1}}={\displaystyle \frac{(1+a_5H_T)^{L_s}(1a_5H_T)^{L_s}}{(1+a_5H_T)^{L_s}+(1a_5H_T)^{L_s}}}=\epsilon _{L_s/2}(a_5H_T)`$ (39)
where $`\epsilon _n(x)`$ is Neuberger’s polar decomposition approximation to $`ϵ(x)`$. Ignoring the term with $`\widehat{A}(m)`$ for the moment, we see that
$`\left[D^{(4)}(1)\right]^1D^{(4)}(m)={\displaystyle \frac{1}{2}}\left[1+m+(1m)\gamma _5\epsilon _{L_s/2}(a_5H_T)\right]=D_{tov}(m).`$ (40)
This is just Neuberger’s polar decomposition approximation to the overlap Dirac operator for auxiliary Hamiltonian $`H_T`$, which we denote by $`D_{tov}`$, “truncated overlap.” For $`L_s\mathrm{}`$ it becomes the exact overlap Dirac operator $`D_{ov}`$ . Alternatively, we can write
$`{\displaystyle \frac{T^{L_s}1}{T^{L_s}+1}}=\mathrm{tanh}\left({\displaystyle \frac{L_s}{2}}\mathrm{log}|T|\right)`$ (41)
and eq. (38), still ignoring the term with $`\widehat{A}(m)`$, becomes the effective 4-d Dirac operator that Neuberger derived for domain wall fermions in . Here we have used that $`L_s`$ is even and written the formula in terms of the absolute value $`|T|`$ to indicate that everything remains well defined when an eigenvalue of $`T`$ becomes negative.
We now can use $`\widehat{A}(m)`$ to project out low-lying eigenvectors, $`v_i`$, of the auxiliary Hamiltonian $`H_T`$ for which $`\epsilon _{L_s/2}`$ with finite $`L_s`$ is not a sufficiently accurate approximation to $`ϵ(x)`$. Let
$`H_Tv_i=\lambda _iv_i,Tv_i=T_iv_i,\widehat{P}_i=v_iv_i^{},`$ (42)
where, from eq. (25),
$`T_i={\displaystyle \frac{1a_5\lambda _i}{1+a_5\lambda _i}}.`$ (43)
The projection can then be achieved by setting
$`\widehat{A}(m)=(1m)\gamma _5Q_{}{\displaystyle \underset{i}{}}g_i\widehat{P}_i,`$ (44)
with
$`g_i`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left[\left(T_i^{L_s}1\right)+\left(T_i^{L_s}+1\right)ϵ(a_5\lambda _i)\right].`$ (45)
Note that $`\widehat{A}(m)`$ vanishes for $`m=1`$, as required for the pseudo-fermions. With this $`\widehat{A}(m)`$ we find, instead of (40),
$`\left[D^{(4)}(1)\right]^1D^{(4)}(m)=D_{ov}(m)=`$ (46)
$`{\displaystyle \frac{1}{2}}\left\{1+m+(1m)\gamma _5\left[\epsilon _{L_s/2}(a_5H_T)\left(1{\displaystyle \underset{i}{}}\widehat{P}_i\right)+{\displaystyle \underset{i}{}}ϵ(a_5\lambda _i)\widehat{P}_i\right]\right\}.`$ (47)
For Boriçi’s variant of the domain wall action, where $`H_T=H_w`$ commutes with $`Q_\pm `$ the “projection operator” $`\widehat{A}(m)`$ can be simplified to
$`\widehat{A}(m)=(1m)\gamma _5{\displaystyle \underset{i}{}}f_i\widehat{P}_i,`$ (48)
with
$`f_i`$ $`=`$ $`{\displaystyle \frac{1}{2}}(a_5\lambda _i1)\left[\left(T_i^{L_s}1\right)+\left(T_i^{L_s}+1\right)ϵ(a_5\lambda _i)\right]`$ (49)
$`=`$ $`\{\begin{array}{cc}(a_5\lambda _i1)\hfill & \text{for }\lambda _i>0\hfill \\ \frac{(1+a_5\lambda _i)^{L_s}}{(1a_5\lambda _i)^{L_s1}}\hfill & \text{for }\lambda _i<0.\hfill \end{array}`$ (50)
### B Propagator
Our next step is to relate the 4-d overlap propagator to the 5-d propagator of the corresponding domain wall fermion action. In all steps below $`\widehat{A}(m)`$ for the eigenvalue projections is included. From (47) we find, obviously,
$`D_{ov}^1(m)=\left[D^{(4)}(m)\right]^1D^{(4)}(1).`$ (51)
To connect this to the 5-d theory, we consider, motivated by the fact that the light 4-d fermion is $`q=(𝒫^1\mathrm{\Psi })_1=\chi _1`$,
$`X`$ $`=`$ $`\left\{𝒫^1\left[D^{(5)}(m)\right]^1D^{(5)}(1)𝒫\right\}_{11}`$ (52)
$`=`$ $`{\displaystyle \frac{1}{Z}}{\displaystyle \underset{i}{}d\mathrm{\Psi }_id\overline{\mathrm{\Psi }}_i\underset{k}{}(𝒫^1\mathrm{\Psi })_1\overline{\mathrm{\Psi }}_k\left[D^{(5)}(1)𝒫\right]_{k1}\mathrm{e}^{S^{(5)}}}`$ (53)
$`=`$ $`{\displaystyle \frac{1}{Z^{}}}{\displaystyle \underset{i}{}d\chi _id\overline{\chi }_i\underset{k}{}\chi _1\overline{\chi }_kQ_{}^1\gamma _5\left[D^{(5)}(1)𝒫\right]_{k1}\mathrm{e}^{S^{(5)}}}`$ (54)
$`=`$ $`{\displaystyle \frac{1}{Z^{}}}{\displaystyle \underset{i}{}d\chi _id\overline{\chi }_i\underset{k}{}\chi _1\overline{\chi }_k\left[D_\chi ^{(5)}(1)\right]_{k1}\mathrm{e}^{S^{(5)}}}.`$ (55)
Here we have used the 5-d $`D_\chi ^{(5)}=Q_{}^1\gamma _5D^{(5)}𝒫`$ introduced in eq. (28). Now, we integrate out, in succession, $`\chi _{L_s},\overline{\chi }_{L_s}`$, $`\chi _{L_s1},\overline{\chi }_{L_s1}`$, $`\mathrm{}`$, $`\chi _2,\overline{\chi }_2`$. From the transformations (31) we see that $`\overline{\chi }_k\overline{\chi }_1T^{k+1}`$ in the process. Therefore, we obtain
$`X={\displaystyle \frac{1}{Z^{}}}{\displaystyle 𝑑\chi _1𝑑\overline{\chi }_1\chi _1\overline{\chi }_1\underset{k}{}T^{k+1}\left[D_\chi ^{(5)}(1)\right]_{k1}\mathrm{e}^{\overline{\chi }_1D^{(4)}(m)\chi _1}}.`$ (56)
But from (28) we see that
$`{\displaystyle \underset{k}{}}T^{k+1}\left[D_\chi ^{(5)}(1)\right]_{k1}=(P_{}P_+)T^{L_s}(P_+P_{})=D^{(4)}(1),`$ (57)
where we used the fact that $`\widehat{A}(1)=0`$. Thus we finally obtain
$`X=\left\{𝒫^1\left[D^{(5)}(m)\right]^1D^{(5)}(1)𝒫\right\}_{11}=\left[D^{(4)}(m)\right]^1D^{(4)}(1)=D_{ov}^1(m).`$ (58)
To solve
$`D_{ov}(m)\psi =b`$ (59)
we introduce $`\stackrel{~}{b}=(b,0,\mathrm{},0)^T`$ and solve
$`D^{(5)}(m)\varphi =D^{(5)}(1)𝒫\stackrel{~}{b}.`$ (60)
$`\psi `$ is then obtained as
$`\psi =\left(𝒫^1\varphi \right)_1.`$ (61)
The physical propagator has a contact term subtracted , and we arrive at the general form for both auxiliary Hamiltonians considered here, with and without projection, and also valid for finite $`L_s`$,
$`\stackrel{~}{D}_{ov}^{}{}_{}{}^{1}(m)={\displaystyle \frac{1}{1m}}\left[D_{ov}^1(m)1\right]={\displaystyle \frac{1}{1m}}\left[\left\{𝒫^1\left[D^{(5)}(m)\right]^1D^{(5)}(1)𝒫\right\}_{11}1\right].`$ (62)
### C Relation of the 5-d and 4-d operators
So far, we have related $`D_{ov}^1(m)`$ and $`detD_{ov}(m)`$, or their truncated versions, to 5-d operators, see (37), (40) and (47) and finally (58). E.g. for eigenvalue calculations it would be useful to establish a similar connection for $`D_{ov}(m)`$ or $`D_{tov}(m)`$. Consider $`D_\chi ^{(5)}(m)`$ introduced in (28). We find
$`\left[D_\chi ^{(5)}(m)D_\chi ^{(5)}(1)\right]_{ij}=`$ (63)
$`\left[(1m)P_+Q_{}^1\gamma _5\widehat{A}(m)\right]\delta _{i1}\delta _{j1}T^1(1m)P_{}\delta _{iL_s}\delta _{j1}.`$ (64)
Thus
$`\left[D_\chi ^{(5)}(1)\right]^1D_\chi ^{(5)}(m)=`$ (65)
$`1+\left[D_\chi ^{(5)}(1)\right]_{i1}^1\left[(1m)P_+Q_{}^1\gamma _5\widehat{A}(m)\right]\delta _{j1}(1m)\left[D_\chi ^{(5)}(1)\right]_{iL_s}^1T^1P_{}\delta _{j1}.`$ (66)
So only the first column (of 4-d blocks) is different from the 5-d unit matrix. Let’s call the entries in the first column $`X_1,X_2,\mathrm{},X_{L_s}`$. Similarly, its inverse is
$`\left[D_\chi ^{(5)}(m)\right]^1D_\chi ^{(5)}(1)=`$ (67)
$`1\left[D_\chi ^{(5)}(m)\right]_{i1}^1\left[(1m)P_+Q_{}^1\gamma _5\widehat{A}(m)\right]\delta _{j1}+(1m)\left[D_\chi ^{(5)}(m)\right]_{iL_s}^1T^1P_{}\delta _{j1}.`$ (68)
Again, only the first column is different from the 5-d unit matrix. Let’s call the entries in the first column $`Y_1,Y_2,\mathrm{},Y_{L_s}`$. Since these 5-d matrices are the inverse of each other, their product is the 5-d unit matrix, i.e.
$`Y_1X_1=1`$ $``$ $`Y_1=(X_1)^1`$ (69)
$`Y_2X_1+X_2=0`$ $``$ $`Y_2=X_2(X_1)^1`$ (70)
$`\mathrm{}`$ $`\mathrm{}`$ (71)
$`Y_jX_1+X_j=0`$ $``$ $`Y_j=X_j(X_1)^1\mathrm{for}j=1,\mathrm{},L_s.`$ (72)
The first equation, $`Y_1=(X_1)^1`$ establishes the relation we are looking for
$`D_{ov}(m)=\left\{𝒫^1\left[D^{(5)}(1)\right]^1D^{(5)}(m)𝒫\right\}_{11}.`$ (73)
Note that the components of the 5-d propagator used for the inverse of the overlap Dirac operator in eq. (58) do not correspond to the physical quark propagator $`q\overline{q}`$ as obtained from domain wall fermions. The reason is that the overlap propagator still needs a subtraction of a contact term and a multiplicative normalization . We consider here the standard domain wall action without projection, and denote the corresponding 4-d Dirac operator $`D_{tov}(m;H_T)`$. To make the connection between the subtracted 4-d propagator and $`q\overline{q}`$ explicit, we write the 1 for the subtraction as
$`1=\left\{𝒫^1\left[D_{DW}^{(5)}(m)\right]^1D_{DW}^{(5)}(m)𝒫\right\}_{11}`$ (74)
and thus
$`D_{tov}^1(m;H_T)1=\left\{𝒫^1\left[D_{DW}^{(5)}(m)\right]^1\left[D_{DW}^{(5)}(1)D_{DW}^{(5)}(m)\right]𝒫\right\}_{11}.`$ (75)
From the matrix for $`D_{DW}^{(5)}`$, eq. (15) with the replacement $`D_{}1`$ and setting $`\widehat{A}(m)=0`$, we see that
$`\left[D_{DW}^{(5)}(1)D_{DW}^{(5)}(m)\right]_{ij}=(1m)\left[P_{}\delta _{iL_s}\delta _{j1}+P_+\delta _{i1}\delta _{jL_s}\right].`$ (76)
Therefore we obtain
$`D_{tov}^1(m;H_T)1=(1m)\left\{𝒫^1[D_{DW}^{(5)}(m)]^1𝒥𝒫\right\}_{11},`$ (77)
where $`𝒥_{ij}=\delta _{i,L_s+1j}`$ is the inversion operator of the fifth direction. Now from (10) we see that the physical fermion degrees are given in terms of the domain wall boundary fermions as
$`q=(𝒫^1\mathrm{\Psi })_1,\overline{q}=(\overline{\mathrm{\Psi }}𝒥𝒫)_1.`$ (78)
Hence we find
$`q\overline{q}=\left\{𝒫^1[D_{DW}^{(5)}(m)]^1𝒥𝒫\right\}_{11}={\displaystyle \frac{1}{1m}}\left[D_{tov}^1(m;H_T)1\right].`$ (79)
The usual domain wall physical fermion propagator automatically contains both the subtraction and multiplicative normalization of the overlap fermion propagator.
The hermitian conjugate of the 5-d operator $`D_{DW^{}}^{(5)}`$, eq. (15), of Boriçi’s variant is different than that of the usual DWF operator, $`D_{DW}^{(5)}`$, without projection. In particular it does not have the (generalized) $`\gamma _5`$ hermiticity, $`D_{DW}^{(5)}=\gamma _5𝒥D_{DW}^{(5)}\gamma _5𝒥`$, with $`𝒥`$ the inversion operator of the fifth direction. Instead we find
$`D_{DW^{}}^{(5)}=\left(\begin{array}{cccccccc}D_+^{}P_{}\widehat{A}^{}& P_+D_{}^{}& 0& 0& \mathrm{}& 0& 0& mP_{}D_{}^{}\\ P_{}D_{}^{}& D_+^{}& P_+D_{}^{}& 0& \mathrm{}& 0& 0& 0\\ 0& P_{}D_{}^{}& D_+^{}& P_+D_{}^{}& \mathrm{}& 0& 0& 0\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ 0& 0& 0& 0& \mathrm{}& P_{}D_{}^{}& D_+^{}& P_+D_{}^{}\\ mP_+D_{}^{}P_+\widehat{A}^{}& 0& 0& 0& \mathrm{}& 0& P_{}D_{}^{}& D_+^{}\end{array}\right)`$ (80)
where
$`D_\pm ^{}=a_5D_w^{}(M)\pm 1,\widehat{A}^{}(m)=(1m){\displaystyle \underset{i}{}}f_i\widehat{P}_i\gamma _5.`$ (81)
### D Eigenvalues of the pseudo-fermion operators
Here, we want to investigate a little more the properties of the pseudo-fermions, needed to cancel the bulk contributions in the 5-d domain wall fermion approaches.
For both the standard domain wall fermion action (7) and for Boriçi’s variant (14) we find for the pseudo-fermion matrix, including the boundary conditions and recalling that the projection operator $`\widehat{A}(1)`$ vanishes,
$`D^{(5)}(1)D^{(5)}(1)=\left(\begin{array}{cccccccc}X& Y& 0& 0& \mathrm{}& 0& 0& Y\\ Y& X& Y& 0& \mathrm{}& 0& 0& 0\\ 0& Y& X& Y& \mathrm{}& 0& 0& 0\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ 0& 0& 0& 0& \mathrm{}& Y& X& Y\\ Y& 0& 0& 0& \mathrm{}& 0& Y& X\end{array}\right)`$ (82)
with $`X`$ and $`Y`$ 4-d hermitian matrices,
$`Y`$ $`=`$ $`P_+D_+D_+^{}P_{}={\displaystyle \frac{1}{2}}\left(a_5D_w+a_5D_w^{}+2\right)`$ (83)
$`X`$ $`=`$ $`D_+D_+^{}+1=a_5^2D_wD_w^{}2Y`$ (84)
for the standard domain wall action, and
$`Y`$ $`=`$ $`D_+P_+D_{}^{}+D_{}P_{}D_+^{}=a_5^2D_wD_w^{}1`$ (85)
$`X`$ $`=`$ $`D_+D_+^{}+D_{}D_{}^{}=2a_5^2D_wD_w^{}+2=4a_5^2D_wD_w^{}2Y`$ (86)
for Boriçi’s variant. Similar relations can be found in Ref. .
Now, let $`S`$ be the shift (or translation) operator in the 5-th direction with anti-periodic boundary condition:
$`S=\left(\begin{array}{cccccc}0& 1& 0& \mathrm{}& 0& 0\\ 0& 0& 1& \mathrm{}& 0& 0\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ 0& 0& 0& \mathrm{}& 0& 1\\ 1& 0& 0& \mathrm{}& 0& 0\end{array}\right)`$ (87)
It is easy to see that $`[S,D^{(5)}(1)D^{(5)}(1)]=0`$, and so $`S`$ and $`D^{(5)}(1)D^{(5)}(1)`$ can be diagonalized simultaneously. The eigenvalues of $`S`$ are $`\mathrm{exp}\{i\pi (2k+1)/L_s\}`$ for $`k=0,\mathrm{},L_s1`$, because of the anti-periodic boundary conditions. The corresponding eigenvectors are
$`w(k)^T=`$ (88)
$`(v_4,v_4\mathrm{exp}\{i\pi (2k+1)/L_s\},v_4\mathrm{exp}\{i\pi 2(2k+1)/L_s\},\mathrm{},v_4\mathrm{exp}\{i\pi (L_s1)(2k+1)/L_s\})^T`$ (89)
with $`v_4`$ some 4-d vector. This also has to be an eigenvector of $`D^{(5)}(1)D^{(5)}(1)`$. From (82), and using (84) and (86) we obtain the eigenvalue equation for $`v_4`$
$`\left\{Ja_5^2D_wD_w^{}2Y\left[1\mathrm{cos}\left({\displaystyle \frac{\pi }{L_s}}(2k+1)\right)\right]\right\}v_4=\lambda (k)v_4`$ (90)
Here $`J=1`$ for domain wall fermions, and $`J=4`$ for Boriçi’s variant, and we have indicated the dependence of $`\lambda `$ on $`k`$, the momentum in the 5-th direction. Since the left-hand-side of (90) remains unchanged under $`kL_sk1`$ we conclude that $`\lambda (L_sk1)=\lambda (k)`$. Therefore, the eigenvalues of $`D^{(5)}(1)D^{(5)}(1)`$ are all (at least) doubly degenerate.
Furthermore, for large $`L_s`$ and small $`k`$ the pseudo-fermion domain wall eigenvalues closely track the Wilson eigenvalues, the eigenvalues of $`D_wD_w^{}`$, since then $`2Y\left[1\mathrm{cos}\left(\frac{\pi }{L_s}(2k+1)\right)\right]Y\pi ^2(2k+1)^2/L_s^2`$ in (90) is only a small perturbation. In particular, at zero crossings of the Wilson Dirac operator, a pair of pseudo-fermion domain wall eigenvalues will go to zero for $`L_s\mathrm{}`$.
From the domain wall Dirac operator $`D_{DW}^{(5)}(m)`$ without the projection matrix $`\widehat{A}(m)`$, we can make a hermitian version, $`H_{DW}^{(5)}(m)=D_{DW}^{(5)}(m)𝒥\gamma _5`$, with $`𝒥_{ij}=\delta _{i,L_s+1j}`$ the inversion operator of the fifth direction. From eq. (15) with $`D_{}1`$ and $`\widehat{A}=0`$ it becomes
$`H_{DW}^{(5)}(m)=\left(\begin{array}{cccccccc}mP_+& 0& 0& 0& \mathrm{}& 0& P_{}& D_+\gamma _5\\ 0& 0& 0& 0& \mathrm{}& P_{}& D_+\gamma _5& P_+\\ 0& 0& 0& 0& \mathrm{}& D_+\gamma _5& P_+& 0\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ P_{}& D_+\gamma _5& P_+& 0& \mathrm{}& 0& 0& 0\\ D_+\gamma _5& P_+& 0& 0& \mathrm{}& 0& 0& mP_{}\end{array}\right)`$ (91)
Since $`D_{DW}^{(5)}(m)D_{DW}^{(5)}(m)=[H_{DW}^{(5)}(m)]^2`$, the vectors $`w(k)`$ and $`w(L_sk1)`$ of eq. (88) are degenerate eigenvectors of $`[H_{DW}^{(5)}(1)]^2`$. To get the eigenvalues of $`H_{DW}^{(5)}(1)`$ we need to compute the $`2\times 2`$ matrix $`w(k_1)|H_{DW}^{(5)}(1)|w(k_2)`$ with $`k_1,k_2=k,L_sk1`$. We find
$`w(k_1)|H_{DW}^{(5)}(1)|w(k_2)=`$ (92)
$`{\displaystyle \underset{j=1}{\overset{L_s}{}}}\mathrm{e}^{\frac{i\pi }{L_s}(2k_1+1)(j1)}\mathrm{e}^{\frac{i\pi }{L_s}(2k_2+1)(L_sj)}v_4|P_{}\mathrm{e}^{\frac{i\pi }{L_s}(2k_2+1)}+D_+\gamma _5P_+\mathrm{e}^{\frac{i\pi }{L_s}(2k_2+1)}|v_4`$ (93)
The sum over $`j`$ gives the constraint $`k_2=L_sk_11`$, and we thus find the $`2\times 2`$ matrix to be of the form
$`\left(\begin{array}{cc}0& b\\ b^{}& 0\end{array}\right)`$ (94)
and therefore the eigenvalues of the hermitian pseudo-fermion domain wall operator come in $`\pm `$ pairs.
### E Preconditioning
Even-odd preconditioning is a common technique to reduce the condition number of a matrix and can be implemented in both versions of the domain wall actions discussed above. With projection, the matrices have a more complicated form reminiscent of preconditioning for the clover fermion action.
Write the matrix $`D^{(5)}`$ as a two by two block matrix
$`D^{(5)}=\left(\begin{array}{cc}𝒜_{EE}& _{EO}\\ _{OE}& 𝒜_{OO}\end{array}\right)`$ (95)
where the upper case characters $`E`$ and $`O`$ label even and odd (checkerboard) sites, respectively, in five dimensions. The matrix $`D^{(5)}`$ can be brought into an even-odd block diagonal form using the lower and upper block triangular matrices
$`L=\left(\begin{array}{cc}I_{EE}& 0\\ _{OE}𝒜_{EE}^1& I_{OO}\end{array}\right)U=\left(\begin{array}{cc}𝒜_{EE}& _{EO}\\ 0& I_{OO}\end{array}\right)`$ (96)
with the transformation
$`\stackrel{~}{D}^{(5)}=L^1D^{(5)}U^1=\left(\begin{array}{cc}I_{EE}& 0\\ 0& 𝒜_{OO}_{OE}𝒜_{EE}^1_{EO}\end{array}\right).`$ (97)
This decomposition always exists as long as $`𝒜_{EE}`$ is non-singular.
The structure of $`𝒜_{EE}`$ is deduced from the structure of $`D_w(M)`$ and the projection matrix $`\widehat{A}`$ that occur in both the standard domain wall and Boriçi variants, $`D_{DW}^{(5)}`$ and $`D_{DW^{}}^{(5)}`$, respectively. In Boriçi’s variant the nearest neighbor coupling in the fifth direction contains also hoppings in the spatial directions leading to a complicated even-odd structure. This case will not be discussed further here. For preconditioning of the standard domain wall action with projection, after a rescaling of the fermion fields to bring $`D_{DW}^{(5)}`$ into the “kappa”-form, i.e., having 1 in the diagonal except for the projection matrix $`\widehat{A}`$, we have
$`𝒜_{EE}=\left(\begin{array}{ccccc}1\widehat{A}_{ee}P_{}& 0& \mathrm{}& 0& \widehat{A}_{eo}P_+\\ 0& 1& \mathrm{}& 0& 0\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ 0& 0& \mathrm{}& 1& 0\\ 0& 0& \mathrm{}& 0& 1\end{array}\right).`$ (98)
Here $`\widehat{A}`$ is rescaled by $`2\kappa `$, with $`\kappa =1/(2(a_5(4+M)+1))`$, and $`e`$ and $`o`$ refer to the even and odd sub-lattice in the 4-d checker-boarding.
For even-odd preconditioning, we need
$`𝒜_{EE}^1=\left(\begin{array}{ccccc}(1\widehat{A}_{ee}P_{})^1& 0& \mathrm{}& 0& (1\widehat{A}_{ee}P_{})^1\widehat{A}_{eo}P_+\\ 0& 1& \mathrm{}& 0& 0\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ 0& 0& \mathrm{}& 1& 0\\ 0& 0& \mathrm{}& 0& 1\end{array}\right).`$ (99)
Therefore, for every application of the projected domain wall Dirac operator, the inversion $`(1\widehat{A}_{ee}P_{})^1\varphi `$ for some 4-d even sub-lattice fermion $`\varphi `$ has to be performed once — the two occurrences in $`𝒜_{EE}^1`$ can be combined into the application on one 4-d even sub-lattice fermion field. In the limit of the fermion mass $`m1`$, the matrices $`𝒜_{EE}^1`$ and $`𝒜_{OO}`$ become the unit matrix. For small $`m`$, the matrix $`(1\widehat{A}_{ee}P_{})`$ should be well conditioned, and the cost of inversion is the number of iterations involving applications of $`D_w`$ on half of the four volume. This “inner” iteration cost is a multiplicative overhead on the cost of inverting $`\stackrel{~}{D}^{(5)}`$.
### F Dynamical fermions
Implementing an HMC algorithm for dynamical domain wall fermion simulations is a straightforward modification of dynamical Wilson fermion simulation code. Implementing projection does not cause any fundamental difficulties. What is needed, in addition to the usual case, for the computation of the fermion contribution to the force is $`\widehat{A}(m;H(U))/U`$, where we indicated the implicit dependence on $`U`$ of the projection operator $`\widehat{A}`$. Using the chain rule, what we need is
$`{\displaystyle \frac{\lambda _i}{U}}\mathrm{and}{\displaystyle \frac{\widehat{P}_i}{U}}={\displaystyle \frac{v_i}{U}}v_i^{}+v_i{\displaystyle \frac{v_i^{}}{U}},`$ (100)
for $`H(U)v_i=\lambda _iv_i`$ with $`v_i^{}v_i=1`$. We easily find
$`{\displaystyle \frac{\lambda _i}{U}}=v_i^{}{\displaystyle \frac{H(U)}{U}}v_i.`$ (101)
and
$`{\displaystyle \frac{v_i}{U}}=\left(H\lambda _i\right)^1\left[v_i^{}{\displaystyle \frac{H(U)}{U}}v_i{\displaystyle \frac{H(U)}{U}}\right]v_i.`$ (102)
Therefore, a different 4-d inversion is needed for every projected eigenvector $`v_i`$. For the standard domain wall fermion action, where $`H=H_T`$, additional 4-d inversions are needed to compute $`\frac{H_T(U)}{U}`$.
We note that the modifications involving projection for the force term may not be necessary in practice. The unprojected action can be used to construct the guiding Hamiltonian for HMC with the projected action used for the computation of the initial and final energies in the Metropolis accept/reject step. A critical factor in the choice will be the acceptance rate which will suffer if the guiding Hamiltonian is not accurate enough without projection.
## III Results
### A Smooth gauge field
We start our analysis by first considering the case of spectral flow on a smooth SU(2) instanton configuration. We use the instanton construction from ref. , namely a single $`8^4`$ instanton with $`\rho =1.5`$ and anti-periodic boundary conditions. In this section we will mainly focus on the standard domain wall operator. Similar results are found for the Boriçi variant.
What we are interested in is first determining how various domain-wall like actions reproduce topology on the lattice. Without projection, we note that the index of the massless four dimensional (4-d) operator is
$`Q=\mathrm{Tr}(\gamma _5D_{tov})={\displaystyle \frac{1}{2}}\mathrm{Tr}\left(\epsilon _{L_s/2}(a_5H(M))\right)`$ (103)
with the auxiliary Hamiltonian $`H`$ given by either $`H_T`$ or $`H_w`$ — the two cases under consideration. For infinite $`L_s`$, $`\epsilon _{L_s/2}`$ becomes $`ϵ`$ and $`Q`$ is a measure of the discrepancy of positive and negative states of $`H`$. A simple way to determine $`Q`$ is to start from some $`M`$ where we know that $`Q=0`$ and increase it to the desired value while counting the change of the number of positive states. We emphasize that this procedure works as long as the spectral flow is smooth — namely, as long as the eigenvalues of $`H`$ are well behaved. While $`H_w`$ is always well behaved, we note that this is not the case for the domain wall $`H_T(M)`$ for $`M>2`$. In particular, the spectral flow method will fail for $`M`$ in the “doubler” region where there are multiple species of chiral fermions. We will address this point in more detail later.
In Fig. 1 we show the spectral flow of the lowest 10 eigenvalues of the hermitian Wilson-Dirac operator $`H_w(M)`$, the lowest 10 eigenvalues of the hermitian domain wall Dirac operator $`𝒥\gamma _5D_{DW}^{(5)}(0)`$ with $`a_5=1`$ for various fifth dimensional extent values $`L_s`$, and the 10 lowest eigenvalues of the hermitian overlap-Dirac operator (with projection, i.e. with chiral symmetry), all as function of the “domain wall height $`M`$”. We recall from eq. (73)
$`D_{ov}(0)=\left\{𝒫^1\left[D^{(5)}(1)\right]^1D^{(5)}(0)𝒫\right\}_{11}`$ (104)
that when $`D_{ov}(0)`$ has zero-eigenvalues we can expect $`D^{(5)}(0)`$ to have zero-eigenvalues as long as $`D^{(5)}(1)`$ is not singular. We see in Fig. 1 that there are exact zeros of $`D_{ov}(0)`$ for $`M`$ beyond the crossing, and the spectrum changes discontinuously at the crossing . For finite $`L_s`$, the $`D_{DW}^{(5)}(0)`$ eigenvalues decrease quickly in $`L_s`$. For $`M`$ beyond the crossing they quickly go to zero and in fact appear exponential in $`L_s`$, but near the crossing, when the eigenvalues of $`H_w`$ are small, the decrease is slowed.
This slow rate of convergence in $`L_s`$, when the eigenvalues of $`H_w`$ are small, is a prime source of difficulty in recent numerical calculations involving domain wall fermions. The reason for the slow rate of convergence is clear when the form of the 4-d operator is examined,
$`D_{tov}(m)={\displaystyle \frac{1}{2}}\left[1+m+(1m)\gamma _5\epsilon _{L_s/2}(a_5H)\right].`$ (105)
When $`\epsilon _{L_s/2}(a_5H)`$ deviates from one, for example for the smallest and largest eigenvalues of $`H`$, there is violation of chiral symmetry. Fig. 2 shows the domain wall spectral flow when the $`5`$ smallest (closest to zero) eigenvectors of $`H_T`$ are projected out as prescribed by eqs. (44),(45) and (47). With projection turned on, $`D_{DW}^{(5)}(0)`$ no longer has (generalized) $`\gamma _5`$ hermiticity, so the eigenvalues of $`+\sqrt{(D_{DW}^{(5)}(0))^{}D_{DW}^{(5)}(0)}`$ are determined. We see that just beyond the crossing there is a discontinuous drop of the eigenvalues to zero. In fact fairly small eigenvalues are seen even for $`L_s=4`$. Also shown are the lowest eigenvalues of $`H_T`$. We see they agree quite well with the eigenvalues of $`H_w`$ up to a factor of two as predicted by eq. (26).
All computations of eigenvalues and eigenvectors where done with the Ritz method . For computing eigenvectors of $`H_T(M)`$, in each Ritz iteration, the application of $`H_T(M)^{}H_T(M)`$ on a vector is needed. From eq. (26) for $`H_T`$, the expression for $`H_T^{}H_T`$ can be combined into
$`H_T^{}H_T=\gamma _5D_w{\displaystyle \frac{1}{(2+a_5D_w^{})(2+a_5D_w)}}\gamma _5D_w,`$ (106)
so a single inversion of a hermitian positive definite operator is needed for each Ritz iteration and Conjugate Gradient for Normal Equations (CGNE) was used. Typically 30 to 60 iterations of CGNE were needed for each inversion to achieve an accuracy of $`10^7`$.
It should come as no surprise that zero eigenvalues of the domain wall operator can be obtained when projection is turned on since this is the same mechanism by which they are obtained for the overlap-Dirac operator . What was key to the latter case was enough projection coupled with a judicious choice of the fifth dimension spacing $`a_5`$. Since $`a_5`$ only serves as a multiplicative scaling of $`H_w`$, $`a_5`$ can be chosen so that the lowest unprojected and highest eigenvalues of $`H_w`$ lie within the range of approximation of $`ϵ`$. For the overlap-Dirac operator in Ref. , an optimal rational approximation was chosen; however, this merely increased the useful range of the approximation. When the condition is met for a valid approximation of $`ϵ`$, there are no other cutoff effects associated with having a finite $`a_5`$. For example, observables of the effective 4-d theory will not have slightly different 4-d lattice spacing dependencies for slightly different choices of $`a_5`$. However, for the conventional domain wall action, $`H_T`$ intrinsically depends on $`a_5`$ as seen in eq. (26). Therefore, even for infinite extent in the fifth direction, where $`\epsilon (a_5H_T)`$ converges to $`ϵ(a_5H_T)=ϵ(H_T)`$, choosing different $`a_5`$ results in different 4-d lattice spacing dependence for physical observables. This fact limits the usefulness of adjusting $`a_5`$ with the goal of imposing arbitrarily precise chiral symmetry for (small) finite fifth dimension as will be shown later.
In Fig. 3 we examine more closely the smallest eigenvalue of the domain wall operator and the deviation of $`\epsilon (a_5H_T)`$ from one for $`a_5=1`$. We see that without projection the eigenvalues decrease approximately exponentially in $`L_s`$ with a slow rate varying with $`M`$, while with projection they drop to about $`10^4`$ after the crossing. This is the accuracy to which the eigenvalues were computed in single precision and is “zero” here. For increasing $`M`$, the zero eigenvalues slowly increase.
Some clue to the slow increase can be seen from the right panel in Fig. 3 which shows the deviation from unity of $`\epsilon (\lambda _i)`$ where $`\lambda _i`$ is the first, fifth and largest eigenvalue of $`H_T`$. The deviation for the fifth and the largest eigenvalue gives a measure of the accuracy of the approximation since the first five eigenvalues are projected out of $`H_T`$. The deviation of the first eigenvalue indicates how much the standard unprojected case deviates. For $`L_s=4`$ the deviation of the largest eigenvalue is clearly visible, and for $`L_s8`$ the deviation is off the bottom of the graph until $`M`$ approaches $`2`$. For $`L_s=8`$ a typical deviation of the fifth eigenvalue is about 1% while nonetheless a domain wall eigenvalue of $`10^4`$ is obtained. However, as $`M`$ is increased, the largest eigenvalue is not well approximated and the domain wall eigenvalue deviates away from zero due to the resulting chiral symmetry breaking.
We will conclude this portion of the analysis with a final look at the role of the pseudo-fermion term in eq. (73) (repeated again in eq. (104)). It was stated that zero eigenvalues of the 5-d $`D^{(5)}(0)`$ are seen when the 4-d $`D_{ov}(0)`$ has zero eigenvalues. However, as shown in Sec. II D, the pseudo-fermion term has a zero eigenvalue in the infinite fifth dimensional extent limit at the zero crossings of $`H_w`$. We show in Fig. 4 the lowest eigenvalue of the hermitian version of the unprojected $`D^{(5)}(0)`$ and the pseudo-fermion operator $`D^{(5)}(1)`$ (recall that the pseudo-fermion operator is unaffected by projection). The pseudo-fermion eigenvalues track the regular eigenvalues up to the crossing, then move away from zero again. As $`L_s`$ increases, the pseudo-fermion eigenvalues appear to go exponentially to zero. We see then by studying eq. (104) that the pseudo-fermions cancel to some extent in the ratio $`(D^{(5)}(1))^1D^{(5)}(0)`$.
For the resultant unprojected 4-d operator $`D_{tov}(0)`$ (finite $`L_s`$), the eigenvalues are larger than the corresponding 5-d ones and the 4-d eigenvalues decrease across a broad region in the mass $`M`$. These eigenvalues are the actual ones occurring in 4-d observables like the pion propagator and the chiral condensate. More to the point, the eigenvalues of the 5-d operator $`D^{(5)}(m)`$ are in general not directly physically relevant to 4-d observables. Stated differently, the eigenvectors or some piece of the eigenvectors of the $`D^{(5)}(m)`$ are not eigenvectors of $`D_{tov}(m)`$ because of the basis changing matrices $`𝒫`$ and the projection onto the first 5-d slice. However, eq. (104) gives us the connection between the 5-d and 4-d matrices and in fact provides a new algorithm to compute the relevant 4-d eigenvalues from only the 5-d operators. It has some attractive features for practical use because of its simplicity and direct connection, but it is not really efficient for the reason that $`D^{(5)}(1)`$ has near zero eigenvalues at $`H_w`$ zero crossings and hence has bad condition numbers, at least in the large $`L_s`$ limit.
We proceed to another test of chiral symmetry given by the generalized Gellmann-Oakes-Renner relation (called GMOR by many authors). In the form used here and with our normalization conventions , it states that
$`mb|(\gamma _5\stackrel{~}{D}_{ov}^1(m))^2|b=b|\stackrel{~}{D}_{ov}^1(m)|b\gamma _5|b=\pm |b,`$ (107)
with the physical (subtracted, see eq. (62)) fermion propagator
$`\stackrel{~}{D}_{tov}^{}{}_{}{}^{1}(m)={\displaystyle \frac{1}{1m}}\left[D_{tov}^1(m)1\right]`$ (108)
and $`L_s`$ taken to the infinite fifth dimensional limit. Averaging over several (chiral) Gaussian random vectors $`|b`$, eq. (107) becomes a stochastic estimate for $`m\chi _\pi =\overline{\psi }\psi `$, which is the familiar Gellmann-Oakes-Renner relation. Eq. (107) holds configuration by configuration for any chiral state $`|b`$. For finite fifth dimensional extent $`L_s`$, the relation is broken and a useful measure is the ratio
$`R=m{\displaystyle \frac{b|(\gamma _5\stackrel{~}{D}_{tov}^1(m))^2|b}{b|\stackrel{~}{D}_{tov}^1(m)|b}}.`$ (109)
This ratio has been studied for domain wall fermions where in quenched SU(3) configurations at $`\beta =5.85`$ and $`5.7`$ significant deviations of the ratio from one are seen for fermion masses on the order $`m=0.01`$ and smaller. We argue that the larger violations seen at $`\beta =5.7`$ are due to the larger ensemble averaged density of zero eigenvalues of $`H_w(M)`$ at $`\beta =5.7`$ compared to $`\beta =5.85`$ .
We show in Fig. 5 a plot of the ratio $`R`$ for the domain wall operator as a function of the fermion mass at a fixed $`M=0.5`$ and $`a_5=1`$ with the instanton background used in Fig. 1. The mass $`M=0.5`$ is just before the zero crossing of $`H_w(M)`$. Without projection, the deviation is large even for $`L_s=32`$. With $`5`$ eigenvalues projected, the deviation from unity is negligible for $`L_s=16`$. For $`L_s=8`$, the violation is large with $`5`$ eigenvectors projected out and is not noticeably changed with $`20`$ eigenvectors projected out. This is explained by noticing in the right panel of Fig. 3 that for $`L_s=8`$ the fifth eigenvalue shows a deviation of $`1|\epsilon _4(\lambda _5)|0.01`$. What is not shown is that this value changes little at the 20th eigenvalue. What we are seeing is that by the 20th eigenvalue we are entering into a relatively dense band of eigenvalues of $`H_T`$ and they change slowly with increasing eigenvalue. Therefore, little is gained by projecting out more eigenvectors. By going to $`L_s=16`$, $`1|\epsilon _8(\lambda _5)|5\times 10^5`$ is much reduced compared to $`L_s=8`$, and hence all the volume modes contributing to $`R`$ have deviation less (or much less) than the worst case resulting in a small overall deviation of $`R`$ from unity. Therefore, projection is quite effective at restoring chiral symmetry, and we see that the GMOR test given by eq. (109) is a quite sensitive measure of chiral symmetry breaking since it detected a deviation in the $`ϵ`$ approximation on order of 1%.
We finish up the discussion of the smooth field case by considering the case of $`M>2`$ (the doubler region). We continue the discussion initiated in Ref. on the difference of the overlap Dirac operator $`D_{ov}(m;H_w)`$ and the 4-d version of the standard domain wall operator $`D_{ov}(m;H_T)`$.
We consider the free propagator with momentum near one of the corners of the Brillouin zone, i.e., with $`\overline{p}_\mu =\mathrm{sin}(p_\mu )1`$ for all $`\mu `$. Set $`M(p)=MB(p)=M_\mu 2\mathrm{sin}^2(p_\mu /2)=M2n+𝒪(\overline{p}_{\mu }^{}{}_{}{}^{2})`$, where $`n`$ is the number of momentum components near $`\pi `$. Then $`D_w(M)=i_\mu \gamma _\mu \overline{p}_\mu M(p)`$ and we find for the inverse free overlap propagator, up to terms of higher order in $`\overline{p}_\mu `$ and $`m`$,
$`\stackrel{~}{D}_{ov}^1(p,m)=\{\begin{array}{cc}2M(p)\times \frac{i_\mu \gamma _\mu \overline{p}_\mu +2mM(p)}{_\mu \overline{p}_{\mu }^{}{}_{}{}^{2}+(2mM(p))^2}\hfill & \text{for }M(p)>0\hfill \\ \frac{i_\mu \gamma _\mu \overline{p}_\mu }{2|M(p)|}\hfill & \text{for }M(p)<0.\hfill \end{array}`$ (110)
This corresponds to a light free fermion propagator only for $`M(p)M2n>0`$. For $`2<M<4`$, for example, the origin and the corners of the Brillouin zone with one momentum component close to $`\pi `$ give light fermions, while for $`0<M<2`$ only the origin gives a light fermion.
For free domain wall fermions, recalling the form of $`H_T`$ from eq. (26),
$`H_T(M)=\gamma _5D_w(M){\displaystyle \frac{1}{2+a_5D_w(M)}}=\gamma _5D_w(M){\displaystyle \frac{1}{a_5D_w(\frac{2}{a_5}M)}}.`$ (111)
we find
$`D_T(M)=\gamma _5H_T(M)=C(p)\left[i{\displaystyle \underset{\mu }{}}\gamma _\mu \overline{p}_\mu {\displaystyle \frac{1}{2}}M(p)\left(2a_5M(p)\right)\right]`$ (112)
with $`C(p)`$ positive, and hence unimportant in $`ϵ(H_T)`$. We see that for free domain wall fermions $`\frac{1}{2}M(p)\left(2a_5M(p)\right)`$ plays the role of $`M(p)`$ for overlap fermions. Hence, from eq. (110) we see that we need $`M(p)\left(2a_5M(p)\right)>0`$ to have a light domain wall fermion. These are exactly the conditions for a normalizable domain wall zero mode , $`1<1a_5M(p)<1`$. The first inequality here is an additional condition for domain wall fermions, not present for overlap fermions. For $`2<M<4`$ it excludes the origin of the Brillouin zone from giving a light fermion, as it did for overlap fermions.
Now consider the interacting case. We know that zero crossings of $`H_w(M)=0`$ also correspond to zero crossings of the domain wall $`H_T(M)`$ for all $`a_5`$. However, eq. (111) suggests that $`H_T(M)`$ can have poles. Zero eigenvalues of $`H_w(M)`$ at some $`M_0`$ are also zero eigenvalues of $`D_w(M_0)`$. From (111) we see that to each such zero eigenvalue is associated a pole of $`H_T(M)`$ at $`M=M_0+2/a_5`$. In other words, all zeros below $`M=2`$ are replicated as poles at $`M+2/a_5`$.
For the single instanton background shown in Fig. 1, there is one downward crossing in the Wilson spectral flow at about $`M=0.55`$. After $`M=0.55`$, the domain wall and overlap Dirac operators have index $`1`$ and hence a single zero mode. For this one downward crossing, there are 4 corresponding upward crossings at $`M=2.17`$ and $`M=2.30`$. After this last crossing, the index of the Dirac operator is $`+3`$ and hence there are 3 zero modes. There are no other zeros of $`H_w(M)`$ and hence $`H_T(M)`$ until $`M4`$. The largest eigenvalue of $`H_T(M)`$ has a behavior as follows: Starting at $`M=0`$, there is a near pairing of positive and negative eigenvalues of the largest eigenvalue $`\lambda _{\mathrm{max}}`$; however, a negative eigenvalue is slightly largest in magnitude. As $`M`$ increases, so does $`|\lambda _{\mathrm{max}}|`$. At $`M=1.249`$ $`\lambda _{\mathrm{max}}`$ reaches $`1`$, giving a zero eigenvalue for $`T^1`$. This corresponds to $`\stackrel{~}{B}`$ in the transfer matrix $`T`$, eqs. (23) and (24), having a zero eigenvalue. In the free field case the domain wall mass at which $`T^1`$ has a zero eigenvalue (namely $`M=1`$) is identified as the mass where the domain walls (each end of the fifth dimension) have the least coupling for fixed $`L_s`$, i.e. this is the optimal domain wall mass. As $`M`$ increases to the pole position of $`M=2.55`$, the eigenvalue $`\lambda _{\mathrm{max}}`$ diverges. The pole is the mechanism by which the topology changes for as $`M`$ increases beyond the pole mass the largest eigenvalue flips sign and decreases from positive infinity. There is a net increase in the number of positive states of $`H_T(M)`$ occurring in $`ϵ(H_T(M))`$ for the 4-d domain wall Dirac operator $`D_{ov}(0;H_T(M))`$ and the index becomes $`+4`$ with four zero modes.
In general then for the standard domain wall action in the smooth field case with $`a_5=1`$, for every crossing of $`H_w(M_0)`$ for $`M_0<2`$, there are four opposite crossings around $`M2`$ and a pole of $`H_T(M)`$ at $`M=2+M_0`$. If at some mass $`M<2`$ there is an index $`Q(M)`$ then $`Q(2+M)=4Q(M)`$. For the overlap Dirac operator $`D_{ov}(0;H_w(M))`$, however, $`Q(2+M)=3Q(M)`$.
Varying the fifth dimension lattice spacing $`a_5`$ gives some freedom in improving the chiral symmetry properties for the induced 4-d Dirac operators at finite $`L_s`$ given in eq. (105). If at some positive $`\lambda _{}<1`$ we have $`1\epsilon _{L_s/2}(a_5\lambda _{})<\delta `$ where $`\delta `$ is some prescribed accuracy, then the range of accuracy is $`|\lambda |`$ in $`[\lambda _{},1/(a_5^2\lambda _{})]`$. For the overlap Dirac operator with $`H_w(M)`$, the flexibility of rescaling by a gauge field dependent $`a_5`$ coupled with projection helped insure that all unprojected eigenvalues of $`H_w(M)`$ where within some prescribed range of accuracy of $`|\epsilon _{L_s/2}(a_5\lambda _i)|`$ approximating $`1`$.
For the standard domain wall operator, $`H_T(M)`$ depends implicitly on $`a_5`$, and choosing different values for $`a_5`$ leads to different cut-off effects, just as choosing different values for $`M`$ does. Therefore, both $`M`$ and $`a_5`$ need to be kept fixed in a given simulation. We saw in Fig. 3 that even for moderate $`L_s`$ the largest eigenvalue $`\lambda _{\mathrm{max}}`$ of $`H_T(M)`$ was usually well within the good approximation region of $`ϵ(a_5\lambda )`$. Thus we would like to use an $`a_5>1`$ to fully use the region of good approximation to $`ϵ`$. However, this can mix in doubler contributions, since for each crossing at $`M_0<2`$ there is the doubler pole of $`H_T(M)`$ at $`M=2/a_5+M_0`$, and as $`a_5`$ increases the pole moves down to lower domain wall mass $`M`$. However, as can be seen in the right panel of Fig. 3 the deviation for the largest eigenvalue of $`H_T`$ becomes sizeable well before the pole at $`M=2/a_5+M_0`$ is reached. This region of large deviation moves down with increasing $`a_5`$ and begins to conflict with our goal of increasing $`a_5`$ to fully utilize the region of good approximation to $`ϵ`$. Since one needs to choose $`M`$ sufficiently larger than the $`M_c`$ around which the physically important crossings, due to large instantons, occur, the choice of $`a_5`$ is quite limited. This is a fundamental limitation of the standard domain wall operator.
### B Quenched gauge field
We now move on to the interesting case of a thermalized quenched SU(3) gauge field configuration. Fig. 6 shows the spectral flow of the hermitian domain wall and Wilson operators for a quenched SU(3) $`\beta =5.85`$, $`8^3\times 32`$ configuration. Near $`M=1.0`$, the Wilson flow has two crossings down and two crossings up that result in a zero index for $`M>1.017`$. There is a fifth crossing at about $`M_0=1.86`$ resulting in an index of $`1`$ for $`M`$ greater than this crossing mass. In the region around $`M=1.0`$ where there is a non-zero index, the domain wall eigenvalues do not go close to zero. Well into the region $`M>1.017`$ where the index is zero there are fairly small domain wall eigenvalues, but no dramatic drop of the lowest eigenvalue is observed for $`M>M_0`$. The overlap Dirac operator, on the other hand, has zero modes whenever the index is non-zero. Fig. 7 shows the spectral flow of the projected domain wall operator on the same configuration as in the Fig. 6, together with the Wilson flow and the low-lying eigenvalues of $`2H_T`$. Near $`M=1.0`$ and after the Wilson crossing around $`1.85`$ there are small (zero) eigenvalues of the projected domain wall operator.
In the left panel of Fig. 8 is shown the spectral flow of the lowest eigenvalue of the standard domain wall operator for $`L_s=10`$ and $`30`$ and the Boriçi variant for $`L_s=10`$ all for $`a_5=1`$. We see in Fig. 8 that a small (near) zero eigenvalue is obtained for $`L_s=30`$ with the standard domain wall $`D_{DW}^{(5)}(0)`$. For $`L_s=10`$, a drop is seen in the lowest eigenvalue for $`M>M_0`$, down to the level of the $`L_s=30`$ eigenvalue. For the Boriçi variant $`D_{DW^{}}^{(5)}(0)`$, a relatively small eigenvalue is seen even for $`L_s=10`$.
The explanation for the small eigenvalues of the Boriçi variant even at $`L_s=10`$ can be seen in the right panel of Fig. 8 which shows the accuracy of the $`|\epsilon _{L_s/2}(\lambda _i)|`$ approximating $`1`$. The lowest eigenvalues of $`H_T(M)`$ are about half those of $`H_w(M)`$. Hence for the domain wall fermion action with $`L_s=10`$ and $`M=1.65`$ we find $`1|\epsilon _5(\lambda _{20})|0.25`$ while for the Boriçi variant using $`H_w(M)`$ the deviation is about $`0.3\%`$. The largest eigenvalue of $`H_T(1.65)`$ is about $`1.13`$ hence the deviation is below the range in the graph even for $`L_s=10`$. For the Boriçi variant, the largest eigenvalue of $`H_w(1.65)`$ is $`5.82`$ and there is about $`6\%`$ deviation. While the condition number is smaller for $`H_T(1.65)`$ than $`H_w(1.65)`$, the eigenvalues are placed more optimally around one for $`H_w`$ which results in overall smaller deviations of $`|\epsilon _{L_s/2}(\lambda _i)|`$ from $`1`$. For this configuration, there is not much need for rescaling the Wilson $`H_w`$ eigenvalues by $`a_5`$ since $`a_5=1`$ is close to optimal. For the standard domain wall action, a larger $`a_5`$ would lower the deviations of $`|\epsilon _{L_s/2}(a_5\lambda _i)|`$ from 1 for this configuration. However, as explained in the previous sub-section, choosing an $`a_5>1`$ might not be desirable, since it could lead to large deviations (due to the proximity of poles) for other configurations in a simulation.
Fig. 8 shows that small eigenvalues of the domain wall fermion actions $`D_{DW}^{(5)}`$ and $`D_{DW^{}}^{(5)}`$ can be easily achieved using projection and the cost of the projection is justified. To achieve the same small eigenvalue without projection would require $`L_s30`$.
In Fig. 9 is shown the ratio $`R`$ defined from the generalized Gellmann-Oakes-Renner relation (GMOR) in (109) using $`M=1.65`$, $`a_5=1`$ and the same configuration as in Fig. 8. For the standard domain wall action, not much is gained from projection for $`L_s=10`$ and violations of the relation are large. For $`L_s=30`$, projection significantly improves chiral symmetry resulting in small violations of the relation while there are large deviations for even this large $`L_s`$ without projection. The Boriçi variant also shows similar improvement with increasing $`L_s`$.
The condition number of a fermionic operator is a useful measure of its convergence properties in solving linear systems of equations, and can be used to construct a bound on the number of iterations for convergence in some iterative methods. The condition number defined here is $`\lambda _{\mathrm{max}}/\lambda _{\mathrm{min}}`$ where $`\lambda _i0`$ is defined from $`D^{}Dv_i=\lambda _i^2v_i`$. The condition numbers of the various fermion operators are shown in Fig. 10. The parameters are chosen the same as for the GMOR test in Fig. 9. The condition number of the overlap Dirac operator $`D_{ov}(m;H_w)`$ is much smaller than the 5-d methods. The condition number of the standard domain wall operator $`D_{DW}^{(5)}`$ is larger for $`L_s=30`$ compared to $`L_s=10`$ while the Boriçi variant is roughly the same for both $`L_s`$ values, but higher than that of $`D_{DW}^{(5)}`$. Projection also does not affect the condition numbers much basically because for this configuration at mass $`M=1.65`$, there should be no zero modes and the smallest eigenvalue of $`D^{(5)}`$ is not strongly affected by projection as can be seen in Fig. 8. The reason the 4-d operator has much smaller condition number is due to the largest eigenvalue which for $`D_{ov}(m)`$ is very close to $`1`$, while for $`D^{(5)}`$ it is roughly $`10`$.
Performance tests and comparisons of the various fermion actions are shown in Fig. 11. The left panel shows the iteration count for the benchmark Conjugate-Gradient for Normal Equations (CGNE) algorithm to achieve a residual accuracy of $`10^5`$ normalized by the source norm for the linear system $`D^{}D\varphi =D^{}b`$. The configuration parameters are those of the GMOR test in Fig. 9. The average and the standard error of the data (not the mean) is derived from the $`12`$ source color and spin inversions necessary to compute a fermion propagator. For the 4-d case (exact chiral symmetry not necessary), we note that a single inversion could be used for all the fermion masses $`m`$ used in Fig. 11 with the convergence governed by the time required for the smallest fermion mass. The 4-d overlap Dirac operator $`D_{ov}(m)`$ has a low iteration count requiring about $`320`$ CG iterations for convergence at $`m=0`$. Without preconditioning, the standard domain wall operator $`D_{DW}^{(5)}(m)`$ requires roughly 1000 iterations for convergence at $`L_s=10`$ and there is a strong $`L_s`$ dependence as expected from the condition number. Preconditioning reduces the condition numbers by about a factor of three. The Boriçi variant requires correspondingly more iterations for convergence. All these results are consistent with a linear scaling in the condition number.
It has been pointed out that both CGNE and Conjugate-Residual are optimal algorithms for the overlap-Dirac operator with exact chiral symmetry. In principle the Conjugate-Residual method would be more desirable since only one $`D_{ov}(m)`$ application is needed per iteration. For CGNE one expects two applications of $`D_{ov}(m)`$, but this can be rewritten as only one application of $`ϵ(H_w)`$ on a vector , so the work is the same per iteration for both methods. However, it was found that Conjugate-Residual had clearly worse convergence in practice compared to CGNE and was not used.
The right panel of Fig. 11 shows the actual cost for convergence measured in total number of $`D_w`$ applications for each linear system solution. This includes the “inner” CG portion for the 4-d overlap Dirac operator. For each “outer” CG iteration, a single application of $`ϵ(H_w)`$ is needed on a vector . This inner CG iteration in turn involves multiplications of a vector with $`H_w^2`$. For the standard domain wall operator, there are $`2L_s`$ applications of $`D_w`$ needed per iteration. For the Boriçi variant, there are $`3L_s`$ applications of $`D_w`$ needed as can be seen from the form of $`D_{DW^{}}^{(5)}(m)`$ and $`(D_{DW^{}}^{(5)}(m))^{}`$ in eqs. (15) and (80).
The cost of the methods can be roughly characterized as follows: the Boriçi variant requires about twice as many $`D_w`$ applications for a given $`L_s`$ compared to the standard domain wall operator with and without projection. This can be expected from the slightly larger condition number (mainly from the largest eigenvalue) and the 50% overhead in $`D_w`$ applications. Preconditioning saves the domain wall operator roughly a factor of two to three for both $`L_s`$ values with and without projection. This decrease is mainly due to the reduction of the largest eigenvalue. The $`L_s=30`$ calculations require about three to four times as many $`D_w`$ applications as the $`L_s=10`$ calculations with and without projection. However, $`L_s=30`$ and projection is needed to pass the GMOR test successfully.
The inversion of $`(1\widehat{A}_{ee}P_{})`$ in eq. (99) requires only about three iterations of a Minimal Residual algorithm to reach an accuracy of $`10^7`$. This multiplicative overhead in the number of “outer” CG iterations, but not proportional to $`L_s`$, is included in the total cost of the preconditioned domain wall operator and results in about an 18% overhead for $`L_s=10`$ and a 5% overhead for $`L_s=30`$ — there is little $`L_s`$ dependence in the “inner” CG. However, for these parameters and configuration (topological index is zero) the projected preconditioned domain wall requires fewer CGNE iterations for convergence, and is less costly in $`D_w`$ applications for $`L_s=30`$ than the unprojected version.
For $`L_s=30`$ with preconditioning and projection, the standard domain wall operator is about three times less costly in $`D_w`$ applications than the 4-d overlap Dirac operator. However, given that with the 4-d overlap Dirac operator all the masses for this spectroscopic calculation can be computed simultaneously for a given color and spin source, the cost benefit of the two approaches is comparable.
## IV Discussion
The five dimensional domain wall fermion action provides a means whereby an effective chiral theory may be obtained. Upon integrating out extra modes (which may be interpreted as flavors) from this extra fifth dimension, a chiral theory is obtained in four dimensions. For finite fifth dimensional extent $`L_s`$, the form of this action in four dimensions is
$`D_{tov}(m;H_T)={\displaystyle \frac{1}{2}}\left[1+m+(1m)\gamma _5\epsilon _{L_s/2}(a_5H_T)\right]`$ (113)
with the auxiliary Hamiltonian $`H_T`$. When $`L_s\mathrm{}`$, $`\epsilon _{L_s/2}(a_5H_T)ϵ(H_T)`$ for $`a_5>0`$ and there is a chiral symmetry in the action. For finite $`L_s`$, chiral symmetry breaking is induced because $`\epsilon _{L_s/2}(a_5H_T)`$ deviates from one. These deviations occur whenever the range of eigenvalues of $`H_T`$ are outside the range of approximation of $`\epsilon _{L_s/2}(x)`$ to $`ϵ(x)`$. Deviations can occur for both large and small eigenvalues of $`H_T`$.
Section II shows how the domain wall operator and a variant can be straightforwardly modified to allow arbitrarily precise chiral symmetry at finite values of $`L_s`$. The domain wall operator as shown in eq. (15) has only two new entries corresponding to projection of eigenvalues of $`H_T`$ or $`H_w`$. When the $`\chi `$ basis is chosen for the fermions fields in the action, eq. (28), the projection term occurs only in the origin along the fifth dimensional line (the flavor index). With this special position, the projection term is unmodified by the successive integrations of the extra flavor fermion fields. After extraction of the pseudo-fermion term, the expected form of the 4-d fermion action is obtained with the projection term acting as (part) of the spectral decomposition of $`ϵ(a_5H)`$. In principle then one can see that once in the $`\chi `$ basis the projection term could account for the entire spectrum of $`H`$ and the fifth dimensional extent could be set to one!
All the normal relations involving the fermion propagator go through with projection, and we have the relation connecting the 5-d and 4-d operators in eq. (73) and the propagator in eq. (62) valid with and without projection.
The main purpose of this paper is to show how projection can be used to achieve exact chiral symmetry at finite fifth dimensional extent. The method is not limited to quenched calculations and methods were outlined in Section II F for use in dynamical fermion calculations. Even-odd preconditioning is practical and for moderate $`L_s`$ is a small overhead. Tests of chiral symmetry were made for a smooth gauged field background composed of a single instanton, and for a quenched SU(3) gauge field background. Small (zero) eigenvalues of the Dirac operator can be achieved given careful control of the range of the approximation of the eigenvalues of the auxiliary Hamiltonian and choice of $`L_s`$. A stringent simple test is the generalized Gellmann-Oakes-Renner relation (GMOR) from eq. (107) and projection appears essential to pass this test.
However, for practical simulations the requirement of no induced chiral symmetry breaking may not be necessary or easy to achieve. In this case, the amount of projection and the fifth dimensional length can be tuned. This may well be necessary anyway since as shown for quenched backgrounds, there is evidence for a non-zero density of zero eigenvalues $`\rho (0;H)`$ of the auxiliary Hamiltonian $`H_w`$ (and hence $`H_T`$) for a large range of $`\beta `$. If one wants to go to very large lattice volume, the number of eigenvalues for projection to achieve some given accuracy for $`\epsilon _{L_s/2}(a_5H)`$ grows like the 4-d lattice volume. For example, this cost is an additional overhead for a spectroscopy calculation which also grows like the 4-d lattice volume. One way to lessen the overhead is to use a weaker coupling since $`\rho (0;H)`$ decreases very rapidly with the coupling which of course may require a larger lattice volume to hold the physical volume fixed. To keep calculation costs down, $`L_s`$ can be chosen suitably small and measures of induced chiral symmetry breaking like PCAC can be used. No effort was made to quantify this for given parameters like $`\beta `$, however, some projection is expected to reduce the induced chiral symmetry breaking from finite $`L_s`$.
For a given calculation, one would like to know what ultimately is the most efficient method to use. Namely, is the 5-d or the direct 4-d method the most efficient? From the present work, when one enforces chiral symmetry there is not really a huge difference in the cost of any of the methods tested. Quite likely then the answer depends on what one wants to do and on ones taste. For eigenvalue calculations of spectral quantities, the 4-d eigenvalues are needed. For spectroscopy or dynamical fermion calculations, either 5-d or 4-d methods can be used. For fixed cutoff dependence, the 5-d domain wall operator is preferable to the 4-d methods using $`ϵ(a_5H_T)`$ – the nontrivial cost of applying $`H_T`$ within an inner CG iteration inside of another inner CG iteration is prohibitive. For 4-d eigenvalue calculations the application of $`D_{ov}(m;H_T)`$ from $`D_{DW}^{(5)}(m)`$ in eq. (73) is probably most efficient. For the overlap Dirac case using $`D_{ov}(m;H_w)`$, the Boriçi variant of $`D_{DW^{}}^{(5)}(m)`$ does not appear competitive for spectroscopy calculations. A multi-grid implementation in the fifth direction might help , but has not been tested here. Other 5-d methods should be tried and tested .
For a direct comparison of the standard domain wall action $`D_{DW}^{(5)}(m)`$ and the overlap operator $`D_{ov}(m)`$, one first must answer just how important is chiral symmetry? Projection is essential for realizing good chiral symmetry properties for small quark masses and moderate $`L_s`$. Clearly computing eigenvectors using $`H_w`$ is much more efficient than computing eigenvectors of $`H_T`$. If one enforces chiral symmetry the un-preconditioned domain wall operators is about the same cost in the number of $`D_w`$ applications as the overlap operator. Preconditioning saves about a factor of three in cost, however the methods are comparable in cost for a spectroscopy calculation when multiple fermion masses are needed. A benefit of the 5-d method is that it gives a simple handle $`L_s`$ in which to reduce the overall cost while incurring some acceptably small amount of induced chiral symmetry breaking.
RGE was supported by DOE contract DE-AC05-84ER40150 under which the Southeastern Universities Research Association (SURA) operates the Thomas Jefferson National Accelerator Facility (TJNAF). UMH was supported in part by DOE contract DE-FG05-96ER40979. Computations were performed on FSU’s QCDSP, operated at TJNAF, and on the HPC workstation cluster at TJNAF and the workstation cluster at CSIT. We thank Rajamani Narayanan for discussion, and UMH would like to thank Artan Boriçi for a useful conversation.
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# Weak lensing surveys and the intrinsic correlation of galaxy ellipticities
## 1. Introduction
The large–scale mass distribution in the Universe is expected, through gravitational lensing, to imprint itself on the pattern of ellipticities measured from background galaxies (e.g., Blandford et al. 1991; Miralda-Escudé 1991). The angular correlations of such ellipticities, or the variance in ellipticities averaged in angular cells (amongst other statistics) can be compared to expectations for different cosmological models, and in principle can discriminate between them (e.g., Blandford et al. 1991; Miralda-Escudé 1991; Kaiser 1992; Jain & Seljak 1997). Since lensing is induced by the foreground mass distribution only, it provides the most direct method of studying the structure of mass in the Universe on large scales. Accordingly, detecting the shear signal induced by large scale structure has been the subject of a great deal of observational effort. Recently, four separate groups have reported detections of this cosmic shear, at levels comparable to that expected from currently popular models of structure formation (Van Waerbeke et al. 2000 \[hereafter VW\]; Bacon, Refregier & Ellis 2000; Wittman et al. 2000; Kaiser, Wilson & Luppino 2000 \[hereafter KWL\]).
In attributing the observed correlations of ellipticities, or the shear variance, to large–scale structure, an important assumption is that the sample of background galaxies used contains no intrinsic correlation of ellipticities. If such a correlation were present in the sample of lensed sources, it could be attributed to lensing, and may enhance the detected signal. The argument for discounting this possibility is that a pair of galaxies separated by a small distance on the sky are nonetheless on average separated by a large distance along the line of sight. If a particular pair of galaxies are separated by a large distance, there is no good theoretical reason to expect their ellipticites to be intrinsicly correlated.
However, the angular correlation of ellipticites predicted to be due to lensing is quite small — on the order of $`10^4`$ out to scales of several arcminutes. To detect this correlation against the random variations of galaxy ellipticities, observers take deep images and use large samples of background galaxies, perhaps $`10^5`$ per square degree. Some of these galaxies can be expected to be close not only in projection, but in real space as well. To what degree should we expect correlations in the actual ellipticities of nearby galaxies, and how much would the projection of such correlations add to any observed lensing signal?
That such intrinsic correlations in galaxy ellipticities may exist is not implausible. If elongation by local tidal fields contributed significantly to galaxy ellipticities, then nearby galaxies could be expected to sample the tidal field in the same fashion, producing similar elongations. Alternately, if some elongation originating from a galaxy’s last merger were to survive for a time comparable to the characteristic merger timescale for its environment, then one might expect galaxies to be preferentially aligned along the local large–scale structure, and thus similarly to each other. On larger scales, such elongation in cosmic structures appears to be present. For instance, the effects of large scale structure on the shapes of galaxy clusters have been the subject of much study. Cosmological N–body simulations have suggested that clusters tend to be oriented towards neighboring clusters or in directions defined by adjoining filaments and the merging subclusters which drain along them. (Dekel, West & Aarseth 1984; West, Dekel & Oemler 1989; West, Villumsen & Dekel 1991; van Haarlem & van de Weygaert 1993; Splinter et al. 1997; de Theije, van Kampen & Slijkhuis 1998). Observations have typically indicated the presence of such alignments, either towards nearby clusters (Binggeli 1982; Flin 1987; West 1989a,b; Rhee, van Haarlem & Katgert 1992; Plionis 1994) or towards nearby large–scale structure in the galaxy distribution (Argyres et al. 1986; Lambas, Groth & Peebles 1988); although not all studies support the presence of such alignments (Struble & Peebles 1985; Ulmer, McMillan & Kowalski 1989; Fong, Stevenson & Shanks 1990). The existence of correlations in the alignment of large–scale structures appears quite possible; perhaps similar intrinsic correlations in alignment exist on galactic scales.
Theoretical expectations for the degree of correlation of intrinsic ellipticities can in principle be derived. The local gravitational shear can be expected to either align the intrinsic angular momentum of nearby galaxies (Lee & Pen 2000), or to similarly deform neighboring, non–rotating galaxies through tidal distortion (Ciotti & Dutta 1994). Thus, the statistics of the local tidal field can be related to the statistics of galaxy angular momenta (Catelan & Theuns 1996a,b; Catelan & Theuns 1997; Sugerman, Summers & Kamionkowski 1999); and therefore to the intrinsic correlations in galaxy ellipticities (Coutts 1996; Lee & Pen 2000; Catelan et al. 2000, in preparation; Mackey and White 2000, in preparation).
There have been numerous attempts to detect intrinsic correlations in galaxy alignments using low redshift samples; the picture painted by this work is unclear, as we can see from the following sample. Flin (1988) considered a sample of 118 galaxies in the Perseus supercluster and found that the spin axes of these galaxies were aligned with the supercluster plane. Muriel & Lambas (1992) reported a correlation of alignments seen with spirals taken from the ESO catalog and analyzed in three dimensions; when only projected data was considered, the correlation was no longer present. Garrido et al. (1993) analyzed a sample covering a large area of sky in the northern hemisphere and claimed to find no evidence for correlations in alignment except within the Coma supercluster. Han, Gould & Sackett (1995) examined the spins of 60 galaxies in the Ursa Major filament and found no evidence for any alignment of spins. Cabanela & Aldering (1998) considered galaxy shapes extracted from a survey of Perseus–Pisces conducted using an automated plate scanner; statistically significant and color dependent correlations of galaxy ellipticities were found. On the other hand, Cabanela & Dickey (1999) used HI observations to determine the spins of 54 galaxies in the Perseus–Pisces supercluster; and found no evidence for preferential alignments of spin vectors. At this time evidence favors an orientation alignment between cD galaxies and the major axis of their parent cluster; but the presence or absence of any other galaxy shape correlations remains undetermined.
In this paper, we use Nbody simulations to make theoretical predictions for the correlation of intrinsic galaxy ellipticities. For simplicity, we work with directly with the projected ellipticities of the simulated dark matter halos, without making assuming any model for the way galaxies form within them. The significance of our results will therefore be entirely dependent on whether galaxy ellipticities behave significantly differently from their halos, a problem we leave to future gasdynamical simulations.
The layout of the paper is as follows. In §2, we describe the N–body dataset used, our halo catalog, and our measurement of projected ellipticities from the halos. We measure the three dimensional correlation functions of projected ellipticities in §3, and in §4 we decribe the construction of simulated surveys from our halo catalogs, with a geometry designed to mimic weak lensing observations. In §5, we project the three dimensional correlation functions into angular statistics, including the shear variance. We also compute these angular measures directly from our simulated surveys (as a consistency check). We compare our results to current observational data in §6, before discussing and summarizing our results in §7.
## 2. Simulated halos
### 2.1. Nbody simulations
Our requirements are that the simulation volume be large enough to capture the large scale tidal field that may cause correlations to arise between halo shapes, while at the same time having enough mass resolution to follow the formation of galaxy sized halos with a reasonable number of particles. We use outputs from Nbody simulations run by the Virgo Consortium (see e.g., the author list of Jenkins et al. 1998 for Virgo members) and which they have generously made public. The simulations are part of a set of different models, although we only use one here, the currently favoured cosmological constant-dominated cold dark matter ($`\mathrm{\Lambda }`$CDM) model. The parameters of this model are as follows, $`\mathrm{\Omega }_m=0.3`$ (the matter density at $`z=0`$), $`\mathrm{\Omega }_\mathrm{\Lambda }=0.7`$ (the contribution of $`\mathrm{\Lambda }`$ to the critical density at z=0), $`\mathrm{\Gamma }=0.21`$, which is the shape parameter of the initial power spectrum. The normalization was set so that $`\sigma _8=0.9`$, (the rms matter fluctuations in $`8h^1\mathrm{Mpc}`$ spheres extrapolated to $`z=0`$).
We use two different simulations, one run with a box size of $`141.3h^1\mathrm{Mpc}`$, and another in a box of size $`240h^1\mathrm{Mpc}`$, both with $`256^3`$ particles. The mass per particle in the former case is $`1.4\times 10^{10}\mathrm{h}^1\mathrm{M}_{}`$, and in the second it is 5 times larger. The Nbody code used was an adaptive particle-particle particle-mesh code (Couchman, Thomas and Pearce), and the gravitational softening length was $`30h^1\mathrm{kpc}`$ (for the smaller box). For more details, the reader is referred to the Virgo papers, including Kauffmann et al. (1999) for the smaller box simulation. Unless stated otherwise, all our analysis will be carried out using the higher resolution simulation, with the other being used as a check of the effects of lower resolution. We use only the $`z=1`$ output in each case, as we will be interested in comparing to lensing observations where the peak of the galaxy distribution is expected to lie close to this redshift.
### 2.2. Halo finding
Our analysis will be restricted to dark matter halos, and we will make no attempt to identify luminous galaxies in the simulations. In order to keep interpretation of our results simple, we will focus mostly on halos picked from the particle distribution using one easy to use groupfinder, the friends-of-friends (FOF) algorithm (e.g., Davis et al. 1985). A quick comparison will also be made with results from the publically available HOP groupfinding code (Eisenstein and Hut 1998).
In FOF, the one free parameter is the linking length as a fraction of the mean interparticle separation, $`b`$. This in a crude way governs the overdensity at the edge of the halo. We present some results using either $`b=0.1`$, or $`b=0.2`$. The latter is often used so that particles will be chosen from regions roughly inside the virial overdensity theshold of halos. With the former we will be picking out the denser (by approximately a factor of $`8`$) central parts of halos.
The HOP groupfinder works in a different way, using a smoothed overdensity field calculated at the positions of the particles. A group merging stage places subgroups together which fall inside a certain density contrast. The HOP algorithm has 6 free parameters; in this paper we use the fiducial values recommended by Eisenstein and Hut (1998).
In Figure 1, we plot the particles which ended up inside halos for each of the three groupfinding choices (FOF<sub>0.1</sub>, FOF<sub>0.2</sub>, and HOP). Only a small fraction of the box is shown, centred on an overdense region (the density inside the volume plotted, which measures $`15\times 15\times 40h^3\mathrm{Mpc}^3`$ is 1.6 times the cosmic mean). A lower mass cut of 20 particles ($`2.8\times 10^{11}\mathrm{h}^1\mathrm{M}_{}`$) has been applied to the halo list, which results in there being 20569 FOF<sub>0.1</sub> halos in the whole simulation volume, 42024 FOF<sub>0.2</sub> halos, and 18023 HOP halos. Of course, the FOF<sub>0.1</sub> halos will tend to coincide with the central regions of the other halos, so that if we define the total mass of a halo to be its virial mass, then this will be larger for the FOF<sub>0.1</sub> objects than the number of particles which we plot in Fig 1. The FOF<sub>0.1</sub> halos are interesting because they will offer a good check of the robustness of our results, when we ask below whether the ellipcities of the dense central regions are correlated in the same way as the outer parts of the halos.
### 2.3. Halo ellipticities
Although it is possible to measure three dimensional ellipticities for halos, we will work here only with quantities which are projected onto the plane of the sky, for ease of comparison with observational data.
We estimate the ellipcity components of each halo by measuring the second moments of the projected mass distribution in a way analogous to that used for the surface brightness distribution of galaxies (Valdes et al. 1983, Miralda-Escudé 1991):
$$e_1=\frac{I_{xx}I_{yy}}{I_{xx}+I_{yy}},e_2=\frac{2I_{xy}}{I_{xx}+I_{yy}},$$
(1)
where
$$I_{xx}=\frac{_{i=1}^N(x_i\overline{x})^2}{N},I_{yy}=\frac{_{i=1}^N(y_i\overline{y})^2}{N},$$
(2)
and
$$I_{xy}=\frac{_{i=1}^N(x_i\overline{x})(y_i\overline{y})}{N}.$$
(3)
Here, $`\overline{x}`$ and $`\overline{y}`$ are the coordinates of the center of mass of the halo in question, which contains $`N`$ particles of equal mass. In the observational case, these moments are often calculated using the surface brightness weighted by a function chosen to maximize the S/N of the measurement (see e.g., VW). In our case, as we have mentioned above, we will probe the sensitivity of our measurements to variations of this sort by comparing halos chosen with a large FOF linking length to those using a small value which effectively give a non-zero weight only to the dense central regions.
When working with pairwise statistics, such as the ellipticity correlation function (see §3), it will be advantageous to redefine $`e_1`$ and $`e_2`$ for each pair of halos, with the $`x`$-axis defined to be line joining the halo centres and the $`y`$-axis a line perpendicular to it.
The quantity $`(e_1,e_2)`$ defines a pseudovector, of length $`e=\sqrt{e_1^2+e_2^2}`$. A positive $`e_1`$ component indicates a stretching along the $`x`$axis, and a negative component a stretching along the $`y`$axis. The $`e_2`$ component is likewise a measure of the stretching along axes at $`45\mathrm{deg}`$ to the $`x`$ and $`y`$ axes. If we have an ellipse of ellipticity $`e`$ ($`e(1q^2)/(1+q^2)`$, where $`q`$ is the ratio of minor and major axis lengths), then
$$e_1=e\mathrm{cos}(2\beta ),e_2=e\mathrm{sin}(2\beta ),$$
(4)
where $`\beta `$ is the angle between the major axis and the $`x`$axis.
In the bottom three panels of Fig. 1, we plot the ellipticities of the halos chosen by the three groupfinders. For each halo, we show a bar with length proportional to $`e`$, oriented along the direction of the major axis. It is difficult to pick out by eye any tendency for halos to be aligned. One worry which we might have had concerns the tendency of close halos to be joined together by thin “necks” of particles and for the groupfinder to count such pairs as one halo. If these halos are distributed along filaments, then we might expect an artificial alignment to arise. This does not seem to be an obvious problem here. For example, if we look at halos picked out by the FOF<sub>0.1</sub> and FOF<sub>0.2</sub> groupfinders, the ellipticities and directions appear to be fairly similar even though different parts of the halo are being used to calculate them. We will check the statistical tendencies for alignments in the next section. The HOP halos do in several cases appear to be made from several small halos joined together. Presumably this could be alleviated by tuning the free parameters which govern the method, but we have not attempted to do this.
In Fig. 2, we show the probability distribution of $`e`$ values for the halos (again with $`>20`$ member particles) picked out by the three groupfinders. They are all fairly similar, with a mean $`e`$ around $`0.4`$, slightly higher than the $`0.3`$ typical of real galaxies (e.g., Wittman et al. 2000).
## 3. Ellipticity correlations in three dimensions
Observationally, ellipticity correlations are measured as a function of angular separation, on the plane of the sky (at least for weak lensing surveys). With distance information (for example, using redshifts), it would be possible to measure them as a function of separation in three dimensions, and this is what it is most natural to do using our Nbody simulations. In this section, we will do this, and later convert the measurements to angular correlations which can be compared to weak lensing survey results. This conversion will be done in two different ways. The first is an analytical projection of our three dimensional results using Limber’s equation (Limber, 1959). The second is a direct measurement from simulated surveys made by projecting the halo distributions in the box, and applying a radial selection function. In the present section, although we will be dealing with halo separations in three dimensions, it is worth bearing in mind that we restrict ourselves to quantities which can be measured directly observationally (albeit with redshifts), so that the ellipticities we will be correlating are projected ellipticities (defined in Equation 1).
Two ellipticity correlations can be defined (following Miralda-Escudé 1991) as
$$c_{11}(r)=e_1(𝐱)e_1(𝐱+𝐫),$$
(5)
and
$$c_{22}(r)=e_2(𝐱)e_2(𝐱+𝐫),$$
(6)
together with a cross-correlation,
$$c_{12}(r)=e_1(𝐱)e_2(𝐱+𝐫),$$
(7)
where $`𝐫`$ is the three dimensional vector of length $`r`$ joining each pair of halos. Here we have defined the ellipticity components, $`e_1`$ and $`e_2`$ with respect to axes which are the projection into the $`xy`$ plane of the line joining each pair of halos and a line orthogonal to it which lies also in the x-y plane.
In general, these functions may be anisotropic, and depend on the separation between pairs both across and along the line of sight. We will plot their full dependance on these quantities later. In our tests of different halo finders and simulation resolution, we will work with $`c_{11}(r)`$ only (although we have checked that our conclusions are valid for the other functions also). This is because $`c_{11}`$ has no strong angular dependence, so that plotting it as a solely a function of $`r`$ does not hide any crucial aspects.
We show $`c_{11}(r)`$ in Fig. 3, for the three different groupfinders. In the top panel, we give results for all halos with $`>20`$ member particles. The error bars are the equivalent of Poisson errors, calculated by randomly rotating the halo major axes and the calculating $`c_{11}(r)`$ for this randomized distribution (see e.g., VW). The dispersion between results for several such randomizations (we use 50) gives the error bar. In order to make best use of all the information contained in the simulation, we have averaged results from the three orthogonal projections of the halo particles.
The first thing to notice is that there is a good detection of a signal, even for the largest bin that we plot. The widest pairs we use have a separation of 0.25 times the box side-length, or $`35h^1\mathrm{Mpc}`$. This is not necessarily that surprising, given that large scale correlations in the matter distribution exist on such scales and beyond.
The FOF<sub>0.1</sub> and FOF<sub>0.2</sub> halos give very similar results. As mentioned earlier, this is a good test of the robustness of the correlation signal and a good argument that there are not systematic problems. It is telling us that the inner and outer parts of halos are responding in the same way, a stability which points to the correlation being a real effect. The HOP results are larger by a factor of $`2`$, and on small scales seem to deviate even more. We have already seen from Fig. 1 that with the free parameters set in the way we have chosen, this groupfinder seems to produce groups which are clumps of subgroups which the FOF algorithm keeps separate. This is probably an indication that for a more reasonable result, we would need to tune some of the six free parameters that govern HOP.
Some of the differences between $`c_{11}(r)`$ results might be due to the fact that the same absolute mass cut (20 particles) for all three groupfinders will result in different halos being chosen. For example, the FOF<sub>0.1</sub> halos will be much rarer than FOF<sub>0.2</sub>. A better way to compare results is to set the space density of halos to be equal. We have done this the bottom panel of Fig.3, where we show $`c_{11}(r)`$ for the most massive 10000 halos in the simulation volume. The lower mass limits for the different groupfinders in this case are 44, 90 and 85 particles ($`6.2\times 10^{11}\mathrm{h}^1\mathrm{M}_{}`$, $`1.3\times 10^{12}\mathrm{h}^1\mathrm{M}_{}`$ and $`1.2\times 10^{12}\mathrm{h}^1\mathrm{M}_{}`$), for FOF<sub>0.1</sub>, FOF<sub>0.2</sub> and HOP, respectively. The results are indeed more similar now, except for the HOP results at the smallest separations. We have also tried applying an upper mass cut, keeping only halos containing less than a certain number of particles. We find similar results (not plotted) to those for all halos.
For the rest of the paper, we have decided to use groups chosen by the FOF<sub>0.1</sub> groupfinder. This is because we are most interested in the inner parts of halos, which is where galaxies will tend to lie. Also, when halos are chosen to have the same space density (Fig. 3b), these halos have a slightly lower correlation than the others, so that we are being conservative.
One can ask about the effect of simulation resolution on our results, whether the simulations have converged and if not, whether lower resolutions yield higher or lower correlations. In Figure 4, we compare halos taken from the larger, low resolution simulation (240 $`h^1\mathrm{Mpc}`$ boxsize) to our previous results. We again use the FOF<sub>0.1</sub> groupfinder, and in order to compare halos with the same mass, we use halos containing $`>20`$ particles in the low resolution case and $`>100`$ in the high resolution case (these halos do in fact have the same space density, indicating that the simulations have converged as far as the mass function is concerned).
We can see that although the ellipticity correlations have not converged, the higher resolution simulation is more correlated. This may be because structure has been allowed to form on smaller scales, whereas if the the structure is unresolved, some of the tidal field which causes correlations to arise is missing. In any case, we can argue that the fact that higher resolution increases the correlated signal means that our results are likely to be conservative. Of course this situation is far from ideal, and one would really like to see if even higher resolution results do in fact converge. Such simulations are beyond the scope of this paper, however. This is particularily true because it is important not to sacrifice a cosmologically interesting volume for the sake of higher resolution, since long range tidal forces may play a significant role in setting up the correlations.
In Fig. 5, we show $`c_{11}`$ and $`c_{22}`$ as a function of halo pair separation across the line of sight ($`\sigma `$) and along the line of sight ($`\pi `$). This presentation is similar to that used in plotting redshift distortions of the matter correlation function (e.g., Kaiser 1987). Only $`c_{22}(\sigma ,\pi )`$ appears to be severely anisotropic. The plot seems to be showing us that projected halos tilt in opposite directions (-ve correlation) when they are side-by-side, but in the same direction when they lie one behind the other. The $`c_{11}(\sigma ,\pi )`$ function is nearly isotropic, and positive, so that pairs of halos have a tendency to stretch along the projection of the line joining them whatever their spatial orientation.
If $`c_{22}`$ is measured averaged in radial bins, $`r`$ (where $`r=\sqrt{\sigma ^2+\pi ^2}`$) then the signal from the bottom right partly cancels that in the top left and the result is a very small correlation function. On the other hand, a projection of $`c_{22}`$ on the plane of the sky (which we will do later when calculating the angular correlations) is equivalent to a projection along the $`\pi `$ axis. Because $`\pi `$ cannot be negative (it is defined as the modulus of the pair separation along the line of sight), then the contribution of the positive $`c_{22}`$ signal at large $`\pi `$ will result in a significant angular signal.
We do not plot the cross-correlation, $`c_{12}(\sigma ,\pi )`$ here, as we find it to be consistent with zero everywhere. We also note that in the observational case, plots like these (and the previous plots which were radially averaged) will be affected by redshift distortions of the pair separations.
## 4. Simulated surveys
One way of making predictions for the ellipticity correlations that can be compared directly to current observations is by using simulated surveys with a similar geometry to real surveys. In this section we describe our creation of such datasets from the Nbody halo catalogs. In §5 we measure angular statistics from the simulated surveys and compare to analytic projection of the three dimensional statistics of §3.
The redshift distribution of faint background galaxies that are used in current weak lensing studies is not known accurately. Most of the observational groups assume that the peak of the redshift histogram lies around $`z=1`$, so we will use this value in setting up our simulated surveys. We note that depending on the actual redshift of the galaxies, the relative contribution and angular scale of the lensing correlations and any intrinsic correlation will vary, something which can undoubtedly used to tell them apart. In our mock surveys, we will also assume that the galaxy population (we refer to halos as “galaxies” when describing the mock surveys) stays constant in comoving coordinates, so that we only make use of the $`z=1`$ simulation output. We leave the study of the effects of redshift evolution to future work. Also, this approach will allow us to compare directly with the analytic projection of the statistics.
Following what has now become a standard technique (e.g., Jain, Seljak & White 2000, Croft et al. 2000) we stack simulation boxes one behind the other, from $`z=0`$ to $`z=2`$. Each simulation box is subjected to a random recentering, and possibility of reflection about one of the three centres. We place each with one of the three axes (randomly chosen) pointing along the line of sight. This procedure is adopted to avoid periodic repetition of structures. When projecting the box contents, we work in comoving coordinates, and set the angular size of the simulated catalog to be one box width at $`z=1`$, where our $`z`$ histogram will peak. At $`z=1`$, an angular scale of $`1\mathrm{deg}.`$ therefore corresponds to $`40h^1\mathrm{Mpc}`$. We make use of the small angles involved by projecting onto a flat plane, rather than a curved sky. This facilitates things because the halo ellipticities have been defined in the same fashion.
The $`N(z)`$ form we choose for our mock surveys is a Gaussian distribution $`N(z)=e^{(z1)^2/(2\sigma _z^2)}`$ (for simplicity’s sake), where we have centered on $`z=1`$. We apply the selection function $`\psi (r)`$ which leads to such an $`N(z)`$ for a uniform distribution of galaxies, where $`\psi (r)`$ is the relative probability of including randomly chosen galaxy at a comoving distance $`r`$ (see e.g., Peebles 1993). We try 3 values for $`\sigma _z`$, $`0.4,0.2`$ and $`0.1`$. The widest of these is designed to mimic roughly the present observational catalogs which use no colour selection. The narrower redshift distributions could arise if some sort of photometric redshifts are used (narrower distributions still are possible, see e.g., Hogg et al. 1998).
It is necessary to project 26 box lengths to reach $`z=2`$. For galaxies past the peak in the $`N(z)`$, the mock survey is wider than the simulation box. At this point we make use of the periodic boundaries and make each layer of boxes periodic across the line of sight. This was also done by e.g., Seljak, Burwell & Pen 2000) and should not cause any problems. In any case, we will compare with the analytic projection of the statistics to make sure.
We use only the FOF<sub>0.1</sub> halos (with $`>20`$ particles) to make our simulated surveys. This results in a rather small number of galaxies per square degree, $`18000`$ (for $`\sigma _z=0.4`$), compared to e.g., $`10^5`$ for VW . The $`N(z)`$ for three surveys made using different value of $`\sigma _z`$ is plotted in Fig. 6. We set the selection function to zero for $`z<0.3`$ and $`z>2`$.
In Fig.7, we show the $`x`$ and $`y`$ components of the ellipticities on the plane of the sky, where we have binned the individual $`e_1`$ and $`e_2`$ values in square cells of side-length 2 arcmins. We plot as a line, the average ellipticity in each cell in a manner analogous to the individual halo ellipticities plotted in Fig 1. We also show, as a grayscale, the angular area density of halos smoothed with a Gaussian filter of FWHM 2 arcmins. Fig. 7a represents a survey with a wider $`z`$ distribution ($`\sigma _z=0.4`$) than Fig. 7b. The obvious difference between the two is that the wider distribution has washed out both the density inhomogeneities and the absolute value of the residual ellipticities. To the eye, there does not seem to be any obvious correlation between the ellipcitity vectors, although as we will see in the next section, there are detectable statistical correlations and a non-zero shear variance. We defer any exploration of the relation between the projected halo density and ellipticities to future work.
## 5. Angular statistics
### 5.1. Angular ellipticity correlations
The ellipticity correlations measured in three dimensions in §3 can be projected for direct comparison with results from angular galaxy surveys. This can be done simply using a modification of Limber’s equation (Limber 1953, Peebles 1993), so that the angular correlation function of $`e_1`$ components of pairs is
$$C_{11}(\theta )=\frac{(r_1r_2)^2dr_1dr_2\psi _1\psi _2(1+\xi (r_{12})c_{11}(\sigma ,\pi )}{[r^2𝑑r\psi ]^2+(r_1r_2)^2𝑑r_1𝑑r_2\psi _1\psi _2\xi (r_{12})},$$
(8)
where $`r_{12}`$ comoving is the separation between points at distance $`r_1`$ and $`r_2`$ from the observer. $`\sigma `$ and $`\pi `$ are the separations across and along the line of sight, which we calculate in the small angle limit, so that
$$\sigma =(\frac{r_1+r_2}{2})\theta ,\pi =(r_1r_2)(1\theta ^2/2),$$
(9)
and $`r_{12}=\sqrt{\sigma ^2+\pi ^2}`$. The selection functions at $`r_1`$ and $`r_2`$ are given by $`\psi _1`$ and $`\psi _2`$, and $`\xi (r)`$ is the halo autocorrelation function. The relations for the $`e_2`$ component correlation, $`C_{22}(\theta )`$ and the cross-correlation, $`C_{12}(\theta )`$ are analogous.
Using this formalism, we calculate $`C_{11}(\theta )`$, $`C_{22}(\theta )`$ and $`C_{12}(\theta )`$ for the three different selection functions also used in our simulated surveys, which yield a Gaussian $`N(z)`$. The results are plotted in Fig.8. We can see that for the first two of these a distinct correlation signal results. The behaviour as a function of redshift width of the galaxy distribution is as we might predict, with a wider distribution yielding less angular correlation. Also very interesting is the fact that $`C_{11}(\theta )`$ is everywhere positive, whereas $`C_{22}(\theta )`$ dips below zero on the largest scales. The $`C_{12}`$ function is much smaller, oscillating about zero, due to noise. These functional forms are very similar to the theoretical expectation for the weak lensing signal (see e.g., Miralda-Escudé 1991, Kaiser 1992). A simple explanation of this similarity is that the shear field which is responsible for the weak lensing correlations is mathematically similar to the shear field responsible for generating galaxy intrinsic shapes and spins from the surrounding large-scale structure, something which can be used to make fully analytic predictions of the effect, in the context of linear theory (Kamionkowski, private communication). We will return briefly to this in §7.
We have put points on Fig 8 which come from directly estimating the correlation functions from our simulated surveys of §4. We use as estimators the angular analogs of equations 5-7. The error bars have been computed from the error on the mean measured from 10 simulated surveys set up with different random seeds. As we use the same underlying Nbody simulation for all, this will result in an underestimate of the cosmic variance. This is adequate for our purposes, because we are comparing to the Limber’s equation projection of the three dimensional statistics which is not likely to be exact anyway, as it involves numerical intergration of a noisy function. This is particularily true of $`C_{22}(\sigma ,\pi )`$. That said, the agreement between the two sets of angular statistics is good, so that we can be fairly confident that there are no serious errors in the conversion of three dimensional statistics to angular ones.
### 5.2. Shear variance
The variance of the galaxy ellipticities binned into cells on the plane of the sky is a simple and interesting statistic to calculate. An estimator for the shear in a cell centred on angular position $`\varphi _𝐢`$ is (following VW, although note that we assume that the shear, $`\gamma =e/2`$, in the limit of intrinsically round galaxies, so that our definition is slightly different) is
$$E[\gamma ^2(\varphi _𝐢)]=[\frac{1}{2N}\underset{j=1}{\overset{N}{}}e_{1,j}]^2+[\frac{1}{2N}\underset{j=1}{\overset{N}{}}e_{2,j}]^2,$$
(10)
where there are $`N`$ galaxies in the cell and $`e_{1,j}`$ and $`e_{2,j}`$ are the ellipticity components of the $`j`$th galaxy, which have been defined with respect to some fixed axes. The shear variance is then $`E[\gamma ^2(\varphi _𝐢)]=(\sigma _e^2/4N)+\gamma ^2`$. Here $`(\sigma _e^2/4N)`$ is an ellipticity shot noise term, arising from the random intrinsic ellipticities of the finite number of galaxies involved, and $`\gamma ^2`$ is the quantity we are interested in.
In terms of the angular ellipticity correlation functions, $`\gamma ^2`$ is given by a double integral over the cell being used (as in the relation between $`\xi `$ and counts-in-cells variance, e.g., Peebles 1993), so that
$$\gamma ^2=\frac{1}{4A^2}C_{11}(\theta )+C_{22}(\theta )d^2A,$$
(11)
where $`A`$ is the cell area. We note that this integral will tend to be dominated by the largest $`\theta `$ scale, so that sensitivity to systematic errors on small scales will be relatively low.
We carry out this integration using our angular ellipticity correlations from Fig. 8. For easy comparison with our simulated survey results (see below), we use square cells of side length $`L`$. In Figure 9, we show the $`\gamma ^2`$ curves that result from the integral, as a function of angular scale $`\theta =L/\sqrt{\pi }`$. This $`x`$-axis scaling was chosen because some obseravtional groups (e.g., VW) calculate $`\gamma ^2`$ in round cells of radius $`\theta `$.
From the plot, we can see that $`\gamma ^2`$ rises as we move from large to small scales. The rms shear, $`\gamma ^2^{1/2}`$ at $`\theta =2`$ arcmins is $`0.6\%`$ for a wide redshift distribution ($`\sigma _z=0.4`$). This is significantly below the theoretical prediction for the rms shear due to weak lensing in different cosmological models, which can vary from $`1\%`$ to $`>3\%`$ (see e.g., VW). Because the contributions to this quantity add in quadrature, the intrinisic shear is likely only to matter for low amplitude models. Of course, as we have only simulated one (relatively high amplitude) model, it could be that the intrinisic shear for low amplitude cosmologies is also lower. For a narrower redshift distribution with $`\sigma _z=0.1`$, our value of $`\gamma ^2^{1/2}`$ doubles. On large scales we also see that there is still signal out to $`\theta =30`$ arcmins.
For comparison, we again show points representing the value of $`\gamma ^2`$ measured from the simulated surveys. We have used the estimator of equation 10, with square cells of side length $`L`$ ($`=\sqrt{\pi }\theta `$). To remove the shot noise term, we use the procedure advocated by VW. This involves randomly rotating the major axes of all the projected halos and calculating $`\gamma ^2`$, which will be a noisy version of the shot noise term. We carry out 1000 such randomizations and use the average $`\gamma ^2`$ measured from them as our $`\sigma _e^2/4N`$, which we subtract. The results are shown in Fig. 9, with error bars which are again the error on the mean from 10 simulated surveys. As with the ellipticity correlation functions, the agreement is good.
## 6. Comparison with observations
As mentioned in §1, there are now several different groups (VW; Bacon, Refregier & Ellis 2000; Wittman et al. 2000; KWL) producing results from large observational campaigns to measure weak lensing by large-scale structure. These different surveys (at the time the published results were submitted) involved from 0.5 to 1.8 square degrees of imaging data, and produced measurements of ellipticity correlations and shear on scales from $`0.6`$ to $`30`$ arcmins. In this section, we will briefly compare these current measurements with the predictions of §5.
Two of these groups (VW, and Wittman et al. 2000) and have published ellipticity correlation functions, using slightly different notation from ours. What we call $`C_{11}(\theta )`$, $`C_{22}(\theta )`$, are called $`e_t(0)e_t(\theta )`$, $`e_r(0)e_r(\theta )`$ by VW, and $`e_1e_1`$, $`e_2e_2`$ by Wittman et al. (2000). The cross correlation $`C_{12}(\theta )`$ is defined analogously.
We show these observational results in Fig. 10, together with those from the simulated surveys (from Fig. 8). The measurements from the two groups are mostly at different scales, but agree within the errors in the small degree of overlap around $`\theta =3.5`$ arcmins. The functional form of our intrinsic correlation follows roughly the observational results. The amplitude is however much lower, by a factor of roughly between 5 and 10 for most of the points, for the most relevant redshift distribution (a wide one, with $`\sigma _z=0.4`$). The Wittman $`\mathrm{et}\mathrm{al}.`$ points appear to be even higher than this on large scales for $`C_{22}`$. We note that in the Wittmann $`\mathrm{et}\mathrm{al}.`$ paper, these points are shown to also be rather higher compared to theoretical weak lensing predictions for $`\mathrm{\Lambda }`$CDM than the $`C_{11}`$ points. For a narrower redshift distribution, which may be achieved in the future with photometric redshifts, the intrinsic correlations are closer to the current observational data. The cross correlation, $`C_{12}`$, while noisy, is small, and consistent with zero in all cases.
If we move to the shear variance, $`\gamma ^2`$, we can plot results from all four groups. We have the plotted these in the same format as Figure 2 of KWL, as Fig.11. As Wittman et al. did not specifically publish $`\gamma ^2`$ values, we have converted their $`C_{11}`$ measurements using the formulae in Kaiser (1992) (as was done in KWL, to which the reader is referred for details), assuming a spectral index of mass fluctuations of $`n=1`$. One difference between this plot and that in KWL is that we plot the results as a function of $`\theta `$, the radius of a top-hat sphere. To do this, we have converted the square box side-length $`L`$ from KWL to the radius of a circle with the same area ($`\theta =L/\sqrt{\pi }`$; this is the same as was done in plotting Fig. 8) . The groups all have results which appear to be consistent within their errors, which is impressive and reassuring.
Turning to the simulation results, we concentrate first on the wide redshift distribution ($`\sigma _z=0.4`$), which is most relevant to the current observations. The intrinsic $`\gamma ^2`$ we predict is a small fraction of the measured value on all scales, except possibly for the largest scale point from KWL, which, as those authors point out is an interesting null detection. The intrinisic $`\gamma ^2`$ is also at a smaller level than the $`1\sigma `$ errors for the observational measurements. It is therefore not really a factor when interpreting the current results, but may become a bit more of an issue when the surveyed areas increase and the statistical errors become smaller. If we now turn to the other two curves, we can see that for the narrow redshift distributions, the intrinsic $`\gamma ^2`$ begins to approach the current measurements, at least on large scales. On small scales, although the errors are large (Fig. 9), there appears to be relatively less signal.
## 7. Summary and Discussion
We have used Nbody simulations to make predictions for intrinsic correlations of galaxy ellipticities, under the assumption that galaxy shapes follow the shapes of their dark matter halos. Measurements of the ellipticity correlation functions in three-dimensions give a distinctive signal, which we measure with relatively small uncertainties on scales from $`0.530h^1\mathrm{Mpc}`$. These correlations vary by less than a factor of $`2`$ for different halo finding techniques and different simulation resolutions. We project these three dimensional correlations into angular statistics, including the shear variance. We have done this both analytically, using a modified Limber’s equation, and by making direct measurements from simulated surveys constructed by projecting the simulation boxes. We find that the amplitude of the angular statistics depends strongly on the redshift width of the galaxy distribution. With widths appropriate to present day surveys, we find that the intrinsic correlations we predict are around $`1020\%`$ of the currently measured signal, somewhat smaller than the $`1\sigma `$ errors on the measurements.
Since the area of the sky surveyed for weak lensing is increasing rapidly, the intrinisic correlation may become detectable from these deep and wide surveys in the future. In any case, it seems to be worth bearing in mind that there could be this sort of contamination. In particular, one possible way of extracting more information from lensing which has received attention is the use of photometric redshift information, to break down the background galaxy distribution into a number of “screens”. This would enable tomography to be carried out (e.g., Hu 1999). We have seen however that the intrinsic correlation may be quite large for these narrow redshift bins, so that it might become comparable to the weak lensing signal (Fig. 11).
Of course, the extra information available in the form of photometric redshifts is likely to be very useful for deciding whether there is an intrinsic component, and if it exists, to separate it from the lensing signal. For example the cross-correlation of ellipticities (or co-variance of the shear) between different redshift bins can be compared to the correlation within bins, with only the later responding to intrinisic correlations. Something along these lines has already been carried out by KWL, albeit with two colour bands which both give wide redshift distributions (but with one deeper than the other). These authors find a higher shear signal for the deeper redshift sample, which is consistent with lensing, but in the wrong direction for intrinsic correlations. For the cross-correlation between samples, they do find slightly anomalous results, however.
Another way of trying to measure any intrinsic ellipticity correlations would be to stick to the local universe, and to measure the three dimensional correlation functions (§3) from a redshift survey. If there really is a signal like that plotted in Fig. 3, then this could be measurable from a relatively small survey (by todays standards), with a few thousand galaxies. Even without redshifts, one might expect to find a measurable intrinsic signal from a relatively nearby angular sample of galaxies, like the APM survey (Maddox et al. 1990), or Sloan Digital Sky Survey (Gunn & Weinberg 1995).
If there are in fact some measurable correlations between real galaxy ellipticities, then this can be understood in the framework of structure formation by gravitational instability, with the ellipticities being linked to the angular momenta of galaxies, which are in turn set up by tidal torques from the shear in the initial density field (e.g., Peebles 1969, Barnes & Efstathiou 1987, Catelan and Theuns 1996a,b). This may explain why the ellipticity correlation functions we measure have similar functional forms to those caused by weak-lensing: both are responding to a cosmic shear field. Detection of correlated ellipticities, if they exist, may be useful for the study of galaxy formation (e.g., Sugerman et al. 2000), or even cosmology (Lee & Pen 2000).
It is also likely that the signal due to the intrinsic correlation will give qualitatively and measurably different results from weak lensing for some statistics we have not considered here. For example, the probability distribution of the lensing convergence is predicted to have a measurable skewness, something which can be used to determine $`\mathrm{\Omega }`$ (Bernardeau et al. 1997). Measurements of this parameter from our simulated surveys by L. Van Waerbeke (private communication) yield a null result, the convergence pdf being consistent with a Gaussian distribution. The intrinisic correlations do not therefore appear to interfere with our ability to do cosmology in this way, and should not act as more than an additional source of noise (albeit correlated) when reconstructed mass maps are made.
On the simulation side, one important issue is the fact that our results have apparently not converged with resolution. Although we find that the higher resolution of two simulations gives more intrinsic correlations, it is possible that given even higher resolution, things will begin to go the other way. Clearly this needs to be tested somehow in the future. Also, perhaps most important of all, we have assumed a very simple relationship between projected halo ellipticities and projected galaxy ellipticities. It is possible that adding gas dynamics and star formation to simulations will result in their being no significant correlation between the two. The tests which we have carried out which have most bearing on this are the use of two sets of different friends of friends groups, which respond to ellipticities either of the whole halo, or just the dense central region. As we find results for the two which are very similar, this is at least some evidence that the intrinsic correlation may be fundamental.
As this paper was being completed, we became aware of similar work by Heavens et al. (2000). These authors use the angular momentum of Nbody halos (also from Virgo simulations, but only at the lower of the two resolutions) to predict the intrinsic correlation of spiral galaxy ellipticities. They reach final results which are broadly similar (although they find much more noise), and also conclude that while these effects are likely to be minor for present surveys, they may become important with small redshift widths (in their case for shallower surveys). We also became aware of analytic work on a similar theme by Catelan et al. (in preparation), and Mackey and White (in preparation).
We thank Ludo Van Waerbeke, Martin White, Volker Springel, Marc Kamionkowski, and Jordi Miralda-Escudé for useful discussions, and Nick Kaiser for help in interpreting some observational results. We also thank George Efstathiou for supplying us with the FOF groupfinder and the Virgo consortium for making their simulation data public. RACC acknowledges support from NASA Astrophysical Theory Grant NAG5-3820.
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# First order phase transition in a nonequilibrium growth process
## I Introduction
Recently Hinrichsen *et. al.* have introduced a discrete model for nonequilibrium surface growth that displays a first order unbinding transition in a region of the parameter space where there is phase coexistence between a growing and a bound surface. They used a restricted solid on solid dynamics, and distinguished the probability of adsorption of particles at the substrate from that at higher layers. This in order to take into account short range interactions between the substrate and the surface.
In the same spirit we study a continuum growth model for an interface interacting with a substrate through a short range attractive potential. The interface separates a solid from a vapor phase, and is driven by a difference in chemical potential due, for example, to a flux of incoming particles, as in Molecular Beam Epitaxy. As the driving force increases the interface eventually detaches from the substrate and grows indefinitely according to KPZ dynamics. We will focus on the nature of the transition from a bound to a moving phase. We resort to numerical simulations and to a mean field approach first introduced in for the KPZ equation.
A similar unbinding process of a surface from a wall has been studied by Muñoz and Hwa . They considered a KPZ equation to which they added strictly attractive or repulsive interactions, and identified a second order unbinding transition for which they calculated some critical exponents.
We will show how the short range attractive potential changes the nature of the phase transition, which becomes first order for an entire range of values of the parameters of the model. The second order transition is recovered only for larger values of the noise, or in the case of a very short ranged potential, but the critical exponent is still different from that of the simple repulsive wall case.
Our model is an example of continuous nonequilibrium system with noise. In the context of noise induced phase transition this kind of systems has been studied by Van den Broeck *et al.* , who showed that multiplicative noise is essential for the transition to take place, since it can trigger instabilities in the short time dynamics of the system. On the other hand systems with multiplicative noise are characterized by a transition to an absorbing phase , a state from which the system cannot escape, characterized by a null value of the order parameter. The phase diagram of this class of systems has been extensively studied in recent years . Most of these works were concerned with second order transitions, but an example of first order has been studied by Müller *et al.* with the help of mean field techniques and numerical simulations.
## II The Model
With the term *surface* we identify the interface between the solid and vapor phases of some substance. In what follows we want to give a statistical physics description of the interactions that may take place between the vapor particles as they are depositing on some kind of substrate.
First of all a note on terminology: an interaction is called attractive *for the surface* when there is some force that repels incoming particles, so that the surface is not growing; On the contrary, attractive interaction between the substrate and the particles will let the surface grow and eventually detach from it.
In our model the system is described by a height field $`h(x)`$ defined on a $`d`$-dimensional continuous substrate, which evolves according to a Langevin equation of the KPZ type. The latter is chosen in order to have a system out of equilibrium.
The Langevin equation for the field $`h(x)`$ is obtained in the usual way:
$$_th(t)=D\left[^2h(h)^2\right]+r\frac{\delta V[h]}{\delta h}+\sigma \eta ,$$
(1)
where $`D`$ is some diffusion coefficient, $`r`$ is the driving force due to the incoming particles, and the nonlinear term $`(h)^2`$ characterizes the nonequilibrium KPZ-like dynamics, and allows for the growth along the direction normal to tilted areas of the surface. This term breaks the up-down symmetry of the system and therefore defines a preferential direction for growth (see *e. g.* ).
The field $`\eta (x,t)`$ is a white Gaussian noise with the following properties:
$`\eta (x,t)=0`$ (2)
$`\eta (x,t)\eta (x^{},t^{})=\delta (xx^{})\delta (tt^{}).`$ (3)
We model the interaction of the substrate with the surface introducing a potential $`V[h]`$ which has the shape of the Morse potential (see Fig. 1):
$$V_M[h]=A(1e^{\beta h}e^{\beta h_0})^2,$$
(4)
where $`A`$ is the depth of the well, $`\beta `$ the hardness of the repulsive wall that mimics the substrate, and $`h_0`$ the position of the local minimum, which can be set to $`0`$ without loss of generality.
Note that the driving force $`r`$ can be included in the potential as a linear term in the height field: $`rh(x)`$. The parameter $`r`$ is the essential one: if it is negative the shape of $`V`$ is a well with an exponential wall on the small-$`h`$ side, whereas if it is positive the minimum close to the repelling wall is only a local one (Fig. 1). Particles incoming from above in this case have to trespass a potential barrier (*i. e.* the local maximum of $`V`$) in order to arrive close to the substrate. This is equivalent to a repulsion for particles, or an attraction for the surface. Physically this term has a precise meaning, since it is due to the difference in chemical potential between the solid and vapor phases, and determines a drive of the interface toward higher or lower values, depending on its sign.
In the dynamics of growing surfaces this term is usually eliminated via the transformation $`hhrt`$, which correspond to viewing the system from a moving reference frame. As one can see the potential in Eq. (4) is not invariant for such a transformation, nor for the more general $`hh+\delta h`$. Notice that the KPZ equation is invariant for both these transformations. In the present case the consequence is that we can uniquely define the mean position of the surface with respect to the substrate, and its state of motion. It is therefore meaningful to study the unbinding of the surface from the substrate.
By inserting the Morse potential one obtains:
$$_th(t)=r+2\beta A(e^{2\beta h}e^{\beta h})+D\left[^2h(h)^2\right]+\sigma \eta .$$
(5)
By performing a Hopf-Cole transformation $`\psi =e^h`$ we map Eq. (5) into a Langevin equation with multiplicative noise; this suggests us the order parameter to be used in the search for a phase transition, namely the spatial average of $`\psi `$, $`m\psi `$:
$$_t\psi (t)=(r+\sigma ^2/2)\psi 2\beta \left[\psi ^{2\beta +1}\psi ^{\beta +1}\right]+^2\psi +\sigma \psi \eta .$$
(6)
Notice that we have used the Ito calculus in doing the transformation, and we have been able to set $`A`$ and $`D`$ equal to one via a suitable rescaling of space and time.
Qualitatively we can have two possible situations in the long time limit of the stationary solution of the equation. Either $`m0`$, which corresponds to a diverging value of $`h`$, *i. e.* to a moving (growing) surface; or $`m0`$, corresponding to the surface being bound to the substrate at some average height. The transition between the two phases is controlled by the parameter $`r`$. Namely, as one can see from Eq. (5), for large negative values of $`r`$ the surface is pushed down against the substrate, whereas for large positive values the surface is pulled away from it. It is the balance between the force $`r`$ and the repulsion from the exponential wall $`e^{2\beta h}`$ that determines the equilibrium position of the surface. Notice finally that in the language of multiplicative noise systems the state with $`m=0`$ is the absorbing one, since the surface can never invert its average velocity and go back toward the substrate.
An equation similar to (6) has been studied by Muñoz and Hwa in in the context of multiplicative noise processes. The system had a soft lower repulsive wall represented by a term $`p\psi ^{p+1}`$ in the Langevin equation, the hard wall limit corresponding to $`p\mathrm{}`$:
$$_t\psi =^2\psi +r\psi p\psi ^{p+1}+\psi \sigma \eta .$$
(7)
The authors found a second order phase transition between a bound and a moving phase. From scaling arguments, the exponent that controls the vanishing of the order parameter close to the transition is found to be $`\beta _r=(2z)/(2z2)`$, where $`z`$ is the KPZ dynamical exponent. Numerical simulations confirmed this result.
## III Mean Field Approximation
To proceed further with the analytical solution of the model we notice that dimensional analysis performed on Eq. 6 along the lines of shows that the upper critical dimension for the depinning transition is 2. This means that mean field results should be correct in 3 dimensions, and provide a reasonable estimate (up to logarithmic corrections) in two dimensions, *i. e.* for real surfaces.
First of all we perform a discretization in space, by inserting a (hyper)cubic lattice, and define the variables $`\psi (𝐱`$$`,t)\psi (x_i,t)\psi _i(t)`$. The discretized Laplacian is:
$$^2\psi _i=\frac{1}{d}\underset{jnn(i)}{}(\psi _j\psi _i),$$
(8)
where the sum runs over the nearest neighbor of the site $`i`$. The mean field approximation we are going to use consists in the following expression for the discretized Laplacian:
$$\frac{1}{d}\underset{jnn(i)}{}(\psi _j\psi _i)\psi \psi _i.$$
(9)
This is equivalent to say that the number of dimensions is very large, or that the number of neighbors (coordination of the lattice) is infinite, or again that the interaction is infinite range.
Equation (6) then becomes:
$$_t\psi _i(t)=(r+\sigma ^2/21)\psi _i2\beta \left[\psi _i^{2\beta +1}\psi _i^{\beta +1}\right]+\psi \sigma \psi _i\eta _i,$$
(10)
from which, going back to the field $`h`$, the discretized version of Eq. (5) becomes:
$$_th_i(t)=(r+1)+2\beta e^{\beta h_i}(e^{\beta h_i}1)e^he^{h_i}+\sigma \eta _i.$$
(11)
### The Phase Transition
We obtain the probability distribution of the stationary solution of Eq. (11) by solving the associated Fokker-Planck equation. The stationary solution is:
$`P_{st}(h)=𝒩e^{\frac{2V_{eff}(h)}{\sigma ^2}},`$where $`𝒩`$ is the normalization factor, and the effective potential is in the present case:
$$V_{eff}(h)=(r+1)h+e^{2\beta h}2e^{\beta h}+e^he^h.$$
(12)
This expression for the probability distribution depends on the value of the order parameter, which must be determined self-consistently. We have to solve therefore the following mean-field equation:
$`m=e^h`$ $`=`$ $`{\displaystyle e^hP_{st}(h)𝑑h}F(m)={\displaystyle \frac{g_J(m)}{g_{J1}(m)}}`$ (13)
$`g_j(m)`$ $``$ $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑he^{(j+1)h}e^{\frac{2}{\sigma ^2}e^{2\beta h}}e^{\frac{4}{\sigma ^2}e^{\beta h}}e^{\frac{2m}{\sigma ^2}e^h},`$ (14)
where $`J=2(r+1)/\sigma ^2`$.
Plotting $`m`$ and $`F(m)`$ on the same graph, it is easy to see that Eq. (13) can have one, two or three solutions in $`m`$, according to the value of $`r`$. The possible situations are represented in Fig. 2. The stability of the solutions can be investigated making use of the results of Shiino . Namely, the solution can be stable only when $`F(m)`$ crosses the line $`m`$ from above. In our case we have that for $`r`$ larger than some $`r_{c_2}`$ the only solution (therefore a stable one) is $`m=0`$; for $`r`$ below some value which will be determined in a moment, there are only 2 solutions, the stable one being $`m0`$; finally for intermediate values of $`r`$ there are three solutions, two of which look stable. Notice that $`m=0`$ remains always a solution; it becomes stable only when $`r`$ reaches the value for which $`F^{}(m=0)=1`$. This value is found analytically to be: $`r_c=\frac{\sigma ^2}{2}`$. This number is also the lowest value of $`r`$ for which three solutions exist.
The situation just described represents a typical case of first order phase transition, and is the same as found in . The value of the order parameter $`m`$ remains different from zero as long as this is the only stable solution. Then for $`r>r_c`$ there is a coexistence region, where two phases are stable, and the final state of the system is determined by the initial conditions. Finally there is the value $`r=r_{c_2}`$ above which the order parameter can be only zero. If we remind the definition of $`m`$ we understand that this phase transition corresponds to the unbinding of the surface from the substrate, $`e^h0`$ meaning $`h\mathrm{}`$. We remark that with the term *coexistence* we do not mean that both phases are present at the same time, rather that depending on the initial conditions both are possible stationary state of the system. In Fig. 3 is shown an example of this phenomenon. In the coexistence region an interface which is initially close to the substrate will remain bounded to it whereas if the interface is initially far it will move further away from the substrate with a velocity $`v>0`$.
### The Moving Phase
It is also possible to determine the velocity of the surface when it is away from the substrate. To do this we define a field $`\varphi _i(t)`$ for the fluctuations of $`h`$ around its mean position: $`\varphi _i(t)=h_i(t)vt`$. Here $`v`$ is the mean velocity of the surface obtained from Eq. (11):
$`v=_th_i=r+1e^he^h,`$where we could neglect the exponential terms in the potential, since they are vanishingly small in the moving phase. The field $`\varphi _i(t)`$ obeys in this phase to the following Langevin equation:
$`_t\varphi _i=_th_iv=(r+1v)e^he^{h_i}+\sigma \eta _i.`$The stationary probability distribution from the associated Fokker-Planck equation is as usual (see *e. g.* ):
$$P_{st}(\varphi )=𝒩e^{\frac{2V_{eff}(\varphi )}{\sigma ^2}}=𝒩e^{\alpha \varphi \gamma e^\varphi },$$
(15)
with $`\alpha =\frac{2}{\sigma ^2}(r+1v)`$ and $`\gamma =\frac{2}{\sigma ^2}e^\varphi `$. The self-consistency relation for $`e^h`$ is:
$`e^h=\gamma {\displaystyle \frac{1}{\alpha 1}},`$from which we have that $`v=r\frac{\sigma ^2}{2}`$. The condition $`v>0`$ shows that $`r_c=\sigma ^2/2`$, as anticipated. In the next section we will see that for large values of the noise strength $`\sigma `$ and of $`\beta `$, the unbinding transition becomes second order. In that part of the phase diagram we can define in the vicinity of the transition a velocity critical exponent $`\theta `$ : $`v(rr_c)^\theta `$. From the previous discussion we see that $`\theta =1`$.
## IV Numerical Findings
To study further the phase diagram of the model in mean field approximation we solve numerically Eq. (13). In particular we have determined the value $`r_c^{}`$ for which its two non-zero solutions coalesce with the null one. Analytically this point corresponds to a negative second derivative of $`F(m)`$ at the critical point $`r_c`$. In Figs. 4 and 5 are shown two typical cuts of the phase diagram at fixed values of the control parameters $`\sigma `$ and $`\beta `$ respectively. One can see that as the value of one of the parameters increases while the other remains fixed, the coexistence region becomes smaller, and eventually disappears for some $`\beta ^{}(\sigma )`$ or $`\sigma ^{}(\beta )`$. At the same time the unbinding transition becomes second order, thus defining a line of tricritical points in the phase diagram ($`\sigma `$,$`\beta `$) of the system. A numerical estimate of the exponent $`\beta _r`$ governing the transition (see *e. g.* ) reveals that, at least for large values of $`\beta `$, we have:
$$m|rr_c|^{\beta _r}|rr_c|^{\sigma ^2/2}$$
(16)
This can be checked analytically by taking the limit for $`\beta \mathrm{}`$ at the critical point in Eq. (13). The condition to have a second order transition is that $`F^{\prime \prime }(m)<0`$ at $`r_c`$; in this case there can be no coexistence region, since $`r_{c_2}=r_c`$. With the same notation of Eq. (13) we have:
$`F^{^{}}(m)`$ $`=`$ $`{\displaystyle \frac{2}{\sigma ^2}}\left(1{\displaystyle \frac{g_Jg_{J2}}{g_{J1}^2}}\right)`$ (17)
$`F^{^{\prime \prime }}(m)`$ $`=`$ $`\left({\displaystyle \frac{2}{\sigma ^2}}\right)^2\left[{\displaystyle \frac{g_{J2}}{g_{J1}}}{\displaystyle \frac{g_Jg_{J3}}{g_{J1}^2}}+2{\displaystyle \frac{g_Jg_{J2}^2}{g_{J1}^3}}\right].`$ (18)
Then we make the change of variable $`x=2m/\sigma ^2e^h`$, to get:
$`g_j(m)`$ $`=`$ $`\left({\displaystyle \frac{2m}{\sigma ^2}}\right)^{j+1}{\displaystyle _{\frac{2m}{\sigma ^2}}^{\mathrm{}}}𝑑xx^{(j+2)}e^x`$
$`=`$ $`\left({\displaystyle \frac{2m}{\sigma ^2}}\right)^{j+1}\left[\mathrm{\Gamma }(j1){\displaystyle _0^{\frac{2m}{\sigma ^2}}}𝑑xx^{(j+2)}e^x\right].`$
Since we are looking for a second order transition, we need to estimate the small $`m`$ limit of $`g_j(m)`$. Therefore we can expand $`e^x`$ in the integral, and solve it explicitly to any order in $`m`$. At the critical point $`r_c=\sigma ^2/2`$ we find that:
$`F^{^{\prime \prime }}(m)={\displaystyle \frac{\frac{2}{\sigma ^2}+1}{\mathrm{\Gamma }(\frac{2}{\sigma ^2}+1)}}\left({\displaystyle \frac{2m}{\sigma ^2}}\right)^{\frac{2}{\sigma ^2}1}<0,`$which means that the transition is always second order at large $`\beta `$, and that the $`m`$-expansion of $`F(m)`$ close to the critical point can be written as:
$`F(m)=(1+\delta r)m+Cm^{\frac{2}{\sigma ^2}+1},`$where $`\delta r=r_cr`$, and $`C`$ is some constant. Solving the mean field equation $`F(m)m=0`$ gives:
$$m\left(\delta r\right)^{\frac{\sigma ^2}{2}},$$
(19)
in agreement with the numerical value.
We note that the latter calculation can be repeated for the model of Eq. (7), to give the same result, thus implying that in the mean field limit of a system with a hard wall the phase transition is governed by a continuous exponent. To our knowledge this finding was not previously reported. A different value for the $`\beta _r`$ exponent, namely $`\beta _r=1/p`$, is given in , where the mean field approximation Eq. (9) is also compared with different results from field theoretical calculations ($`\beta _r=1`$). However it is not clear whether the hard wall limit can be obtained from the aforementioned calculations.
Another feature of the phase diagram is that in the limit of small $`\sigma `$ $`r_{c2}\beta /2`$ , but we were unable to derive this result analytically from Eq. 13.
We check if there is any functional relation between $`\beta ^{}`$ and $`\sigma ^{}`$. Figure 6 shows a log-log plot of the tricritical points together with the linear fitting. The value of the slope is in this case about $`4`$, which means that $`\beta ^{}\sigma ^4`$, or that $`\beta ^{}r_c^2`$. Notice however that this linear fitting is a fairly poor approximation, and there is no reason why we should expect a priori any simple functional relation between these two quantities.
We have also checked with simulations on a finite number of sites that these results are the same as those obtained from Eq. (11). With the help of these simulations we could see that hysteresis associated with the first order transition and the coexistence region does indeed occur.
## V Results in Finite Dimension
To perform numerical simulation of Eq. (5) in finite dimensions we introduced as usual a (hyper)cubic lattice on the space of coordinates, and used for the discretized Laplacian the expression in Eq. (8). Going back to the $`h`$ field we obtained the following Langevin equation on a lattice:
$$_th_i=r+1+2\beta e^{\beta h_i}(e^{\beta h_i}1)e^{h_i}\frac{1}{d}\underset{jnn(i)}{}e^{h_j}+\sigma \eta _i.$$
(20)
To have a better accuracy we chose to apply the Heun method for solving stochastic differential equations (see *e. g.* ). We run the simulations long enough for the system to reach its stationary state, and then measured the order parameter at intervals larger than some estimated autocorrelation time.
The results from three-dimensional simulations are in quantitative agreement with the mean field analysis, thus confirming that the upper critical dimension for this class of systems is $`2`$.
It is now interesting to see how this analysis is modified in the two-dimensional case, the physical one, when in the renormalization group sense only logarithmic corrections are expected. We simulated systems on a square lattice with linear size $`L`$ up to 60. A cut of the phase space is shown in Fig. 7. It is easily seen how the main features of the phase diagram remain unaltered in 2 dimensions; namely we can still identify a first order phase transition with a coexistence region where the system displays hysteresis (Fig. 3), and a crossover to a second order transition for large values of the noise.
Similar results are obtained even in 1 dimension. It must be noticed how the coexistence region becomes smaller as one goes to low dimensionality, but is still present.
## VI Discussion and Conclusions
In a sense the combined effect of the attracting potential and of the noise is similar to that of a superposition of quenched an thermal noise (*cfr.* ). Namely in both situations we have a force opposing the growth that depends on the height of the surface, and a noise term that randomly modulates in time and space such interaction. This is why we can use the quenched noise terminology and define *e. g.* the velocity exponent. Of course it is the attractive nature of the Morse potential that gives the transition its first order character, and we expect similar results for other short range attractive potential as well.
We have shown how a simple modification of the nonlinear part of the KPZ equation can describe the dynamics of a surface that interacts with the substrate via a short range interaction. The mean field analysis has revealed that a first order unbinding transition takes place for a large range of the control parameters, and that in a whole region of the phase diagram the system’s stationary position depends on the initial conditions (hysteresis). Furthermore a crossover to a second order transition takes place for large values of the control parameters, and the exponent controlling the vanishing of the order parameter depends continuously on the strength of the noise.
The same situation is recovered in simulations above and below the upper critical dimension $`d_c=2`$
#### Acknowledgments
We acknowledge stimulating discussions with M. A. Muñoz, H. Hinrichsen, D. Mukamel, A. De Martino.
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# Untitled Document
##
During the last few years the new large - scale neutrino telescopes NT-200 /1/ and
AMANDA /2/ have been taken into operation. Two deep underwater arrays, NESTOR/3/ and ANTARES /4/, are under construction.
The Cherenkov light from high energy muons, electromagnetic and hadronic showers can be recorded at distances of 20 – 100 m depending upon the light absorption of water or ice. Impulse laser light sources are widely used in these arrays for calibration. Light pulses from such sources can be recorded over even larger distances.
At such distances, the difference between group and phase velocity of light in water or ice is essential. As far as I know, this fact hasn’t been mentioned in the literature from the time of the first neutrino telescope projects in the middle of the seventies /5,6/ and hasn’t been taken into account for the data analysis and MC calculations.
The velocity of light pulses in matter is given by the group velocity /7,8/ :
$$V_{gr}=d\omega /dk$$
$`(1)`$
where $`\omega `$ is the frequency, $`k`$ is the wave vector.
The connection between $`\omega `$ and $`k`$ can be written in the following form ( dispersion equation):
$$\omega =ck/n(\omega )$$
$`(2)`$
where $`c`$ is the velocity of light in vacuum, $`n`$ is the refraction index.
As it is known, refraction indices given in handbooks like /9/ correspond to the phase velocity due to the method they are measured /10/ :
$$n(\omega )=c/V_{ph}(\omega )$$
$`(3)`$
where $`V_{ph}`$ is the phase velocity.
It is reasonable to accept a group refraction index for the group velocity similar to the phase index of refraction:
$$n_{gr}(\omega )=c/V_{gr}(\omega )$$
$`(4)`$
where $`V_{gr}`$ is the group velocity.
Using (2) and replacing the $`\omega `$ dependence of $`n`$ by a wavelength dependence, the relation between $`n_{gr}`$ and $`n`$ can be derived:
$$n_{gr}(\lambda )=n(\lambda )\lambda dn/d\lambda $$
$`(5)`$
where $`\lambda `$ is the wavelength.
$`dn/d\lambda `$ is negative for water over the visible light range, and so $`n_{gr}`$ is larger than $`n`$. Fig.1 shows the wavelength dependence of $`n`$ ( curve 1) for destilled water at 20<sup>0</sup>C /9/, and the same for $`n_{gr}`$ (curve 2).
The density of water influences the refraction index at large depths of neutrino telescopes. But this correction is rather small due to the comparatively small compression coefficient of water. For the conditions of NT-200 ( depth 1100 m, water temperature 4<sup>0</sup>C) the variation of the refraction index is about 0.002 only. This is more than 10 times smaller than the difference between $`n_{gr}`$ and $`n`$. This correction of $`n`$ would reach 0.01 for the arrays to be placed in the ocean depths of 4 -5 km, but even in this case the depth correction is significantly smaller than that due to the difference between group and phase velocities. The curve 3 on fig.1 shows the dependence of $`n_{gr}`$ on $`\lambda `$ for the conditions of the telescope NT-200.
All curves at fig.1 were calculated for destilled water. The light absorption of natural water differs from that of destilled water, and so there is a difference in the image part of their dielectric penetration. As Kramers-Kroning equation connects the real part of the dielectric penetration with it’s imaginary part /7/, the real refraction index differs, in principle, from that of destilled water. Such a correction would be non - negligible only if there are the intensive lines of absorption in the visible light absorption spectrum.
Fig.2 shows the dependencies of delays on the distance between light source and receiver which emerge from the difference between the phase and group velocities for different wavelengths.
For arrays with angular resolution 1-2<sup>0</sup> replacement of the phase velocity by the group velocity probably would not lead to essential changes in track reconstruction. For the projects which claim an angular accuracy 0.1-0.2<sup>0</sup> and absorption length of $`50`$m /4/ the use of group velocity in track reconstruction procedures seems to be absolutely necessary.
## Acknowledgments
I would like to express thanks to G.V.Domogatsky, J-A.M. Djilkibaev, S.I.Klimushin, B.K.Lubsandorzhiev, E.A.Osipova, V.V.Prosin, C.Spiering and V.B.Tsvetkov for useful discussion.
## References
1. I.A.Belolaptikov et al.,
The Baikal Underwater Neutrino Telescope: Design, Performance, and First Results.
Astroparticle Physics 7(1997) 263-282.
2.E.Andres et al.,
The AMANDA Neutrino Telescope: Principle of Operation and First Results.
Astroparticle Physics 13 (2000) 1-20.
3.B.Monteleoni et al.,
NESTOR a deep sea physics laboratory for the Mediterranean
Proc. of the 17th International Conference on Neutrino Physics and Astrophysics, (Neutrino 96)
Helsinki, Finland, June 13-19, 1996.
4.Ph.Amram et al.,
The ANTARES Project
Nucl.Physics B (Proc. Suppl.) 75A (1999) 415.
5.V.S.Berezinsky and G.T.Zatsepin
Sov.Phys.Usp. 20 (1977) 361.
6.DUMAND-75. Proceeding of the 1975 Dumand Summer Study,Bellingham,
edited by P.Kotzer.Published by Western Washington State College.
7.L.D.Landau , E.M.Lifshiz. Electrodynamics of Continuous Media.
Moscow, Nauka, 1982, p 345.
8.A.Sommerfeld, OPTIK , Wiesbaden 1950.
9.Physical Quantities , Handbook,ed. I.S.Grigorev and E.Z.Meylihoba,
Energoatomizdat, Moscow 1991, p. 790.
10. A.I.Stogarov and N.F.Timofeeva,
Trudi GOI, 1963, v 31, N 160 (in Russian).
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# Casimir Energy of a Semi-Circular Infinite Cylinder
## I Introduction
When calculating the ground state energy of a quantum field (the Casimir energy) the main problem is to single out a finite part of the vacuum energy which is initially divergent. Usually for this purpose a subtraction procedure is used with preliminary regularization of the divergent expressions (for example, by introducing ultraviolet cutoff). However in quantum field theory treated with allowance for nontrivial boundary conditions or in the space-time with curvature, a complete renormalization procedure is not formulated explicitly. Therefore, for any specific problem the subtraction procedure should be invented anew. As a result, one succeeds in calculating the Casimir energy only in the problems with known spectra or at least implicitly known spectra. Practically it implies the boundary conditions of high symmetry (parallel plates, sphere, cylinder).
In studies of the Casimir energy it is wildly used the zeta function technique which is also referred to as the zeta regularization or zeta renormalization. In fact, the use of the zeta functions, as well as other regularizations, gives only regularized quantities for ground state energy, for effective potential and so on. The necessity to renormalize the expressions obtained in this way certainly remains. However, in some problems the zeta technique gives at once a finite result. Usually the latter is considered to be a renormalized physical answer though generally it is not the case.
When using the zeta regularization in one or another problem, it is desirable to know beforehand whether the finite result can be obtained in this way. In order to answer this question the general analysis of the divergences in the problem at hand should be accomplished. This can be done by calculating the heat kernel coefficients depending on the geometry of the manifold under consideration. For a large class of situations these coefficients have been obtained. However there is a number of problems (for example, boundaries with edges or corners) for which no general results regarding the heat trace are known.
In this situation it is undoubtedly worth carrying out, in the framework of the zeta function technique, the calculations of the Casimir energy for new configurations, both the cases being interesting with finite result and with pole contributions left in the final expression for the vacuum energy.
In the present paper we address the calculation of the Casimir energy for boundaries with edges, more precisely, the vacuum energy of electromagnetic field will be calculated for a semi-circular cylindrical shell by making use of the relevant zeta functions. This shell is obtained by crossing an infinite circular cylindrical shell by a plane passing through the symmetry axes of the cylinder. All the surfaces, including the infinite cutting plane, are assumed to be perfectly conducting. Obviously it is sufficient to consider only a half of this configuration (left or right) which we shall refer to as a semi-circular cylindrical shell or, for sake of shortening, as a semi-circular cylinder. The internal boundary value problem for this configuration is nothing else as a semi-cylindrical waveguide. In the theory of waveguides it is well known that a semi-circular waveguide has the same eigenfrequencies as the cylindrical one but without degeneracy (without doubling) and safe for one frequency series (see below). Notwithstanding the very close spectra, the zeta function technique does not give a finite result for a semi-circular cylinder unlike for a circular one. First the Casimir energy of an infinite perfectly conducting cylindrical shell has been calculated in Ref. by introducing ultraviolet cutoff and recently this result was derived by zeta function technique (see also Refs. ). As far as we know the asymmetric boundaries such as a semi-circular cylinder have not been considered in the Casimir problem.
The paper is organized as follows. In Sec. II the electromagnetic spectra are considered in details for cylindrical and semi-cylindrical shells. The general solution of the Maxwell equations for boundary conditions chosen is expressed in terms of two scalar functions, longitudinal components of the electric and magnetic Hertz vectors. These scalar functions are the eigenfunctions of the two-dimensional transverse Laplace operator and obey the Dirichlet and Neumann boundary conditions on the conducting surfaces. In Sec. III the spectral zeta function is constructed for the Dirichlet boundary value problem. To this end, the technique is used which has been elaborated before for representing the spectral zeta function, with given eigenfrequency equations, in terms of contour integral. When carrying out the analytic continuation of the zeta function into the physical region, the uniform asymptotic expansion for the modified Bessel functions is used. In the same way, in Sec. IV the zeta function is constructed for a scalar field obeying the Neumann boundary conditions given on the surface of a semi-circular cylindrical shell. The Section V is concerned with the complete zeta function for electromagnetic field with boundary conditions on the semi-circular cylinder. Transition to the relevant two-dimensional problem is also considered here. In the Conclusion (Sec. VI) the results obtained are summarized, and the origin of the pole singularities of the zeta functions at hand and their relation to the respective boundary value problem are briefly discussed.
## II Eigenmodes of Electromagnetic Field for Circular and Semi-Circular Cylinders
The construction of the solutions to the Maxwell equations with boundary conditions given on closed surfaces proves to be nontrivial problem. Mainly it is due to the vector character of the electromagnetic field . In the case of cylindrical symmetry the electric $`𝐄`$ and magnetic $`𝐇`$ fields are expressed in terms of the electric ($`𝚷^{}`$) and magnetic ($`𝚷^{\prime \prime }`$) Hertz vectors having only one non-zero component
$`𝚷^{}`$ $`=`$ $`𝐞_z\mathrm{\Phi }(r,\phi )e^{\pm ik_z^{}z},`$ (1)
$`𝚷^{\prime \prime }`$ $`=`$ $`𝐞_z\mathrm{\Psi }(r,\phi )e^{\pm ik_z^{\prime \prime }z}.`$ (2)
Here the cylindrical coordinate system $`r,\phi ,z`$ is used with $`z`$ axes directed along the cylinder axes. The common time–dependent factor $`e^{i\omega t}`$ is dropped. The scalar functions $`\mathrm{\Phi }(r,\phi )`$ and $`\mathrm{\Psi }(r,\phi )`$ are the eigenfunctions of the two–dimensional transverse Laplace operator and meet, respectively, the Dirichlet and Neumann conditions on the boundary $`\mathrm{\Gamma }`$
$$(\mathbf{}_{}^2+\gamma ^{{}_{}{}^{}\mathrm{\hspace{0.17em}2}})\mathrm{\Phi }(r,\phi )=0,\mathrm{\Phi }(r,\phi )|_\mathrm{\Gamma }=0,$$
(3)
$$(\mathbf{}_{}^2+\gamma ^{\prime \prime 2})\mathrm{\Psi }(r,\phi )=0,\frac{\mathrm{\Psi }(r,\phi )}{n}|_\mathrm{\Gamma }=0,$$
(4)
where $`\mathbf{}_{}^2`$ is the transverse part of the Laplace operator
$$\mathbf{}_{}^2=\frac{^2}{r^2}+\frac{1}{r}\frac{}{r}+\frac{1}{r^2}\frac{^2}{\phi ^2}$$
(5)
and
$$\gamma ^{{}_{}{}^{}\mathrm{\hspace{0.17em}2}}=\omega ^2k_z^{{}_{}{}^{}\mathrm{\hspace{0.17em}2}},\gamma ^{{}_{}{}^{\prime \prime }\mathrm{\hspace{0.17em}2}}=\omega ^2k_z^{{}_{}{}^{\prime \prime }\mathrm{\hspace{0.17em}2}}.$$
(6)
First we consider a cylindrical shell. In this case the functions $`\mathrm{\Phi }(r,\phi )`$ and $`\mathrm{\Psi }(r,\phi )`$ should be $`2\pi `$-periodic in angular variable $`\phi `$. As a result the Dirichlet boundary value problem (3) has the following unnormalized eigenfunctions ($`E`$-modes)
$$\mathrm{\Phi }_{nm}(r,\phi )=\genfrac{}{}{0pt}{}{\mathrm{sin}}{\mathrm{cos}}(n\phi )\{\begin{array}{cc}J_n(\gamma _{nm}^{}r),r<a,\hfill & \\ H_n^{(1)}(\overline{\gamma }_{nm}^{}r),r>a,\hfill & \end{array}$$
(7)
where $`a`$ is the cylinder radius, $`J_n(x)`$ are the Bessel functions, $`H^{(1)}(x)`$ are the Hankel functions of the first kind, and $`\gamma _{nm}^{}`$, $`\overline{\gamma }_{nm}^{}`$ stand for the roots of the frequency equations
$`J_n(\gamma _{nm}^{}a)=0,H_n^{(1)}(\overline{\gamma }_{nm}^{}a)=0,`$ (8)
$`n=0,1,2,\mathrm{},m=1,2,\mathrm{}.`$ (9)
For the Neumann boundary value problem (4) we have the $`H`$-modes
$$\mathrm{\Psi }_{nm}(r,\phi )=\genfrac{}{}{0pt}{}{\mathrm{sin}}{\mathrm{cos}}(n\phi )\{\begin{array}{cc}J_n(\gamma _{nm}^{\prime \prime }r),r<a,\hfill & \\ H_n^{(1)}(\overline{\gamma }{}_{}{}^{\prime \prime }{}_{nm}{}^{}r),r>a,\hfill & \end{array}$$
(10)
where $`\gamma _{nm}^{\prime \prime }`$ and $`\overline{\gamma }_{nm}^{\prime \prime }`$ are the roots of the equations
$`{\displaystyle \frac{d}{dr}}J_n(\gamma _{nm}^{\prime \prime }r)|_{r=a}=0,{\displaystyle \frac{d}{dr}}H^{(1)}(\overline{\gamma }_{nm}^{^{\prime \prime }}r)|_{r=a}=0,`$ (11)
$`n=0,1,2,\mathrm{},m=1,2,\mathrm{}.`$ (12)
As usual, it is assumed that for $`r>a`$ the eigenfunctions should satisfy the radiation condition.
It is important to note that each root
$$\gamma _{nm}^{},\overline{\gamma }_{nm}^{^{}},\gamma _{nm}^{^{\prime \prime }},\overline{\gamma }_{nm}^{^{\prime \prime }},n1,m1$$
(13)
is doubly degenerate since, according to Eqs. (7), (10), there are two eigenfunctions which are proportional to either $`\mathrm{sin}(n\phi )`$ or $`\mathrm{cos}(n\phi )`$. The frequencies with $`n=0`$
$$\gamma _{0m}^{},\overline{\gamma }_{0m}^{^{}},\gamma _{0m}^{^{\prime \prime }},\overline{\gamma }_{0m}^{^{\prime \prime }},m=1,2,\mathrm{}$$
(14)
are independent on $`\phi `$, and the degeneracy disappears.
For given Hertz vectors $`𝚷^{}`$ and $`𝚷^{\prime \prime }`$ the electric and magnetic fields are constructed by the formulas
$`𝐄`$ $`=`$ $`\mathbf{}\times \mathbf{}\times 𝚷^{},𝐇=i\omega \mathbf{}\times 𝚷^{}(E\text{-modes}),`$ (15)
$`𝐄`$ $`=`$ $`i\omega \mathbf{}\times 𝚷^{\prime \prime },𝐇=\mathbf{}\times \mathbf{}\times 𝚷^{\prime \prime }(H\text{-modes}).`$ (16)
It has been proved that the superposition of these modes gives the general solution to the Maxwell equations in the problem under consideration. An essential merit of using the Hertz polarization vectors is that in this approach the necessity to satisfy the gauge conditions does not arise.
Let us consider a waveguide which is obtained by cutting the infinite cylindrical shell by a plane passing through the symmetry axes of the cylinder (see Fig. 1). All the surfaces are assumed to be perfectly conducting. In this case the boundary value problems (3) and (4) for the Hertz electric ($`𝚷^{}`$) and magnetic ($`𝚷^{\prime \prime }`$) vectors have the following eigenfunctions
$`\mathrm{\Phi }_{nm}(r,\phi )=\mathrm{sin}(n\phi )\{\begin{array}{cc}J_n(\gamma _{nm}^{}r),r<a,\hfill & \\ H_n^{(1)}(\overline{\gamma }_{nm}^{}r),r>a,\hfill & \end{array}`$ (17)
$`n=1,2,\mathrm{},m=1,2,\mathrm{}`$ (18)
and
$`\mathrm{\Psi }_{nm}(r,\phi )=\mathrm{cos}(n\phi )\{\begin{array}{cc}J_n(\gamma _{nm}^{\prime \prime }r),r<a,\hfill & \\ H_n^{(1)}(\overline{\gamma }_{nm}^{\prime \prime }r),r>a,\hfill & \end{array}`$ (19)
$`n=0,1,2,\mathrm{},m=1,2,\mathrm{}.`$ (20)
The frequencies $`\gamma _{nm}^{}`$, $`\overline{\gamma }_{nm}^{}`$, $`\gamma _{nm}^{\prime \prime }`$, and $`\overline{\gamma }_{nm}^{\prime \prime }`$ are determined by the same equations (8) and (11). However the new spectral problem has two essential distinctions: i) the frequencies (13) are now nondegenerate, and ii) two series of eigenfrequencies
$$\gamma _{0m}^{},\overline{\gamma }_{0m}^{},m=1,2,\mathrm{}$$
(21)
are absent. At first sight one could expect that such a change of the spectrum cannot influence drastically on the ultraviolet behaviour of the relevant spectral density. However, as it will be shown below, the zeta function for a semi-circular cylinder, unlike for a circular one, does not provide a finite answer for the Casimir energy in the problem in question.
In view of all above-mentioned the zeta function for electromagnetic field obeying the boundary conditions on the surface of the semi-circular cylinder is the sum of two zeta functions for scalar massless fields satisfying the Dirichlet and Neumann conditions on the lateral of this cylinder.
## III Zeta Function for Dirichlet boundary value problem
First we consider the Dirichlet boundary conditions. We shall proceed from the following representation for the zeta function in terms of a contour integral for given frequency equations (8) with $`n=1,2,\mathrm{}`$
$$\zeta _{\text{cyl}}^\text{D}(s)=\frac{1}{2\pi i}_{\mathrm{}}^{\mathrm{}}\frac{dk_z}{2\pi }\underset{n=1}{\overset{\mathrm{}}{}}_C𝑑\gamma (\gamma ^2+k_z^2)^{s/2}\frac{d}{d\gamma }\mathrm{ln}\frac{J_n(\gamma a)H_n^{(1)}(\gamma a)}{J_n(\mathrm{})H_n^{(1)}(\mathrm{})}.$$
(22)
The contour $`C`$ consists of the imaginary axis $`(i\mathrm{},i\mathrm{})`$ and a semi-circle of an infinite radius in the right half-plane of a complex variable $`\gamma `$. The details of obtaining this integral representation can be found in Refs. . Contribution into Eq. (22) of integration along a semi-circle of infinite radius vanishes. Therefore upon integration over $`k_z`$ this formula acquires the form
$$\zeta ^\text{D}(s)=C(s)\underset{n=1}{\overset{\mathrm{}}{}}_0^{\mathrm{}}𝑑yy^{1s}\frac{d}{dy}\mathrm{ln}\left[2yI_n(y)K_n(y)\right]$$
(23)
with
$$C(s)=\frac{a^{s1}}{2\sqrt{\pi }\mathrm{\Gamma }\left(\frac{s}{2}\right)\mathrm{\Gamma }\left(\frac{3s}{2}\right)}.$$
(24)
In order to accomplish the analytic continuation of (23) into the physical region including the point $`s=1`$, we shall use the uniform asymptotic expansion for the modified Bessel functions
$`\mathrm{ln}\left[2ynI_n(ny)K_n(ny)\right]`$ $`=`$ $`\mathrm{ln}(yt)+{\displaystyle \frac{t^2}{8n^2}}(16t^2+5t^4)`$ (26)
$`+{\displaystyle \frac{t^4}{64n^4}}(13284t^2+1062t^41356t^6+565t^8)+O(n^6),`$
where $`t=1/\sqrt{1+y^2}`$. Following the usual procedure applied in the analogous calculations, we add and subtract in the integrand in Eq. (23) the first two terms of the asymptotic expansion (24). After that we combine all the terms there in the following way
$`\zeta _{\text{cyl}}^\text{D}(s)`$ $`=`$ $`C(s)\left[Z_1(s)+Z_2(s)+Z_3(s)\right],`$ (27)
$`Z_1(s)`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}n^{1s}{\displaystyle _0^{\mathrm{}}}𝑑yy^{1s}{\displaystyle \frac{d}{dy}}\mathrm{ln}\left({\displaystyle \frac{y^2}{1+y^2}}\right),`$ (28)
$`Z_2(s)`$ $`=`$ $`{\displaystyle \frac{1}{8}}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}n^{1s}{\displaystyle _0^{\mathrm{}}}𝑑yy^{1s}{\displaystyle \frac{d}{dy}}\left[t^2(16t^2+5t^4)\right],`$ (29)
$`Z_3(s)`$ $`=`$ $`{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}n^{1s}{\displaystyle _0^{\mathrm{}}}dyy^{1s}{\displaystyle \frac{d}{dy}}[\mathrm{ln}(2ynI_n(yn)K_n(ny))`$ (31)
$`\mathrm{ln}{\displaystyle \frac{y}{\sqrt{1+y^2}}}{\displaystyle \frac{t^2(16t^2+5t^4)}{8n^2}}].`$
Analytic continuation of the function $`Z_1(s)`$ into vicinity of the point $`s=1`$ can be accomplished in the same way as it has been done in Ref. . Therefore we write here only the final result of this continuation
$$Z_1(s)=\frac{1}{2}\zeta (s1)\mathrm{\Gamma }\left(\frac{3s}{2}\right)\underset{m=1}{\overset{\mathrm{}}{}}\frac{\mathrm{\Gamma }\left(m{\displaystyle \frac{1s}{2}}\right)}{m\mathrm{\Gamma }(m)}.$$
(32)
The integral in Eq. (28) converges when $`1<\text{Re }s<3`$, and the sum over $`n`$ is finite for Re$`s>0`$. Thus, the regions, where the integral and the sum exist, overlap, and this formula can be used for constructing the analytic continuation needed. For this aim we substitute the sum by the Riemann zeta function
$$\underset{n=1}{\overset{\mathrm{}}{}}n^{1s}=\zeta (s+1)$$
(33)
and define the integral as an analytic function by making use of the formula
$$_0^{\mathrm{}}𝑑yy^{1s}\frac{d}{dy}t^{2(\rho 1)}=(1\rho )\frac{\mathrm{\Gamma }\left(\frac{3s}{2}\right)\mathrm{\Gamma }\left(\rho \frac{3s}{2}\right)}{\mathrm{\Gamma }(\rho )},32\text{ Re }\rho <\text{ Re }s<3.$$
(34)
In view of the poles of the gamma functions on the right-hand side of this relation, the integral on the left-hand side of it is well defined, as a function of the complex variable $`s`$, only in the region indicated in Eq. (34). Doing the analytic continuation of this integral we define it outside this region also by this equation, keeping in mind that the gamma functions involved should be treated as the analytic functions over all the plane of the complex variable $`s`$ safe for the known poles. This gives
$$Z_2(s)=\frac{1}{8}\zeta (s+1)\mathrm{\Gamma }\left(\frac{3s}{2}\right)\mathrm{\Gamma }\left(\frac{1+s}{2}\right)\left[1+3(1+s)\frac{5}{8}(3+s)(1+s)\right].$$
(35)
In order to investigate the convergence of the integral entering in Eq. (29) it makes sense to substitute in the integrand the logarithmic function by expansion (24). After that it is easy to be convinced that the integral under consideration converges when $`3<\text{Re }s<3`$. The sum over $`n`$ in this formula is finite for $`\text{Re }s>2`$. Hence, the function $`Z_3(s)`$ is an analytic function without singularities in the domain $`2<\text{Re }s<3`$. It is quiet enough for our purpose, and the analytic continuation is unnecessary.
Summarizing we conclude that Eqs. (24), (27), (31), (32), and (35) afford the analytic continuation needed and define the zeta function $`\zeta ^\text{D}(s)`$ as an analytic function in the region including the point $`s=1`$.
Now we are able to calculate the value of the zeta function $`\zeta ^\text{D}(s)`$ at the point $`s=1`$. For the coefficient $`C(s)`$ in Eq. (24) we have
$$C(1)=\frac{1}{4\pi a^2}.$$
(36)
From Eq. (32) it follows that
$$Z_1(1)=\frac{1}{2}\underset{s1}{lim}\zeta (s1)\left[\mathrm{\Gamma }\left(\frac{1+s}{2}\right)+\underset{m=2}{\overset{\mathrm{}}{}}\frac{1}{m(m1)}\right].$$
(37)
With allowance for the relations
$$\mathrm{\Gamma }(x)=\frac{1}{x}\gamma +O(x),\underset{m=2}{\overset{\mathrm{}}{}}\frac{1}{m(m1)}=1,\zeta (2)=0,$$
(38)
where $`\gamma `$ is the Euler constant, $`\gamma =0.577215\mathrm{}`$, one derives
$`Z_1(1)`$ $`=`$ $`\underset{s1}{lim}{\displaystyle \frac{1}{2}}\left[\zeta (2)+\zeta ^{}(2)(s+1)+O\left((s+1)^2\right)\right]\left[{\displaystyle \frac{2}{s+1}}\gamma +O(s+1)\right]`$ (39)
$`=`$ $`\zeta ^{}(2)=0.030448.`$ (40)
Using the values of the Riemann zeta function and its derivative at the origin
$`\zeta (0)={\displaystyle \frac{1}{2}},\zeta ^{}(0)={\displaystyle \frac{1}{2}}\mathrm{ln}(2\pi )`$
and taking into account the behaviour of the gamma function near zero (see Eq. (38)) we deduce from Eq. (35)
$`Z_2(1)`$ $`=`$ $`{\displaystyle \frac{1}{8}}\underset{s1}{lim}\left[\zeta (0)+\zeta ^{}(0)(s+1)+O\left((s+1)^2\right)\right]`$ (42)
$`\times \left[{\displaystyle \frac{2}{s+1}}\gamma +𝒪\left(s+1\right)\right]\left[1+{\displaystyle \frac{7}{4}}(s+1)\right]`$
$`=`$ $`{\displaystyle \frac{7}{32}}{\displaystyle \frac{\gamma }{16}}+{\displaystyle \frac{1}{8}}\mathrm{ln}(2\pi )+{\displaystyle \frac{1}{8}}{\displaystyle \frac{1}{s+1}}|_{s1}.`$ (43)
When calculating $`Z_3(1)`$ we shall use Eq. (31) for sevral first values of $`n,nn_0`$ and for $`n>n_0`$ we substitute the asymptotic expansion (26) into (31) with the result
$`Z_3^{\text{as}}(s)`$ $`=`$ $`{\displaystyle \frac{1}{64}}\left({\displaystyle \underset{n=n_0+1}{\overset{\mathrm{}}{}}}n^{3s}\right){\displaystyle _0^{\mathrm{}}}𝑑yy^{1s}{\displaystyle \frac{d}{dy}}\left[t^4(13284t^2+1062t^41356t^6+565t^8)\right]`$ (44)
$`=`$ $`{\displaystyle \frac{1}{64}}\left({\displaystyle \underset{n=n_0+1}{\overset{\mathrm{}}{}}}n^{3s}\right)\mathrm{\Gamma }\left({\displaystyle \frac{3s}{2}}\right)[13\mathrm{\Gamma }\left({\displaystyle \frac{3+s}{2}}\right)+142\mathrm{\Gamma }\left({\displaystyle \frac{5+s}{2}}\right){\displaystyle \frac{532}{3}}\mathrm{\Gamma }\left({\displaystyle \frac{7+s}{2}}\right)`$ (46)
$`+{\displaystyle \frac{113}{2}}\mathrm{\Gamma }\left({\displaystyle \frac{9+s}{2}}\right){\displaystyle \frac{113}{24}}\mathrm{\Gamma }\left({\displaystyle \frac{11+s}{2}}\right)].`$
The value of $`n_0`$ should be chosen so as to provide the accuracy needed. This algorithm with $`n_0=6`$ gives for $`Z_3(1)`$
$$Z_3(1)=0.022806$$
(47)
Summing up Eqs. (40), (43), and (47) we obtain
$`\zeta ^\text{D}(1)`$ $`=`$ $`{\displaystyle \frac{1}{4\pi a^2}}\left({\displaystyle \frac{7}{32}}+0.022806{\displaystyle \frac{\gamma }{16}}+{\displaystyle \frac{1}{8}}\mathrm{ln}(2\pi )+\zeta ^{}(2)+{\displaystyle \frac{1}{8}}{\displaystyle \frac{1}{s+1}}|_{s1}\right)`$ (48)
$`=`$ $`{\displaystyle \frac{1}{a^2}}\left(0.0005230.009947{\displaystyle \frac{1}{s+1}}|_{s1}\right).`$ (49)
Thus the zeta function $`\zeta ^\text{D}(s)`$ has a pole at the point $`s=1`$, therefore it does not give the finite (renormalized) value for the respective Casimir energy
$$E^\text{D}=\frac{1}{2}\zeta ^\text{D}(1).$$
(50)
It implies that further renormalization is required.
## IV Zeta Function for Neumann Boundary Value Problem
When constructing the zeta function for the boundary value problem (4) with $`\mathrm{\Gamma }`$ being a semi-circular infinite cylinder, we shall again proceed from the frequency equations (now from Eq. (11). It should be taken into account that all these roots are not degenerate. Therefore we can write analogously to Eq. (22)
$$\zeta ^\text{N}(s)=\frac{1}{2\pi i}_{\mathrm{}}^{\mathrm{}}\frac{dk_z}{2\pi }\underset{n=0}{\overset{\mathrm{}}{}}_C𝑑\gamma (\gamma ^2+k_z^2)^{s/2}\frac{d}{d\gamma }\mathrm{ln}\frac{J_n^{}(\gamma a)H_n^{(1)^{}}(\gamma a)}{J_n^{}(\mathrm{})H_n^{(1)^{}}(\mathrm{})}.$$
(51)
The contour $`C`$ is the same as in Eq. (22) and the prime on the Bessel and Hankel functions denotes differentiation with respect to the entire argument.
The product of the derivatives of the modified Bessel functions $`I_n^{}(z)K_n^{}(z)`$ has the following asymptotics when $`n`$ is fixed and $`|z|`$ is large
$$I_n^{}(z)K_n^{}(z)=\frac{1}{2z}\left[1+\frac{4n^23}{2(2z)^2}+\frac{(4n^21)(4n^245)}{8(2z)^4}+O(z^6)\right].$$
(52)
Taking this into account in calculation of the denominator in Eq. (51), we obtain for $`\zeta ^\text{N}(s)`$ upon integration over $`k_z`$
$$\zeta ^\text{N}(s)=C(s)\underset{n=0}{\overset{\mathrm{}}{}}_0^{\mathrm{}}𝑑yy^{1s}\frac{d}{dy}\mathrm{ln}\left[2yI_n^{}(y)K_n^{}(y)\right]$$
(53)
with the same function $`C(s)`$ as in Eq. (24).
Further we shall use the uniform asymptotic expansion for the derivatives of the Bessel functions
$`\mathrm{ln}\left[2ynI_n^{}(ny)K_n^{}(ny)\right]`$ $`=`$ $`\mathrm{ln}(yt)+{\displaystyle \frac{t^2}{8n^2}}(3+10t^27t^4)+{\displaystyle \frac{t^4}{n^4}}({\displaystyle \frac{27}{64}}`$ (55)
$`+{\displaystyle \frac{109}{16}}t^2{\displaystyle \frac{733}{32}}t^4+{\displaystyle \frac{441}{16}}t^6{\displaystyle \frac{707}{64}}t^8)+𝒪(n^6).`$
In order to render the integral in the term with $`n=0`$ in Eq. (53) convergent we add and subtract the second term from the asymptotics (55). For $`n1`$ in Eq. (53) we add and subtract in respective intagrands the first two terms of the asymptotic expansion (55). After that we combine all the terms in the following way
$`\zeta ^\text{N}(s)`$ $`=`$ $`C(s)\left[V_0(s)+V_1(s)+V_2(s)+V_3(s)\right],`$ (56)
$`V_0(s)`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}𝑑yy^{1s}{\displaystyle \frac{d}{dy}}\left\{\mathrm{ln}[2yI_0^{}(y)K_0^{}(y)]{\displaystyle \frac{t^2}{8}}(3+10t^27t^4)\right\},`$ (57)
$`V_1(s)`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}n^{1s}{\displaystyle _0^{\mathrm{}}}𝑑yy^{1s}{\displaystyle \frac{d}{dy}}\mathrm{ln}\left({\displaystyle \frac{y^2}{1+y^2}}\right)=Z_1(s),`$ (58)
$`V_2(s)`$ $`=`$ $`{\displaystyle \frac{1}{8}}\left({\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}n^{1s}+1\right){\displaystyle _0^{\mathrm{}}}𝑑yy^{1s}{\displaystyle \frac{d}{dy}}\left[t^2(3+10t^27t^4)\right],`$ (59)
$`V_3(s)`$ $`=`$ $`{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}n^{1s}{\displaystyle _0^{\mathrm{}}}dyy^{1s}{\displaystyle \frac{d}{dy}}\{\mathrm{ln}[2ynI_n^{}(ny)K_n^{}(ny)]`$ (61)
$`+\mathrm{ln}(yt){\displaystyle \frac{t^2}{8n^2}}(3+10t^27t^4)\}.`$
Taking into account the behaviour of the product $`I_0^{}(y)K_0^{}(y)`$ at the origin and at infinity
$`2yI_0^{}(y)K_0^{}(y)=y+{\displaystyle \frac{1}{8}}(1+4y4\mathrm{ln}2+\mathrm{ln}y)y^3+O(y^5\mathrm{ln}y),`$
$$2yI_0^{}(y)K_0^{}(y)=1\frac{3}{8y^2}+\frac{45}{128y^4}+O(y^6)$$
(62)
it is easy to show that Eq. (57) defines $`V_0(s)`$ as an analytic function in the region $`3<\text{Re }s<1`$. Under this condition the integration by parts can be done here
$$V_0(s)=(1s)_0^{\mathrm{}}𝑑yy^s\left\{\mathrm{ln}[2yI_0^{}(y)K_0^{}(y)]\frac{t^2}{8}(3+10t^27t^4)\right\}.$$
(63)
The function $`V_1(s)`$ differs only in sign of the function $`Z_1(s)`$ from the proceeding Section. The integral in Eq. (58) is convergent when $`1<\text{Re }s<3`$. The sum over $`n`$ in this formula is finite when Re $`s>0`$. Thus the regions, where the integral and the sum exist, overlap and this formula can be used for constructing the analytic continuation needed by making use of the substitutions (33) and (34). Substituting the sum in Eq. (59) by the Riemann zeta function and doing the integration according to Eq. (34) one obtains
$$V_2(s)=\frac{1}{8}[\zeta (1+s)+1]\mathrm{\Gamma }\left(\frac{3s}{2}\right)\mathrm{\Gamma }\left(\frac{1+s}{2}\right)\left[35(1+s)+\frac{7}{8}(1+s)(3+s)\right].$$
(64)
The convergence of the integral in Eq. (61) can be determined in the same line as it has been done for the function $`Z_3(s)`$ in the preceding Section. This integral converges when $`3<\text{Re }s<3`$, and the sum encountered here is finite for $`\text{Re }s>2`$. Hence there is no need to do analytic continuation for $`V_3(s)`$.
Finally the zeta function $`\zeta ^\text{N}(s)`$ for the massless scalar field obeying the Neumann boundary conditions on a semi-circular cylinder is determined explicitly by Eqs. (56), (58), (63), and (64) in a finite domain of the complex plane $`s`$ containing the closed interval of the real axis $`1\text{Re }s0`$.
Now we turn to the calculation of the value of the function $`\zeta ^\text{N}(s)`$ at the point $`s=1`$. Integration in Eq. (63) gives
$`V_0(1)`$ $`=`$ $`2{\displaystyle _0^{\mathrm{}}}𝑑yy\left\{\mathrm{ln}[2yI_0^{}(y)K_0^{}(y)]+{\displaystyle \frac{3}{8}}t^2\right\}+{\displaystyle \frac{13}{16}}`$ (65)
$`=`$ $`20.475215+0.8123=1.76393.`$ (66)
From Eqs. (58) and (40) it follows that
$$V_1(1)=Z_1(1)=\zeta ^{}(2)=0.03044.$$
(67)
Developing the functions $`\zeta (1+s)`$ and $`\mathrm{\Gamma }((1+s)/2)`$ in Eq. (64) near the point $`s=1`$ one obtains
$`V_2(1)`$ $`=`$ $`{\displaystyle \frac{1}{8}}\left[\zeta (0)+\zeta ^{}(0)(s+1)+1+O\left((s+1)^2\right)\right]\left[{\displaystyle \frac{2}{1+s}}\gamma +O(s+1)\right]`$ (69)
$`\times [3{\displaystyle \frac{13}{4}}(s+1)]={\displaystyle \frac{13}{32}}{\displaystyle \frac{3}{16}}\gamma +{\displaystyle \frac{3}{4}}\zeta ^{}(0)+{\displaystyle \frac{3}{8}}{\displaystyle \frac{1}{s+1}}|_{s1}.`$
When calculating $`V_3(s)`$ for $`s=1`$ numerically we cannot use the method applied in the preceding section because it requires now to take into account the next terms in the uniform asymptotic expansion (55). Instead of this we calculate numerically the first 30 terms in the sum (61) with the result
$$V_3(1)=0.04366.$$
(70)
Substituting in Eq. (61) the logarithm by its uniform asymptotic expansion (55) we derive a rough estimation for $`V_3(s)`$ without numerical integration
$`V_3^{\text{as}}(s)`$ $`=`$ $`\zeta (3+s){\displaystyle _0^{\mathrm{}}}𝑑yy^{1s}{\displaystyle \frac{d}{dy}}\left[t^4\left({\displaystyle \frac{27}{64}}+{\displaystyle \frac{109}{16}}t^2{\displaystyle \frac{733}{32}}t^4+{\displaystyle \frac{441}{16}}t^6{\displaystyle \frac{707}{64}}t^8\right)\right]`$ (71)
$`=`$ $`\zeta (3+s)\mathrm{\Gamma }\left({\displaystyle \frac{3s}{2}}\right)[{\displaystyle \frac{27}{64}}\mathrm{\Gamma }\left({\displaystyle \frac{3+s}{2}}\right){\displaystyle \frac{109}{32}}\mathrm{\Gamma }\left({\displaystyle \frac{5+s}{2}}\right)+{\displaystyle \frac{733}{192}}\mathrm{\Gamma }\left({\displaystyle \frac{7+s}{2}}\right)`$ (73)
$`{\displaystyle \frac{441}{384}}\mathrm{\Gamma }\left({\displaystyle \frac{9+s}{2}}\right)+{\displaystyle \frac{707}{7680}}\mathrm{\Gamma }\left({\displaystyle \frac{11+s}{2}}\right)].`$
For $`s=1`$ it gives
$$V_3^{\text{as}}(1)=\frac{839}{2^635}\zeta (2)=\frac{839}{960}\frac{\pi ^2}{6}=1.43760,$$
(74)
that is very far from Eq. (70) having only the right sign.
Summing up $`V_i`$, $`i=0,1,2,3`$ with allowance for Eq. (36) we arrive at the final result
$`\zeta ^\text{N}(1)`$ $`=`$ $`{\displaystyle \frac{1}{4\pi a^2}}[{\displaystyle \frac{13}{32}}+0.95043\zeta ^{}(2){\displaystyle \frac{3}{16}}\gamma `$ (76)
$`{\displaystyle \frac{3}{8}}\mathrm{ln}(2\pi )0.04366+{\displaystyle \frac{3}{8}}{\displaystyle \frac{1}{s+1}}|_{s1}]`$
$`=`$ $`{\displaystyle \frac{1}{a^2}}\left(0.043450.0298{\displaystyle \frac{1}{s+1}}|_{s1}\right).`$ (77)
Thus both the zeta functions for Dirichlet and Neumann boundary conditions have the pole at the point $`s=1`$. Hence an additional renormalization is needed in order for a finite physical value of the relevant Casimir energies to be obtained.
## V Vacuum Energy of Electromagnetic Field with Boundary Conditions on a Semi-Circular Cylinder
Analysis of the spectral problem for the electromagnetic field with boundary conditions on a semi-circular cylinder (see Sec. II) implies that the zeta function for this field is the sum of two zeta functions calculated in the preceding Sections
$$\zeta ^{\text{EM}}(s)=\zeta ^\text{D}(s)+\zeta ^\text{N}(s).$$
(78)
Substitution of Eqs. (49) and (73) into Eq. (78) gives
$`\zeta ^{\text{EM}}(1)`$ $`=`$ $`{\displaystyle \frac{1}{4\pi a^2}}\left[{\displaystyle \frac{1}{4}}+0.95043{\displaystyle \frac{\gamma }{4}}{\displaystyle \frac{1}{4}}\mathrm{ln}(2\pi )0.04366+{\displaystyle \frac{1}{2}}{\displaystyle \frac{1}{s+1}}|_{s1}\right]`$ (79)
$`=`$ $`{\displaystyle \frac{1}{a^2}}\left(0.044010.03978{\displaystyle \frac{1}{s+1}}|_{s1}\right).`$ (80)
In both the zeta functions $`\zeta ^\text{D}(s)`$ and $`\zeta ^\text{N}(s)`$ the pole terms have the same sign. As a result the pole contribution in the sum (78) retains. Thus, the situation here proves to be analogous to that when calculating, in the framework of zeta technique, the vacuum energy for spheres in spaces of even dimensions.
As was noted above, we have derived the exact expressions for the zeta functions in question which determine these functions as analytic functions of the complex variable $`s`$ in a finite region of the plane $`s`$ containing the closed interval of the real axis $`1\text{Re }s0`$. It enables one to construct in a straightforward way the spectral zeta functions for relevant boundary value problem on the plane by making use of the relation
$$\zeta _{\text{s-cir}}(s)=2\sqrt{\pi }\frac{\mathrm{\Gamma }\left({\displaystyle \frac{s+1}{2}}\right)}{\mathrm{\Gamma }\left({\displaystyle \frac{s}{2}}\right)}\zeta _{\text{s-cyl}}(s),$$
(81)
where $`\zeta _{\text{s-cir}}`$ is the Dirichlet or the Neumann zeta function for a semi-circle, and $`\zeta _{\text{s-cyl}}`$ is the respective zeta function for semi-circular cylinder. We shall use this relation for calculating the values $`\zeta _{\text{s-cir}}^\text{D}(1)`$ and $`\zeta _{\text{s-cir}}^\text{N}(1)`$ which determine the vacuum energy of the massless scalar fields defined on the half-plane and obeying, respectively, the Dirichlet or Neumann boundary conditions on a semi-circle (see Fig. 1).
For $`\zeta _{\text{s-cir}}^\text{D}(1)`$ we get from Eqs. (81). (26), and (23)
$$\zeta _{\text{s-cir}}^\text{D}(1)=\frac{1}{\pi a}\underset{i=1}{\overset{3}{}}Z_i(0).$$
(82)
When $`s=0`$ integration in Eq. (28) can be done explicitly with the result
$$Z_1(0)=\zeta (1)_0^{\mathrm{}}𝑑y\mathrm{ln}\frac{y}{\sqrt{1+y^2}}=\frac{1}{12}\left(\frac{\pi }{2}\right)=\frac{\pi }{24}.$$
(83)
From Eq. (35) it follows that
$$Z_2(0)=\frac{\pi }{128}\left(\frac{1}{s}|_{s0}+\gamma \right).$$
(84)
Numerical integration in Eq. (31) with $`s=0`$ gives
$$Z_3(0)=0.00304.$$
(85)
Summing up Eqs. (83), (84), and (85) we arrive at the result
$`\zeta _{\text{s-cir}}^\text{D}(1)`$ $`=`$ $`{\displaystyle \frac{1}{a}}\left({\displaystyle \frac{1}{24}}{\displaystyle \frac{\gamma }{128}}+0.00097{\displaystyle \frac{1}{128}}{\displaystyle \frac{1}{s}}|_{s0}\right)`$ (86)
$`=`$ $`{\displaystyle \frac{1}{a}}\left(0.038127{\displaystyle \frac{1}{128}}{\displaystyle \frac{1}{s}}|_{s0}\right).`$ (87)
Following the same way one can write
$$\zeta _{\text{s-cir}}^\text{N}(1)=\frac{1}{\pi a}\underset{n=0}{\overset{3}{}}V_i(0).$$
(88)
Using Eq. (63) one gets
$`V_0(0)`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}𝑑y\left\{\mathrm{ln}\left[2yI_0^{}(y)K_0^{}(y)\right]+{\displaystyle \frac{3}{8}}t^2\right\}{\displaystyle \frac{\pi }{64}}`$ (89)
$`=`$ $`0.475175{\displaystyle \frac{\pi }{64}}.`$ (90)
From Eqs. (58) and (83) it follows that
$$V_1(0)=Z_1(0)=\frac{\pi }{24}.$$
(91)
Equation (64) gives
$$V_2(0)=\frac{5\pi }{128}\left(1+\gamma +\frac{1}{s}|_{s0}\right).$$
(92)
For $`V_3(0)`$ numerical integration in Eq. (61) with $`s=0`$ gives
$$V_3(0)=0.005659.$$
(93)
Finally, we have
$`\zeta _{\text{s-cir}}^\text{N}(1)`$ $`=`$ $`{\displaystyle \frac{1}{a}}\left[0.15132{\displaystyle \frac{5}{128}}\left(\gamma +{\displaystyle \frac{5}{3}}\right)+0.00180{\displaystyle \frac{5}{128}}{\displaystyle \frac{1}{s}}|_{s0}\right]`$ (94)
$`=`$ $`{\displaystyle \frac{1}{a}}\left(0.2371030.0124{\displaystyle \frac{1}{s}}|_{s0}\right).`$ (95)
Both the functions $`\zeta _{\text{s-cir}}^\text{D}(s)`$ and $`\zeta _{\text{s-cir}}^\text{N}(s)`$ have the pole at the point $`s=1`$ with the coefficients of the same (negative) sign. For electromagnetic field defined on a plane the boundary conditions reduce to the Neumann conditions. Hence the relevant zeta function is $`\zeta _{\text{s-cir}}^\text{N}(s)`$.
## VI Conclusion
In the paper the spectral zeta functions are constructed for massless scalar fields obeying the Dirichlet and Neumann boundary conditions on a semi-circular infinite cylinder. Proceeding from this, the zeta function for electromagnetic field is also derived for such a configuration. In all three cases, the final expressions for the relevant Casimir energy contains the pole contribution. Hence for obtaining the physical result an additional renormalization is needed.
It is essential that for the zeta functions $`\zeta (s)`$ the exact formulas are derived which determine these functions in a finite region of the complex variable $`s`$ but not at the vicinity of one point $`s=1`$. This allowed one to get in a straightforward way the zeta functions for the two dimensional (plane) version of the boundary value problem at hand, i.e. the zeta functions for scalar fields defined on a half-plane and obeying the Dirichlet and Neumann boundary conditions on a semi-circle. In this case the final expression for the vacuum energy contains the pole contributions also.
Notwithstanding the spectrum of a semi-circular cylinder is very close to the spectrum of circular one, the zeta function technique does not give a finite value for vacuum energy in the first case and does for the second configuration. In a recent paper the divergences found in our consideration are attributed to the existence of edges or corners in the boundaries under investigation.
Closing, it is worth noting that, as far as we know, such boundary conditions with asymmetric geometry (semi-circular cylinder) has been considered in the Casimir problem for the first time.
###### Acknowledgements.
Research is supported by the fund MURST ex 40% and 60%, art. 65 D.P.R. 382/80. This work has been accomplished during the visit of V.V.N. to the Salerno University. It is a pleasure for him to thank Professor G. Scarpetta, Drs. G. Lambiase and A. Feoli for warm hospitality. The financial supports of IIASS and ISTC (Project # 840) are acknowledged. G.L. thanks the UE fellowship, P.O.M. 1994/1999, for financial support.
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# Limits on Neutrino Masses from Large-Scale Structure
## 1 What Can Cosmology Tell Us About Neutrino Masses?
There is now evidence for neutrino oscillations from SuperKamiokande fukudaetal98 ($`\mathrm{\Delta }m^210^3`$eV<sup>2</sup>), the solar neutrino problem bahcallks98 ($`\mathrm{\Delta }m^210^5`$eV<sup>2</sup>)<sup>1</sup><sup>1</sup>1for vacuum oscillations, $`\mathrm{\Delta }m^210^{10}`$eV<sup>2</sup> and LSND athanassopoulosetal96 ($`0.1`$ eV$`{}_{}{}^{2}\mathrm{\Delta }m^220`$eV<sup>2</sup>). While these oscillations only test the difference in squared masses, they give evidence that the mass of at least one (and probably at least three) neutrino species is non-zero. Laboratory limits from tritium beta decay rule out the possibility of an electron neutrino more massive than 4.4 eV belesevetal95 . Present cosmological bounds on the masses of other neutrino species are stricter than those from laboratory experiments; a 45 eV neutrino would lead to $`\mathrm{\Omega }_\nu =1`$ kolbt90 , so for a universe at less than critical density the neutrinos must all be lighter than this. The exception to this is if a neutrino is so massive ($`>1`$ MeV) that it was non-relativistic during freeze-out, i.e. Cold Dark Matter (CDM).
Each model of structure formation predicts Cosmic Microwave Background (CMB) anisotropy and large-scale structure inhomogeneities. Massive neutrinos lead to slightly different predictions for CMB anisotropies and significantly different predictions for large-scale structure. The CMB radiation power spectrum is given by $`C_{\mathrm{}}=_kC_{\mathrm{}}(k)P_p(k)`$ and the matter power spectrum by $`P(k)=T^2(k)P_p(k)`$, where $`C_{\mathrm{}}(k)`$ and $`T(k)`$ are transfer functions predicted by a given model of structure formation and $`P_p(k)`$ is the primordial power spectrum of density perturbations that originated in the early universe. Massive neutrinos alter these transfer functions, so for a given $`P_p(k)`$ we can test a model with massive neutrinos by comparing its predictions to the observed $`C_{\mathrm{}}`$ and $`P(k)`$. The effect is to exponentially damp the matter transfer function, $`T(k)`$, on scales smaller than the neutrino free-streaming scalehuet98 ,
$$k0.026\left(\frac{m_\nu }{1eV}\right)^{\frac{1}{2}}\mathrm{\Omega }_m^{1/2}hMpc^1,$$
(1)
which is equal to the horizon size when the neutrinos become non-relativistic. For mostly Cold Dark Matter and a fraction of Hot Dark Matter (massive neutrinos), the damping is no longer exponential but still quite significant. The effect on the CMB is more subtle; relativistic neutrinos increase the radiation density before decoupling, which affects the shape of the acoustic peaks in the CMB angular power spectrum.
## 2 Approach
We assume here that $`\mathrm{\Lambda }`$CDM is the correct model of structure formation and that the primordial power spectrum is well-described by a power-law. We start with a version of $`\mathrm{\Lambda }`$CDM which is in good agreement with observations of Type Ia supernovae, the cluster baryon fraction, the primordial deuterium abundance, and Hubble’s constant, with $`\mathrm{\Omega }_m=0.4`$, $`\mathrm{\Omega }_b=0.04`$, and $`h=0.7`$.<sup>2</sup><sup>2</sup>2Smaller values of $`\mathrm{\Omega }_m`$ will lead to tighter limits on $`\mathrm{\Omega }_\nu `$. However, this $`\mathrm{\Lambda }`$CDM model is not a great fit to large-scale structure data gawisers98 , so we investigate whether adding a Hot Dark Matter component will improve the fit.
Our data compilation includes COBE and smaller-scale CMB anisotropy detections, measurements of $`\sigma _8`$ from galaxy clusters vianal96 ; bahcallfc97 , a measurement of the matter power spectrum from peculiar velocities kolattd97 , redshift-space matter power spectra from SSRS2+CfA2 dacostaetal94 , LCRS linetal96 , PSCz sutherlandetal99 and APM Clusters tadrosed98 , and a real-space matter power spectrum derived from the APM angular correlation function eisensteinz99 . Using all available CMB and large-scale structure data gives us information on intermediate scales, which helps to differentiate between variations in the primordial power spectrum and the reduction in small-scale power caused by massive neutrinos. We analyze this data compilation using the methods of Gawiser & Silk gawisers98 .
Even given all of this data, we need to assume something about the primordial power spectrum. We have tried using Harrison-Zel’dovich (scale-invariant, $`P_p(k)=Ak`$) and scale-free ($`P_p(k)=Ak^n`$) primordial power spectra. For inflationary models, the most reasonable parameterization is $`\mathrm{log}P_p(k)=\mathrm{log}A+n\mathrm{log}k+\alpha (\mathrm{log}k)^2+\mathrm{}`$ with successive terms decreasing in importance in the slow-roll regime. Unfortunately, an unconstrained primordial power spectrum $`P_p(k)`$ can easily mimic the effect that massive neutrinos have on the matter transfer function $`T(k)`$, making it nearly impossible to limit the neutrino mass.
## 3 Results
The results presented here were first determined by Gawisergawiser99 . We find that as HDM is added, the combined fit to CMB and large-scale structure deteriorates. This occurs because adding HDM reduces the power on physical scales shorter than the neutrino free-streaming length, which degrades the fit to large-scale structure data (see Figure 1). For $`\mathrm{\Omega }_\nu =0.05`$, a blue tilt of the primordial power spectrum ($`n=1.3`$) is necessary to counteract the damping of small-scale perturbations by free-streaming of the massive neutrinos. Even with this best-fit value of $`n`$, the fit to the data is worse than with no HDM, because CMB observations disfavor such a high value of $`n`$. For a higher HDM fraction, an even higher value of $`n`$ is preferred ($`n=1.5`$ for the $`\mathrm{\Omega }_\nu =0.1`$ model of Figure 1), leading to an even worse fit to the CMB data.
Our limits on the neutrino mass are based upon an attempt to search the reasonable parameter space around this fiducial model to produce the best fit possible to the data for a given neutrino mass. Since disagreement with CMB data is the main problem once a blue tilt is considered, we have tried to alleviate this by adding a significant tensor component or early reionization. Each of these effects reduces the small-scale CMB power relative to COBE scales, which helps to reconcile $`n>1`$ with the CMB data. However, no parameter combination helps enough to make $`\mathrm{\Lambda }`$CHDM a better fit than the fiducial $`\mathrm{\Lambda }`$CDM model, and this allows us to set an upper limit on the sum of the neutrino masses, $`\mathrm{\Sigma }m_\nu =94\mathrm{\Omega }_\nu h^2`$ eV. An upper limit on $`\mathrm{\Omega }_\nu `$ leads to an upper limit on the mass of the most massive neutrino. If the mass is split between multiple nearly-equal-mass neutrinos, the limit on the sum of the masses is tighter because, for example, two 1 eV neutrinos depress the power spectrum more than one 2 eV neutrino because they both become non-relativistic later.
If $`\mathrm{\Lambda }`$CDM is right, and $`H_0`$ is about 65 $`h^1`$Mpc and $`n=1`$, then $`\mathrm{\Omega }_\nu 0.05,m_\nu 2`$ eV.
If $`n`$ can vary, $`\mathrm{\Omega }_\nu 0.1,m_\nu 4`$ eV.
If $`P_p(k)`$ is not a power-law (non-inflationary or a complicated two-field inflationary model), then all bets are off.
This is compatible with the recent claim by Croft et al. crofthd99 that the Lyman $`\alpha `$ forest power spectrum plus COBE limits the neutrino mass to 5 eV or less. Our method appears more robust as the Lyman $`\alpha `$ forest power spectrum has an uncertain normalization and covers a narrower range of scales than our large-scale structure compilation; moreover, the origin of the Lyman $`\alpha `$ forest is not yet well understood. Fukugita et al. fukugitals00 use only COBE and $`\sigma _8`$ and assume $`n=1`$, yielding $`m_\nu 3`$ eV. This is highly model-dependent because the primordial power spectrum is nearly degenerate with neutrino free-streaming when structure formation is only probed at two different spatial scales. $`\mathrm{\Lambda }`$CHDM has also been explored by Valdarnini et al. valdarninikn98 and Primack & Gross primackg98 with different analysis methods and significantly smaller data compilations.
Hu et al. huet98 discuss the future prospects of including the Sloan Digital Sky Survey (SDSS) $`P(k)`$ in a method similar to that used here; they expect to detect or limit $`m_\nu `$ down to 0.5 eV. Cooray cooray99 gives an expected future limit from measuring $`P(k)`$ with surveys of weak gravitational lensing of $`m_\nu 3`$eV.
## 4 Conclusions
The currently-favored $`\mathrm{\Lambda }`$CDM model does not prefer the addition of a Hot Dark Matter component. This leads to an upper limit on the mass of the most massive neutrino of 4eV if a power-law primordial power spectrum is assumed. This limit is comparable to tritium beta-decay limits on the electron neutrino mass, and it should improve considerably with data from SDSS and the MAP satellite. Our results are already in conflict with the portion of the parameter space compatible with LSND athanassopoulosetal96 that requires a mass difference greater than $`16`$ eV<sup>2</sup>.
In evaluating these results and other work on this subject, the reader is encouraged to check what assumptions have been made about cosmological models, the selection and normalization of data, and the primordial power spectrum. Restrictive assumptions can lead to tight but meaningless limits on the neutrino mass.
## Acknowledgments
I would like to thank Joe Silk for initially suggesting an investigation of adding Hot Dark Matter to $`\mathrm{\Lambda }`$CDM and for his continuing collaboration in this research.
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# Dynamics of the formation of an event horizon
## 1 INTRODUCTION
The formation of a black-hole event horizon has attracted a great deal of attention on the part of physicists for a long time. An enormous amount of material has been written on this subject (see, for example, Ref. ); nevertheless, the treatment of this problem within the general theory of relativity has created more questions than known solutions.
One of the main questions concerning this problem is still the reciprocal influence of accreting matter on a black hole. The motion of test particles in the field of a black hole has been considered hitherto for the most part, but they, as we know, do not exhibit a reciprocal influence, which can be enormous when a falling particle achieves the speed of light as it crosses the event horizon.
1. In this paper we consider the special, but physically real case of spherically symmetric accretion on a central body without allowance for rotation. The following notation is adopted: the speed of light $`c`$ and the gravitational constant $`G`$ are set equal to unity. In these units the gravitational radius for a given mass $`M`$ is $`r_g=2M`$, i.e., the radius of the event horizon in free space for the same mass concentrated at the center.
Let us devise a likely model for the evolution of the system. We assume that our system is a cooling massive star having a radius $`R_0`$ and a gravitational radius $`R_{G_0}=2M(R_0)`$, where $`R_{G_0}<R_0`$. The matter comprising this body is initially at rest (“dust” with the equation of state $`P=\alpha \epsilon `$, where $`P`$ is the pressure in the matter, $`\epsilon `$ is the energy density, and $`\alpha `$ is a constant). In the next moment the matter begins to fall freely.\[1)\]<sup>)</sup> If it is assumed that the gravitational fields are not excessively strong and that the dust density\[2)\]<sup>)</sup> is fairly small in the initial moment, a force field with a finite energy is needed to retain it in the initial moment. After this field is removed, the dust leaves the system and ceases to interact with it after a time of the order of the size of the system, i.e., after a time much shorter than the time during which the dust manages to partially settle and the gravitational fields increase dramatically. Thus, this model is physically consistent.
What subsequently happens to the system? The dust begins to fall toward the center of the body, increasing its mean density and the gravitational radius $`r_g(r)`$ for the mass $`M(r)`$ at a certain radius $`r`$. If we would neglect the reciprocal influence of the pressure of the moving matter on the dynamics of the system and on its gravitational field, then after all of the matter has unavoidably fallen and the inequality
$$r_g(r)=2M(r)r$$
$`(1)`$
holds at one of the points $`r`$ of the system, an event horizon would form at that point according to Schwarzchild’s solution for a gravitational field in a vacuum, i.e., the velocity of the falling matter relative to the $`r=\mathrm{const}`$ surfaces would reach the speed of light (see below). Is this what actually happens? The achievement of the speed of light by the matter causes a change in the sign of the interval and is therefore an invariant event, which does not depend on the choice of the reference frame.
An attempt to solve this problem in a reference frame which is stationary at infinity leads at once to a contradiction. In fact, the analytically exact, nonstationary model, in principle, cannot be studied. If, however, a simplification is made and it is assumed that the system is quasistationary at a certain moment in time in the range from one radius to a certain radius known to be large, but still far smaller than the dimensions of the system, then it can be stated, at the very least, that the components $`g_{tt}(r)`$ and $`g_{rr}(t)`$ of the metric have singularities (zeros and poles) in this reference frame.\[3)\]<sup>)</sup> When the parameters of the system are chosen so that there would be a region in space where the inequality (1) is sure to be satisfied, it becomes clear that the metric does not have singularities, regardless of whether the inequality (1) is satisfied.
This can be shown by assuming that if a singularity appears at a certain point $`r_0`$, the component of the metric near it can be represented in the form
$$g_{ii}(r)\mathrm{const}(rr_0)^{y_i},$$
where $`y_i`$ is a certain number. When such a metric is substituted into the equations, it is found that they do not have a solution for any $`y_i0`$.
This apparently indicates that the singularities and thus the horizon of the rapidly moving matter are eliminated (the right-hand side of the Einstein equations, which is equal to zero in a vacuum, becomes singular in the presence of ultrarelativistic falling matter when the radial component of the three-velocity tends to unity and the radial and temporal components of the four-velocity tend to infinity; this is also the reason for the elimination of the singularities of the metric).
However, in reality all this stems from the inapplicability of the quasistationary approximation in the case of strong gravitational fields. It is inapplicable because the passage of time in the system is highly nonuniform due to the nonuniformity of the component $`g_{tt}(r)`$ of the metric. This causes the picture, which appears to be stationary far from the center, to become highly nonstationary to an observer approaching the symmetry center of the system.
Nevertheless, this does not remove the question posed: do an horizon and a black hole appear in the real nonstationary case?
2. An answer to the question posed can be found by selecting a comoving reference frame. The problem was solved in this frame in Ref. 1 (Sec. 103) in the special case of $`\alpha =0`$ (see below). Matter is at rest in the reference frame chosen, and its motion can be evaluated only from the variation of the “circumferential” or photometric distances $`r`$, which are related to the center of the system and are defined as the circumferences of the respective circles around the center: $`2\pi r`$. When the radius $`r`$ is defined as such, it is convenient to represent the metric in the form
$$ds^2=e^\nu dt^2e^\lambda dR^2r^2(d\theta ^2+\mathrm{sin}^2\theta d\phi ^2).$$
$`(2)`$
Here $`R`$ is the coordinate of a dust particle in the comoving reference frame or its index, and $`e^\nu `$, $`e^\lambda `$, and $`r`$ are functions of $`R`$ and $`t`$. It is noteworthy that at zero pressure, i.e., when $`\alpha =0`$, we have $`\nu =0`$, i.e., the reference frame is simultaneously synchronous.
To solve the problem posed we write out the Einstein equation in the comoving reference frame:
$$r^2e^\lambda (1+r\nu ^{}/r^{})e^\nu (2r\ddot{r}+\dot{r}^2r\dot{r}\dot{\nu })=1+8\pi \alpha r^2\epsilon ,$$
$`(3a)`$
$$2\dot{\mu }^{}+\dot{\mu }\mu ^{}\dot{\lambda }\mu ^{}\nu ^{}\dot{\mu }=0,$$
$`(3b)`$
$$\left(\lambda +2\mu +\frac{2}{1+\alpha }\mathrm{ln}\epsilon \right)=0,$$
$`(3c)`$
$$\left(\nu +\frac{2\alpha }{1+\alpha }\mathrm{ln}\epsilon \right)^{}=0.$$
$`(3d)`$
Here $`\mu =2\mathrm{ln}r`$, a prime denotes differentiation with respect to $`R`$, and a dot denotes differentiation with respect to $`t`$. Equations (3) were derived in Ref. \[Eqs. (2), (5), and (6) of problem 5 in Sec. 100\].
It follows from (3d) that
$$\nu =\frac{2\alpha }{1+\alpha }\mathrm{ln}\epsilon +f^{}(t)$$
and that by transforming the time $`t`$ in the interval element (2) the function $`f^{}(t)`$ can be set equal to $`[2\alpha /(1+\alpha )]\mathrm{ln}\epsilon _{}`$, where $`\epsilon _{}`$ is a constant with the dimensions of energy density, which expresses the measurement scale of $`\epsilon `$. Then
$$\nu =\frac{2\alpha }{1+\alpha }\mathrm{ln}\frac{\epsilon }{\epsilon _{}}.$$
$`(4)`$
We next assign the indices $`R`$ to the dust particles so that $`r=R`$ in the initial moment. Under such initial conditions $`r^{}(R,t)`$ corresponds to $`(n_0/n)^{1/3}`$, where $`n(R,t)`$ is the concentration of dust particles and $`n_0`$ is its value at the initial moment.
Let us now ascertain the conditions which must be imposed on the initial distribution of the dust. The most important among them is that the inequality (1) need not hold within the matter at the initial moment. It means that there is no horizon in all space in the initial moment. It thus imposes an upper limit on the initial density of the dust and on the initial dimensions of the system. More specifically, if the initial density distribution of the dust is set equal to $`\epsilon _0(R)`$, then, according to (1), the maximum radius of the body $`R_{\mathrm{max}}`$ is uniquely specified by the expression
$$R_{\mathrm{max}}=2\underset{0}{\overset{R_{\mathrm{max}}}{}}4\pi \epsilon _0(R)R^2𝑑R.$$
$`(5)`$
Then, it follows from (3c) and (4) that
$$\frac{}{t}[\alpha (\lambda +2\mu )\nu ]=0$$
or
$$\nu =\alpha [\lambda +2\mu +f^{}(R)],$$
$`(6)`$
where $`f^{}(R)`$ is an arbitrary function that depends on the initial conditions.
3. Let us now find the initial values for all the variables in our problem. We have already assigned these values for $`r`$ and $`\epsilon `$. From (4) it follows that
$$\nu _0=\frac{2\alpha }{1+\alpha }\mathrm{ln}\frac{\epsilon _0}{\epsilon _{}}.$$
$`(7)`$
To find the initial value of $`\lambda `$ we take advantage of the fact that the problem has already been solved for $`\alpha =0`$, and we can therefore utilize the familiar expression for $`\lambda _0|_{\alpha =0}`$ from Ref. (Sec. 103.6):
$$\lambda _0(R)=\mathrm{ln}[1S(R)],$$
$`(8)`$
where for $`\alpha =0`$ we have
$$S(R)=2M(R)/R,$$
$`(9)`$
and $`M(R)`$ is the mass within the radius $`R`$ at the initial moment.
The expression for $`S(R)`$ for an arbitrary value of $`\alpha `$ is the same. It can be obtained from Eq. (4) in problem 5 of Sec. 100 in Ref. , where the Einstein equations in matter in the comoving reference frame were found for a centrally symmetric system. We write out this equation:
$$e^\lambda \left[\mu ^{\prime \prime }+\frac{3}{4}\mu ^2\frac{\mu ^{}\lambda ^{}}{2}\right]+\frac{1}{r^2}+\frac{1}{2e^\nu }\left[\dot{\lambda }\dot{\mu }+(\dot{\mu })^2/2\right]=8\pi \epsilon .$$
$`(10)`$
Expressing $`\mu `$ in terms of $`r`$ ($`\mu =\mathrm{ln}r^2`$) and combining similar terms, we can bring this expression into the form
$$8\pi r^{}\epsilon r^2=[r(r^2e^\lambda 1)]^{}+\frac{r^{}}{e^\nu }\left[\dot{\lambda }r\dot{r}+(\dot{r})^2\right].$$
$`(11)`$
Taking into account the expression (8), as well as the fact that, according to the expression (100.23) in Ref. , the equality
$$2M(r)=\underset{0}{\overset{r}{}}8\pi \epsilon (\stackrel{~}{r},t)\stackrel{~}{r}^2𝑑\stackrel{~}{r}|_{t=\mathrm{const}}$$
holds for the initial moment in time, when $`\dot{r}=0`$ and $`r^{}=1`$, we obtain the expression (9) for $`S(R)`$ after preliminarily integrating (11) over $`R`$ from 0 to $`R`$.
Substituting the expression (8) into (6) and taking into account (7), we find that
$$f^{}(R)=\frac{2}{1+\alpha }\mathrm{ln}\frac{\epsilon _0}{\epsilon _{}}+\mathrm{ln}[1S(R)]\mathrm{ln}R^4.$$
$`(12)`$
4. Now, plugging (6) into (3b) and dividing everything by $`\dot{\mu }\mu ^{}`$, we obtain the expression
$$\frac{1}{\dot{r}}(2\mathrm{ln}\mu ^{}+\mu \lambda )^{}=\frac{\nu ^{}}{r^{}}=\alpha \frac{[\lambda +2\mu +f^{}(R)]^{}}{r^{}}.$$
$`(13)`$
Taking into account that $`e^\nu (2r\ddot{r}+\dot{r}^2r\dot{r}\dot{\nu })=(e^\nu r\dot{r}^2)^{}/\dot{r}`$ and introducing the notation
$$U(R,t)=(\dot{r})^2,Q(R,T)=r^2e^\lambda ,$$
$`(14)`$
we see that Eq. (3a) can be written as an equation for $`U`$:
$$\frac{\dot{U}}{\dot{r}}+aU=\sigma ,$$
$`(15)`$
where
$$a(R,t)=\frac{1}{r}\left(1\frac{r\dot{\nu }}{\dot{r}}\right),\sigma (R,t)=\frac{1}{r}\left[Q\left(1+\frac{r\nu ^{}}{r^{}}\right)18\pi \alpha r^2\epsilon \right]e^\nu .$$
This equation has a solution which satisfies the initial conditions:
$$U(R,t)=\frac{1}{\gamma ^{}(R,t)}\underset{0}{\overset{t}{}}\gamma ^{}(R,\stackrel{~}{t})\sigma (R,\stackrel{~}{t})\stackrel{~}{r}𝑑\stackrel{~}{t},\gamma ^{}(R,t)=\mathrm{exp}\left[\underset{0}{\overset{t}{}}a(R,\stackrel{~}{t})\stackrel{~}{r}𝑑\stackrel{~}{t}\right].$$
$`(16)`$
Finding $`U`$, we can obtain an expression for the square of the velocity of the matter relative to the $`r=\mathrm{const}`$ surfaces from the form of the metric (2) (see Appendix 1):
$$V^2(R,t)=Ue^\nu e^\lambda /r^2.$$
$`(17)`$
The expression for $`\gamma ^{}`$ can easily be found:
$$\gamma ^{}(R,t)=C(R)re^\nu ,$$
where the multiplier $`C(R)`$ for $`\gamma ^{}(R,t)`$, which does not depend on $`t`$, can be taken out of the integral sign in (16) and canceled; therefore, it can be set equal to unity. Then
$$\gamma ^{}\sigma =r^2e^\lambda \left(1+\frac{r\nu ^{}}{r^{}}\right)18\pi \alpha r^2\epsilon .$$
Alternatively, taking into account that Eq. (13) can now be rewritten as
$$\frac{(2\mathrm{ln}\mu ^{}+\mu \lambda )^{}}{\dot{r}}=\frac{(\mathrm{ln}Q)^{}}{\dot{r}}=\frac{\nu ^{}}{r^{}},$$
$`(18)`$
we obtain
$$\gamma ^{}\sigma =\frac{(r(Q1))^{}}{\dot{r}}8\pi \alpha r^2\epsilon .$$
$`(19)`$
Then (16) is rewritten in the form
$$U=\frac{e^\nu }{r}[r(Q1)R(Q_01)+2\alpha m(R,t)],$$
$`(20)`$
where we have introduced the notation
$$m(R,t)=\underset{t}{\overset{0}{}}4\pi \stackrel{~}{\epsilon }\stackrel{~}{r}^2\stackrel{~}{r}𝑑\stackrel{~}{t}=\underset{r}{\overset{R}{}}4\pi \stackrel{~}{\epsilon }\stackrel{~}{r}^2\stackrel{~}{r}𝑑\stackrel{~}{r}|_{R=\mathrm{cosnt}}.$$
$`(21)`$
5. For $`\alpha =0`$, taking into account (6), (9), (18), and (20), we can easily obtain an analytically exact expression for $`U`$ and $`V`$:
$$U_{\alpha =0}=S(R)\left(\frac{R}{r}1\right)=\frac{2M(R)}{r}S(R).$$
$`(22)`$
Substituting this expression into (17), for the velocity we obtain
$$V_{\alpha =0}^2=\frac{2M(R)/rS(R)}{1S(R)}.$$
$`(23)`$
Hence $`V_{\alpha =0}=1`$ when $`r=r_0=2M(R)`$. This coincides with the results in Sec. 100 of Ref. , where the problem has already been solved for this case.
6. Finally, let us consider the location of the horizon\[4)\]<sup>)</sup> in the presence of a nonzero pressure. For this purpose we plug the expressions (20) and (8) for $`e^{\lambda _0}`$ into formula (17). After some relatively simple transformations, we ultimately obtain
$$r=\frac{2M(R)+2\alpha m(R,t)}{1Q(1V^2)}.$$
$`(24)`$
As will be shown in Appendix 1, the horizon appears at the point and at the time where the velocity of the falling matter relative to the $`r=\mathrm{const}`$ surfaces reaches unity, i.e., where $`V=1`$. In addition, the speed of light relative to the falling matter at this site is also, as always, equal to unity.
Hence, according to (24) and (9), the horizon radius $`r_{\mathrm{hor}}`$ is given by the formula
$$r_{\mathrm{hor}}=2M(R)+2\alpha m(R,t_{\mathrm{hor}}).$$
$`(25)`$
Thus, the horizon is displaced to a larger radius in comparison to the value in a vacuum $`r_{0_{\mathrm{hor}}}=2M(R)`$ by $`2\alpha m(R,t_{\mathrm{hor}})`$. In this case the quantity $`m(R,t)`$ has the meaning of the mass which would accumulate if we would join layers of dust with the initial radius $`R`$ and the thickness $`d\stackrel{~}{r}(\stackrel{~}{t})`$ to one another up to the radius $`r(R,t)`$ at the moment when this $`d\stackrel{~}{r}`$ layer passes through the joining point.
7. Regarding the possible values of $`\alpha `$ we note that $`\alpha =0`$ corresponds to dustlike matter without interactions between the particles. The results obtained for them are the same \[see (23)\] as the results for test particles in a central field of mass $`M`$ (see Sec. 101 in Ref. ). However, of course, such an equation of state of matter cannot correspond to reality near the horizon. It is reasonable to assume that the ultrarelativistic equation of state of matter, in which $`\alpha =1/3`$, holds near the horizon. Therefore, the location of the horizon should probably be sought with just such a value of $`\alpha `$.
8. When $`\alpha 0`$, it would appear that the falling matter should be slowed under the action of the pressure gradient, and the horizon should therefore form later, i.e., be displaced toward smaller values of $`r`$, but, as we have just shown, it is displaced toward larger values of $`r`$ by $`2\alpha m(R,t_{\mathrm{hor}})`$. What is the reason for this contradiction? It can be seen from the initial equations (3) that the reason should be sought in Eq. (3a). For this purpose we explore Eqs. (3a) and (4) in the initial moment for the case of $`\alpha 1`$. In that moment $`\dot{r}=0`$ and $`r^{}=1`$; therefore, we write
$$[1S(R)]\left(12ar\frac{\epsilon ^{}}{\epsilon }\right)2r\stackrel{~}{r}\left(\frac{\epsilon }{\epsilon _{}}\right)^{2\alpha }1+8\pi \alpha r^2\epsilon .$$
Since
$$(\epsilon /\epsilon _{})^{2\alpha }1+2\alpha \mathrm{ln}\frac{\epsilon }{\epsilon _{}},S(R)=\frac{2M(R)}{R},r=R,$$
then, after performing some relatively simple transformations, in the linear approximation with respect to $`\alpha `$ we obtain
$$\stackrel{~}{r}=\frac{GM(R)}{r^2}\left[12\alpha \mathrm{ln}\frac{\epsilon }{\epsilon _{}}\right]\frac{P}{\rho }\left[1\frac{2GM(r)}{rc^2}\right]4\pi \alpha Gr\rho ,$$
$`(26)`$
where $`\rho (R,r)=\epsilon (R,r)/c^2`$ is the density of the matter. Here, for the sake of clarity we use the ordinary (Gaussian) system of units with $`G1`$ and $`c1`$. It can be seen from (26) that the first term corresponds to the ordinary Newtonian force of gravity, and the second term corresponds to the interaction force between the particles, i.e, the pressure gradient (just this force is the cause of the slowing of the fall of the matter in the first stage). The remaining terms do not appear in the equation of motion in the Newtonian approximation (the corrections in square brackets are also neglected in that case), but, as we have already seen, the last term begins to dominate over the second term at high energies; therefore, a shift of the horizon toward larger radii appears. Thus, the contradiction has been resolved. Physically this corresponds to the “gravity of pressure” in the general theory of relativity, which surpasses the gradient terms at high energies.
9. The analysis performed allows us to draw the following conclusions.
First, a shift of the horizon toward a larger radius in comparison to the Schwarzchild radius due to the “gravity of pressure” has been discovered. We stress that this effect is purely dynamic and is not observed in the static case (after all the matter has fallen).
Second, according to the results in Appendix 2, the evolution of the entire system at a constant value of $`\alpha `$ is completely specified by the energy density distribution profile in the initial moment, i.e., for example, by the normalized density distribution of the matter and by the value of the parameter $`S`$ at an arbitrary point on this distribution.
If the evolution of only one spherical layer of matter with the index $`R`$ must be described, it is completely specified by three dimensionless parameters in the initial moment in that layer and, in this sense, does not depend on the initial distribution of the matter in the system below and above that layer. However, this in no way signifies the independence of the spherical layers in the general case, since just these three parameters, as will be seen from Appendix 2, govern the interaction of the layers. Consequently, integration of the system leads to a complete family of self-similar solutions.
Third, according to Appendix 2, a local extremum appears on the $`V(R)|_{t=\mathrm{const}}`$ curve for a specific choice of initial parameters, and when $`V=1`$, it leads to the formation of a second apparent horizon in the system (an analog of the second horizon in the Reissner–Nordström and Kerr–Newman solutions for an electrically charged rotating static black hole; for an interpretation of these solutions, see, for example, Refs. and ).
## 1
We have hitherto used the term horizon to refer to a trapping surface, or an apparent horizon, as it is called in the literature.
Let us ascertain the difference between an event horizon and an apparent horizon in greater detail in an example. We assume that we already have a stationary black hole of mass $`M`$ and that there is an apparent horizon at $`r=2M`$. Now we assume that another chunk of matter with a mass $`\delta M`$ falls into our black hole. After it falls, the radius of the apparent horizon increases to $`2(M+\delta M)`$. Thus, if an observer is placed between these radii before the additional chunk of matter falls, he would then be outside the black hole, but after the chunk of matter falls he would be inside it. The concept of an event horizon is global and is determined by the entire course of evolution of the black hole or, stated differently, by all the mass which falls into it at any time.
The existence of an apparent horizon, which specifies a black hole locally, is sufficient for the existence of a black hole. As follows from our arguments, in the spherically symmetric case the two horizons ultimately coincide and form a static black hole described by Schwarzchild’s solution. Therefore, we shall henceforth use the term horizon to refer to the apparent horizon.
Let us prove that the horizon in a system with spherical symmetry forms at the moment when a falling particle with a nonzero rest mass achieves the speed of light relative to the $`r=\mathrm{const}`$ surfaces at the same point. For this purpose we write the law of motion for the particle in the form
$$r(R,T)=R\underset{0}{\overset{t}{}}\sqrt{U(R,\stackrel{~}{t})}𝑑\stackrel{~}{t}.$$
$`(27)`$
We now assume that we are located on a dust particle with the index $`R_{\mathrm{}}`$ and we are tracking a dust particle with the index $`R_p`$, which sends us a light beam passing through the radii $`r_p(R_p,t)`$, from the large radius $`r_{\mathrm{}}(t)`$. The criterion for determining that the dust particle has not yet reached the horizon is the fact that we still see light from it, i.e., the light propagating still crosses the radii $`r>r_p`$. Therefore, the criterion for determining that the dust particle has reached the horizon is an event in which the light propagating from $`R_p`$ can no longer cross the radii $`r>r_p`$. Let us express this criterion mathematically.
In Fig. 1 the vertical straight line $`abcde`$ denotes the world line of an $`R_p`$ dust particle in the coordinates $`R`$ and $`t`$ of the comoving reference frame from the moment of rest ($`a`$) to the center of the system ($`e`$) at $`r=0`$. In this case of solid curves passing through points $`e`$, $`d`$, $`c`$, and $`b`$ denote, respectively, lines of constant values of $`r(R,t)`$ for $`r=0`$, $`r<r_{\mathrm{hor}}`$, $`r=r_{\mathrm{hor}}`$, and $`r>r_{\mathrm{hor}}`$. The dashed lines emerging from these points denote the cones within which light emitted by the $`R_p`$ dust particle can propagate (light cannot propagate outside these cones). Therefore, according to the criterion indicated above, the horizon forms at the point where the cone is tangent to the $`r=\mathrm{const}`$ line. In the figure this line is designated as $`r=r_{\mathrm{hor}}`$, and it passes through point $`c`$. For clarity, Fig. 1 shows that the light cone intersects lines with $`r>r_p`$ at point $`b`$; therefore, there is still no horizon at that point. This figure also shows that at point $`d`$ the light cone is located entirely above the $`r=\mathrm{const}`$ curve passing through point $`d`$. Consequently, this light cone intersects only lines with $`r<r_p`$, and therefore point $`d`$ is already located below the horizon.
Let us examine the expression (27) on one of the $`r=\mathrm{const}`$ curves and take its complete differential on that curve:
$$0=dR\sqrt{U}dt\frac{1}{2}\underset{0}{\overset{t}{}}\frac{U^{}(R,\stackrel{~}{t})}{\sqrt{U(R,\stackrel{~}{t})}}𝑑R𝑑\stackrel{~}{t},$$
or
$$\sqrt{U}\frac{dt}{dR}|_{r=\mathrm{const}}=1\frac{1}{2}\underset{0}{\overset{t}{}}\frac{U^{}(R,\stackrel{~}{t})}{\sqrt{U(R,\stackrel{~}{t})}}𝑑\stackrel{~}{t}.$$
$`(28)`$
Next, differentiating (27) with respect to $`R`$, we obtain the following expression for $`r^{}`$:
$$r^{}=1\frac{1}{2}\underset{0}{\overset{t}{}}\frac{U^{}}{\sqrt{U}}𝑑\stackrel{~}{t},$$
with consideration of which from (28) we find
$$\frac{dt}{dR}|_{r=\mathrm{const}}=\frac{r^{}}{\sqrt{U}}.$$
$`(29)`$
Thus, we have found an expression for the slope of an $`r=\mathrm{const}`$ curve relative to the $`R`$ axis.
To find the slope of a light cone, by definition, for light we have $`ds^2=0`$. Hence, from (2) it follows that
$$\frac{dt}{dR}|_{\mathrm{ligth}}=\sqrt{e^{\lambda \nu }}.$$
$`(30)`$
According to the foregoing statements, the criterion for the absence of a horizon is the condition
$$\frac{dt}{dR}|_{\mathrm{ligth}}<\frac{dt}{dR}|_{r=\mathrm{const}}.$$
$`(31)`$
Substituting the expressions (29) and (30) therein and taking into account (17), we obtain this criterion in the form
$$|V|=\frac{\sqrt{Ue^{\lambda \nu }}}{r^{}}<1.$$
$`(32)`$
Here, according to (29), the rate of motion of the matter relative to the $`r=\mathrm{const}`$ lines has the form
$$|V|=\frac{dl}{d\tau }|_{r=\mathrm{const}}=\sqrt{e^{\lambda \nu }}\frac{dt}{dR}|_{r=\mathrm{const}}.$$
Thus, the assertion that a horizon forms at the moment when the matter achieves the velocity $`V=1`$ relative to the $`r=\mathrm{const}`$ surfaces has been proved. The horizon surface separates regions in which $`r`$ is space-similar and time-similar.
## 2
To solve the equations describing collapse, we first bring them into dimensionless forms. For this purpose it is convenient to introduce the following notation:
$$x=r/R,\gamma =\frac{\rho _0(R)}{\rho }=\frac{8\pi \epsilon _0(R)R^2}{3S(R)}.$$
In this Appendix we find the ranges of permissible values of $`\gamma `$ and $`S`$, investigate the character of collapse at these values of the parameters, and obtain numerical solutions for $`V^2`$. We must first of all know the form of the function $`r^{}(x)`$. Differentiating (27), we obtain\[5)\]<sup>)</sup>
$$r^{}(x)=1+\frac{1}{2}R\underset{1}{\overset{x}{}}[\mathrm{ln}U(R,\stackrel{~}{x})]^{}𝑑\stackrel{~}{x}.$$
$`(33)`$
Unfortunately, an analytically exact expression for $`r^{}`$ can be found only in the case of $`\alpha =0`$, the character of collapse can be assessed exactly only at that value of $`\alpha `$. However, the main features of that character, as will be seen below, remain the same as in the case of $`\alpha 0`$. Therefore, let us first investigate the case of $`\alpha =0`$.
Thus, we should find $`r^{}(R,x)`$. According to the expression (22) for $`U`$, we obtain
$$\mathrm{ln}[U(R,x)]=\mathrm{ln}[S(R)]+\mathrm{ln}\left(\frac{1}{x}1\right).$$
Introducing the notation $`y=r^{}x`$ and taking into account that $`x^{}=y/R`$, we have
$$(\mathrm{ln}U)^{}=\frac{S^{}}{S}\frac{y}{Rx(1x)}.$$
The substitution of this expression into (33) gives
$$y(x)+x1=\frac{1}{2}\underset{1}{\overset{x}{}}\left[\frac{RS^{}}{S}\frac{y(\stackrel{~}{x})}{\stackrel{~}{x}(1\stackrel{~}{x})}\right]𝑑\stackrel{~}{x}.$$
$`(34)`$
Differentiating (34) with respect to $`x`$, we obtain
$$\frac{y(x)}{x}+a^{}(x)y(x)=\sigma ^{}(R),$$
$`(35)`$
where we have introduced the notation
$$a^{}(x)=\frac{1}{2x(ax)},\sigma ^{}(R)=\frac{RS^{}}{2S}1=\frac{3}{2}+\frac{3}{2}\gamma .$$
As can be seen, Eq. (35) coincides in form with Eq. (15), and the initial conditions, $`y|_{t=0}=0`$, are the same; therefore, the method used to solve it is similar. The solution has the form
$$y(x)=r^{}x=\sigma ^{}\left[\sqrt{\frac{1x}{x}}\mathrm{arctan}\sqrt{\frac{1x}{x}}(1x)\right].$$
$`(36)`$
Let us find the domain of $`r^{}`$. First, the condition for compression of the matter has the form $`r^{}1`$. Second, the condition that dust layers with different $`R`$ do not intersect\[6)\]<sup>)</sup> has the form $`r^{}>0`$. Thus,
$$0<r^{}1.$$
$`(37)`$
We assume that the $`V^2(R)`$ curve for $`t=t_m=\mathrm{const}`$ has a local extremum, and we presume (to fix ideas) that it is a maximum. Then the horizon appears specifically at the local maximum, i.e., the point $`R=R_{\mathrm{extr}}`$. We now find the condition for a maximum. First, at that point we should have $`V^2(R_{\mathrm{extr}},x)=V_{\mathrm{extr}}^2`$. Second, since it is the first point at which the velocity of the matter achieves the value $`V_{\mathrm{extr}}`$ and the rate of collapse increases with time, in the vicinity of this point we should have $`V^2<V_{\mathrm{extr}}^2`$, or
$$\frac{V^2(R,t_m)}{R}>0,R<R_{\mathrm{extr}},$$
$`(38)`$
$$\frac{V^2(R,t_m)}{R}<0,R>R_{\mathrm{extr}}.$$
If it turns out that (38) holds with opposite inequality signs, there will be a local minimum on the $`V^2(R)`$ curve at the point $`R_{\mathrm{extr}}`$ at the moment when the velocity $`V_{\mathrm{extr}}`$ is achieved at that point, i.e., the matter will achieve the velocity $`V_{\mathrm{extr}}`$ last at that point.
The condition for an extremum is written in the form
$$\frac{V^2(R,x)}{R}=0,$$
where, according to (23),
$$V^2(R,x)|_{\alpha =0}=\frac{11/x}{1a/S(R)}.$$
Differentiating this expression with respect to $`R`$, we obtain
$$\frac{V^2}{R}|_{t=t_m}=\frac{S/R}{1S}\left[y+\frac{13\gamma }{1S}\right],$$
$`(39)`$
where it has been taken into account that $`x^{}=y/R`$ and $`S^{}/S=(3\gamma 1)/R`$. Then, with allowance for the fact that $`0<x1`$, $`0<\gamma 1`$, $`1<y0`$, and $`0<S<1`$, the condition (38) can be rewritten in the form
$$y>\frac{3\gamma 1}{1S},R<R_{\mathrm{extr}},$$
$`(40)`$
$$y<\frac{3\gamma 1}{1S},R>R_{\mathrm{extr}}.$$
If we introduce the notation $`z=\sqrt{(1x)/x}`$ and take into account that, according to formula (23), $`z=V\sqrt{1/S1}`$, from (40) we obtain
$$\frac{3}{2}(1\gamma )\left[z\mathrm{arctan}(z)\frac{z^2}{1+z^2}\right]\frac{3\gamma 1}{1S}>0,R<R_{\mathrm{extr}},$$
$$\frac{3}{2}(1\gamma )\left[z\mathrm{arctan}(z)\frac{z^2}{1+z^2}\right]\frac{3\gamma 1}{1S}<0,R>R_{\mathrm{extr}},$$
or for the extremum point we can write
$$(1\gamma )\left[z\mathrm{arctan}(z)\frac{z^2}{1+z^2}+\frac{2}{1S}\right]\frac{4/3}{1S}=0.$$
$`(41)`$
This formula can be used to construct the plot of $`\gamma (S,V_{\mathrm{extr}})`$ separating positive and negative values of the derivative $`(V^2)^{}`$ and to determine the character of the extremum. The corresponding curves for various values of $`V_{\mathrm{extr}}`$ are shown in Fig. 2. The regions where $`V^{}>0`$ are located above and to the right of them, and the regions where $`V^{}<0`$ are located below and to the left of them.
It is seen from Fig. 2 that there can be (for a definite choice of the distribution profile of the matter in the initial moment and of the parameter $`S`$ at a certain point $`R_{}`$) two values of $`R`$, at which $`V=1`$ at a certain moment in time, and, therefore, the appearance of a second horizon in the system is possible.
The appearance of a second horizon is not news in the physics of black holes (see, for example, the Reissner–Nordström or Kerr–Newman solution in Ref. ).
The results obtained in this Appendix apply to the case of the absence of pressure, although the case of $`\alpha =1/3`$ is of experimental interest. Therefore, we used formulas (4), (6), (8), (9), (12), (17), (18), (20), and (33) to introduce new dimensionless variables ($`\widehat{\nu }=\nu \nu _0`$, $`\widehat{\lambda }=\lambda \lambda _0`$, $`\widehat{U}=Ue^{\nu _0}`$, and $`\widehat{\epsilon }=\epsilon /\epsilon _0`$) and equations for them. The initial conditions for them take the form
$$\widehat{\nu }_0=\widehat{\lambda }_0=\widehat{U}_0=0,r_0^{}=\widehat{\epsilon }_0=1.$$
Designating the new coordinates as $`x=r/R`$ and $`\xi =R/R_{}`$ ($`R_{}=\mathrm{const}`$) and introducing the parameters\[7)\]<sup>)</sup>
$$h=\xi _\xi \nu _0,h_{\mathrm{cr}}=S\frac{1+3\alpha \gamma }{1S},\eta =h/h_{\mathrm{cr}},$$
we obtain equations for the new variables in the form
$$e^{\widehat{\nu }/\alpha }=x^4e^{\widehat{\lambda }},$$
$$\widehat{\epsilon }=(e^{\widehat{\nu }/\alpha })^{(1+\alpha )/2},$$
$$\mathrm{ln}(r^2e^{\widehat{\lambda }})=\underset{1}{\overset{x}{}}\left[\xi \frac{_\xi \widehat{\nu }}{r^{}}+\frac{h}{r^{}}\right]𝑑\stackrel{~}{x},$$
$`(42)`$
$$\widehat{U}=e^{\widehat{\nu }}\left[r^2e^{\widehat{\lambda }}(1S)1+\frac{S}{x}\alpha \frac{3\gamma S}{x}\underset{1}{\overset{x}{}}\widehat{\epsilon }\stackrel{~}{x}^2𝑑\stackrel{~}{x}\right],$$
$$r^{}=1+\frac{1}{2}\underset{1}{\overset{x}{}}\xi _\xi (\mathrm{ln}\widehat{U})d\stackrel{~}{x}\frac{h}{2}(1x).$$
In the new variables the velocity is
$$V^2=\frac{\widehat{U}e^{\widehat{\lambda }\widehat{\nu }}}{(r^{})^2(1S)}.$$
Hence $`x_{\mathrm{hor}}=S+2\alpha \stackrel{~}{m}`$, where
$$\stackrel{~}{m}=\frac{m}{R}=\epsilon _0R^2\underset{x}{\overset{1}{}}4\pi \widehat{\epsilon }\stackrel{~}{x}^2𝑑\stackrel{~}{x},\epsilon _0R^2=\frac{3\gamma S}{8\pi }.$$
$`(43)`$
We note that this formula and formula (36) can be used to find the corrections $`\delta r_{\mathrm{hor}}`$ in (25) to the displacement of the horizon in the linear approximation with respect to $`\alpha `$, since, according to Eq. (103.11) from Ref. for $`\alpha =0`$, we have
$$8\pi \epsilon r^2=\frac{2M^{}}{r^{}}=\frac{8\pi \epsilon _0R^2}{r^{}},$$
or
$$\widehat{\epsilon }x^2=1/r^{}.$$
This expression can be substituted into (43) and a quadrature expression can be obtained for the correction sought.
In addition, we numerically integrated the equations for the case of $`\alpha 0`$ using a difference scheme, and the results for various values of $`\alpha `$ are presented in Fig. 2. As it should be, according to (25), the plots of $`V^2(x)`$ are displaced upward and to the right as $`\alpha `$ is increased from $`\alpha =0`$ to $`\alpha =1/3`$.
The numerical calculations confirm that the analytical results of this Appendix remain valid for the real equation of state of matter: $`P=\alpha \epsilon `$. Figure 2 shows plots of $`\gamma (S)`$, $`\xi (S)`$, and $`\eta (S)`$ for the special case of a Gaussian density distribution: $`\epsilon _0(\xi )/\epsilon _0(0)=\mathrm{exp}(3\xi ^2)`$. Comparing this figure with Fig. 2, we can see that the $`\gamma (S)`$ curve in Fig. 2 crosses the $`\gamma (S)`$ curves in Fig. 2 in the downward direction roughly at the point $`S0.92`$, if we proceed from $`\xi =0`$ to $`\xi =1`$. As can be seen in Fig. 2, the point $`S0.92`$ corresponds to $`\xi 0.85`$ and $`\eta 0.1`$; therefore, since the region where $`V^{}>0`$ is located above and to the right of the curves in Fig. 2 and the region where $`V^{}<0`$ is located below and to the left of these curves, the point $`\xi 0.85`$ should be a local maximum on the $`V(\xi )`$ curve for a constant value of $`t`$.
This analytical result is confirmed by a numerical calculation of $`V(\xi )|_{t=\mathrm{const}}`$ curves, whose results are shown in Fig. LABEL:5 with the predicted maxima.
To conclude this Appendix we would like to say a few words regarding the initial characteristics and distribution of the matter.
When the equations of the model were brought into dimensionless form, it was found that the solution for a spherical layer of matter with the index $`R`$ is completely specified by three dimensionless parameters in the initial moment in that layer: $`0<S<1`$, $`0<\gamma <1`$, and $`0<\eta <1`$. This corresponds to assigning the initial conditions for the gravitational potential and two parameters which determine the distribution of the matter and the pressure gradient near the point under consideration. Thus, upon integration we at once find a whole family of self-similar solutions,\[8)\]<sup>)</sup> which can be characterized by these three parameters alone and which contains the dependence on the other layers of matter above and below the radius $`R`$ considered.
We thank N. S. Kardashev, V. L. Ginzburg, B. V. Komberg, V. N. Lukash, and Yu. M. Bruk, as well as all the participants in the seminars of the Division of Theoretical Physics and the Astrocosmic Center of the P. N. Lebedev Physics Institute of the Russian Academy of Sciences for fruitful discussions of this work and for their important comments.
With our sincerest gratitude we recall D. A. Kirzhnits, with whom we formulated the ideas and initial approaches used in the present work.
This work was supported by the Russian Foundation for Basic Research (Grant No. 96-15-96616).
Translated by P. Shelnitz
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# Supersymmetry in Singular Quantum Mechanics
## 1 Introduction
Supersymmetry is a beautiful and, simultaneously, a tantalizing symmetry \[1-7\]. On the one hand, supersymmetry leads to field theories and string theories with exceptional properties \[8-9\]. On the other hand, supersymmetry also predicts degenerate superpartner states which are not observed experimentally and, consequently, one expects that supersymmetry must be spontaneously (dynamically) broken. However, unlike ordinary symmetries, spontaneous breaking of supersymmetry has so far proved extremely difficult in the conventional framework. Consequently, in the context of supersymmetry, one constantly looks for alternate, unconventional methods of breaking of this symmetry \[6-7\]. There is, of course, the breaking of supersymmetry due to instanton effects which is well understood. However, several authors, in recent years have suggested that supersymmetry may be broken in the presence of singular potentials or boundaries in a nonstandard manner \[10-12\]. The examples, where such a breaking has been discussed, are simple quantum mechanical models which nonetheless arise from the non-relativistic limit of some field theories. It is for this reason that, in an earlier paper, we had examined a candidate relativistic $`2+1`$ dimensional field theory to see if the manifestation of such a mechanism was possible in a field theory. However, a careful examination of the theory revealed that supersymmetry prevails at the end although it might appear naively, in the beginning, that supersymmetry would be broken in the nonstandard manner. This prompted us to re-analyze the quantum mechanical models, where this mechanism was demonstrated, more carefully and a systematic and critical examination, once again, reveals that supersymmetry is manifest even in such singular quantum mechanical models which is the main result of this talk.
Since our discussion would be entirely within the context of one dimensional supersymmetric quantum mechanics, let us establish the essential notations here. Given a superpotential, $`W(x)`$, we can define a pair of supersymmetric potentials as
$$V_+=\frac{1}{2}\left(W^2(x)+W^{}(x)\right),V_{}=\frac{1}{2}\left(W^2(x)W^{}(x)\right)$$
(1)
where “prime” denotes differentiation with respect to $`x`$. With $`\mathrm{}=1`$ and $`m=1`$, we can, then, define a pair of Hamiltonians which describe a supersymmetric system as
$$H_+=\frac{1}{2}\frac{d^2}{dx^2}+V_+,H_{}=\frac{1}{2}\frac{d^2}{dx^2}+V_{}$$
(2)
In fact, defining the supercharges as
$$Q=\frac{1}{\sqrt{2}}\left(\frac{d}{dx}+W(x)\right),Q^{}=\frac{1}{\sqrt{2}}\left(\frac{d}{dx}+W(x)\right)$$
(3)
we recognize that we can write the pair of Hamiltonians in eq. (2) also as
$$H_+=Q^{}Q,H_{}=QQ^{}$$
(4)
All the eigenstates of the two Hamiltonians $`H_+`$ and $`H_{}`$ would be degenerate except for the ground state with vanishing energy which would correspond to the state satisfying
$$Q|\psi _+=0,\mathrm{or},Q^{}|\psi _{}=0$$
(5)
For a given superpotential, at most one of the two conditions in eq. (5) can be satisfied (that is, at most, only one of the two conditions in (5) would give a normalizable state). Namely, the ground state with vanishing energy is unpaired and can belong to the spectrum of either $`H_+`$ or $`H_{}`$ depending on which of the conditions leads to a normalizable state. This corresponds to the case of unbroken supersymmetry. If, on the other hand, the superpotential is such that neither of the states in eq. (5) is normalizable, then, supersymmetry is known to be broken by instanton effects . In this talk, we carefully analyze the models \[10-12\] where supersymmetry is thought to be broken because of singular nature of the potentials and show that when carefully analyzed, the systems with singular potentials have manifest supersymmetry.
## 2 Super “Half” Oscillator
Let us consider a particle moving in the harmonic oscillator potential on the “half” line
$$V(x)=\{\begin{array}{ccc}\frac{1}{2}(\omega ^2x^2\omega )\hfill & \mathrm{for}\hfill & x>0\hfill \\ \mathrm{}\hfill & \mathrm{for}\hfill & x<0\hfill \end{array}$$
(6)
The spectrum of this potential is quite clear intuitively. Namely, because of the infinite barrier in the negative axis, we expect the wave function to vanish at the origin leading to the conclusion that, of all the solutions of the oscillator on the full line, only the odd solutions (of course, on the “half” line there is no notion of even and odd) would survive in this case. While this is quite obvious, let us analyze the problem systematically for later purpose.
First, let us note that singular potentials are best studied in a regularized manner because this is the only way that appropriate boundary conditions can be determined correctly. Therefore, let us consider the particle moving in the regularized potential
$$V(x)=\{\begin{array}{ccc}\frac{1}{2}(\omega ^2x^2\omega )\hfill & \mathrm{for}\hfill & x>0\hfill \\ & & \\ \frac{c^2}{2}\hfill & \mathrm{for}\hfill & x<0\hfill \end{array}$$
(7)
with the understanding that the limit $`|c|\mathrm{}`$ is to be taken at the end. The Schrödinger equation can now be solved in the two regions. Since $`|c|\mathrm{}`$ at the end, for any finite energy solution, we have the asymptotically damped solution, for $`x<0`$,
$$\psi ^{(II)}(x)=Ae^{(c^22ϵ)^{\frac{1}{2}}x}$$
(8)
Since the system no longer has reflection symmetry, the solutions, in the region $`x>0`$, cannot be classified into even and odd solutions. Rather, the normalizable (physical) solution would correspond to one which vanishes asymptotically. The solutions of the Schrödinger equation, in the region $`x>0`$, are known as the parabolic cylinder functions and the asymptotically damped physical solution is given by
$$\psi ^{(I)}(x)=BU((\frac{ϵ}{\omega }+\frac{1}{2}),\sqrt{2\omega }x)$$
(9)
It is now straightforward to match the solutions in eqs. (8, 9) and their first derivatives across the boundary at $`x=0`$ and their ratio gives
$$\frac{1}{\sqrt{c^22ϵ}}=\frac{1}{2\sqrt{\omega }}\frac{\mathrm{\Gamma }(\frac{ϵ}{2\omega })}{\mathrm{\Gamma }(\frac{ϵ}{2\omega }+\frac{1}{2})}$$
(10)
It is clear, then, that as $`|c|\mathrm{}`$, this can be satisfied only if
$$\frac{ϵ}{2\omega }+\frac{1}{2}\stackrel{|c|\mathrm{}}{}n,n=0,1,2,\mathrm{}$$
(11)
In other words, when the regularization is removed, the energy levels that survive are the odd ones, namely, (remember that the zero point energy is already subtracted out in (6) or (7)) $`ϵ_n=\omega (2n+1)`$. The corresponding physical wave functions are nontrivial only on the half line $`x>0`$ and have the form
$$\psi _n(x)=B_nU((2n+\frac{3}{2}),\sqrt{2\omega }x)=\stackrel{~}{B}_ne^{\frac{1}{2}\omega x^2}H_{2n+1}(\sqrt{\omega }x)$$
(12)
Namely, only the odd Hermite polynomials survive leading to the fact that the wave function vanishes at $`x=0`$. Thus, we see that the correct boundary condition naturally arises from regularizing the singular potential and studying the problem systematically.
We now turn to the analysis of the supersymmetric oscillator on the half line. One can define a superpotential
$$W(x)=\{\begin{array}{ccc}\omega x& \mathrm{for}\hfill & x>0\hfill \\ \mathrm{}& \mathrm{for}\hfill & x<0\hfill \end{array}$$
(13)
which would, naively, lead to the pair of potentials
$$V_\pm (x)=\{\begin{array}{ccc}\frac{1}{2}(\omega ^2x^2\omega )& \mathrm{for}\hfill & x>0\hfill \\ \mathrm{}& \mathrm{for}\hfill & x<0\hfill \end{array}$$
(14)
Since, this involves singular potentials, we can study it, as before, by regularizing the singular potentials as
$`V_+(x)`$ $`=`$ $`\{\begin{array}{ccc}\frac{1}{2}(\omega ^2x^2\omega )\hfill & \mathrm{for}\hfill & x>0\hfill \\ & & \\ \frac{c_+^2}{2}\hfill & \mathrm{for}\hfill & x<0\hfill \end{array}`$ (18)
$`V_{}(x)`$ $`=`$ $`\{\begin{array}{ccc}\frac{1}{2}(\omega ^2x^2+\omega )\hfill & \mathrm{for}\hfill & x>0\hfill \\ & & \\ \frac{c_{}^2}{2}\hfill & \mathrm{for}\hfill & x<0\hfill \end{array}`$ (22)
with the understanding that $`|c_\pm |\mathrm{}`$ at the end.
The earlier analysis can now be repeated for the pair of potentials in eq. (22). It is straightforward and without going into details, let us simply note the results, namely, that, in this case, we obtain
$`ϵ_{+,n}`$ $`=`$ $`\omega (2n+1)\psi _{+,n}(x)=B_{+,n}e^{\frac{1}{2}\omega x^2}H_{2n+1}(\sqrt{\omega }x)`$
$`ϵ_{,n}`$ $`=`$ $`2\omega (n+1)\psi _{,n}(x)=B_{,n}e^{\frac{1}{2}\omega x^2}H_{2n+1}(\sqrt{\omega }x)`$ (23)
Here $`n=0,1,2,\mathrm{}`$. There are several things to note from this analysis. First, only the odd Hermite polynomials survive as physical solutions since the wave function has to vanish at the origin. This boundary condition arises from a systematic study involving a regularized potential. Second, the energy levels for the supersymmetric pair of Hamiltonians are no longer degenerate. Furthermore, the state with $`ϵ=0`$ no longer belongs to the Hilbert space (since it corresponds to an even Hermite polynomial solution). This leads to the conventional conclusion that supersymmetry is broken in such a case and let us note, in particular, that in such a case, it would appear that the superpartner states do not belong to the physical Hilbert space (Namely, in this case, the supercharge is an odd operator and hence connects even and odd Hermite polynomials. However, the boundary condition selects out only odd Hermite polynomials as belonging to the physical Hilbert space.).
There is absolutely no doubt that supersymmetry is broken in this case. The question that needs to be addressed is whether it is a dynamical property of the system or an artifact of the regularization (and, hence the boundary condition) used. The answer is quite obvious, namely, that supersymmetry is broken mainly because the regularization (and, therefore, the boundary condition) breaks supersymmetry. In other words, for any value of the regularizing parameters, $`c_\pm `$ (even if $`|c_+|=|c_{}|`$), the pair of potentials in eq. (22) do not define a supersymmetric system and hence the regularization itself breaks supersymmetry. Consequently, the breaking of supersymmetry that results when the regularization is removed cannot be trusted as a dynamical effect.
### Regularized Superpotential
Another way to understand this is to note that for a supersymmetric system, it is not the potential that is fundamental. Rather, it is the superpotential which gives the pair of supersymmetric potentials through Riccati type relations. It is natural, therefore, to regularize the superpotential which would automatically lead to a pair of regularized potentials which would be supersymmetric for any value of the regularization parameter. Namely, such a regularization will respect supersymmetry and, with such a regularization, it is, then, meaningful to ask if supersymmetry is broken when the regularization parameter is removed at the end. With this in mind, let us look at the regularized superpotential
$$W(x)=\omega x\theta (x)+c\theta (x)$$
(24)
Here $`c`$ is the regularization parameter and we are supposed to take $`|c|\mathrm{}`$ at the end. Note that the existence of a normalizable ground state, namely, the form of the superpotential in eq. (13) selects out $`c>0`$ (otherwise, the regularization would have broken supersymmetry through instanton effects as we have mentioned earlier).
The regularized superpotential now leads to the pair of regularized supersymmetric potentials
$`V_+(x)`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left[(\omega ^2x^2\omega )\theta (x)+c^2\theta (x)c\delta (x)\right]`$
$`V_{}(x)`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left[(\omega ^2x^2+\omega )\theta (x)+c^2\theta (x)+c\delta (x)\right]`$ (25)
which are supersymmetric for any $`c>0`$. Let us note that the difference here from the earlier case where the potentials were directly regularized (see eq. (22)) lies only in the presence of the $`\delta (x)`$ terms in the potentials. Consequently, the earlier solutions in the regions $`x>0`$ and $`x<0`$ continue to hold. However, the matching conditions are now different because of the delta function terms. Carefully matching the wave function and the discontinuity of the first derivative across $`x=0`$ for each of the wavefunctions and taking their ratio, we obtain the two conditions
$`{\displaystyle \frac{1}{(c^22ϵ_+)^{1/2}c}}`$ $`=`$ $`{\displaystyle \frac{1}{2\sqrt{\omega }}}{\displaystyle \frac{\mathrm{\Gamma }(\frac{ϵ_+}{2\omega })}{\mathrm{\Gamma }(\frac{ϵ_+}{2\omega }+\frac{1}{2})}}`$ (26)
$`{\displaystyle \frac{1}{(c^22ϵ_{})^{1/2}+c}}`$ $`=`$ $`{\displaystyle \frac{1}{2\sqrt{\omega }}}{\displaystyle \frac{\mathrm{\Gamma }(\frac{ϵ_{}}{2\omega }+\frac{1}{2})}{\mathrm{\Gamma }(\frac{ϵ_{}}{2\omega }+1)}}`$ (27)
It is now clear that, as $`c\mathrm{}`$, (26) and (27) give respectively, $`ϵ_{+,n}=2\omega n`$ and $`ϵ_{,n}=2\omega (n+1)`$ with $`n=0,1,2,\mathrm{}`$. The corresponding wave functions, in this case, have the forms
$`\psi _{+,n}(x)`$ $`=`$ $`B_{+,n}e^{\frac{1}{2}\omega x^2}H_{2n}(\sqrt{\omega }x)`$
$`\psi _{,n}(x)`$ $`=`$ $`B_{,n}e^{\frac{1}{2}\omega x^2}H_{2n+1}(\sqrt{\omega }x)`$ (28)
This is indeed quite interesting for it shows that the spectrum of $`H_+`$ contains the ground state with vanishing energy. Furthermore, all the other states of $`H_+`$ and $`H_{}`$ are degenerate in energy corresponding to even and odd Hermite polynomials as one would expect from superpartner states. Consequently, it is quite clear that if the supersymmetric “half” oscillator is defined carefully by regularizing the superpotential, then, supersymmetry is manifest in the limit of removing the regularization. This should be contrasted with the general belief that supersymmetry is broken in this system (which is a consequence of using boundary conditions or, equivalently, of regularizing the potentials in a manner which violates supersymmetry). Of course, we should worry at this point as to how regularization independent our conclusion really is. Namely, our results appear to follow from the matching conditions in the presence of singular delta potential terms and, consequently, it is worth investigating whether our conclusions would continue to hold with an alternate regularization of the superpotential which would not introduce such singular terms to the potentials. We have done this which shows that our result is regularization independent.
## 3 Oscillator with $`\frac{1}{x^2}`$ Potential
In the last section, we showed that, in the presence of one kind of singularity, namely, a boundary, supersymmetry is unbroken. In what follows, we will study another class of supersymmetric models, namely, the supersymmetric oscillator with a $`\frac{1}{x^2}`$ potential, where there is a genuine singularity in the potential not necessarily arising from a boundary. A naive analysis of this model also shows that supersymmetry is broken by such a singular potential (for certain parameter ranges). However, this conclusion can be understood, again, as a consequence of regularizing the potential which, as we have seen before, does not respect supersymmetry. In stead, we will show through a careful analysis that, when the superpotential is regularized, supersymmetry is manifest in this model as well (with a lot of interesting features). In this section, however, we will systematically analyze only the quantum mechanical system corresponding to an oscillator in the presence of a $`\frac{1}{x^2}`$ potential (postponing the discussion of the supersymmetric case to the next section). This system has been analyzed by several people \[16-18\] and the most complete analysis appears to be in ref. . However, we feel that, while the energy levels derived in are correct, the wave functions are not (namely, the extensions of the solutions from the positive to the negative axis are incomplete and the wave functions, of course, become quite crucial when one wants to extend the analysis to a supersymmetric system) and, consequently, we present a careful analysis of this system regularizing the singular potential in a systematic manner. With the supersymmetric system in mind (to follow in the next section), we write the potential for the system as (with $`\mathrm{}=m=\omega =1`$)
$$V(x)=\frac{1}{2}\left[\frac{g(g+1)}{x^2}+x^22g+1\right]$$
(29)
The singular potential is repulsive for $`g>0`$ or $`g<1`$ while it is attractive for $`1<g<0`$. It is also worth noting here that the Schrödinger equation, in this case, is invariant under
$`g`$ $``$ $`(g+1)`$
$`ϵ`$ $``$ $`ϵ+2g+1`$ (30)
This symmetry, of course, would also be reflected in the solutions. Furthermore, the fixed point of this symmetry, namely, $`g=\frac{1}{2}`$ separates the two branches (namely, for every value of $`\lambda `$ there exist two distinct values of $`g`$ corresponding to two distinct branches separated at the branch point) in the parameter space.
### Regularized Potential
The Schrödinger equation can be solved quite easily for $`x>0`$ as was also done in . However, to determine correctly how this wavefunction should be extended to the negative axis, it is more suitable to regularize the potential near the origin and study the problem carefully. Let us consider a potential of the form
$$V(x)=\{\begin{array}{ccc}\frac{1}{2}\left[\frac{g(g+1)}{x^2}+x^22g+1\right]\hfill & \mathrm{for}\hfill & |x|>R\hfill \\ & & \\ \frac{1}{2}\left[\frac{g(g+1)}{R^2}+R^22g+1\right]\hfill & \mathrm{for}\hfill & |x|<R\hfill \end{array}$$
(31)
Namely, we have regularized the potential in a continuous manner preserving the symmetry in eq. (30) with the understanding that the regularization parameter $`R0`$ at the end. With this regularization, the Schrödinger equation has to be analyzed in three distinct regions. However, since the potential has reflection symmetry, we need to analyze the solutions only in the regions $`R<x<R`$ and $`x>R`$.
The potential is a constant in the region $`R<x<R`$ and hence the Schrödinger equation is quite simple here. The solutions can be classified into even and odd ones and take the forms
$$\psi ^{(II)even}(x)=A(R)\mathrm{cosh}\kappa x,\psi ^{(II)odd}(x)=B(R)\mathrm{sinh}\kappa x$$
(32)
where we have defined
$$\kappa =\sqrt{\frac{g(g+1)}{R^2}+R^2(2ϵ+2g1)}\frac{\sqrt{g(g+1)}}{R}$$
(33)
Since $`R`$ is small (and we are to take the vanishing limit at the end), the last equality holds only if $`g0\mathrm{or}1`$ which we will assume. The special values of $`g`$ corresponding to the absence of a singular potential have to be treated separately and we will come back to this at the end of this section. We note here that the normalization constants, $`A`$ and $`B`$, can, in principle depend on the regularization parameter which we have allowed for in writing down the form of the solutions in eq. (32).
The potential is much more complicated in the region $`x>R`$. However, the physical solution can be obtained in terms of confluent hypergeometric functions in the form, (for $`x>0`$)
$`\psi ^{(I)}(x)`$ $`=`$ $`C(R)e^{\frac{1}{2}x^2}[{\displaystyle \frac{\mathrm{\Gamma }(g\frac{1}{2})}{\mathrm{\Gamma }(\frac{1}{2}g\frac{ϵ}{2})}}x^{g+1}M(1{\displaystyle \frac{ϵ}{2}},g+{\displaystyle \frac{3}{2}},x^2)`$ (34)
$`+{\displaystyle \frac{\mathrm{\Gamma }(g+\frac{1}{2})}{\mathrm{\Gamma }(1\frac{ϵ}{2})}}x^gM({\displaystyle \frac{1}{2}}g{\displaystyle \frac{ϵ}{2}},g+{\displaystyle \frac{1}{2}},x^2)]`$
Once again, we have allowed for a dependence of the normalization constant, $`C`$, on the regularization parameter, $`R`$. However, for a nontrivial solution to exist, we require that
$$C(R)\stackrel{R0}{}C0$$
So far, we have the general solutions, in the two regions, where energy is not quantized and which should arise from the matching conditions. Furthermore, we have not bothered to evaluate the solution in the region $`x<R`$ which clearly would be the same as in the region $`x>R`$. However, the matching conditions would determine how we should extend the solutions in the region $`x>R`$ to the region $`x<R`$. Therefore, let us now examine the matching conditions systematically since there are two possible cases.
$`(i)`$ Even Solution
We can match the even solution of the region $`R<x<R`$ and its derivative with those of the region $`x>R`$ at $`x=R`$. Taking the ratio and remembering that $`R`$ is small (which is to be taken to zero at the end), we obtain to the leading order in $`R`$
$$\sqrt{g(g+1)}\mathrm{tanh}\sqrt{g(g+1)}=\frac{(g+1)\frac{\mathrm{\Gamma }(g\frac{1}{2})}{\mathrm{\Gamma }(\frac{1}{2}g\frac{ϵ}{2})}R^{g+1}g\frac{\mathrm{\Gamma }(g+\frac{1}{2})}{\mathrm{\Gamma }(1\frac{ϵ}{2})}R^g}{\frac{\mathrm{\Gamma }(g\frac{1}{2})}{\mathrm{\Gamma }(\frac{1}{2}g\frac{ϵ}{2})}R^{g+1}+\frac{\mathrm{\Gamma }(g+\frac{1}{2})}{\mathrm{\Gamma }(1\frac{ϵ}{2})}R^g}$$
(35)
Since the left hand side is independent of $`R`$, for consistency, the right hand side must also be and this can happen in two different ways.
First, for $`g>\frac{1}{2}`$, it is clear that relation (35) can be satisfied if (we assume from now on that $`n=0,1,2,\mathrm{}`$.)
$$ϵ_n=2(n+1)2f_1(g)R^{2g+1}$$
(36)
with a suitable choice of $`f_1(g)`$.
On the other hand, for $`g<\frac{1}{2}`$, if
$$ϵ_n=(2n2g+1)2f_2(g)R^{2g1}$$
(37)
relation (35) can be satisfied with a suitable choice of $`f_2(g)`$. It is clear that the two possible branches of the solution simply reflect the symmetry in eq. (30).
This analysis shows that when the regularization is removed (namely, $`R0`$), we have an even extension of the solution of the forms $`g>\frac{1}{2}`$, $`ϵ_n=2(n+1)`$ with
$$\psi _n(x)=C_n\frac{\mathrm{\Gamma }(g\frac{1}{2})}{\mathrm{\Gamma }(g\frac{1}{2}n)}e^{\frac{1}{2}x^2}M(n,g+\frac{3}{2},x^2)\{\begin{array}{ccc}x^{g+1}& \mathrm{for}\hfill & x>0\hfill \\ |x|^{g+1}& \mathrm{for}\hfill & x<0\hfill \end{array}$$
(38)
and $`g<\frac{1}{2}`$, $`ϵ_n=2n2g+1`$ with
$$\psi _n(x)=C_n\frac{\mathrm{\Gamma }(g+\frac{1}{2})}{\mathrm{\Gamma }(g+\frac{1}{2}n)}e^{\frac{1}{2}x^2}M(n,g+\frac{1}{2},x^2)\{\begin{array}{ccc}x^g& \mathrm{for}\hfill & x>0\hfill \\ |x|^g& \mathrm{for}\hfill & x<0\hfill \end{array}$$
(39)
$`(ii)`$ Odd Solution In a similar manner, we can determine the odd solutions which have the forms $`g>\frac{1}{2}`$, $`ϵ_n=2(n+1)`$ with
$$\psi _n(x)=C_n\frac{\mathrm{\Gamma }(g\frac{1}{2})}{\mathrm{\Gamma }(g\frac{1}{2}n)}e^{\frac{1}{2}x^2}M(n,g+\frac{3}{2},x^2)\{\begin{array}{ccc}x^{g+1}& \mathrm{for}\hfill & x>0\hfill \\ |x|^{g+1}& \mathrm{for}\hfill & x<0\hfill \end{array}$$
(40)
and $`g<\frac{1}{2}`$, $`ϵ_n=2n2g+1`$ with
$$\psi _n(x)=C_n\frac{\mathrm{\Gamma }(g+\frac{1}{2})}{\mathrm{\Gamma }(g+\frac{1}{2}n)}e^{\frac{1}{2}x^2}M(n,g+\frac{1}{2},x^2)\{\begin{array}{ccc}x^g& \mathrm{for}\hfill & x>0\hfill \\ |x|^g& \mathrm{for}\hfill & x<0\hfill \end{array}$$
(41)
### Understanding of the Result
The conclusion following from this analysis, therefore, is that every energy level of this system is doubly degenerate. Both even and odd extensions of the solution are possible for every value of the energy level. The energy levels which we have obtained are, of course, identical to those obtained in . The crucial difference is in the structure of the wave functions, namely, that both even and odd extensions of the solution are possible for every value of the energy (Incidentally, the solutions we have obtained in terms of confluent hypergeometric functions also coincide with generalized Laguerre polynomials as was obtained in ref. .). It is crucial, therefore, to ask if such a conclusion is physically plausible. To understand this question, let us recapitulate the results from a simple quantum mechanical model which is well studied. Namely, let us look at a particle moving in a potential of the form
$$V(x)=\{\begin{array}{ccc}\gamma \delta (x)& \mathrm{for}\hfill & |x|<a\hfill \\ \mathrm{}& \mathrm{for}\hfill & |x|>a\hfill \end{array}$$
It is well known that the solutions of this system can be classified into even and odd ones with energy levels ($`\mathrm{}=m=1`$)
$$E_n^{even}=\frac{n^2\pi ^2}{2(a+\frac{1}{\gamma })^2},E_n^{odd}=\frac{n^2\pi ^2}{2a^2}$$
The even and the odd solutions, of course, have distinct energy values for any finite strength of the delta potential. However, when $`\gamma \mathrm{}`$, both the even and the odd solutions become degenerate in energy. Namely, a delta potential with an infinite strength leads to a double degeneracy of every energy level corresponding to both even and odd solutions. The connection of this example with the problem we are studying is intuitively clear. Namely, we can think of
$$\frac{g(g+1)}{x^2}=\underset{\eta 0}{lim}\frac{g(g+1)}{x^2+\eta ^2}=\underset{\eta 0}{lim}\left(\frac{\pi g(g+1)}{\eta }\right)\left(\frac{1}{\pi }\frac{\eta }{x^2+\eta ^2}\right)$$
It is clear that for $`g0\mathrm{or}1`$, the singular $`\frac{1}{x^2}`$ potential behaves like a delta potential with an infinite strength and it is quite natural, therefore, that this system has both even and odd solutions degenerate in energy.
It is also clear from this analysis that it is meaningless to take the $`g=0\mathrm{or}1`$ limit from the results obtained so far simply because the characters of the two problems are quite different. As we have argued, for any finite value of $`g`$ not coinciding with those special values, the potential behaves, at the origin, like a delta potential of infinite strength while for the special values, there is no such potential. The two cases are related in a drastically discontinuous manner. As a result, one cannot treat the $`\frac{g(g+1)}{x^2}`$ as a perturbation and obtain the full, correct solution simply because there is nothing perturbative (small) about this potential for any “nontrivial” value of $`g`$. Another way of saying this is to re-emphasize what we have already observed following eq. (33), namely, the character of $`\kappa `$ and, therefore, the matching conditions change depending on whether or not $`g`$ differs from the special values $`0,1`$.
## 4 Supersymmetric Oscillator with $`\frac{1}{x^2}`$ Potential:
The supersymmetric version of the case studied is obtained from a superpotential of the form
$$W(x)=\frac{g}{x}x$$
(42)
In this case, it is easily seen that a normalizable ground state wavefunction exists only for $`g>\frac{1}{2}`$.
The superpotential and, therefore, the potential is singular at the origin. Thus, once again, the proper way to study the spectrum of such a system is by regularizing the superpotential. We introduce the regularized superpotential
$$W(x)=\theta (x|R|)\left(\frac{g}{x}x\right)+\theta (|R|x)\left(\frac{g}{R}R\right)\frac{x}{R}$$
(43)
Here, $`R`$ denotes the regularization parameter which is to be taken to zero at the end. The regularized superpotential is continuous and the resulting pair of potentials take the forms
$`V_+`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left[\theta (x|R|)\left({\displaystyle \frac{g(g1)}{x^2}}+x^22g1\right)+\theta (|R|x)\left(\left({\displaystyle \frac{g}{R^2}}1\right)^2x^2+\left({\displaystyle \frac{g}{R^2}}1\right)\right)\right]`$
$`V_{}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left[\theta (x|R|)\left({\displaystyle \frac{g(g+1)}{x^2}}+x^22g+1\right)+\theta (|R|x)\left(\left({\displaystyle \frac{g}{R^2}}1\right)^2x^2\left({\displaystyle \frac{g}{R^2}}1\right)\right)\right]`$
The solutions of the pair of Hamiltonians can now be studied. Without going into details , let us note that the solutions, in this case, again turn out to be confluent hypergeometric functions. There are several interesting features that arise in this case. For example, it turns out that, in the limit $`R0`$, $`H_+`$ has three sets of normalizable solutions – one even and two odd. The three sets of normalizable solutions of $`H_{}`$ also correspond to one even and two sets of odd solutions. While one of the three sets of solutions correspond to a supersymmetric system, there are additional solutions which apparently have no relation to one another.
The proper understanding of the solutions comes really from recognizing that, given a bosonic system, there is an arbitrariness in supersymmetrizing the system. It is much like the arbitrariness of whether a spin $`\frac{1}{2}`$ particle belongs to a supersymmetric multiplet $`(0,\frac{1}{2})`$ or $`(\frac{1}{2},1)`$. The different solutions really correspond to different possible supersymmetrizations and matching has to be done carefully. When analyzed carefully, it turns out that supersymmetry is manifest in the system. In addition, in this case, the problem can be solved algebraically because of a special symmetry in the problem known as shape invariance. The algebraic solution also coincides with the explicit solutions obtained.
## 5 Conclusion
In this talk, we have discussed systematically two classes of supersymmetric quantum mechanical models - one consisting of a singular boundary and the other with a singular potential. We have shown that, contrary to the conventional understanding \[10-12\], supersymmetry is manifest in these systems. In particular, for a system with a singular potential such as $`\frac{1}{x^2}`$, the solution of the Schrödinger equation leads to several distinct solutions corresponding to distinct supersymmetrizations of the system. Consequently, it becomes quite important to identify the appropriate wavefunctions when supersymmetric properties are being investigated. Finally, we would like to conclude by noting that supersymmetry is known to be robust at short distances (high energies). The singularities discussed in the quantum mechanical models occur at short distances and, therefore, it is intuitively quite clear that they are unlikely to break supersymmetry. Our detailed, systematic analysis only reinforces this.
## Acknowledgments
This work was supported in part by the U.S. Dept. of Energy Grant DE-FG 02-91ER40685 and NSF-INT-9602559.
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# 1 Introduction
## 1 Introduction
It is often claimed that the duality operation is only defined in even dimensional spacetimes and that selfduality is further restricted to twice-odd spacetime dimensional theories. The purpose of this paper is to extend the notion of both duality symmetry as well as selfduality to the odd (2+1) dimensional spacetime in the electromagnetic context. Naturally, the conventional duality symmetry in even dimensions is also contained in our approach. To achieve this we introduce an alternative definition for the duality operation which is valid in any dimensions. Specifically we analyse and explore the impact of the Gauss law kernel’s parity and the dual projection procedure over the duality operation of the Maxwell theory in different spacetime dimensions. The crucial issue concerns the spacetime dependence of the parity property of a generalized curl involved in the resolution of the Gauss law. We show that this property is decisive in determining the proper actions and the corresponding group of symmetry for each specific dimension.
The role of the duality operation in the investigation of concrete physical systems in different areas is by now well recognized. This is a symmetry transformation that is fundamental for investigations in arenas as distinct as quantum field theory, statistical mechanics and string theory. Duality is a general concept relating physical quantities in different regions of the parameter space. It relates a model in a strong coupling regime to its dual version working in a weak coupling regime, providing valuable information in the study of strongly interacting models. The selfduality present in D=4k-2 dimensions has attracted much attention because it seems to play an important role in many theoretical models.
The electromagnetic duality transformation is defined by the Hodge-star operation that involves multiplication by the appropriate $`ϵ`$ symbol. Consider a general (n-1)-form and its field strength
$$F_{k_1\mathrm{}k_n}_{[k_n}A_{k_1\mathrm{}k_{n1}]}.$$
(1)
The dual field is then defined as,
$$^{}F^{k_1\mathrm{}k_n}=\frac{1}{n!}ϵ^{k_1\mathrm{}k_{2n}}F_{k_{n+1}\mathrm{}k_{2n}}$$
(2)
Note that only in $`2n`$ dimensions will the $`n`$-form field be of the same rank as its dual. The action, field equation and Bianchi identity for a source free field are
$`S`$ $`=`$ $`c_n{\displaystyle d^{2n}xF_{k_1\mathrm{}k_n}F^{k_1\mathrm{}k_n}}`$
$`0`$ $`=`$ $`_{k_1}F^{k_1\mathrm{}k_n}`$
$`0`$ $`=`$ $`_{k_1}^{}F^{k_1\mathrm{}k_n}`$ (3)
where $`c_1=1/2`$, $`c_2=1/4`$, etc. The field equation and the Bianchi identity are of the same form so that the duality transformation $`F^{}F`$ is a symmetry at the level of these equations but is not present at the level of the action. The dependence with dimensionality appears to be crucial. The duality symmetry is characterized by a one-parameter continous SO(2) group for D=4k ($`kZ_+`$) while for the D=4k-2 case it is described by a discrete group with just two elements. Most of the analysis of this dimensional dependence take the algebraic point of view where the distinction among different dimensions is manifest by the double dualization operation following from the identities,
$$^{}F=\{\begin{array}{cc}F,\hfill & \text{if }D=4k2\hfill \\ F,\hfill & \text{if }D=4k\text{.}\hfill \end{array}$$
(4)
It was then shown that the duality groups $`G`$ preserving the form of the action were subgroups of those preserving equations of motion and Bianchi identities, obtained by taking the intersection of the former with the group O(2), the symmetry group for the energy-momentum tensor.
Most of the discussion about duality transformations as a symmetry for the actions and the existence of self-duality are based on these concepts. We feel that this algebraic viewpoint is rather restrictive. It is only defined for even dimensional spacetimes leading to separate consequences regarding the duality groups and actions. The usual lore is that only the 4D Maxwell theory and its 4k extensions possess duality as a symmetry but selfduality would only be definable in D=4k-2. On the other hand for the 2D scalar field theory and its 4k-2 extensions duality is not even a well defined concept. A solution for these problems came with the recognition of an internal structure in the space of potentials. The internal space effectively unifies the selfduality concept in the different $`4k2`$ and $`4k`$ dimensions. The dual field is now defined to include an internal index $`(\alpha ,\beta )`$ and an extended dualization is defined as,
$`\stackrel{~}{F}^\alpha `$ $`=`$ $`ϵ^{\alpha \beta }^{}F^\beta ;D=4k`$
$`\stackrel{~}{F}^\alpha `$ $`=`$ $`\sigma _1^{\alpha \beta }^{}F^\beta ;D=4k2`$ (5)
where $`\sigma _k`$ are the usual Pauli matrices and $`ϵ_{\alpha \beta }`$ is the fully antisymmetric $`2\times 2`$ matrix with $`ϵ_{12}=1`$. Now, irrespective of the dimensionality, the double dual operation yields,
$$\stackrel{~}{\stackrel{~}{F}}=F$$
(6)
which generalises (4) . Self and anti-self dual solutions are now well defined in all even $`D=2k`$ dimensions.
With the above background, it is useful to put our work in a proper perspective by observing the following:
* The procedure of dual projection developed here is basically analogous to a canonical transformation, except that the former is performed at the level of the actions while the latter, as is well known, is at the hamiltonian level. Since throughout the paper only first order actions will be considered the equivalence between the dual projection and the canonical transformations becomes manifest.
* The Maxwell action in any even dimension is decomposed, by the dual projection method, into two pieces, one of which carries the $`SO(2)`$ symmetry while the other has the $`Z_2`$ symmetry. By specialising to $`4k`$ or $`4k2`$ dimensions, we find that one of the terms in the action becomes a total derivative which can be ignored. In this way the conventional results of duality symmetry characterising the $`SO(2)`$ group for $`D=4k`$ and the $`Z_2`$ group for $`D=4k2`$ dimensions are reproduced.
* The dual projection in D=4k-2 dimensions leads to a diagonal form of the actions with two pieces manifesting the opposite chiralities. It is then possible to impose a chiral constraint to eliminate either of the pieces. What remains is the action for a chiral boson. This generalises the usual construction of chiral bosons to higher dimensions.
* The case D=2 seems to be a special point. Although it qualifies as a member of the D=4k-2 group, it will be shown that it allows for the realization of both D=4k-2 and D=4k constructions.
* By passing to the momentum space, it is possible to derive results analogous to the D=2 case for higher dimensions. In other words there is a duality transformation among the Fourier modes showing the complementary nature of the symmetries: the $`SO(2)`$ symmetry for $`D=4k2`$ dimensions and the $`Z_2`$ symmetry for the $`D=4k`$ dimensions. It might be mentioned that this nature of duality symmetry was earlier shown by us using different methods.
* The algebraic analysis leaves out the possibility of obtaining duality symmetric electromagnetic Maxwell theory in odd dimensions. For the special case of three dimensions this will be achieved here.
In the next section we shall show that the use of the two key concepts; namely, the dimensional dependence of the parity property of the generalized rotation operator involved in the resolution of the Gauss law and the dual projection method, reproduce the known results of the algebraic analysis, clarifying their physical origins. In the third section we discuss some special instances like the two-dimensional case that possess both $`SO(2)`$ and $`Z_2`$ representations. In the fourth section, which contains the centeral result of this paper, attention is given to the D=(2+1) Maxwell theory that is studied in full details. We disclose the presence of an internal space of potentials where duality is realized as a $`SO(2)`$ rotation and also as a discrete $`Z_2`$ symmetry. The special equivalence between Maxwell theory and the scalar field via Hodge-dualization is discussed from our approach. The last section is reserved to a discussion of our conclusions and perpectives.
## 2 Parity and Dual Projection in Even Dimensional Spacetimes
The main argument of this report is the dimensional dependence of the parity property of the generalized rotation operator and a canonical transformation that we call dual projection. A systematic derivation of selfdual actions for two and four dimensional cases was proposed in using the dual projection procedure. Here we generalise the method to inclued all even dimensions. Subsequently these ideas will be exploited to discuss the consequences in a three dimensional theory. The generalized rotation operator, which is basically a functional curl, is defined as,
$$(ϵ)ϵ_{k_1k_2\mathrm{}k_{D1}}_{k_{D1}}$$
(7)
Clearly the dimensional dependence of this operator’s parity is given by,
$$𝒫(ϵ)=\{\begin{array}{cc}+1,\hfill & \text{if }D=4k\hfill \\ 1,\hfill & \text{if }D=4k2\text{,}\hfill \end{array}$$
(8)
where parity is defined as
$$\mathrm{\Phi }(ϵ\mathrm{\Psi })=𝒫(ϵ)\mathrm{\Psi }(ϵ\mathrm{\Phi })$$
(9)
The consequence of this property is best appeciated after a dual projection of the action. First, the theory is reduced to its first-order form as,
$$S=d^Dx\left[\pi \dot{A}\frac{1}{2}\pi \pi \frac{1}{2}BB+A_0(\pi )\right]$$
(10)
where we used the notation (anti symmetrisation is implied by the brackets)
$$\mathrm{\Phi }\mathrm{\Psi }\mathrm{\Phi }_{[k_1k_2\mathrm{}k_{D1}]}\mathrm{\Psi }_{[k_1k_2\mathrm{}k_{D1}]}$$
(11)
and defined the magnetic field as
$$B=(ϵ)A$$
(12)
In the four dimensional case, $`B_k=ϵ_{kmn}_mA_k`$ is a three-vector while in three dimensions, the magnetic field is a scalar, $`B=ϵ_{km}_mA_k`$. Clearly the two-dimensional instance represents a special situation due to the absence of a Gauss law and will be treated separately in the next section. There does not seem to exist any difficulty in dimensions $`D>4`$. Note that $`A_0`$ in (10) generically denotes the multiplier in any even dimension, that enforces the Gauss constraint. For example, it is just $`A_0`$ in four dimensions while it is $`A_{0i}`$ in six dimensions and so on.
The important point to observe is that the Gauss constraint is trivially solved, in any dimension, using the generalized curl (7),
$$\pi =(ϵ)\varphi $$
(13)
where $`\varphi `$ is a ($`\frac{D}{2}1`$)-form potential. For instance, in D=4 and D=6 which are generic for D=4k and D=4k-2, this solution reads
$`\pi _k`$ $`=`$ $`ϵ_{kmn}_m\varphi _n`$
$`\pi _{km}`$ $`=`$ $`ϵ_{kmnpq}_n\varphi _{pq}`$ (14)
The next step is to perform the canonical transformations,
$`A`$ $`=`$ $`\mathrm{\Phi }^{(+)}+\mathrm{\Phi }^{()}`$
$`\pi `$ $`=`$ $`\eta (ϵ)(\mathrm{\Phi }^{(+)}\mathrm{\Phi }^{()})`$ (15)
with $`\eta =\pm 1`$ defining the signature of the operation. The effect of the dual projection procedure into the first-order Maxwell action is the creation of an internal space of potentials in which the duality symmetry is local and manifest. In terms of the internal space potentials $`\mathrm{\Phi }^{(+)}`$ and $`\mathrm{\Phi }^{()}`$ the action now reads,
$$S=d^Dx\left\{\eta \left[\dot{\mathrm{\Phi }}^{(\alpha )}\sigma _3^{\alpha \beta }B^{(\beta )}+\dot{\mathrm{\Phi }}^{(\alpha )}ϵ^{\alpha \beta }B^{(\beta )}\right]B^{(\beta )}B^{(\beta )}\right\}$$
(16)
where $`B^{(\beta )}=(ϵ\mathrm{\Phi }^{(\beta )})`$ and $`\sigma _3^{(\alpha \beta )}`$ and $`\sigma _2^{(\alpha \beta )}=iϵ^{(\alpha \beta )}`$ are the $`2\times 2`$ Pauli matrices. Notice that while the hamiltonian sector of the first-order action is unique, the symplectic sector is composed by two distinct parts with separate consequences. We can now appreciate the impact of the dimensionality over the symplectic structure of (16) and the role of the parity in selecting the proper action and the corresponding duality group. Parity (or dimensionality) has no influence over the hamiltonian since it only involves quadratic forms. For twice odd dimensions the second term of the symplectic sector is a total derivative and may be discarded. The remaining piece diagonalizes the action providing a generalization of the two-dimensional action describing chiral bosons. The $`Z_2`$ property is manifest by the interchange between the internal space potentials $`\mathrm{\Phi }^{(\pm )}\mathrm{\Phi }^{()}`$ mapping one chirality into the other.
For twice even dimensions, on the other hand, it is the first term that becomes a total derivative. The action does not diagonalize but presents an explicit one-parameter continous SO(2) symmetry. In the D=4 case this action corresponds to duality symmetric Maxwell theory quoted in the literature . The important point to stress is that the derivative operator involved in the dual projection has been determined by the solution of the Gauss constraint. This automactically fixes the dependence of parity with dimensionality and explains its effect over the electromagnetic actions and duality groups. Incidentally, observe that due to its intrinsic diagonal form, the phase space of the chiral boson solution (D=4k-2) may be reduced if we impose a chiral constraint as,
$$\pi =\pm (ϵ)A$$
(17)
Each of these constraints eliminates one of the (internal) chiral potentials thereby leading to an action for chiral bosons. These are the generalisations of the usual actions for chiral bosons in two dimensions. However, the same situation cannot be reached in the twice-even instance due to the special form of the sympletic sector (the hamiltonian poses no obstruction to reduction in either case). In the reduced phase-space the remaining chiral boson carries a representation for half the number of degrees of freedom of the original system. On the other hand, the duality symmetric system mantains the phase space structure intact. Therefore, this system should not be considered as the 4D analog of the 2D chiral boson. Although, due to the possibility of two distinct signatures in the dual projection, there exist either a self-dual or anti self-dual decomposition (but not simultaneously), we believe that this situation should not be confused with the distinct chiralities that appear (simultaneously) in the dual projection of D=4k-2 dimensional systems.
We have therefore reproduced completely the results known from the algebraic approach plus presented a derivation of the appropriate actions displaying the internal potentials for each case. Also worth of mention is the fact that there are two and not one self-dual action, labelled by the signature $`\eta `$ of the dual projection, describing opposite aspect of the self-duality symmetry. The most useful and striking feature in this dual projection procedure is that it is not based on evidently even dimensional concepts and may be extended to the odd dimensional situation.
## 3 Special Examples
In this section we discuss some special dimensions and situations. The 4D Maxwell example deserves detailed analysis since it is the paradigm of the duality symmetry. The other example is the (1+1) dimensional scalar theory. The absence of a gauss constraint leads to a crucial change from the Maxwell case.
### 3.1 The Electromagnetic Duality
Exploiting the ideas elaborated in the previous sections, it is straightforward to implement the selfduality projection in the electromagnetic theory. Let us start with the usual Maxwell action,
$$S=\frac{1}{4}d^4xF_{\mu \nu }F^{\mu \nu }$$
(18)
which is expressed in terms of the electric and magnetic fields as,
$$S=\frac{1}{2}d^4x\left(E_k^2B_k^2\right)$$
(19)
where,
$`E_i`$ $`=`$ $`F_{0i}=_0A_i+_iA_0`$
$`B_i`$ $`=`$ $`ϵ_{ijk}_jA_k`$ (20)
The following duality transformation,
$$E_kB_k;B_k\pm E_k$$
(21)
is known to preserve the invariance of the full set comprising Maxwell’s equations and the Bianchi identities although the action changes its signature. The Maxwell Lagrangean is next recast in a symmetrised first order form that displays an Sp(2,R) symmetry when we treat ($`P_k,A_k`$) as a doublet,
$$=\frac{1}{2}\left(P_k\dot{A_k}\dot{P}_kA_k\right)\frac{1}{2}P_k^2\frac{1}{2}B_k^2+A_0_kP_k$$
(22)
Next a canonical transformation is invoked. There are two possibilities (assigning different signatures for the dual projection) which translate from the old set $`(P_k,A_k)`$ to the new ones $`(A_k^1,A_k^2)`$. It is, however, important to recall that the Maxwell theory has a Gauss constraint that is implemented by the Lagrange multiplier $`A_0`$. The new variables are chosen in two different ways which solve this constraint and implement distinct signatures to the dual projection as,
$`P_k`$ $``$ $`B_k^2;A_kA_k^1`$
$`P_k`$ $``$ $`B_k^1;A_kA_k^2`$ (23)
It is now simple to show that, in terms of the new variables, the original Maxwell action takes the form,
$$S_\pm =\frac{1}{2}d^4x\left(\pm \dot{A}_k^\alpha ϵ^{\alpha \beta }B_k^\beta B_k^\alpha B_k^\alpha \right)$$
(24)
It is duality symmetric under the full $`SO(2)`$.
Let us next introduce the proper and improper $`O(2)`$ rotation matrices as $`R^+(\theta )`$ and $`R^{}(\phi )`$ with determinant $`+1`$ and $`1`$, respectively,
$`R^+\left(\theta \right)=\left(\begin{array}{cc}\mathrm{cos}\theta & \mathrm{sin}\theta \\ \mathrm{sin}\theta & \mathrm{cos}\theta \end{array}\right)`$ (27)
$`R^{}\left(\phi \right)=\left(\begin{array}{cc}\mathrm{sin}\phi & \mathrm{cos}\phi \\ \mathrm{cos}\phi & \mathrm{sin}\phi \end{array}\right)`$ (30)
Note that the matrix that corresponds to improper rotations, $`R^{}(\phi )`$ switches the actions $`S_+`$ and $`S_{}`$ into one another.
Using the basic brackets following from the canonical transformations or from the symplectic structure of the theory,
$$[A_\alpha ^i(x),ϵ^{jkl}^kA_\beta ^l(y)]=\pm i\delta ^{ij}ϵ_{\alpha \beta }\delta (𝐱𝐲)$$
(31)
we can verify that the generators of the $`SO(2)`$ rotations are given by a Chern-Simon like structure,
$$Q^{(\pm )}=\frac{1}{2}d^3xA_k^\alpha B_k^\alpha $$
(32)
so that finite transformations are given by,
$$A_k^\alpha A_{}^{}{}_{k}{}^{\alpha }=e^{iQ\theta }A_k^\alpha e^{iQ\theta }$$
(33)
Let us stress on the fact that there are two distinct structures for the duality symmetric actions. These must correspond to the opposite aspects of some symmetry. By looking at the equations of motion obtained from (24),
$$\dot{A}_k^\alpha =\pm ϵ^{\alpha \beta }_k\times A_k^\beta $$
(34)
it is possible to verify that these are just the self and anti-self dual solutions,
$$F_{\mu \nu }^\alpha =\pm ϵ^{\alpha \beta }^{}F_{\mu \nu }^\beta ;^{}F_{\mu \nu }^\beta =\frac{1}{2}ϵ_{\mu \nu \rho \lambda }F_\beta ^{\rho \lambda }$$
(35)
obtained by setting $`A_0^\alpha =0`$. It may be observed that the opposite aspects of the dual symmetry are contained in the internal space.
To close our arguments, let us now comment on another property, which is related to the existence of two distinct actions (24), by replacing (21) with a new set of transformations,
$`E_\alpha `$ $``$ $`R_{\alpha \beta }^{}(\phi )E_\beta `$
$`B_\alpha `$ $``$ $`R_{\alpha \beta }^{}(\phi )B_\beta `$ (36)
Notice that these transformations preserve the invariance of the hamiltonian following from either $`S_+`$ or $`S_{}`$. The kinetic terms in the action change signatures so that $`S_+`$ swaps to $`S_{}`$. The discretised version of (36) is obtained by setting $`\phi =0`$,
$`E_\alpha `$ $``$ $`\sigma _1^{\alpha \beta }E_\beta `$
$`B_\alpha `$ $``$ $`\sigma _1^{\alpha \beta }B_\beta `$ (37)
It is precisely the $`\sigma _1`$ matrix that reflects the proper into improper rotations,
$$R^+(\theta )\sigma _1=R^{}(\theta )$$
(38)
which illuminates the reason behind the swapping of the actions in this example.
### 3.2 The Scalar Theory in 1+1 Dimensions
The ideas developed in the previous section are now implemented and elaborated in $`1+1`$ dimensions. In particular we show that two distinct dual projections are possible in this case, leading to either $`Z_2`$ or SO(2) group of dualities. Notice first that in D=2 there is no photon and the Maxwell theory trivialises so that the electromagnetic field can be identified with a scalar field. Thus all the results presented here can be regarded as equally valid for the “photon” field. There is no Gauss constraint however so that we are free to choose any operator in the dual projection. Our computations will be presented in a very suggestive notation which illuminates the Maxwellian nature of the problem. The action for the free massless scalar field is given by,
$$S=\frac{1}{2}d^2x\left(_\mu \varphi \right)^2$$
(39)
and the equation of motion reads,
$$\ddot{\varphi }\varphi ^{\prime \prime }=0$$
(40)
where the dot and the prime denote derivatives with respect to time and space components, respectively. Introduce a change of variables using electromagnetic symbols,
$$E=\dot{\varphi };B=\varphi ^{}$$
(41)
Obviously, $`E`$ and $`B`$ are not independent but constrained by the identity,
$$E^{}\dot{B}=0$$
(42)
In these variables the equation of motion and the action are expressed as,
$`\dot{E}B^{}=0`$
$`S={\displaystyle \frac{1}{2}}{\displaystyle d^2x\left(E^2B^2\right)}`$ (43)
so that the transformations,
$$E\pm B;B\pm E$$
(44)
display a duality between the equation of motion and the ‘Bianchi’-like identity (42) but the action changes its signature. Note that there is a relative change in the signatures of the duality transformations (44) with respect to the true electromagnetic duality (21), arising basically from dimensional considerations. This symmetry coresponds to the improper group of rotations.
To illuminate the close connection with the Maxwell formulation, we introduce covariant and contravariant vectors with a Minkowskian metric $`g_{00}=g_{11}=1`$,
$$F_\mu =_\mu \varphi ;F^\mu =^\mu \varphi $$
(45)
whose components are just the ‘electric’ and ‘magnetic’ fields defined earlier,
$$F_\mu =(E,B);F^\mu =(E,B)$$
(46)
Likewise, with the convention $`ϵ_{01}=1`$, the dual field is defined,
$`^{}F_\mu `$ $`=`$ $`ϵ_{\mu \nu }^\nu \varphi =ϵ_{\mu \nu }F^\nu `$ (47)
$`=`$ $`(B,E)`$
The equation of motion and the ‘Bianchi’ identity are now expressed by typical electrodynamical relations,
$`_\mu F^\mu `$ $`=`$ $`0`$
$`_\mu ^{}F^\mu `$ $`=`$ $`0`$ (48)
To expose a duality symmetric action, the basic principle of our approach is adopted. We convert the original second order form (3.2) to its first order version displaying the Sp(2) symmetry and then invoke a canonical transformation to provide an internal index. An auxiliary field is therefore introduced at the first step,
$$=PE\frac{1}{2}P^2\frac{1}{2}B^2$$
(49)
where $`E`$ and $`B`$ have already been defined. The following canonical transformation,
$`B`$ $``$ $`\left(\mathrm{\Phi }^{(+)}+\mathrm{\Phi }^{()}\right)`$
$`P`$ $``$ $`\left(\mathrm{\Phi }^{(+)}\mathrm{\Phi }^{()}\right)`$ (50)
leads to an action with fields taking values in the internal space
$$S=d^2x\left[\left(\mathrm{\Phi }^{(+)}\dot{\mathrm{\Phi }}^{(+)}\mathrm{\Phi }^{()}\dot{\mathrm{\Phi }}^{()}\right)+\left(\mathrm{\Phi }^{(+)}\dot{\mathrm{\Phi }}^{()}\mathrm{\Phi }^{()}\dot{\mathrm{\Phi }}^{(+)}\right)\left(\mathrm{\Phi }^{(\alpha )}\mathrm{\Phi }^{(\alpha )}\right)\right]$$
(51)
As discussed previously, due to the absence of a true Gauss law in this case, we are free of any imposition regarding the choice of the operator $``$ in the dual projection. To display this arbitrariness, we choose, for each group of symmetry transformation,
$$=\{\begin{array}{cc}_x,\hfill & \text{leading to }Z_2\hfill \\ \sqrt{_x^2},\hfill & \text{leading to }SO(2)\text{.}\hfill \end{array}$$
(52)
The first choice is traditional. The odd parity of the operator diagonalizes the action by eliminating the second term in the sympletic sector. The resulting actions,
$`S`$ $`=`$ $`S_++S_{}`$
$`S_\pm `$ $`=`$ $`{\displaystyle d^2x\left(\dot{\mathrm{\Phi }}^\pm _x\mathrm{\Phi }^\pm _x\mathrm{\Phi }^\pm _x\mathrm{\Phi }^\pm \right)}`$ (53)
correspond to the well known right and left chiral boson theories.
To examine the symmetry content it is possible to recast (3.2) in a very suggestive form,
$`S_\pm `$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle d^2x\left[\pm _x\mathrm{\Psi }^\alpha \sigma _1^{\alpha \beta }\dot{\mathrm{\Psi }}^\beta _x\mathrm{\Psi }^\alpha _x\mathrm{\Psi }^\alpha \right]}`$ (54)
$`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle d^2x\left[\pm B_\alpha \sigma _1^{\alpha \beta }E_\beta B_\alpha ^2\right]}`$
where $`\mathrm{\Psi }^\pm =\mathrm{\Phi }^+\pm \mathrm{\Phi }^{}`$. In the second line the action is expressed in terms of the electromagnetic variables. This action is duality symmetric under the transformations of the basic scalar fields,
$$\mathrm{\Psi }_\alpha \sigma _1^{\alpha \beta }\mathrm{\Psi }_\beta $$
(55)
which, in the notation of $`E`$ and $`B`$, is given by,
$`B_\alpha `$ $``$ $`\sigma _1^{\alpha \beta }B_\beta `$
$`E_\alpha `$ $``$ $`\sigma _1^{\alpha \beta }E_\beta `$ (56)
It is quite interesting to observe that, contrary to the 4D electromagnetic theory, the transformation matrix in the $`O(2)`$ internal space of potentials is not the epsilon, but rather a $`\sigma _1`$ Pauli matrix. This result is in agreement with that found from general algebraic arguments which stated that for $`D=4k2`$ dimensions there is a discrete $`\sigma _1`$ symmetry. Observe that (3.2) is a manifestation of the original duality (44) which was also effected by the same operation. It is important to stress that the above transformation is only implementable at the discrete level. Moreover, since it is not connected to the identity, there is no generator for it. In this sense it is observed that duality symmetry is not defined in these twice odd dimensions. To complete the picture, we also mention that the following rotation,
$$\mathrm{\Psi }_\alpha ϵ_{\alpha \beta }\mathrm{\Psi }_\beta $$
(57)
interchanges the actions (54),
$$S_+S_{}$$
(58)
Thus, except for a rearrangement of the the matrices generating the various transformations, most features of the electromagnetic example are perfectly retained. The crucial point of departure is that now all these transformations are only discrete.
The second choice in (52) is new and unexpected since it does not fit into the known dimensional classification. It leads to a continous SO(2) duality transformation, characteristic of the 4k dimensional spacetimes, instead of the discrete $`Z_2`$. The resulting action is,
$$S=d^2x\left[\left(\mathrm{\Phi }^{(+)}\dot{\mathrm{\Phi }}^{()}\mathrm{\Phi }^{()}\dot{\mathrm{\Phi }}^{(+)}\right)\left(\mathrm{\Phi }^{(\alpha )}\mathrm{\Phi }^{(\alpha )}\right)\right]$$
(59)
which is appropriate for the even-parity dual projection. It is easy to obtain the basic symplectic brackets from here as,
$$\{\mathrm{\Phi }^\alpha (x),\mathrm{\Phi }^\beta (y)\}=ϵ^{\alpha \beta }\delta (xy)$$
(60)
Now observe that the action (59) is manifestly invariant under the continuous duality transformations,
$$\mathrm{\Phi }^\alpha R_{\alpha \beta }^+\mathrm{\Phi }^\beta $$
(61)
where $`R_{\alpha \beta }^+`$ is the usual $`SO(2)`$ rotation matrix (27). The generator of the infinitesimal symmetry transformation is given by,
$$Q^\pm =\frac{1}{2}d^2x\mathrm{\Phi }^\alpha \mathrm{\Phi }^\alpha $$
(62)
while the transformation by a finite angle $`\theta `$ is generated by,
$`\mathrm{\Phi }^\alpha \stackrel{~}{\mathrm{\Phi }}^\alpha `$ $`=`$ $`e^{i\theta Q}\mathrm{\Phi }^\alpha e^{i\theta Q}`$ (63)
$`=`$ $`R_{\alpha \beta }^+(\theta )\mathrm{\Phi }^\beta `$
We therefore conclude that the scalar theory in two dimensions manifests all features of duality symmetry pertaining to either twice odd or twice even dimensions. It, therefore, goes beyond the results found by the algebraic approach.
## 4 Dual Projection in (2+1) dimensions
As we have stressed, conventional group theoretical arguments fail to discuss duality in odd dimensional theories. This is possible in our approach. As an example we consider the (2+1) dimensional Maxwell theory. We shall see that the solution of the Gauss constraint leads naturally to a differential operator to be used as the dual projector. In subsection 4.1 the canonical transformation in the dual projection involves an even-parity operation. The resulting action displays a continous SO(2) group of symmetry transformations. This is a new result that could not be disclosed by algebraic methods. However, since in three spacetime dimensions vector fields are duality related to scalars where there is no Gauss law restriction, an odd-parity kernel also exists. The complete diagonalization of the electromagnetic action into two distinct type actions, that would be a prototype of chiral bosons in three dimensions, cannot be done in the coordinate space. On the other hand, if momentum space approach of the dual-projection is adopted, such a structure is then shown to exist if a special combination of the Fourier modes is considered. This issue will be studied in subsection 4.2.
### 4.1 Even Parity Projection
To begin with we see that the solution of the Gauss constraint that takes proper care of the spatial indices must involve a canonical scalar field,
$$\pi _kϵ_{km}_m\varphi $$
(64)
The problem here is that, in contrast to the even dimensional cases, parity is not a good property to look for in the generalized curl operator ($`ϵ`$). However, exploiting the property,
$$^2=\left(ϵ_{km}_m\right)\left(ϵ_{kn}_n\right)$$
(65)
we may find an even solution for the dual projection by the following canonical transformations,
$`\pi _k`$ $`=`$ $`\eta ϵ_{km}_m\left(\varphi ^+\varphi ^{}\right)`$
$`A_k`$ $`=`$ $`{\displaystyle \frac{ϵ_{km}_m}{\sqrt{^2}}}\left(\varphi ^++\varphi ^{}\right)`$ (66)
where $`\eta `$ gives the signature of the dual projection and $`\sqrt{^2}`$ is included for dimensional reasons. Substituting these into the Maxwell action (10) we obtain,
$$S_\eta =d^3x\left(\eta \dot{\varphi }^\alpha ϵ^{\alpha \beta }B^\beta B^\alpha B^\alpha \right)$$
(67)
where $`B^\alpha `$ is a shorthand for
$$B^\alpha =\sqrt{^2}\varphi ^\alpha ;\varphi ^\alpha =\varphi ^+;\varphi ^{}$$
(68)
Not surprisingly, the resulting action is explicitly duality invariant displaying a continuous $`SO(2)`$ symmetry.
Let us next examine the symmetry contents revealed in (67). These actions are manifestly invariant under the continuous SO(2) transformations,
$$\varphi _\alpha R_{\alpha \beta }^+\varphi _\beta $$
(69)
where $`R_{\alpha \beta }^+`$ is the proper $`SO(2)`$ rotation matrix (27). The generator of the infinitesimal symmetry transformation is given by the Chern-Simons form,
$$Q_\eta =\frac{\eta }{2}d^3x\varphi _\alpha B_\alpha $$
(70)
and the finite transformations (69) are generated as,
$`\varphi _\alpha \stackrel{~}{\varphi }_\alpha `$ $`=`$ $`e^{i\theta Q}\varphi _\alpha e^{i\theta Q}`$ (71)
$`=`$ $`R_{\alpha \beta }^+(\theta )\varphi _\beta `$
This result comes by using the basic symplectic brackets obtained from (67),
$$\{\varphi _\alpha (\stackrel{}{x}),B_\beta (\stackrel{}{y})\}=\eta ϵ_{\alpha \beta }\delta (\stackrel{}{x}\stackrel{}{y})$$
(72)
This is the parallel of the usual constructions done in the 4D Maxwell theory and its D=4k extensions to induce a duality symmetry in the action.
It is interesting to check the dual projection procedure when applied to a (2+1) dimensional scalar field theory. Recall that a vector field in 3D has only one degree of freedom and spin zero. The vector and the scalar fields are related by the dualization,
$$^{}F_{\mu \nu }=ϵ_{\mu \nu \lambda }^\lambda \varphi $$
(73)
A simple analysis shows that under this transformation, $`B\dot{\varphi }`$ and $`E_kϵ_{km}_m\varphi `$ meaning that there is an exchange of the potential and sympletic sectors between the models. Clearly the Maxwell Gauss law is automatically satisfied in the scalar representation,
$$.𝐄=0\stackrel{}{}_k\left(ϵ_{km}_m\varphi \right)=0$$
(74)
showing that in the dual point of view the gauge constraint has no dynamical consequences. It is now easy to apply the dual projection to the conventional scalar action, written in a first order form,
$$S=𝑑x\left[\pi \dot{\varphi }\frac{1}{2}\pi ^2\frac{1}{2}\left(_i\varphi \right)^2\right]$$
(75)
by performing the following canonical transformations,
$`\varphi `$ $`=`$ $`\varphi _++\varphi _{}`$
$`\pi `$ $`=`$ $`\eta \sqrt{^2}\left(\varphi _+\varphi _{}\right)`$ (76)
The resulting action reproduces the result (67), as it should. This also shows the equivalence of these theories in the context of dual projection.
### 4.2 Odd Parity Projection
To disclose an odd parity projection we note the fundamental difference from the two dimensional case. In the latter, there is only one space dimension and hence it is straightforward to define the odd projection in the coordinate space. For any space dimension greater than one, this projection becomes ambiguous in the coordinate space. The standard way to bypass this problem is to go over to the momentum space. Let us therefore introduce a two-dimensional basis, $`\left\{\widehat{e}_a(k,x),a=1,2\right\}`$, with $`(k,x)`$ being conjugate variables and the orthonormalization condition given as,
$$𝑑x\widehat{e}_a(k,x)\widehat{e}_b(k^{},x)=\delta _{ab}\delta (k,k^{})$$
(77)
We choose the vectors in the basis to be eigenvectors of the Laplacian, $`^2=`$,
$$^2\widehat{e}_a(k,x)=\omega ^2(k)\widehat{e}_a(k,x)$$
(78)
The action of $``$ over the $`\widehat{e}_a(k,x)`$ basis is
$$\widehat{e}_a(k,x)=\omega (k)M_{ab}\widehat{e}_b(k,x)$$
(79)
that together with definition (78) gives <sup>1</sup><sup>1</sup>1Here we use the matricial notation where $`\left(\stackrel{~}{M}\right)_{ab}=M_{ba}`$.
$$\stackrel{~}{M}M=I$$
(80)
Let us use this basis to represent the elementary fields. Since the Maxwell theory here is equivalent to a spin zero scalar, it suffices to analyze this last case. The canonical scalar and its conjugate momentum have the following expansion,
$`\mathrm{\Phi }(x)`$ $`=`$ $`{\displaystyle 𝑑kq_a(k)\widehat{e}_a(k,x)}`$
$`\mathrm{\Pi }(x)`$ $`=`$ $`{\displaystyle 𝑑kp_a(k)\widehat{e}_a(k,x)}`$ (81)
with $`q_a`$ and $`p_a`$ being the expansion coefficients. It leads to a representation of the action as a two-dimensional oscillator. The phase-space is four-dimensional, representing two degrees of freedom per mode,
$$S=𝑑k\left\{p_a\dot{q}_a\frac{1}{2}p_ap_a\frac{\omega ^2}{2}q_aq_a\right\}$$
(82)
Let us now consider the following canonical transformation,
$`p_a(k)`$ $`=`$ $`\omega (k)ϵ_{ab}\left(\phi _b^{(+)}\phi _b^{()}\right)`$
$`q_a(k)`$ $`=`$ $`\left(\phi _b^{(+)}+\phi _b^{()}\right)`$ (83)
such that (82) gets diagonalized,
$$S=S_++S_{}$$
(84)
where,
$$S_\pm =𝑑k\omega (k)\left(\pm \dot{q}_aϵ_{ab}q_b\omega (k)q_aq_a\right)$$
(85)
This action displays the $`Z_2`$ symmetry since, under the transformation $`\phi _a^\alpha \sigma _1^{\alpha \beta }\phi _a^\beta `$, the two pieces $`S_+`$ and $`S_{}`$ are swapped. Hence the theory shows both $`SO(2)`$ and $`Z_2`$ symmetries, depending upon the nature of the transformation.
It may be useful to point out that the actions $`S_\pm `$ are the analogues actions for chiral bosons. Each of these actions characterises one degree of freedom per mode in phase space or half degree of freedom in configuration space that represent a chiral scalar. Since the Maxwell theory is equivalent to a scalar field theory in three dimensional spacetime, this is valid for the photon field as well. This shows that it is indeed possible to obtain a phase space reduced, diagonal selfdual solution for the Maxwell field in this odd dimensional spacetime.
## 5 Conclusions
In this paper we have developed a new technique for obtaining duality symmetric actions. Different aspects of duality symmetry were discussed. Our technique was based on an operation which was termed as a dual projection. Since the analysis was always carried out for first order systems an equivalence between the lagrangian and hamiltonian approaches was possible. Indeed the dual projection entailed a change of variables which was a canonical transformation in the phase space. The analysis was completely general, which required an appropriate definition of the functional curl used for solving the Gauss law constraint. The conventional results for the construction of duality symmetric actions in any even dimension was contained in a single expression. Zooming in on the particular even dimension, either $`4k`$ or $`4k2`$, immediately displayed the relevant $`SO(2)`$ or $`Z_2`$ symmetry, respectively. This sharp line of distinction was related to the parity of the functional curl. Our analysis, however, went beyond. Using the same techniques it was possible to discuss the property of duality symmetry in odd dimensions, which were not analysed earlier. Since the definition of parity of the functional curl was problematic in these dimensions, it was not straightforward to give a general discussion for arbitrary odd dimensions, as was the case for the even dimensions. Hence we elaborated our methods by concentrating on the three dimensional Maxwell theory.
Historically, the study of duality symmetry began by considering the symmetry among the electric and magnetic fields. This led to an invariance of the equations of motion but not of the actions. However it was felt that these were composites and one should study the symmetry properties in the context of potentials which were regarded as the basic entities. This was achieved by introducing an internal space. The new duality symmetry involving the potentials preserved the invariance of the corresponding actions. Two distinct classes of symmetries, pertaining to either $`4k`$ or $`4k2`$ dimensions were found. Now the potentials can also be regarded as composites that were defined in terms of their fourier modes. Pursuing this line of research and studying the symmetries in terms of these modes we were able to show that, suitably interpreted, the distinction among the duality groups could be obiliterated. In other words, in all even dimensions, both $`SO(2)`$ and $`Z_2`$ symmetry groups could coexist.
The present work goes a step further. We show that it is possible to obtain duality symmetric actions for odd dimensions also. At the level of potentials, the duality symmetry in a three dimensional Maxwell theory was shown to posses the $`SO(2)`$ symmetry. Furthermore, by passing to the Fourier modes the $`Z_2`$ symmetry was also revealed.
Acknowledgements CW would like to thank the S.N. Bose Centre for Basic Sciences for the invitation, kind hospitality and financial support during his stay here. This author is partially supported by CNPq, FUJB and FAPERJ, Brasilian scientific agencies.
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# I Introduction
## I Introduction
Already when introducing the notion of a surface quantity Gibbs implicitly entertained the idea of a phase field $`\varphi `$: any density of an extensive quantity (e.g., the mass density) between two coexisting phases changes gradually (but swiftly) from its value in one phase to its value in the other. The existence of a transition zone, though microscopically of atomic extent (far enough from a second-order phase transition), underlies the very Gibbs definition of surface quantities. In phase transition phenomena, either of first or second order, this notion has been adopted in Landau’s spirit. Because energy is an extensive quantity, too, there is an extra energetic cost associated with the transition region, characterized in the appropriate thermodynamical potential density by a term of the form $`ϵ^{}(\varphi )^2`$, $`ϵ^{}`$ being the stiffness of the transition region.
The notion of a phase field has appeared abundantly in the literature in the context of phase transition phenomena. The transition width diverges for a second-order phase transition at the critical point, and thus it is essential to introduce the transition region. For a first-order transition, such as a liquid-solid interface, confering an importance to the interface thickness may seem quite anecdotic if one is interested in properties which occur on a scale larger than the atomic one; typical examples are dendritic patterns occuring at the scale of a $`\mu m`$. Nevertheless, it is here, where phase-field modeling has become most useful in numerical treatments.
Before phase-field models became popular, it seemed quite natural to treat the surface as a geometric location on which boundary conditions are imposed (e.g., for a moving front the normal velocity is proportional to the jump in the gradient of the temperature or concentration field). This is the so-called sharp interface approach, adopted both in analytical and numerical studies in a variety of contexts of front problems.
There has been an upsurge of interest in the phase-field approach to free-boundary problems more recently, though the method was actually introduced pretty early , as a computational tool to model solidification. Various studies have demonstrated the virtues of this method in moving-boundary problems.
Regarding the way how to use phase-field models, there are two distinctly different philosophies. These may be best discussed considering dendritic growth, where a set of well-established continuum equations exists describing phenomena in terms of a sharp interface. On the basis of this knowledge, a phase-field model can be justified by simply showing that it is asymptotic to the correct sharp-interface description, i.e., that the latter arises as the sharp-interface limit of the phase-field model when the interface width is taken to zero. This is definitely a sufficient condition for the phase-field model to yield a correct description of the continuum limit, providing the interface thickness is taken small enough. Small enough sometimes may mean impractically small. The second approach to phase-field modeling is to guess or derive an appropriate form for the free energy of the two-phase system, including the energy cost of the transition region and to regard this as a physical model in its own right. In this case, one might actually forgo considering the limit of small interface thickness and such a model would even make sense, if the strict limit of vanishing interface width did not correspond to sensible physics.
Of course, a problem arises, if a phase-field model obtained in the second way gives predictions that are different from that of the sharp-interface equations. A follower of the first philosophy would then discard the phase-field model, whereas one of the second might contemplate the possibility that his model contain more physics than the sharp-interface model. In the case of dendritic growth the situation is pretty clear: the sharp interface model gives the right answers. However, this statement cannot be generalized easily, since not all sharp-interface models are as well-founded as that for dendritic growth and because the extreme smallness of the interface width cannot be always guaranteed (it might for example become doubtful for a phase transition that is only weakly first order).
A related issue is the question of thermodynamic consistency, i.e., the derivation of the model in the spirit of Gibbs from a free-energy or entropy functional. It is clear that with a known sharp-interface limit in mind, there is no need at all to obtain a phase-field model this way (which would mean to make it “thermodynamically consistent”) as long as one ensures its asymptotic approach to this limit. In fact, it has turned out that in some cases where both a thermodynamically consistent formulation of a phase-field model and a nonvariational formulation exist, the latter was numerically more efficient and hence preferable on practical grounds.
On the other hand, thermodynamic consistency has its virtues. This can be seen particularly well in the case considered here, the influence of elasticity on the stability of a solid interface. It is quite straightforward to write down the contribution of the elastic energy to the total free energy. Hence, if we have a good idea about the physical origin of the free energy to be considered, the corresponding phase-field model is easily obtained, and it is bound to be right. As a result, one may derive sharp-interface equations in cases where they are not known.
For the Grinfeld instability to be considered here, the sharp-interface equations are well known. Nevertheless, it is of course tremendously satisfying, if they simply pop out of the phase-field equations as the sharp-interface limit. Not only does this provide a natural countercheck of our ansatz for the free energy, but it also gives us a new angle of view at the instability, leading to the prediction of circumstances in which the Grinfeld instability should not occur under anisotropic stress, but might appear with isotropic stress. We shall consider this point in section II.
Let us return to the advantages of the phase-field method. The first virtue of phase fields is pretty obvious: instead of tracking permanently the a priori unknown interface position in the sharp-interface limit, and imposing nontrivial boundary conditions for the discontinuity of the fields, the interface in the phase-field approach is nothing but the location of a rapid variation of the field $`\varphi `$, while the two phases are treated as the same entity. Thus there is no boundary condition to be imposed in the transition region, a fact which greatly facilitates both the numerical implementation and analysis. This is done at some price: one must, in principle, mix disparate length (and thus time) scales: the pattern length and the interfacial width, whose ratio may range over many orders of magnitude. This may render the numerical procedure excessively expensive, a fact which would quickly take us back to the sharp-interface problem, where the small length scales out of the problem.
As discussed recently, the sharp-interface limit that one would like to represent when writing the phase-field equations, only makes physical sense on “outer” length scales much larger than the physical extent $`l_c`$ of the transition region, and thus does not depend structurally on the details of the interface shape on the inner $`l_c`$ scale. The mathematical question, formulated in the framework of phase-field models, of formally recovering the sharp interface description via an asymptotic (multi-scale) expansion in the limit $`l_c0`$ might, from this point of view, seem irrelevant.
It should however be kept in mind that the actual matching conditions are imposed for the limit, where an “inner” variable (defined in the transition region) goes to infinity. This entails a certain amount of liberty in the choice of functions defining the free-energy density, because the precise behavior of these functions on the inner scale does not matter. Hence, the validity of a phase-field model can indeed be judged by simply showing that it asymptotically reproduces the correct sharp-interface description. Whether the additional information encoded in the structure of the phase-field model on the inner scale is physically relevant, is a question to be decided on a case by case basis (as implied by our discussion above).
An example where this is relevant, is provided by the Young condition for the contact angle of a droplet on a substrate. This is a condition on “outer” scales, while the inner scale is rather governed by van der Waals interactions of a thin liquid film with the substrate, leading to some nontrivial corrections of the wetting profile at small atomic length scales. In other words, what matters in the phase-field description is not that the width of the interface be of atomic extent, but rather that it be small in comparison with the scale of the pattern of interest.
This physical argument has been cast in a mathematical form in , where the thin-interface limit was considered, arising from an alternative asymptotic procedure. This has to be contrasted with the sharp-interface limit with the small parameter being the ratio of the interface thickness to the capillary length. Making the width of the interface small only in comparison with the scale of the pattern leads to a rather important enhancement of the computing speed, thus rendering the phase-field approach attractive with regard to numerical efficiency as well. Unfortunately, in most systems the thin-interface limit is not as easily accessible as for dendritic growth in the thermal model . Hence, it is difficult in general to make our argument mathematically rigorous. However, we will give a consideration of length scales in section II that clarifies what are the ratios that have to be kept small for our model to be a good description. Moreover, we have checked that the dependence of the results on the interface width becomes weak for the small values that we use for the latter.
Finally the phase-field approach has the additional virtue of regularizing instabilities, such as the development of cuspy structures, often setting severe limitations in numerical studies. In the elastic system, we are led to the question whether the sharp-interface models still make sense in the limit where they predict finite-time singularities. No such singularities arise in the phase-field model, thus possibly extending the range of validity of the latter beyond that of the former. This will be discussed in more detail in section IV.
In the context of growth phenomena phase-field approaches have been introduced in problems involving temperature or concentration fields. There are myriads of situations, however, where the corresponding transition is monitored, or at least affected, by strain. A typical situation is a solid under uniaxial stress. This leads to the Asaro-Tiller-Grinfeld (ATG) instability (see Ref. for a recent review. A surface corrugation allows to lower the stored elastic energy. Other examples of particular interest include solid-solid transformations, phases in nonequilibrium gels, molecular beam epitaxy, solidification of lava, etc. It is thus highly desirable to develop a phase-field approach including the stress as an active variable. Very recently, two groups, Müller-Grant (MG) and Kassner-Misbah (KM) , have independently developed such an approach and have given a brief account on it.
This paper will present extensive discussions of this question and give new results. We shall also provide a comparison between the MG and KM models, and point out similarities and dissimilarities. As already known from sharp-interface simulations a solid under stress presents a somehow stringent behavior in that no stable steady-state solution seems to exist. This is also the case from analytical studies in the long-wavelength limit. Nozières had shown that the bifurcation from the planar front to the deformed one is subcritical (the analogue of a first-order transition). The study of Spencer and Meiron focused on structures with a given basic wave number in the absence of gravity (where the instability does not have a threshold) and on systems, in which the transport mechanism necessary for the instability to manifest itself is surface diffusion. They find that in the unstable range of wavenumbers (i.e., for wavenumbers below the marginal value, above which surface tension stabilizes the planar interface) there exist finite-amplitude steady-state solutions, if the wavenumber is close enough to marginal. This branch of steady-state solutions terminates by structures developing cusp singularities, despite the stabilizing influence of surface tension. It cannot be overemphasized that this result is not an artifact of their numerics. Indeed, they investigate carefully the effects of numerical fine-graining using a code with spectral accuracy, and their discretization sequence seems to get fine enough to render their extrapolation valid. The evidence for the appearance of true cusps in the sharp-interface continuum model becomes compelling, if one takes into account the work of Chiu and Gao who found analytical solutions developing cusp singularities in finite time. The conclusion of Spencer and Meiron is that for generic initial conditions, including sufficiently small wavenumbers, finite-time singularities will always occur. Moreover, they state that this is also true in the presence of gravity (beyond the threshold) and under fairly general conditions.
For a physical system, finite-time singularities will be prevented by intervening effects that are not considered in the model. This would mean that nonlinear elasticity and plasticity have to be taken into account. For example, the formation of dislocations (plasticity) could blunt cusps again.
Two questions then arise. What kind of structures can be expected before cusp formation and what kind of structures will prevail eventually? Our previous study has shown that the initial cellular pattern may develop into a super-structure, where a groove or several of them accelerate in a spectacular fashion, thus relieving the stresses in their surroundings significantly and causing nearby grooves to recede. Further evolution of the structure was difficult to handle numerically, due to the development of cusps which appear in the numerics as they should according to the analytic results . In the phase-field approach presented here, no cusps can arise. The question is of course legitimate, whether the model, which allows to track the dynamics of the structure much beyond the times where earlier studies had to fail numerically, still gives a faithful description of the physics. Here we take the point of view, that the details of the description in the locations, where stresses become large, may not be correctly captured by the model, since we neither include nonlinear elastic effects nor other effects such as capillary overpressure explicitly . If they have the right sign, these might prevent cusp formation . However, the result of any such effect must be to blunt cusps, which the phase-field model does. As we shall see, it does so in a non-obtrusive way by introducing a cutoff for interface curvature. Moreover, away from the cusps, stresses are low enough for linear elasticity to apply. Hence, we believe that the development of the overall morphology is still correctly described by the phase-field model.
A more detailed justification would point out that nonlinear elasticity will first make itself felt by stresses increasing more slowly as a function of strains than in the linear case and that next plasticity will act to introduce an upper cutoff for stresses, where the material will yield. Now the effect of any resulting modification in the stress-strain relationship on the remaining body can be reproduced by cutting out the piece of material where linear elasticity ceases to hold and by requiring boundary conditions at the edge of the cut-out piece that correspond to the correct stresses. In the case of material yielding near a singularity of the stress field (obtained assuming linear elasticity), these boundary conditions will essentially be that the stresses are close to the yield stress at the boundary. This is mimicked by the phase-field model in which the maximum supportable stress is, for a given geometry, determined by the interface width.
Therefore, we think this is one case, where the phase-field model can do more in the description of the physical system than its sharp-interface limit.
It will emerge that usually one leading groove continues to deepen while neighboring ones recede after the winner has started to relieve the stress that kept them growing. Finally the surface shows a single deep groove evolving in time and becoming a location of a strong stress accumulation, possibly until the fracture threshold is reached. Presumably, before that stage is reached the validity of the model will break down. We shall make some speculation on future directions to elucidate the physical behavior in real systems. For generic initial conditions, we may consider the dynamics a continual coarsening process which initially develops as described in and later is dominated by groove growth and shrinkage.
The paper is organized as follows. In section II we give the continuum equations ordinarily used in the description of the Grinfeld instability. This is mainly done in order to introduce appropriate length and time scales in nondimensionalizing the equations; then we present our phase-field approach and discuss how the interface width has to be chosen in comparison with the other length scales. We demonstrate how the phase-field model can be employed to derive new sharp-interface equations in the presence of body forces breaking rotational invariance.
Section III presents validation results, a comparison of the MG and KM models and describes the main findings of our simulations. Section IV sums up the results and discusses perspectives. The mathematically rigorous asymptotic expansion used to derive the sharp-interface limit has been relegated to an appendix, as the calculation is somewhat lengthy and would interrupt the flow of the text. Since we use the MG model in a slightly different form from that presented originally , we give the connection between the two formulations in a second appendix. Finally, appendix C contains the analytic derivation via conformal mapping of the stresses for a particular interface shape to be compared with the numerics.
## II The Grinfeld instability
### A Sharp-interface equations
A description of the basic ingredients of the Grinfeld instability has been given elsewhere . Therefore, we may restrict ourselves to explaining the physical mechanism and giving the equations.
We wish to describe the behavior of a solid submitted to uniaxial stress, at the surface of which material transport is possible. Consider the example of a solid in contact with its melt. An accidental corrugation of the surface will act to reduce the stress at its tip and increase it in the valleys next to it. That is, the solid can decrease its elastic energy density by growing tips (where the stress is lower) and by increasing the depths of valleys (where it gets rid of material having a higher density of elastic energy due to larger stresses). This tendency is most easily cast into equations by writing down the chemical potential difference $`\mathrm{\Delta }\mu =\mu _s\mu _{\mathrm{}}`$ at the interface (the subscripts refer to the solid and liquid or nonsolid phases, respecively) :
$$\mathrm{\Delta }\mu =\frac{1\nu ^2}{2E\rho _s}(\sigma _{tt}\sigma _{nn})^2+\frac{1}{\rho _s}\gamma \kappa +\frac{\mathrm{\Delta }\rho }{\rho _s}g\zeta (x)$$
(1)
Herein, the first term is of elastic origin; $`\sigma _{tt}`$ and $`\sigma _{nn}`$ are the normal stresses tangential and perpendicular to the interface, $`E`$ is Young’s modulus, $`\nu `$ the Poisson number, and $`\rho _s`$ the density of the solid. The second term describes the stabilizing influence of the surface stiffness $`\gamma `$, taken isotropic here (so it becomes identical to the surface energy). $`\kappa `$ is the curvature of the interface; for simplicity, we consider the two-dimensional case only. Finally, the third term is the contribution of gravity ($`g`$), where $`\mathrm{\Delta }\rho =\rho _s\rho _{\mathrm{}}`$ is the density contrast between the solid and the liquid (or vacuum) and $`\zeta (x)`$ is the interface position, given by its $`z`$ coordinate (the $`z`$ axis is oriented antiparallel to the gravitational force). Equation (1) holds for plane strain. For plane stress, the prefactor $`1\nu ^2`$ has to be dropped.
The dynamics is then described by giving the normal velocity in terms of the chemical potential difference. For a solid in contact with its melt this would simply be
$$v_n=\frac{1}{k}\mathrm{\Delta }\mu ,$$
(2)
where $`k`$ is an inverse mobility with the dimension of a velocity. In the case of a solid in contact with vacuum and surface diffusion as the prevailing transport mechanism we would have $`v_n=D^2\mathrm{\Delta }\mu `$ instead.
Of course, in order to compute $`v_n`$, we must first obtain the stresses entering (1). This involves solving an elastic problem ($`_j\sigma _{ij}=0`$) with a prescribed external stress and boundary conditions on the interface, assuming an appropriate constitutive law. Ordinarily, Hooke’s law for isotropic elastic bodies is assumed \[therefore we have only two elastic constants in (1)\]. Neglecting the capillary overpressure, which usually is a good approximation, we have as boundary conditions at the interface $`\sigma _{nn}=p`$, where $`p`$ is the pressure in the second phase, and $`\sigma _{nt}=0`$, i.e., the shear stress vanishes.
A linear stability analysis of a planar interface under the dynamics given by equations (1) and (2) yields the following dispersion relation ($`\omega `$ is the growth rate, $`q`$ the wave number)
$$\omega =\frac{1}{k\rho _s}\left\{\frac{2\sigma _0^2(1\nu ^2)}{E}q\gamma q^2\mathrm{\Delta }\rho g\right\}.$$
(3)
$`\sigma _0`$ is the uniaxial external stress. Eq. (3) provides us with a critical wave number $`q_c=\sqrt{\mathrm{\Delta }\rho g/\gamma }`$ (an inverse capillary length) and a critical stress $`\sigma _{0c}=\left[\gamma q_cE/(1\nu ^2)\right]^{\frac{1}{2}}`$, below which the planar interface is stable. The wave number of the fastest-growing mode can be inverted to give a length
$$\mathrm{}_1=\frac{\gamma E}{\sigma _0^2(1\nu ^2)}.$$
(4)
Apart from a prefactor of $`2/\pi `$, this length is identical to the so-called Griffith length.
Note that a planar interface will not remain at its original equilibrium position, even when only a subthreshold stress is applied, i.e., when it is “stable”. Our dynamical equations predict that it has a nonzero velocity as long as $`\mathrm{\Delta }\mu `$ is different from zero. Hence it will recede to smaller values of $`\zeta `$. If the density of the solid is bigger than that of the liquid, the chemical potential on the solid side of the receding interface decreases faster than that on the liquid side and there is a new equilibrium position, which can be computed directly from (1) and which evaluates to
$$\mathrm{\Delta }\zeta \mathrm{}_2=\frac{1\nu ^2}{2\mathrm{\Delta }\rho gE}\sigma _0^2.$$
(5)
Equations (4) and (5) provide us with two independent length scales of the problem, the first of which is due to a competition between stress and surface energy, while the second arises from the competition between stress and gravity. For the purpose of nondimensionalizing equations, $`\mathrm{}_1`$ is more appropriate, as this length does not diverge in the limit of vanishing density difference (or gravity). To obtain a natural time scale $`\tau `$, we can replace $`q`$ in either of the two wave-number dependent terms of (3) by $`1/\mathrm{}_1`$. This leads to
$$\tau =\frac{k\rho _s\gamma E^2}{\sigma _0^4(1\nu ^2)^2}.$$
(6)
The nondimensional version of the dispersion relation then reads ($`\stackrel{~}{\omega }=\tau \omega `$, $`\stackrel{~}{q}=\mathrm{}_1q`$):
$$\stackrel{~}{\omega }=2\stackrel{~}{q}\stackrel{~}{q}^2\frac{\mathrm{}_1}{2\mathrm{}_2},$$
(7)
which shows clearly that the problem without gravity (when $`\mathrm{}_2`$ becomes infinite) can be made parameter free, i.e., elastic and other parameters only set the time and length scales; apart from that we should expect the same dynamics for all systems. With gravity, the dynamics is essentially determined by the ratio of the two length scales introduced.
### B Elastic energy and state of reference
Let us now proceed to investigate the contributions to the free energy of the same system. The phase-field model will then consist in writing the free-energy density that takes into consideration the global elastic energy in both phases.
As is usual with elastic problems, it is important to specify the state of reference defining the positions of material particles with respect to which displacements are measured. This is crucial whenever the reference state is not that of an undeformed body but one that is subject to prestraining (which will turn out useful later). In order to make this point clear, and in the hope of helping subsequent discussions, we would like first to dwell on this issue.
If the only energy present is elastic energy and Hooke’s law holds true, the free energy per unit volume can be written as
$$f=\mu u_{ij}u_{ij}+\frac{\lambda }{2}u_{ii}^2,$$
(8)
where summation over double subscripts is implied. $`\lambda `$ and $`\mu `$ are the Lamé constants, $`\mu `$ being better known as the shear modulus. For plane strain, these elastic constants are related to Young’s modulus and the Poisson ratio via $`\mu =E/[2(1+\nu )]`$ and $`\lambda =E\nu /[(1+\nu )(12\nu )]`$. The stress tensor $`\sigma _{ij}`$ is then
$$\sigma _{ij}=\frac{f}{u_{ij}}=2\mu u_{ij}+\lambda u_{kk}\delta _{ij}=2\mu (u_{ij}u_{kk}\frac{\delta _{ij}}{d})+Ku_{kk}\delta _{ij}.$$
(9)
$`K=2\mu +\lambda /d`$ is the bulk modulus, and the last relation has the advantage of making explicit the parts of the stress tensor causing pure shape and pure volume changes, respectively. ($`d`$ is the spatial dimension.) We will nevertheless mainly use the relations containing $`\mu `$ and $`\lambda `$, which are more compact. The implied reference state here is $`u_{ij}=0`$, for which $`\sigma _{ij}=0`$.
However, if we choose a reference state given by a different strain tensor $`u_{ij}^{(0)}`$, setting $`\stackrel{~}{u}_{ij}=u_{ij}u_{ij}^{(0)}`$, then we should not simply replace $`u_{ij}`$ by $`\stackrel{~}{u}_{ij}`$ in (9), as the stress tensor and the elastic energy are, in principle, measurable quantities and should thus be unaffected by a change of strain reference state. Hence, we have to write
$$\sigma _{ij}=2\mu \left(\stackrel{~}{u}_{ij}+u_{ij}^{(0)}\right)+\lambda \left(\stackrel{~}{u}_{kk}+u_{kk}^{(0)}\right)\delta _{ij},$$
(10)
and change definition (8) accordingly, i.e. replace $`u_{ij}`$ by $`\stackrel{~}{u}_{ij}+u_{ij}^{(0)}`$. In this situation, the zero-strain state would not be stress-free. An alternative way to specify a reference state would then consist in giving the stress of the zero-strain state.
In general, the thermodynamic system under consideration will not only contain elastic contributions. Then the equilibrium state, corresponding to a minimum of the free energy, may not be a state of vanishing strain. A trivial example is a solid in equilibrium with its melt, where the equilibrium state in the solid corresponds to the strain produced by the equilibrium pressure $`p`$ of the liquid (the equilibrium stress tensor of the solid is $`p\delta _{ij}`$). The form of the free-energy density accounting for such a situation is not (8) \[which does not exhibit a minimum at $`u_{ij}^{(\mathrm{eq})}0`$\] but
$$f=\mu \left(u_{ij}u_{ij}^{(\mathrm{eq})}\right)\left(u_{ij}u_{ij}^{(\mathrm{eq})}\right)+\frac{\lambda }{2}\left(u_{ii}u_{ii}^{(\mathrm{eq})}\right)^2.$$
(11)
This is manifestly minimum at $`u_{ij}^{(\mathrm{eq})}`$ and the nonzero value of the latter quantity takes into account nonelastic contributions to the free energy. If we now define the stress from the first relation in Eq. (9), i.e., $`\sigma _{ij}=f/u_{ij}`$, it will be nonzero only, if there are forces driving the system away from equilibrium. If we rather define it via the second equality of Eq. (9), meaning that we set $`\sigma _{ij}=2\mu u_{ij}+\lambda u_{kk}\delta _{ij}`$ (which is now different from $`f/u_{ij}`$), it will describe, in addition to these forces the prestress necessary to keep the system in equilibrium. An invariant relation between stresses and strains follows from requiring
$$f=_{u_{ij}^{(\mathrm{eq})}}^{u_{ij}}\left(\sigma _{ij}\sigma _{ij}^{(\mathrm{eq})}\right)𝑑u_{ij},$$
(12)
which leads to
$$\sigma _{ij}\sigma _{ij}^{(\mathrm{eq})}=2\mu \left(u_{ij}u_{ij}^{(\mathrm{eq})}\right)+\lambda \left(u_{kk}u_{kk}^{(\mathrm{eq})}\right)\delta _{ij}.$$
(13)
Once both $`\sigma _{ij}^{(\mathrm{eq})}`$ and $`u_{ij}^{(\mathrm{eq})}`$ are specified, this equation gives us an unambiguous relationship between stresses and strains. Depending on which variables we choose to define the reference state, we obtain the conjugate variables of the same state from (13). If we choose, for instance, a strain-free state as reference, this equation will provide us with the corresponding stress of reference, if we choose a stress-free state of reference, it will yield the strain of reference.
As an example we can look at a case where the equilibrium stress is $`\sigma _{ij}^{(\mathrm{eq})}=p_0\delta _{ij}`$, and ask what should the strain be. Hooke’s law – written in a such a way that the absence of strain implies the absence of stress as well \[see (9)\] – then gives us an equilibrium strain
$$u_{ij}^{(\mathrm{eq})}=\frac{p_0}{2\mu +\lambda d}\delta _{ij}.$$
(14)
Since the free energy must not depend on the choice of reference state, it is clear that it does not matter whether we use $`u_{ij}`$ or $`\stackrel{~}{u}_{ij}`$ in (13), providing that we use the correct values of $`u_{ij}^{(\mathrm{eq})}`$ and $`\stackrel{~}{u}_{ij}^{(\mathrm{eq})}`$, respectively. Suppose that we choose another reference state characterized by the strain tensor $`\stackrel{~}{u}_{ij}`$ in such a way that when the strain is zero, the stress is equal to $`p_{0s}\delta _{ij}`$. A vanishing strain then corresponds to a pre-stressed situation. If the equilibrium stress is again $`p_0\delta _{ij}`$ as above, we must have a new equilibrium strain $`\stackrel{~}{u}_{ij}^{(\mathrm{eq})}`$ obeying, according to our invariant relation (13), $`p_{0s}\delta _{ij}+p_0\delta _{ij}=2\mu \stackrel{~}{u}_{ij}^{(\mathrm{eq})}\lambda \stackrel{~}{u}_{kk}^{(\mathrm{eq})}\delta _{ij}`$. After a simple manipulation we obtain
$$\stackrel{~}{u}_{ij}^{(\mathrm{eq})}=\frac{p_{0s}p_0}{2\mu +\lambda d}\delta _{ij}$$
(15)
Of course, Eq. (14) is a special case of (15) for $`p_{0s}=0`$. The free energy, expressed by $`\stackrel{~}{u}_{ij}`$, then reads
$`f`$ $`=`$ $`\mu (\stackrel{~}{u}_{ij}\stackrel{~}{u}_{ij}^{(\mathrm{eq})})(\stackrel{~}{u}_{ij}\stackrel{~}{u}_{ij}^{(\mathrm{eq})})+{\displaystyle \frac{\lambda }{2}}(\stackrel{~}{u}_{ii}\stackrel{~}{u}_{ii}^{(\mathrm{eq})})^2`$ (16)
$`=`$ $`\mu \stackrel{~}{u}_{ij}\stackrel{~}{u}_{ij}+{\displaystyle \frac{\lambda }{2}}\stackrel{~}{u}_{ii}\stackrel{~}{u}_{jj}+(p_0p_{0s})\stackrel{~}{u}_{ii}+{\displaystyle \frac{d}{2(2\mu +\lambda d)}}(p_0p_{0s})^2`$ (17)
It should be realized that Eqs. (11) and (17) describe exactly the same situation, if corresponding values for the strain fields without and with a tilde are inserted. At this point we have said nothing about applied external stresses or so. However, if we choose, say, vanishing displacement as boundary condition at a planar interface, directed along the $`x`$ direction, then this will correspond to two different physical situations for the two different choices of the state of reference. Let us for simplicity assume $`p_0=0`$. According to (14), we then have $`u_{ij}^{(\mathrm{eq})}=0`$, and (11) implies (8). Setting $`u_{xx}=0`$ in Eq. (9) we obtain, because of the boundary condition $`\sigma _{zz}=0`$ that also $`u_{zz}=0`$, and there is no stress at all. On the other hand, if we set $`\stackrel{~}{u}_{xx}=0`$, we have to use (13) and (15) to obtain the elastic state of the solid. The boundary condition for $`\sigma _{zz}`$ implies $`u_{zz}=p_{0s}/(2\mu +\lambda )`$, which in turn leads to $`\sigma _{xx}=2\mu p_{0s}/(2\mu +\lambda )`$; hence vanishing displacement along our planar interface means a solid that is homogeneously strained in the $`x`$ direction with a prestress $`\sigma _0=\sigma _{xx}`$.
As we shall see later, the latter choice of the state of reference has been made in the phase-field model discussed in , the former (setting $`p_{0s}=0`$) in . These are the most natural choices, although an infinity of (less natural) alternatives is available.
### C The phase-field model
The total (solid+liquid) free energy of the system can be written as
$$F[\varphi ,\{u_{ij}\}]=𝑑V\left[f(\varphi ,\{u_{ij}\})+\frac{1}{2}\mathrm{\Gamma }ϵ^2(\varphi )^2\right]$$
(18)
where $`ϵ`$ is a length parameter controlling the order of magnitude of the transition region described by the phase field. $`\mathrm{\Gamma }=3\gamma /ϵ`$ is the energy density corresponding to the surface energy $`\gamma `$ being distributed over a layer of width $`ϵ`$. (The factor 3 is just a convenient choice, simplifying later derivations.)
If we start from the invariant form (13), we can set up a whole class of phase-field models at once and specify the reference state later. In order to be able to write a single elastic energy expression for the two-phase system, we formally treat the liquid as a shear-free solid (not including hydrodynamics). We will discuss some consequences of this approach later.
A straightforward ansatz for the elastic energy density is then
$`f_{\mathrm{el}}(\varphi ,\{u_{ij}\})`$ $`=`$ $`h(\varphi )f_{\mathrm{sol}}(\{u_{ij}\})+[1h(\varphi )]f_{\mathrm{liq}}(\{u_{ij}\})`$ (19)
$`=`$ $`h(\varphi )\left\{\mu \left(u_{ij}u_{ij,s}^{(\mathrm{eq})}\right)\left(u_{ij}u_{ij,s}^{(\mathrm{eq})}\right)+{\displaystyle \frac{\lambda }{2}}\left(u_{ii}u_{ii,s}^{(\mathrm{eq})}\right)^2\right\}`$ (21)
$`+\left[1h(\varphi )\right]{\displaystyle \frac{\stackrel{~}{\lambda }}{2}}\left(u_{ii}u_{ii,\mathrm{}}^{(\mathrm{eq})}\right)^2,`$
where $`f_{\mathrm{sol}}(\{u_{ij}\})`$ and $`f_{\mathrm{liq}}(\{u_{ij}\})`$ are the densities of elastic energy in the solid and in the liquid, respectively, and where $`h(\varphi )`$ may be interpreted as a “solid fraction”, which must be equal to one in the solid and equal to zero in the liquid. We choose $`h(\varphi )=\varphi ^2(32\varphi )`$ for reasons of convenience: with this choice $`h^{}(\varphi )=0`$ for $`\varphi =0`$ and for $`\varphi =1`$, i.e., in the bulk phases. This leads to the advantage (see appendix A) that the zeroth-order solution of the asymptotic expansion in powers of $`ϵ`$ is valid to all orders in the outer region considered.
Since different reference states may be chosen in the solid and in the liquid, the equilibrium strains carry a subscript $`s`$ or $`\mathrm{}`$, respectively. This would not be necessary here, because the prefactor \[$`h(\varphi )`$ or $`1h(\varphi )`$\] decides whether the equilibrium expression for the strain in the liquid or in the solid has to be taken. However, as soon as we take derivatives with respect to $`\varphi `$, this criterion of distinction becomes ambiguous, so we prefer to make the difference explicit from the outset. $`\stackrel{~}{\lambda }`$ is the bulk modulus of the liquid.
To account for the possibility of a phase transition, we introduce a double well potential
$$f_{\mathrm{dw}}(\varphi )=2\mathrm{\Gamma }g(\varphi ),$$
(22)
where $`g(\varphi )=\varphi ^2(1\varphi )^2`$. The minimum at $`\varphi =1`$ corresponds to the solid phase, the one at $`\varphi =0`$ to the liquid phase. Note that while this potential looks similar to the fourth-order polynomial used in the Landau theory of second-order phase transitions, it is employed in quite a different manner here. The two minima correspond to the two phases and the symmetry of the potential is of secondary importance; in Landau’s approach, symmetry considerations are at the heart of the theory, the symmetric minima describe the same phase, and the second phase corresponds to the unstable maximum in between. Since in our case both phases sit at a minimum, the transition described by the double well potential is of first order. We do not need a sixth-order polynomial as would be necessary in Landau’s theory for first-order phase transitions.
Gravity will be included in essentially the same way as in the sharp-interface equations discussed above; i.e., its effect as a body force in the mechanical equilibrium condition is neglected but its influence on the chemical potential is taken into account. This is a good approximation usually (one can estimate the cross-effect of gravity on the elastic energy to be on the order of $`\rho _sgH/\sigma _01`$, for typical heights $`H`$ of the sample). Then the contribution of gravity to the free-energy density becomes
$`f_{\mathrm{grav}}(\varphi ,z)`$ $`=`$ $`(zz_0)\left\{\rho _sh(\varphi )+\rho _{\mathrm{}}\left[1h(\varphi )\right]\right\}g`$ (23)
$`=`$ $`(zz_0)h(\varphi )\mathrm{\Delta }\rho g+(zz_0)\rho _{\mathrm{}}g,`$ (24)
where we have taken the zero point of this potential energy at $`z=z_0`$. Note that in taking fixed values for $`\mathrm{\Delta }\rho `$ and for $`\rho _s`$, we also neglect the second-order effect caused by density changes due to strain.
Finally, we wish to be able to control the equilibrium position of the interface “by hand” via addition of a constant to the free-energy density of one phase; this phenomenogical contribution to the total free-energy density may be conveniently written as
$$f_\mathrm{c}(\varphi )=h(\varphi )\frac{1\nu ^2}{2E}\sigma _{00}^2=h(\varphi )\frac{2\mu +\lambda }{8\mu (\mu +\lambda )}\sigma _{00}^2$$
(25)
and it is normalized such that setting $`\sigma _{00}=\sigma _0`$, we can keep the equilibrium position of the planar interface at the fixed value $`z_0`$, independent of $`\sigma _0`$. This is useful, for example, if one wishes to assess the relative position of the maxima or minima of an evolving structure with respect to a planar interface at the same external stress. Because of the recession of a planar interface according to (5), such a comparison would otherwise be difficult.
Collecting all contributions, we obtain for the total free-energy density
$`f(\varphi ,\{u_{ij}\},z)`$ $`=`$ $`f_{\mathrm{dw}}(\varphi )+f_{\mathrm{el}}(\varphi ,\{u_{ij}\})+f_{\mathrm{grav}}(\varphi ,z)+f_\mathrm{c}(\varphi )`$ (26)
$`=`$ $`\mathrm{\Gamma }(2g(\varphi )+{\displaystyle \frac{ϵ}{3\gamma }}\{h(\varphi )[\mu (u_{ij}u_{ij,s}^{(\mathrm{eq})})(u_{ij}u_{ij,s}^{(\mathrm{eq})})+{\displaystyle \frac{\lambda }{2}}(u_{ii}u_{ii,s}^{(\mathrm{eq})})^2]`$ (29)
$`+\left[1h(\varphi )\right]{\displaystyle \frac{\stackrel{~}{\lambda }}{2}}\left(u_{ii}u_{ii,\mathrm{}}^{(\mathrm{eq})}\right)^2`$
$`+h(\varphi )[\mathrm{\Delta }\rho g(zz_0){\displaystyle \frac{2\mu +\lambda }{8\mu (\mu +\lambda )}}\sigma _{00}^2]+(zz_0)\rho _{\mathrm{}}g\}).`$
Note that here the terms in braces, in particular the elastic term, have acquired a prefactor $`ϵ`$. This $`ϵ`$ dependence is spurious, as we have taken the prefactor $`\mathrm{\Gamma }1/ϵ`$ in front of everything, and the factor $`ϵ`$ just serves to cancel this. In fact, the only contribution to the free energy that can depend on $`ϵ`$ explicitly is the double well potential, which must ensure that in the limit $`ϵ0`$ the only possible states are the bulk phases and must therefore become infinite for all values of $`\varphi `$ different from 1 or 0. All the other energies can depend on $`ϵ`$ only implicitly via $`h(\varphi )`$, the local solid fraction of the two-phase system.
We then require $`\varphi `$ to satisfy a relaxation equation for a non-conserved order parameter. This equation takes the form
$$\frac{\varphi }{t}=R\frac{\delta F}{\delta \varphi },$$
(30)
and the prefactor $`R`$ should contain the mobility $`1/k`$ defined in (2). The dimension of $`R`$ must be (energy density $`\times `$ time)<sup>-1</sup>, which leads us to choosing $`R=1/(3k\rho _sϵ)`$. This essentially amounts to setting the time scale for the evolution of $`\varphi `$ (which must be related to the width of the transition region, because it is only in this region where $`\varphi `$ has an appreciable dynamics).
We arrive at
$`{\displaystyle \frac{\varphi }{t}}`$ $`=`$ $`{\displaystyle \frac{\gamma }{k\rho _s}}[^2\varphi {\displaystyle \frac{1}{ϵ^2}}(2g^{}(\varphi )+{\displaystyle \frac{ϵ}{3\gamma }}h^{}(\varphi )\{\mu (u_{ij}u_{ij,s}^{(\mathrm{eq})})(u_{ij}u_{ij,s}^{(\mathrm{eq})})+{\displaystyle \frac{\lambda }{2}}(u_{ii}u_{ii,s}^{(\mathrm{eq})})^2`$ (32)
$`{\displaystyle \frac{\stackrel{~}{\lambda }}{2}}(u_{ii}u_{ii,\mathrm{}}^{(\mathrm{eq})})^2+(zz_0)\mathrm{\Delta }\rho g{\displaystyle \frac{2\mu +\lambda }{8\mu (\mu +\lambda )}}\sigma _{00}^2\})]`$
Herein, $`g^{}(\varphi )`$ and $`h^{}(\varphi )`$ are the derivatives of $`h(\varphi )`$ and $`g(\varphi )`$ with respect to their argument. As we have mentioned before, $`h^{}(\varphi )`$ vanishes in the solid as well as in the liquid phases \[see eq. (A27)\].
In writing down an equation for the evolution of the elastic variables, we have to be careful about the fact that the strains $`u_{ij}`$, $`i,j=1,\mathrm{},d`$ are not independent quantities. Therefore, the variational derivatives $`\delta F/\delta u_{ij}`$ are not independent. Instead of introducing Lagrangian multipliers, we can however exploit the fact that the components $`u_i`$, $`i=1,\mathrm{},d`$ of the displacement, related with the strains via Eq. (43), are independent variables. Assuming that the time scales of our problem are large in comparison with sound propagation times, we conclude that the variational derivatives $`\delta F/\delta u_i`$ are equal to zero. This is is an adiabaticity assumption. Hence, we obtain
$`0={\displaystyle \frac{\delta F}{\delta u_i}}`$ $`=`$ $`{\displaystyle \frac{}{x_j}}{\displaystyle \frac{\delta F}{\delta u_{ij}}}`$ (33)
$`=`$ $`{\displaystyle \frac{}{x_j}}\left\{h(\varphi )\left(\sigma _{ij}\sigma _{ij}^{(\mathrm{eq})}\right)[1h(\varphi )]\left(pp^{(\mathrm{eq})}\right)\delta _{ij}\right\}.`$ (34)
This is nothing but a generalized elasticity problem, with the generalized stress tensor given by the quantity in braces.
Before moving to a demonstration of the sharp-interface limit, let us discuss scales. Since the elastic problem (34) is formally linear in the strains, rendering it nondimensional is straightforward and unenlightening. On the other hand, trying to cast (32) into nondimensional form, we realize that besides the length and time scales discussed in subsection II A, we need a third length scale $`\mathrm{}_3=\gamma /K`$, apart from the width of the transition region $`ϵ`$. So the phase-field model contains four length scales altogether.
Normalizing elastic moduli and stresses by the bulk modulus, i.e., setting $`M=\mu /K`$, $`\mathrm{\Lambda }=\lambda /K`$, $`\stackrel{~}{\mathrm{\Lambda }}=\stackrel{~}{\lambda }/K`$, $`\mathrm{\Sigma }_{00}=\sigma _{00}/K`$, we obtain
$`{\displaystyle \frac{\varphi }{\stackrel{~}{t}}}`$ $`=`$ $`\stackrel{~}{}^2\varphi {\displaystyle \frac{\mathrm{}_1^2}{ϵ^2}}(2g^{}(\varphi )+{\displaystyle \frac{ϵ}{3\mathrm{}_3}}h^{}(\varphi )\{M(u_{ij}u_{ij,s}^{(\mathrm{eq})})(u_{ij}u_{ij,s}^{(\mathrm{eq})})+{\displaystyle \frac{\mathrm{\Lambda }}{2}}(u_{ii}u_{ii,s}^{(\mathrm{eq})})^2`$ (36)
$`{\displaystyle \frac{\stackrel{~}{\mathrm{\Lambda }}}{2}}(u_{ii}u_{ii,\mathrm{}}^{(\mathrm{eq})})^2+{\displaystyle \frac{2\mathrm{}_3}{\mathrm{}_2}}(\stackrel{~}{z}\stackrel{~}{z}_0){\displaystyle \frac{2M+\mathrm{\Lambda }}{8M(M+\mathrm{\Lambda })}}\mathrm{\Sigma }_{00}^2\}),`$
where $`\stackrel{~}{t}=t/\tau `$ is the nondimensional time, and $`\stackrel{~}{z}=z/\mathrm{}_1`$, $`\stackrel{~}{}=\mathrm{}_1`$ are nondimensionalized spatial operators. Physically, $`\mathrm{}_3`$ represents an atomic scale. For many materials, $`\gamma /K`$ is on the order of the lattice constant.
For the phase-field model to work properly, we must impose some conditions on the length scale $`ϵ`$. We definitely need $`ϵ/\mathrm{}_11`$ to have a decently sharp interface. Moreover, the $`h^{}(\varphi )`$ term must not become too large in comparison with the $`g^{}(\varphi )`$ one, otherwise the minima of the double well move away from the positions $`\varphi =0`$ and $`\varphi =1`$. This appears to suggest that we also need $`ϵ/\mathrm{}_31`$. We compute some typical values. Using the material parameters of solid He , the system for which the Grinfeld instability has been unambiguously demonstrated by Torii and Balibar , we obtain the estimates $`\mathrm{}_10.1`$ cm, $`\mathrm{}_20.1`$ cm, $`\mathrm{}_310^9`$ cm, and $`\tau 1`$ s. If we had to require $`ϵ\mathrm{}_3`$, we would have a problem with very disparate length scales, as our numerical grid would have to be smaller than $`\mathrm{}_3`$, whereas the length scales governing pattern formation are $`\mathrm{}_1`$ and $`\mathrm{}_2`$. Fortunately, the quantity $`ϵ/3\mathrm{}_3`$ appearing in (36) is multiplied by squared strains, and the $`u_{ij}`$ are on the order of $`10^4`$. Moreover, we have $`2\mathrm{}_3/\mathrm{}_22\times 10^8`$ and the last term in braces can be estimated by $`\frac{1}{4}\mathrm{\Sigma }_{00}^22\times 10^8`$. Therefore, the actual condition for our model to be useful is $`10^8\times ϵ3\mathrm{}_3`$, which is much easier to achieve. In our simulations, we typically had $`10^8\times ϵ/3\mathrm{}_30.1`$.
Equations (32)-(34) constitute the basic phase-field equations for the phase transformation under stress.
To specify our model completely, we have to indicate the equilibrium stresses and strains. Let us assume the following forms for the stress-strain relationships in the two phases,
$`\sigma _{ij}`$ $`=`$ $`p_{0s}\delta _{ij}+2\mu u_{ij}+\lambda u_{kk}\delta _{ij},`$ (37)
$`p`$ $`=`$ $`p_0\mathrm{}\stackrel{~}{\lambda }u_{kk},`$ (38)
and require the equilibrium pressure to be $`p_0`$. For a planar interface, this fixes the normal stress in $`z`$ direction to be $`\sigma _{zz}^{(\mathrm{eq})}=p_0`$. If we assume the equilibrium stress tensor to be isotropic (a very natural assumption in most cases), we have $`\sigma _{ij}^{(\mathrm{eq})}=p_0\delta _{ij}`$, and $`u_{ij,s}^{(\mathrm{eq})}`$ in the solid is given by (15). In the liquid, we have $`u_{ii,\mathrm{}}^{(\mathrm{eq})}=(p_0\mathrm{}p_0)/\stackrel{~}{\lambda }`$. Note that only the displacement divergence $`𝐮=u_{ii}`$ appears in the elastic energy of the liquid. This gives us a degree of freedom (neither $`u_{xx}`$ nor $`u_{zz}`$ are fixed separately in the liquid, only their sum is) that will turn out important later. (Without this degree of freedom, it would not be feasible to treat the liquid as a shear-free solid, as will be discussed in appendix A.)
Inserting these equilibrium values into (32) and (34), we obtain as basic equation of motion for the phase field (introducing the abbreviation $`\stackrel{~}{k}=k\rho _s`$)
$$\frac{\varphi }{t}=\frac{\gamma }{\stackrel{~}{k}}\left\{^2\varphi \frac{1}{ϵ^2}\left[2g^{}(\varphi )+\frac{ϵ}{3\gamma }h^{}(\varphi )\left(\mu u_{ij}u_{ij}+\frac{\lambda \stackrel{~}{\lambda }}{2}u_{ii}^2+\mathrm{\Delta }pu_{ii}+\mathrm{\Delta }W+\mathrm{\Delta }\rho g(zz_0)\right)\right]\right\}$$
(39)
where we have defined further abbreviations
$$\mathrm{\Delta }p=p_0\mathrm{}p_{0s}$$
(40)
and
$$\mathrm{\Delta }W=\frac{1}{2}d\frac{(p_0p_{0s})^2}{2\mu +\lambda d}\frac{1}{2}\frac{(p_0p_0\mathrm{})^2}{\stackrel{~}{\lambda }}\frac{2\mu +\lambda }{8\mu (\mu +\lambda )}\sigma _{00}^2.$$
(41)
The elastic problem can be cast into the suggestive form
$$0=\frac{}{x_j}\left\{h(\varphi )\sigma _{ij}[1h(\varphi )]p\delta _{ij}\right\},$$
(42)
from which it is even more transparent that the expression in braces is nothing but a generalized stress tensor. Note that the phase-field model always guarantees exact mechanical equilibrium with respect to this stress tensor, but that the validity of a linear relationship between strains and generalized stresses is only warranted outside the interface region, where the values of $`\varphi `$ cease to depend on the $`u_{ij}`$. (This means that in the vicinity of sharp groove tips we will automatically have deviations from Hooke’s law, albeit they are not modeled to satisfy a particular nonlinear constitutive relation.)
These equations are to be solved subject to the conditions that the phase field approaches its limiting values in the bulk phases. To make them closed equations, we have to replace $`\sigma _{ij}`$ and $`u_{ij}`$ by the field variables $`u_i`$ using the definition of the strain tensor,
$$u_{ij}=\frac{1}{2}\left(\frac{u_i}{x_j}+\frac{u_j}{x_i}\right),$$
(43)
and Hooke’s law.
It remains now to be shown that this model reproduces the sharp-interface limit when the width of the interface is small. This calculation is given in the appendix. Its central result is formula (A57), which we rewrite here in nondimensional form (for $`\sigma _{00}=0`$):
$$\stackrel{~}{v}_n=\left\{\frac{1}{2}\frac{(\sigma _{tt}\sigma _{nn})^2}{\sigma _0^2}+\stackrel{~}{\kappa }+\frac{\mathrm{}_1}{2\mathrm{}_2}(\stackrel{~}{\zeta }\stackrel{~}{\zeta }_0)\right\}.$$
(44)
The ansatz proposed in is slightly different. One difference mainly concerns the interpretation or philosophy of the approach. In the MG model, the phase field is considered the variable determining the shear modulus. The shear modulus is the macroscopic quantity deciding whether a piece of condensed matter is solid or fluid. Hence, the phase-field order parameter differentiates between liquid and fluid and has a transparent meaning in the context of liquid-solid transitions. Of course, the model can be extended easily to the case of two solids with nonvanishing shear moduli on both sides of the interface. In the KM approach, the traditional and more conventional view is taken that the phase field decides between two phases characterized by their respective free energy densities. That one of these phases is a liquid is of secondary importance, as it were. Again, in principle, it might be another solid. Of course, if the second phase is chosen a liquid, then its shear modulus must vanish. And indeed, this is guaranteed in the current form of both models by construction. For ease of further comparison, we give the phase-field equations of in appendix B and show how they are mapped onto the form (39,42).
In concluding this section, we would like to comment briefly on the consequences of an anisotropic equilibrium strain. Suppose we submit a body consisting of piezoelectric material to a homogeneous electric field. (Alternatively, we could consider some magnetrostrictive material under the influence of a magnetic field.) This body will contract or expand until it reaches a new equilibrium state compatible with the body forces exerted by the field. The new state will have anisotropic strain and, assuming isotropic elastic properties, an anisotropic stress tensor as well. What will the surface dynamics of such a body be, if uniaxial stress is applied in addition, as in the setup of the Grinfeld instability? Of course, the assumption that the equilibrium stress remain constant is an oversimplification now, since the dielectric properties of the solid and its melt will usually differ, hence the electric field would become inhomogeneous as soon as an interfacial shape change occurs. Let us nevertheless assume the simplest possible situation, an anisotropic but constant equilibrium state
$$\sigma _{ij}^{(\mathrm{eq})}=p_0\delta _{ij}+\chi _0\delta _{ix}\delta _{jx}.$$
(45)
Using the stress-strain relationship (37), this can be inserted into our expression for the elastic energy density of the solid, which then becomes
$$f_{\mathrm{sol}}(\{u_{ij}\})=\mu u_{ij}u_{ij}+\frac{\lambda }{2}u_{ii}^2+\mathrm{\Delta }pu_{ii}+\frac{1}{2}\frac{\mathrm{\Delta }p^2}{\mu +\lambda }\chi _0u_{11}\chi _0\frac{\mathrm{\Delta }p}{2(\mu +\lambda )}+\chi _0^2\frac{2\mu +\lambda }{8\mu (\mu +\lambda )},$$
(46)
where we have set $`p_0\mathrm{}=p_0`$ for simplicity.
It is then straightforward to derive the sharp-interface limit for this modified phase-field model. The result reads (on setting $`\sigma _{00}=0`$)
$`v_n`$ $`=`$ $`{\displaystyle \frac{1}{k\rho _s}}\{\gamma \kappa +\mathrm{\Delta }\rho gz+{\displaystyle \frac{1\nu ^2}{2E}}(\sigma _{tt}\sigma _{nn}\chi _0)^2`$ (48)
$`+[{\displaystyle \frac{2(1\nu ^2)}{E}}\chi _0(\sigma _{tt}\sigma _{nn}){\displaystyle \frac{1+\nu }{E}}\chi _0^2]n_x^2+{\displaystyle \frac{2\nu (1+\nu )}{E}}\chi _0^2n_x^4\},`$
where $`n_x`$ is the component of the interface normal in $`x`$ direction. Rotational invariance is broken.
We are not aware of any previous mention of this equation in the literature, nor do we think this interesting case has been treated. In fact, what we have demonstrated here, is how the phase-field model can be used to derive hitherto unknown sharp-interface equations in a transparent way.
It is clear from (48) that an isotropic stress tensor, i.e., $`\sigma _{tt}=\sigma _{nn}`$ does not necessarily entail a stable planar interface, whereas setting $`\sigma _{tt}\sigma _{nn}=\chi _0`$, i.e., providing an anisotropic stress tensor, we will have a linearly stable planar front solution with interface position $`z=0`$. This is easily seen from the fact that the terms containing $`n_x^2`$ and $`n_x^4`$ do not contribute in a linear stability analysis, because $`n_x`$ is directly proportional to the perturbation and hence its square and fourth power have to be dropped. Note also that the symmetry of the dynamics with respect to a replacement of $`\sigma _{tt}\sigma _{nn}`$ by its negative value does not hold anymore in this situation.
While this equation opens a new line of research, we will refrain here from pursuing this topic any further.
## III Numerical results
### A Validation of the model
In order to verify that our phase-field description leads to a quantitatively correct description of the instability, at least before cusps set in, we have performed a number of numerical tests. Based on a simple finite-difference scheme, the numerical implementation is set up in a rectangular geometry.
The bottom half of the rectangle is filled with solid, the top with liquid. This is realized by setting the phase field $`\varphi `$ equal to a tanh-like function taking the value one in the bottom region and zero in the top region of the geometry. $`\varphi `$ is kept at these values one and zero exactly at the bottom and top lines of the numerical grid, respectively. Periodic boundary conditions are applied at the lateral boundaries. The initial interface is set by an appropriate modulation of the region where $`\varphi `$ crosses the value $`\frac{1}{2}`$ and was in most cases taken to be sinusoidal or flat with a random perturbation.
The boundary and initial conditions for the fields $`u_x`$ and $`u_z`$ are chosen differently for the KM and MG models as will be described now.
Within the KM model, where we assume strains to vanish at equilibrium (hence $`\mathrm{\Delta }p=0`$), we took the $`x`$ derivatives of both displacement fields periodic in the $`x`$ direction in our initial simulations, in order to obtain periodic strains. Later, we switched to simpler helical boundary conditions for $`u_x`$, i.e., we took $`u_x(L,z)=u_x(0,z)+Lu_{xx,0}`$, where $`L`$ is the length of the rectangle along the $`x`$ direction, and periodic boundary conditions for $`u_z`$. This change in boundary conditions did not affect results in any essential way. All the simulations of the KM model discussed here were carried out with these boundary conditions (whereas those in were done with periodic $`x`$ derivatives). At the bottom of the system, the values of the fields are fixed to values corresponding to a homogeneously strained solid; at the top, $`u_x`$ is fixed and the derivative $`_zu_z`$ chosen such that the condition $`𝐮=0`$ is satisfied. $`u_x`$ is initialized as a linear function $`u_x=xu_{xx,0}`$ and the inital $`u_z`$ is determined via integration of Eq. (A41).
For simulations of the MG model (or rather its variant considered here), the fields were all taken periodic in the $`x`$ direction, whereas the boundary conditions in the $`z`$ direction were as in the KM model. We did not yet attempt to use spectral methods for the solution which would require periodicity in the $`z`$ direction as well (achievable by simply reflecting the system at its bottom, and including the image into the numerical box ). Initialization was done by setting $`u_x=0`$ everywhere and computing $`u_z`$ from (A41) again.
The elastic equations were solved by successive overrelaxation, the time integration was performed by a formally second-order accurate midpoint scheme. Since we did not update the elastic fields at the half time step, the formal accuracy was not attained. The most time-consuming part of the simulation was the relaxation scheme and a way to overcome its restrictions has been given in as is discussed in appendix B. Since it requires an approximation to the solution of the elastic problem even at the analytic level ($`\mu /K`$ has to be small), we did not implement it in our two-dimensional simulations. We intend to compare the quality of this approximation to the solution of the full problem, before employing it in a 3D simulation, where its use is essential for reasons of computational efficiency. Most of our computations were done with the material parameters of Helium to facilitate comparison with experiments Therefore, whenever we do not indicate different choices, our parameters were chosen as described in . Times and lengths given without units are in seconds and centimetres, respectively. Since our nondimensional time unit is about one second and the nondimensional length scale about 0.1 cm, this simply corresponds to using 10$`\mathrm{}_1`$ instead of $`\mathrm{}_1`$ as the basic length scale.
One of our numerical tests consisted in reproducing the instability threshold to within 2% accuracy, another one in verifying the subcritical nature of the bifurcation, first demonstrated analytically by Nozières . A short discussion of the last feature has been given in , so we will not elaborate on it here . We consider a few more tests, however.
Figure 1 gives the dispersion relation determined for three values of the external stress and compares it with the analytic result from linear stability theory. The KM model was used here as it gave more accurate results at finite $`ϵ`$.
To obtain the dispersion relation, we simply followed the dynamics of a system initialized with a small-amplitude cosine profile for a number of different wavelengths, and computed the amplitudes of the evolving structure for a series of times. Then the amplitudes were fitted to an exponential function which provided the growth rate of the interface. Taking $`\sigma _{00}`$ equal to the applied stress $`\sigma _0`$, we fixed the average position of the interface. Moreover, the amplitudes were computed in two different ways, both of which are not influenced by the average interface position. The first method was simply to take the square root of the spatial variance of $`\zeta (x)`$; as a second measure for the amplitude we took the modulus of the Fourier component corresponding to the wavelength chosen. On the figure, these two methods give essentially indistinguishable results within the size of the symbols. System sizes used were the wavelength $`\lambda _f`$ of the fastest-growing mode and a number of rational multiples and fractions thereof (ranging from $`\frac{1}{4}\lambda _f`$ to $`3\lambda _f`$). Since we kept the number of numerical grid points the same for all the systems at $`\lambda _f`$, the mesh size had to be varied. The interface thickness $`ϵ`$ was in general kept above $`\frac{3}{2}`$ of the mesh size, which gives a resolution of 5 points for the region where the phase field varies between 10% and 90% of its maximum value. For smaller values of $`ϵ`$, locking effects to be discussed shortly became conspicuous .
The agreement between analytic results and numerically determined points is satisfactory both above ($`\sigma _0=2.8\times 10^4`$ dynes/cm<sup>2</sup>) and below ($`\sigma _0=2.4\times 10^4`$ dynes/cm<sup>2</sup>) the instability threshold. Two points are worth mentioning. First at $`q30`$ cm<sup>-1</sup>, there are two symbols each for the growth rates corresponding to the two larger stresses. These were given to roughly indicate the possible error in the numerical result when the growth rate has a large negative value. Points below $`q20`$ cm<sup>-1</sup> did not show a comparable error. The two different values were obtained by fitting with the initial and the final half of the data points, respectively. We ascribe the difference to the fact that the amplitude of the interface becomes smaller than its width $`ϵ`$, a situation in which the phase-field description is no longer reliable. For example, the final planar interface is not located exactly at $`z=0`$, about which the initial cosine was centered, albeit the deviation is smaller than the interface width. Second, the overall agreement is surprisingly good in view of the fact that a phase-field model is not particularly well-suited to the determination of a dispersion relation at all. For in order to approach the limit of an infinitesimal perturbation of a planar interface one should choose very small amplitudes, but they must not be smaller than the interface width $`ϵ`$. Reduction of $`ϵ`$ is possible in principle but soon leads to prohibitive computation times. With a sharp-interface model that we investigated in parallel , it was no problem to take amplitudes of $`10^4`$ and to obtain nice exponential growth or decay during long time intervals, whereas here we were restricted to starting amplitudes on the order of 0.05 or larger.
Our next test consists in investigating the dynamics of a planar interface with both the KM and MG models. From (1, 2) we obtain the equation of motion
$$\dot{\zeta }=\frac{1}{k\rho _s}\left(\frac{(1\nu ^2)}{2E}\sigma _0^2+\mathrm{\Delta }\rho g\zeta \right),$$
(49)
which is, given the initial condition $`\zeta (0)=0`$, solved by
$$\zeta (t)=\left(1e^{\mathrm{\Delta }\rho gt/k\rho _s}\right)\frac{(1\nu ^2)}{2\mathrm{\Delta }\rho gE}\sigma _0^2.$$
(50)
This analytic result is compared with simulations of the two models in Fig. 2.
What is cleared up by the figure is that even with a well-resolved interface width (we have $`ϵ=4h`$) the MG model is slightly off the analytic final position, whereas the KM one converges well towards it. With larger values of the numerical mesh size, convergence of the former model gets even worse. For $`h=0.007`$, $`ϵ=0.011`$ the KM model still agrees reasonably well with the analytic curve while the MG one is off by about 10% for $`t=4`$. Both models show deviations from exponential behavior with this set of parameters due to metastability effects of the discrete set of interface points. This problem, which is particularly critical for interface pieces parallel to one of the coordinate axes, has been discussed in detail in . At small interface velocities, the sum of the energies of the discrete points of the phase field in the double well potential may vary at successive time steps (whereas the energy of a continuous field is degenerate under arbitrary translations). Therefore, the interface is slowed down, if the energy of its discretization increases due to motion and accelerated if it decreases. For sufficiently small driving force, the interface may stop moving at all, i.e., lock into some favorable position. Apparently, the MG model is more susceptible to these effects than the KM one.
The ultimate reason for the different behavior of the two models is that they are only asymptotically equivalent, i.e., they describe the same system only in the limit $`ϵ0^+`$. For any finite $`ϵ`$, the equations obeyed by the phase field and the displacements are not the same in the two models. One can observe this directly by comparing the different terms contributing to, e.g., $`_t\varphi `$. In the MG model, the term $`\mathrm{\Delta }pu_{ii}`$ of eq. (39) is frequently the largest interface term affecting $`_t\varphi `$, whereas this term is equal to zero in the KM model. Moreover, the sum of all terms multiplying $`h^{}(\varphi )`$ is not the same in both equations.
The difference can also be seen in comparing a numerical simulation of the the sharp-interface model (1,2) itself, using an integral equation approach, with the phase-field models. We will report on details of this alternative approach elsewhere . Figure 3 shows the interface evolving in the phase-field calculation for two different values of $`ϵ`$ and compares them with the sharp-interface result starting from the same initial condition, after the same time interval.
Again the KM model fares slightly better in the comparison for the same value of $`ϵ`$. In the groove, however, both models deviate from the sharp-inteface result but approach it more closely for the smaller interface thickness. The sharp-interface model produces a more strongly pointed groove, as expected. It should be emphasized that this simulation is not far from the limit of the temporal validity range of the sharp-interface model. This limit is signalled either by a crash of the program due to the singular behavior of the bottom of the groove or by the appearance of a spurious steady state, which can be achieved by overstabilization of the interface.
Why the KM model agrees better with the sharp-interface model for a given value of the interface width is a difficult question, to which we cannot offer any deep answer. Also we cannot exclude that for a different choice of the functions $`h(\varphi )`$ and $`g(\varphi )`$, the MG model would be superior. It should be kept in mind that the functions employed in are not the same as those used here (see App. B).
Normally, we would thus prefer to use the KM model for calculations to be presented. However, since we wish to make sure that effects of translational symmetry breaking are not due to our using a model in which periodicity is imperfectly implemented, we will use the MG variant in the following simulations. The differences between the two models are small, after all. Also the MG model has the advantage to be more easily treated using pseudospectral methods based on Fourier series due to its periodic boundary conditions, with a gain in accuracy that might allow to offset its apparent disadvantage.
The conclusion from Fig. 3 is that the phase-field models give decent agreement with a sharp-interface calculation in regions where the curvature is not too large. Whereas the sharp-interface computation cannot be meaningfully continued by very much beyond the time shown in the figure ($`t=0.25`$), the phase-field models both have no problem in continuing the simulation to times well beyond $`t=1`$.
As anticipated above, we take the point of view that a real solid cannot develop exact cusps, because plastic effects such as the generation of dislocations will intervene. These will relieve stresses and thus prevent infinite densities of the elastic energy. The phase-field model does the same thing and we shall see below that it does so by introducing a cutoff to the curvature. More quantitative modeling would require to explicitly take into account models of nonlinear elasticity or plasticity, which is beyond the scope of this article. Nevertheless, as we can see from Fig. 3, the behavior far from the sharp tip of the groove is described reasonably well by the phase-field model for both values of $`ϵ`$ and is almost independent of the interface width. Therefore, we believe that the phase-field approach correctly reproduces the qualitative behavior of a situation in which plastic effects occur only in the minima of the grooves.
Results obtained under this hypothesis will be discussed in the next section.
### B Dynamics of extended systems
When simulating periodic structures, one realizes that for small supercritical stresses, where the system takes a long time to develop deep grooves, one often observes symmetry breaking and one of the grooves getting ahead of the others. This symmetry breaking must be triggered by numerical noise from roundoff or truncation errors. For high stresses, where the system develops grooves reaching the system bottom within a relatively short lapse of time, this does not happen. Figure 4 gives an example of a structure grown at about three times the critical stress. The interface is plotted at constant time increments ($`\mathrm{\Delta }t=0.05`$). A shift in the chemical potential of the solid has been made to keep the position of a planar front fixed.
We see that the structure remains periodic in the time interval considered and that three equally deep grooves evolve. Note the peculiar shape of the cells. From flat tips there emerge slightly curved slopes on the side of the cells. Then there is a sharp bend downward into the deep groove. The appearance of this bend renders it plausible that the time of formation of a cusp in the sharp-interface description has already passed and from then on the dynamics should be governed by the curvature bound. In the final stage of this dynamics all grooves move at constant velocity. Figure 5 gives the curvatures of the interfaces displayed in Fig. 4 and demonstrates that the radius of curvature at the bottom of the grooves remains constant and is close to $`ϵ`$, which was equal to 0.02 in this simulation.
The curvature was calculated from the contour line defining the interface position. Since the representation of this line ($`\varphi =0.5`$) was constructed by determining its intersection points with the squares of the numerical grid, the discretization points were unevenly spaced (two intersections with grid lines parallel to the $`x`$ and $`y`$ axes can be arbitrarily close to each other, the next may be as far away as $`\sqrt{2}h`$). Therefore, our curvature results are pretty noisy, even after application of a smoothing procedure. A superior method for their determination would be to use the full representation or the phase field instead of just the contour line information. Nevertheless, they clearly indicate the approximate constancy of the curvature in the groove tips.
Since we have the stresses at our disposal, too, we can calculate the final velocity of the grooves. Figure 6 gives a contour plot of the stresses $`\sigma _{xx}`$, corresponding to the final time of the simulation from Fig. 4, $`t=0.45`$. The interface is drawn as as solid line, the contour lines are broken lines in different styles. What we have plotted here, is not a generalized stress tensor component, as defined by (A5), but simply the stress in the solid. Therefore, the contour lines for stresses far in the liquid are meaningless (in the dynamic equations, they are multiplied by $`h(\varphi )0`$), although they become important when entering the interface region. From the figure, we estimate a maximum value of $`\sigma _{tt}2\times 10^5`$ in the bottom of the groove (and a similar value is obtained from the corresponding figure for $`\sigma _{zz}`$). Inserting this, the value of the curvature and the position $`\zeta `$ of the groove bottom into (1), we obtain for the velocity $`v_n=1.6`$. Assuming that the interface grew at this velocity from the outset, we obtain for its final position the value $`\zeta =0.7`$. The data show that it is actually at $`\zeta =0.72`$, which is easily explicable by the inaccuracy of our estimate of the maximum stress. From the contour plot we do not obtain more than a rough figure as stresses vary rapidly in the interface region.
In a straight and narrow crack, the stress scales with the square root of the distance from its tip . Therefore, a reduction of the tip radius by a factor of two will increase both the stress term and the curvature term of (1) by a factor of two as well. As long as the gravity term in (1) is negligible (which, incidentally, it is not in the simulation of Fig. 4, its contribution is about as large as that of the curvature for the last curve), this means that the velocity of the groove will roughly double when $`ϵ`$ is halved. This trend has been confirmed in the simulations, although the observed ratio is slightly smaller than the predicted one, but then our grooves do not yet really have an extremely small width compared with their length.
The next three figures show a simulation at a stress roughly 20% above the critical value. Our numerical box contains six wavelengths of the pattern initially. One of the grooves has however been made by 2% deeper than its neighbours. Contrary to the situation in Fig. 4, no prestress was applied, so a planar interface would move downward to a new equilibrium position. This kind of motion is superimposed on the shape-changing dynamics and serves nicely in separating the curves on the plot.
The temporal dynamics can be divided into several stages. At first, the sinusoidal pattern changes its shape in the way already discussed by Nozières : the tips become flat, the grooves pointed. After some time, the interface becomes similar to a cycloid but with different depths of the grooves. Also, the dynamics almost comes to a halt. Below, we shall discuss the similarity with a cycloid in more detail (see Fig. 10). It holds up to $`t=2.5`$ approximately, which is the time of the lowest curve in Fig. 7. At this point the apparent periodicity of the pattern has doubled. (Of course, strictly speaking this periodicity has been broken from the outset by our making one groove a little deeper. But this was only to avoid its being broken by numerical noise in an uncontrolled manner, i.e., to introduce a well-defined perturbation.)
The groove that was ahead initially, wins the competition for the elastic field; the losing grooves fall back and even close again. This is shown in Fig. 8, displaying the temporal continuation of Fig. 7. In the initial structure of Fig. 8 (the solid line that is shallowest in the big grooves), the smaller grooves are deeper than in the final one (the dashed line which is deepest in the big grooves but shallowest in the small ones). At the end of the period of time depicted in Fig. 8, there are three clear survivors and three losers of the competition.
Finally, as shown in Fig. 9, only one groove survives. Its velocity is almost constant over a range of times. Eventually it slows down and grows sideways towards the end, which may have to do with the fact that it gets too close to the bottom of the numerical box (which is at $`\zeta =1`$). Also gravity has a decisive decelerating effect here.
What we observe, then, is a coarsening process that seems to proceed via imperfect period doubling transitions. Because our system has only six grooves, we cannot explicitly see more than the first period doubling here. These transitions are local in the following sense. Not all grooves surviving the first period doubling get ahead of the others simultaneously. Rather what happens is that first the winning groove gets ahead of its nearest neighbours, screening each of them off the stress field on one side a little. This causes these neighbours to grow more slowly, making them screen off their next neighbours on the other side less. So these get ahead of their neighbour grooves, and so on. The perturbation made by one groove moves through the array in an alternating fashion. In an infinite system, one could imagine a series of “near” period doublings propagating through the array. These morphology changes are not exact period doubling transitions, because there is no restabilization of a structure with doubled periodicity. The system remains dynamic (but see the discussion on gravity below) which means that the foremost groove does not get slower than its competitors, which would be necessary for length adjustment.
The first of these period doublings may be discussed analytically in some detail. Consider the shape of the interface close to the last time of Fig. 7. It can be modeled approximately by a curve that we would like to call a “double cycloid”. A parameter representation of this curve is given by
$`x`$ $`=`$ $`\xi A\mathrm{sin}k\xi B\mathrm{sin}2k\xi `$ (51)
$`z`$ $`=`$ $`A\mathrm{cos}k\xi B\mathrm{cos}2k\xi `$ (52)
Figure 10 compares a double cycloid with the interface at $`t=2`$. The wavenumber $`2k`$ (= 9.425) is given by the basic periodicity of the initial interface (before it is perturbed), the amplitudes $`A`$ and $`B`$ have been fitted “by eye” and the double cycloid has been shifted using translational invariance in the $`x`$ direction. (Its position in the $`z`$ direction can also be adjusted, which corresponds to a particular choice of the initial chemical potential of the solid.)
Since we made only one of the grooves deeper than the others, the agreement of the groove minima is not quite perfect, as we can adjust only the depths of this groove and its nearest neighbours by an appropriate choice of the two constants $`A`$ and $`B`$. Had we taken an initial perturbation of periodicity length $`2\pi /k`$ instead of a local one in the simulation, a much better agreement would have been obtained. The purpose of this comparison, however, is not to claim that the interface shape goes precisely to a double cycloid but only to show that it may be well approximated by such a curve, which can be considered a cycloid (with amplitude $`B`$) modified by a small perturbation of twice its wavelength. In our fit shown in Fig. 10, we have $`AB/10`$.
Our key observation is then that we can solve the sharp-interface elastic problem for a double cycloid exactly in an extension of the work of Gao et al. , using a conformal mapping technique. This solution is given in appendix C. In what follows, we will neglect the gravity term, a procedure that we justify later. The evaluation of the nondimensional velocity via Eq. (44) for the double cycloid yields, in the bottoms of the grooves \[see appendix, Eq. (C41)\]
$$\stackrel{~}{v}_n=\frac{1}{2\left[1k\left(A(1)^m+2B\right)\right]^2}\left\{\left(1+2Bk+\frac{1+Bk}{1Bk}Ak(1)^m\right)^2\alpha k\left[A(1)^m+4B\right]\right\},$$
(53)
where $`\alpha =2k/q_f`$ is the ratio of the actual wavenumber of the basic cycloid and the wavenumber of the fastest-growing mode. The formula with odd $`m`$ holds for the minima with depth $`2Bk+Ak`$, that with even $`m`$ for those with depth $`2BkAk`$. A condition for the solution to hold is that there are no self-crossings of the curve, therefore we must require $`Ak+2Bk<1`$. Let us now assume that $`AB`$, i.e. that the pattern actually is a slightly perturbed cycloid (where the perturbation has twice the basic wavelength $`\pi /k`$). Then the denominator in Eq. (53) in front of the braces goes to zero for even $`m`$ as $`2Bk`$ approaches the value 1. This is the finite-time singularity, already identified by Gao al. . The velocity goes to $`\mathrm{}`$, if the braces remain positive, which they do for small enough $`\alpha `$, i.e., when the wavelength is large enough. For small $`A`$, we can expand (53). This gives
$`\stackrel{~}{v}_n`$ $`=`$ $`{\displaystyle \frac{1}{2\left[1k\left(A(1)^m+2B\right)\right]^2}}\{(12Bk)^2+4(1\alpha )Bk`$ (55)
$`+Ak(1)^m[2{\displaystyle \frac{(1+Bk)(1+2Bk)}{1Bk}}\alpha ]\},`$
a formula that shows that the marginal value of $`\alpha `$ is 1. Thus, for wavelengths larger than that corresponding to the fastest-growing mode ($`\alpha 1`$), the velocity will diverge in the deepest minima, leading to cusps in the sharp-interface limit. We could leave the gravity term out of this consideration, because it never diverges for finite $`\zeta `$.
Now assuming we are at or slightly above the wavelength of the fastest-growing mode, we can see from (55) that for $`(12Bk)^2<Ak`$ the velocity is positive in the secondary minima corresponding to odd $`m`$ . This means resolidification and closure of the corresponding grooves.
Suppose for a moment that $`A=0`$. Then the system with a sharp interface will evolve towards a cusped cycloid, i.e., $`2Bk`$ will increase towards 1. But this means that eventually a point will be reached where $`12Bk`$ is small enough that any perturbation will be larger than $`(12Bk)^2`$. In this case, our equations state that (for $`1\alpha 1`$) the tip perturbed in this way in the right direction (i.e., the perturbation must reduce the depth of the groove) will recede again, its velocity will become positive. A groove tip that is perturbed in the other direction will approach the cusp singularity even faster and reduce the speed of its neighbours. Of course, not all perturbations are periodic; what happens when only a local perturbation is applied, can be seen from the simulation.
What the analytic calculation shows, then, is that the first period-doubling bifurcation happens before the cusp singularity is reached, if the periodicity of the system is equal to the wavelength of the fastest-growing mode. Whereas the bifurcation to a set of alternatingly receding and advancing grooves may happen for any wavelength larger than this one, whether it happens before or after the predictable time of cusp formation will in general depend on the strength of the perturbations present in the system. In the simulation of Figs. 7-9, the periodicity of the unperturbed system is $`\frac{2}{3}`$, the wavelength of the fastest-growing mode is 0.5.
Ordinarily, the time when the finite-time singularity appears in the sharp-interface system will be too short for the losing grooves to have appreciably retracted. In our phase-field model, there are no finite-time singularities, so the evolution can continue. It is then highly plausible that further period doublings occur, even though we have no analytic model for these. But on general grounds, we expect screening of neighbouring grooves to become more effective as all grooves get deeper. Hence the process should repeat, even at wavelengths far from, but above, $`\lambda _f`$.
The difference between the cases of a wavelength close to that of the fastest-growing mode and one far above it is that in the former case, the first period doubling will happen before the time $`t^{}`$, at which cusps form in the sharp-interface limit, whereas in the latter case, it will happen afterwards. This case is, in fact, realized in Fig. 4, where the wavelength of the fastest-growing mode is about one tenth of the periodicity length. From Fig. 6, we can infer that the translational symmetry with respect to the basic wavelength $`\lambda =\frac{2}{3}`$ has already been broken by numerical noise (the stress pattern does not show this periodicity in the upper half of the picture, this symmetry breaking will slowly propagate into the lower half where everything still appears periodic).
Another interesting conclusion from formula (53) is that for $`\alpha >1`$, i.e., for systems with small enough wavelength, stable steady states may be possible, because then surface tension may succeed in overwhelming the effects of stress. For $`\alpha >1`$ and $`A=0`$, the formula predicts that a cycloid becomes stationary in its minima before the appearance of cusps. We hope to report on this aspect in the future.
Finally, let us have a look at a system with a random initial condition. Figure 11 shows the evolution starting from an interface resulting from uniformly random perturbations of a planar front. We see that first some 10 waves develop, which is already a coarsened structure, as the wavenumber of the fastest-growing linear mode would correspond to about 24 waves fitting into the system. However, the initial amplitude is too small for this wavelength to become clearly visible. Some time later, there are much fewer features and eventually, only two grooves remain.
Whether one of the two will die off in the end is not clear, since this is a simulation with gravity. Hence the largest groove is bound to stop at some time, because the stress and curvature terms remain constant once all other grooves are sufficiently small, but the gravity term continues to increase. If the second-largest groove still has a positive velocity when the first stops, it will not reverse its growth direction, but only grow to a point where its velocity becomes zero.
An example, where the final state actually consists of two grooves, is shown in Fig. 12. Here the applied stress is smaller than in Fig. 11, so the pattern actually does come to a halt within the numerical box, after a long time ($`t60`$). Note that during most of the period where two grooves are dominant, one of them is ahead. Once it stops, the second approaches and in the end it has the same length as the first, to numerical accuracy. In the case of two periodically repeated grooves, this is to be expected for symmetry reasons. With three or more grooves, it is also conceivable that not all of them are the same length in the steady state.
We think that in the absence of gravity, the situation in this strongly nonlinear region is very similar to the evolution of a Saffmann-Taylor finger in a Laplacian field. The Lamé equations determining the displacement field are scale invariant just as the Laplace equation (and in fact, Eq. (44) is scale invariant for $`\mathrm{}_2=\mathrm{}`$). Once a strongly nonlinear state has been reached, none of the length scales discussed in section II can play a role anymore, since they only govern the local behavior of the growth pattern. The long-range elastic field will determine the factor $`\sigma _{tt}\sigma _{nn}`$ of the destabilizing term in (44) and this factor will be the larger, the fewer competitors of a groove have grown to the same depth. This will lead to smaller grooves not growing anymore. This situation bears strong similarities to the growth of thermal cracks described in . The main difference is that there a loser in the competition will simply stop growing. In our case, it will even shrink again, for the crystal can not only melt but also freeze again, and whether it will do so is simply determined by the chemical potential difference (1).
An analogous behavior is found in the side branching activity of a dendrite in the region about 20 to 50 tip radii behind the tip . There coarsening is observed, too, which also proceeds via imperfect period doubling. If this dynamics can be described in terms of a series of nonequilibrium phase transitions at all, these would have to be considered first-order transitions because of the discussed locality aspect. There is no diverging length scale in a single transition.
We expect that the dynamics of large systems can be described by scaling laws similar to those given previously for the growth of needles in a Laplacian field . The fact that “needles” can shrink again in the elastic problem should modify the long-time behavior of the needle density, which must pass through a maximum and then go to zero as a function of time, for any needle length.
To some extent, this expectation is supported by the coarsening scenario described in in which extended systems without gravity are studied using random initial conditions. They measured the Fourier transform of the height-height correlation function $`S(q,t)`$ and observed dynamical scaling. For early times, they observed a strong similarity between this behavior and early-stage soinodal decomposition in long-range systems. For later times, when the linear theory no longer describes the data, coarsening is evident: The location of the peak $`q_{\mathrm{max}}`$ of $`S(q,t)`$ moves to smaller wavenumbers, as the peak height increases and sharpens. The peak height follows $`S(q_{\mathrm{max}},t)t^{\alpha +1}`$, where $`\alpha 2`$, while the peak width sharpens with time as $`wt^\gamma `$, where $`\gamma 0.5`$. The former dependence is due to the total interface length increasing linearly with time for any unstable wavenumber. The latter dependence is due to competitive ordering between different wavenumbers, analogous to phase ordering. Within the accuracy of their study, they find that the structure factor shows scale invariance: $`S(q,t)/S(q_{\mathrm{max}},t)=S^{}(q^{})`$, where the scaled wave number $`q^{}=(qq_{\mathrm{max}})/w`$. Fitting to $`S^{}(q^{})^\delta `$ and $`S^{}(1/q^{})^\psi `$, for small and large $`q^{}`$ respectively, gives $`\delta 12`$, and $`\psi 56`$.
It is however difficult to assess to which time regimes these results correspond when compared with the present simulations, because the freedom to rescale parameters has been used extensively in . Since the vanishing of grooves does not seem to be a dominant mechanism of coarsening in their simulations, it is likely that the time windows considered in and here have little overlap and that the stage of needle-like growth of the grooves is never reached in . Of course, the concept of short and long times is ambiguous in the absence of gravity due to the scale invariance of (44) for $`\mathrm{}_2=\mathrm{}`$. However, one can compare the depths of grooves with their lateral distances to decide whether growth is best described as competition of wavenumbers (a Fourier space concept) or as competition of needles (a real-space concept). In this context, it is also important to realize that the validity of the simulations in is restricted to small values of $`\mu /K`$ with the parameter $`\eta _0`$ (see app. B) being $`O(1)`$, hence these simulations are quantitative only for external stresses $`\sigma _0=O(\mu )`$. For these stresses, the state describable as a forest of needles is only obtained after a long time \[of order $`(K/\mu )^4`$\]. Therefore, the scaling exponents obtained in might not be relevant to the scenario discussed here, which makes an analytic treatment along the lines of even more desirable, because it could provide these exponents.
In an infinite system what we have depicted here is probably just the continuation of the coarsening scenario described in .
It should also be pointed out that for sufficiently wide systems, i.e., in particular for infinite ones, this dynamics may be an intermediate state only. Once a groove becomes sufficiently long, stresses along its side may become large enough to provoke a Grinfeld instability of the “side walls” of this crack-like structure, as has been shown by Brener and Marchenko . Whether or not this happens, depends on how efficiently the stresses are relaxed along the grooves, on the speed of the grooves, on the perturbation amplitude, etc. This secondary instability might completely change the scaling behavior, possibly leading to tip splitting of the grooves and tree-like structures. So far, we have not seen anything of this kind in our simulations.
In a system of finite lateral extent (or periodicity) we think coarsening will in the absence of gravity generally lead to the disappearance of all grooves except one which will grow at constant velocity. If gravity is present, several grooves can survive and in a sufficiently deep system they will stop once they reach a depth where the gravity term compensates the stress one.
## IV Conclusions
In this article, we have constructed a class of phase-field models from a free-energy functional including the elastic energy density. A salient feature of the model is that the liquid is treated as a shear-free solid, which is to be contrasted with phase-field models taking into account hydrodynamic effects in solidification, where the solid is usually treated as a liquid of infinite viscosity . Our approach implies the artificial introduction of coherence conditions at the interface which is however counterbalanced by the fact that the only relevant elastic variable in the liquid is $`𝐮`$. A whole class of models is obtained instead of a single one as a consequence of the freedom of choice for the state of reference used in measuring displacements. We compared the two most natural choices and found them to yield slightly different numerical results despite their asymptotic equivalence.
Investigating a large number of laterally small and extended systems, we believe to be able now to describe the generic dynamic behavior. For systems smaller than the wavelength of the fastest-growing mode of linear stability theory but larger than that of the marginal mode (where surface tension stabilizes the planar interface), stable steady-state strucures are possible even in the sharp-interface limit. This is similar to the findings by Spencer and Meiron for the case of transport via surface diffusion, even though we think the whole picture is more complex than what they described . Here, we did not show detailed results on small systems but focused on extended systems.
The case without gravity is particularly simple, as the equation of motion can then be made parameter free \[Eq. (44) without the last term\]. Initially, an interface may grow periodically but as soon as perturbations break the periodicity, coarsening will proceed via approximate period doubling transitions. Supposing randomized perturbations, the interface will, after a sufficient lapse of time, not look much different from one started with random initial conditions (compare Figs. 9 and 11). If the system is of finite lateral extent (but infinitely deep), only a single groove will survive growing at constant velocity, determined by the final constant stress and constant surface tension terms near its tip. The final velocity will scale with the system width $`L`$ and the radius of curvature $`ϵ`$ as $`v_nL/ϵ`$. All the other grooves will eventually retract, i.e., they will not even survive keeping a finite depth, which is different from the behavior of cracks . For wide systems, stresses near the groove tips may become large enough to trigger a secondary instability which would considerably modify the system behavior and allow the appearance of complex crack morphologies. It is however possible that this will arise only in the case of finite perturbations, as grooves may grow too fast in this situation for the instability to develop before it is “advected” (relative to the groove tip) into a region of very small stresses along the groove. If the system is laterally infinite, the system state will first follow dynamical scaling as studied in and should then cross over to the scaling dynamics described here with the number of grooves continually decreasing according to a power law, possibly with logarithmic corrections. Alternatively, the coarsening scenarios observed in and in this article might be governed by the same scaling laws, with their difference being only apparent. The emphasis of was on the scaling laws governing coarsening, that of the present study is on the mechanism of coarsening. Obviously, this situation calls for large-scale simulations in order to determine the scaling exponents in cases where the mechanism presented here is definitely at work already.
If the mentioned secondary instability becomes important, the identified state of competing grooves will be only of intermediate nature. Of course, all our considerations hold only as long as linear elasticity remains valid in the bulk, nonlinear elastic effects may alter the scenario.
With gravity included (which was not considered in ), there are some modifications. First, it is now possible for a planar interface to be stable (apart from a vertical translation). Once the threshold of the instability has been exceeded, the behavior will be similar to the case without gravity. However, we predict that it is possible for several grooves to survive in a finite system and that they will eventually stop growing, because the stress does not increase beyond a certain magnitude due to the lateral system width, whereas the gravity term increases as long as a groove gets deeper. That several grooves may survive has to do with the fact that now we have length scales in Eq. (44). More simple-mindedly we can immediately see that once the biggest groove stops, the second-largest will not retract, if it still has a downward velocity at that moment. Once the second stops, we can repeat the argument for the third, and so on. The final state will consist of a number of grooves, probably of different lengths and disordered. (For large-width systems, the aforementioned secondary instability may again complicate the picture.) In laterally infinite systems, it seems likely that a scaling state will prevail, possibly with a modified scaling exponent. Because now both stresses at the tips of the largest grooves and the gravity terms continue to grow, but they will both grow linearly with the length of the grooves. If initially the stress was large enough to overcompensate the gravity term, it will presumably stay like that. It is, however, not excluded that starting from specific initial conditions a system can be stabilized by gravity in the end.
Acknowledgments
This work was supported by the German Research Society (Deutsche Forschungsgemeinschaft) under grant Ka 672/4-2, which is gratefully acknowledged. Moreover, we acknowledge a PROCOPE grant for travel exchanges by the DAAD (German academic exchange service), grant no. 9619897, and the A.P.A.P.E. (corresponding French organisation), grant no. 97176, as well as financial support by the TMR network “Pattern formation, noise and spatio-temporal chaos in complex systems”.
## A Derivation of the sharp-interface limit
In deriving the sharp-interface limit, we will restrict ourselves to the two-dimensional case as we did in discussing the sharp-interface equations (where we used only two stress components and only one curvature). The generalization to three dimensions is, however, straightforward. We may use (39) and (42) as outer equations, to be used in the regions where the gradient of the phase field is small. For convenience, we set $`z_0=0`$.
To obtain the inner equations, we transform to a local system of (orthogonal) curvilinear coordinates comoving with the interface, with one coordinate axis parallel to $`\varphi `$; the corresponding coordinate will be called $`r`$, while the second will be conveniently expressed by the arclength $`s`$ along the interface . We introduce a stretched variable setting $`r=\stackrel{~}{ϵ}\rho `$. It is then easy to see that a distinguished limit of (39), leading to a nontrivial inner equation that allows to satisfy the boundary conditions, is obtained by setting $`\stackrel{~}{ϵ}=ϵ`$. In saying this, we have assumed that the stresses and strains behave properly under rescaling, i.e., do not diverge. Designating by capital letters the values of the fields in the inner domain (where the gradient of the phase field is large), we then have the inner equations
$`{\displaystyle \frac{v}{ϵ}}_\rho \mathrm{\Phi }`$ $`=`$ $`{\displaystyle \frac{\gamma }{\stackrel{~}{k}}}\{{\displaystyle \frac{1}{ϵ^2}}_\rho ^2\mathrm{\Phi }+{\displaystyle \frac{\kappa }{ϵ}}_\rho \mathrm{\Phi }+_s^2\mathrm{\Phi }{\displaystyle \frac{1}{ϵ^2}}[2g^{}(\mathrm{\Phi })`$ (A2)
$`+{\displaystyle \frac{ϵ}{3\gamma }}h^{}(\mathrm{\Phi })(\mu U_{ij}U_{ij}+{\displaystyle \frac{\lambda \stackrel{~}{\lambda }}{2}}U_{ii}^2+\mathrm{\Delta }pU_{ii}+\mathrm{\Delta }W+\mathrm{\Delta }\rho gz)]\},`$
$`0`$ $`=`$ $`{\displaystyle \frac{1}{ϵ}}_\rho \stackrel{~}{\mathrm{\Sigma }}_{\rho \rho }+_s\stackrel{~}{\mathrm{\Sigma }}_{\rho s}+\kappa \left(\stackrel{~}{\mathrm{\Sigma }}_{\rho \rho }\stackrel{~}{\mathrm{\Sigma }}_{ss}\right),`$ (A3)
$`0`$ $`=`$ $`{\displaystyle \frac{1}{ϵ}}_\rho \stackrel{~}{\mathrm{\Sigma }}_{s\rho }+_s\stackrel{~}{\mathrm{\Sigma }}_{ss}+\kappa \left(\stackrel{~}{\mathrm{\Sigma }}_{\rho s}+\stackrel{~}{\mathrm{\Sigma }}_{s\rho }\right),`$ (A4)
where
$$\stackrel{~}{\mathrm{\Sigma }}_{\alpha \beta }=h(\mathrm{\Phi })\mathrm{\Sigma }_{\alpha \beta }(1h(\mathrm{\Phi }))P\delta _{\alpha \beta }$$
(A5)
is the generalized stress tensor of the two-phase system. In order to obtain (A2A4), we have used
$``$ $`=`$ $`{\displaystyle \frac{1}{ϵ}}𝐧_\rho +𝐭_s,`$ (A6)
$`^2`$ $`=`$ $`{\displaystyle \frac{1}{ϵ^2}}_\rho ^2+{\displaystyle \frac{\kappa }{ϵ}}_\rho +_s^2.`$ (A7)
Derivatives such as $`_xu_x`$ can then be expressed in invariant form as $`𝐞_x(𝐞_x)𝐮`$, which leads to the following relations for the strain tensor components in the new coordinates:
$`U_{\rho \rho }`$ $`=`$ $`{\displaystyle \frac{1}{ϵ}}_\rho U_\rho ,`$ (A8)
$`U_{ss}`$ $`=`$ $`_sU_s+\kappa U_\rho ,`$ (A9)
$`U_{\rho s}`$ $`=`$ $`U_{s\rho }={\displaystyle \frac{1}{2}}\left(_sU_\rho +{\displaystyle \frac{1}{ϵ}}_\rho U_s\kappa U_s\right).`$ (A10)
The next step consists in solving the outer and inner equations via an asymptotic analysis that leads to a globally valid approximation for small interface thickness $`ϵ`$ which approaches the sharp-interface equations as $`ϵ0^+`$. To this end we expand both outer and inner fields in powers of $`ϵ`$:
$`\varphi (x,z,t)`$ $`=`$ $`\varphi _0(x,z,t)+ϵ\varphi _1(x,z,t)+\mathrm{},`$ (A11)
$`u_{ij}(x,z,t)`$ $`=`$ $`u_{ij}^{(0)}(x,z,t)+ϵu_{ij}^{(1)}(x,z,t)+\mathrm{},`$ (A12)
and
$`\varphi (x,z,t)`$ $`=`$ $`\mathrm{\Phi }(\rho ,s,t)=\mathrm{\Phi }_0(\rho ,s,t)+ϵ\mathrm{\Phi }_1(\rho ,s,t)+\mathrm{},`$ (A13)
$`u_{ij}(x,z,t)`$ $`=`$ $`U_{ij}(\rho ,s,t)=U_{ij}^{(0)}(\rho ,s,t)+ϵU_{ij}^{(1)}(\rho ,s,t)+\mathrm{},`$ (A14)
where, due to the transformation properties of tensors, we can think of the subscripts $`i`$, $`j`$ as running either over the values ($`x`$,$`z`$) or ($`r`$,$`s`$) and ($`\rho `$,$`s`$), respectively. Our basic field equations are, however, equations not for the strains but for the displacement fields. Thus, the expansion of the $`U_{ij}(\rho ,s,t)`$ induces one for the displacement components:
$`u_r`$ $`=`$ $`U_\rho (\rho ,s,t)=U_\rho ^{(0)}(\rho ,s,t)+ϵU_\rho ^{(1)}(\rho ,s,t)+\mathrm{},`$ (A15)
$`u_s`$ $`=`$ $`U_s(\rho ,s,t)=U_s^{(0)}(\rho ,s,t)+ϵU_s^{(1)}(\rho ,s,t)+\mathrm{}.`$ (A16)
Now the physical requirement that both $`u_r`$ and $`u_s`$ remain finite in the limit $`ϵ0^+`$ allows us to conclude from (A8) and (A10) that neither $`U_\rho ^{(0)}`$ nor $`U_s^{(0)}`$ can depend on $`\rho `$, hence
$`U_\rho ^{(0)}`$ $`=`$ $`U_\rho ^{(0)}(s,t),`$ (A17)
$`U_s^{(0)}`$ $`=`$ $`U_s^{(0)}(s,t).`$ (A18)
Furthermore, we have matching conditions for $`1\rho ϵ^1`$ that can be obtained from the inner and outer expansions by equating equal powers of $`ϵ`$ (and taking into account that the variable $`r`$ is itself $`ϵ`$ dependent):
$`\mathrm{\Phi }_0(\rho ,s,t)`$ $``$ $`\varphi _0(r,s,t)|_{r=\pm 0},\rho \pm \mathrm{},`$ (A19)
$`\mathrm{\Phi }_1(\rho ,s,t)`$ $``$ $`[\varphi _1(r,s,t)+\rho _r\varphi _0(r,s,t)]|_{r=\pm 0},\rho \pm \mathrm{},`$ (A20)
$`U_{\alpha \beta }^{(0)}(\rho ,s,t)`$ $``$ $`u_{\alpha \beta }^{(0)}(r,s,t)|_{r=\pm 0},\rho \pm \mathrm{},`$ (A21)
where we use the $``$ symbol in the sense of asymptotic equality, i.e., $`f(x)g(x)`$, $`xx_0`$ is equivalent to $`lim_{xx_0}f(x)/g(x)=1`$, and for two series in $`xx_0`$ we require this relation for each corresponding pair of terms.
The relations induced by (A21) for the displacements are more complicated. We just give two examples. Because each derivative with respect to $`r`$ comes with a factor $`1/ϵ`$ when transformed into a derivative w.r.t. $`\rho `$, we have $`U_{\rho \rho }^{(0)}=_\rho U_\rho ^{(1)}`$ and hence
$$\underset{\rho \pm \mathrm{}}{lim}_\rho U_\rho ^{(1)}(\rho ,s,t)=_ru_r^{(0)}(r,s,t)|_{r=\pm 0}.$$
(A22)
Our second example is even more instructive. We write
$`\underset{\rho \pm \mathrm{}}{lim}U_{ss}^{(0)}(\rho ,s,t)`$ $`=`$ $`\underset{\rho \pm \mathrm{}}{lim}\left(_sU_s^{(0)}(s,t)+\kappa U_\rho ^{(0)}(s,t)\right)`$ (A23)
$`=`$ $`_sU_s^{(0)}(s,t)+\kappa U_\rho ^{(0)}(s,t)`$ (A24)
$`=`$ $`[_su_s^{(0)}(r,s,t)+\kappa u_\rho ^{(0)}(r,s,t)]|_{r=\pm 0},`$ (A25)
which shows that the linear combination $`u_s^{(0)}+\kappa u_\rho ^{(0)}`$ must be continuous across the interface.
Finally, we need the expansions of all functions of $`\varphi `$ in powers of $`ϵ`$, e.g.:
$$h(\varphi )=h(\varphi _0)+ϵh^{}(\varphi _0)\varphi _1+ϵ^2\left(h^{}(\varphi _0)\varphi _2+\frac{1}{2}h^{\prime \prime }(\varphi _0)\varphi _1^2\right)+\mathrm{},$$
(A26)
and we will use the obvious abbreviations $`h_0`$, $`h_0^{}`$, etc., for functions of $`\varphi _0`$. Let us note a few useful relations in passing:
$`h^{}(\varphi )`$ $`=`$ $`6\varphi (1\varphi ),`$ (A27)
$`g(\varphi )`$ $`=`$ $`\varphi ^2(1\varphi )^2=\left({\displaystyle \frac{1}{6}}h^{}(\varphi )\right)^2,`$ (A28)
$`g^{}(\varphi )`$ $`=`$ $`2\varphi (1\varphi )(12\varphi )={\displaystyle \frac{1}{18}}h^{}(\varphi )h^{\prime \prime }(\varphi ),`$ (A29)
We have now collected all the prerequisites to perform the asymptotic analysis providing the sharp-interface limit.
First, we note that the outer solution to lowest order \[order $`ϵ^2`$ of (39)\] is simply given by $`g^{}(\varphi _0)=0`$, which yields the solutions $`\varphi _0=0`$, $`\varphi _0=1`$, and $`\varphi _0=\frac{1}{2}`$ \[see (A29)\]. The last of these is unstable and also not compatible with the boundary conditions in typical numerical setups. We assume $`\varphi _0=0`$ for $`r>0`$, corresponding to the liquid phase, and $`\varphi _0=1`$ for $`r<0`$, corresponding to the solid phase. Equation (A27) tells us that $`h^{}(\varphi _0)=0`$, and hence these solutions are valid at all orders of $`ϵ`$. Using $`h(\varphi _0)=0`$ in the liquid and $`h(\varphi _0)=1`$ in the solid, we immediately see that (42) turns into the mechanical equilibrium condition for the liquid and solid, respectively:
$`_ip`$ $`=`$ $`0\text{(liquid)},`$ (A30)
$`_j\sigma _{ij}`$ $`=`$ $`0\text{(solid)}.`$ (A31)
This is again true at all orders of $`ϵ`$, and we can write the zeroth-order piece of the result in the form:
$`p^{(0)}`$ $`=`$ $`p_0\mathrm{}\stackrel{~}{\lambda }u_{kk}^{(0)}=p_0=\mathrm{const}.,`$ (A32)
$`\sigma _{ij}^{(0)}`$ $`=`$ $`p_{0s}\delta _{ij}+2\mu u_{ij}^{(0)}+\lambda u_{kk}^{(0)}\delta _{ij}.`$ (A33)
Later, we will look at two reference states in particular. One is the “natural” choice $`p_{0s}=p_0\mathrm{}`$, i.e., the unstrained state is hydrostatic and corresponds to the same pressure in the liquid and in the solid. If moreover, this pressure is chosen equal to the equilibrium pressure $`p_0`$, then we have $`u_{kk}^{(0)}0`$ in the liquid at equilibrium. The second choice corresponds to assuming a finite difference $`\mathrm{\Delta }p=p_0\mathrm{}p_{0s}`$ while keeping $`p_0\mathrm{}=p_0`$. This means that zero strain corresponds to a prestressed solid, with a stress tensor $`\sigma _{ij}=p_0\delta _{ij}+\mathrm{\Delta }p\delta _{ij}`$, i.e., the deviation from equilibrium is the isotropic tensor $`\mathrm{\Delta }p\delta _{ij}`$. Both approaches can be exploited numerically.
We now consider the inner solution. The lowest order of (A2) gives
$$_\rho ^2\mathrm{\Phi }_02g^{}(\mathrm{\Phi }_0)=0.$$
(A34)
This equation can be solved by standard methods. Multiplying by $`_\rho \mathrm{\Phi }_0`$, we immediately obtain a first integral, written down here for further reference,
$$_\rho \mathrm{\Phi }_0=2\mathrm{\Phi }_0(1\mathrm{\Phi }_0)=\frac{1}{3}h_0^{},$$
(A35)
which is solved by $`\mathrm{\Phi }_0=\frac{1}{2}(1\mathrm{tanh}\rho )`$, and this solution satisfies the matching conditions (A19).
From the strain equations (A3), (A4) we obtain to lowest order
$`_\rho \stackrel{~}{\mathrm{\Sigma }}_{\rho \rho }^{(0)}`$ $`=`$ $`0,`$ (A36)
$`_\rho \stackrel{~}{\mathrm{\Sigma }}_{s\rho }^{(0)}`$ $`=`$ $`0,`$ (A37)
which on integration from $`\rho =\mathrm{}`$ to $`\rho =\mathrm{}`$ together with the matching conditions (A21) yields
$`\sigma _{rr}^{(0)}|_{r=0^{}}`$ $`=`$ $`\stackrel{~}{\mathrm{\Sigma }}_{\rho \rho }^{(0)}(\mathrm{})=\stackrel{~}{\mathrm{\Sigma }}_{\rho \rho }^{(0)}(\mathrm{})=P=p_0,`$ (A38)
$`\sigma _{sr}^{(0)}|_{r=0^{}}`$ $`=`$ $`\stackrel{~}{\mathrm{\Sigma }}_{s\rho }^{(0)}(\mathrm{})=\stackrel{~}{\mathrm{\Sigma }}_{s\rho }^{(0)}(\mathrm{})=0.`$ (A39)
The limiting values of $`\stackrel{~}{\mathrm{\Sigma }}_{\rho \rho }^{(0)}`$ and $`\stackrel{~}{\mathrm{\Sigma }}_{s\rho }^{(0)}`$ can be gathered from (A5). Obviously, these two equations constitute the condition of mechanical equilibrium at the interface, as $`\sigma _{rr}^{(0)}`$ and $`\sigma _{sr}^{(0)}`$ are the normal and shear components of the stress tensor of the outer solution there.
However, the strain equations provide more information than just mechanical equilibrium on the outer scale. We write (A36) explicitly in terms of the strains and integrate indefinitely with respect to $`\rho `$, which yields
$$h_0\left[(2\mu +\lambda \stackrel{~}{\lambda })U_{\rho \rho }^{(0)}+(\lambda \stackrel{~}{\lambda })U_{ss}^{(0)}+\mathrm{\Delta }p\right]+\stackrel{~}{\lambda }(U_{\rho \rho }^{(0)}+U_{ss}^{(0)})=f(s,t),$$
(A40)
where $`f(s,t)`$ is a function of integration, to be determined from the matching conditions. This is straightforward and yields $`f(s,t)=p_0\mathrm{}p_0`$. Moreover, we know that the only spatial dependence of $`U_{ss}^{(0)}`$ is that on $`s`$ \[see (A25) and (A18)\], which suggests to solve for $`U_{\rho \rho }^{(0)}`$ in terms of $`U_{ss}^{(0)}`$. The result is
$`U_{\rho \rho }^{(0)}`$ $`=`$ $`{\displaystyle \frac{1}{(2\mu +\lambda \stackrel{~}{\lambda })h_0+\stackrel{~}{\lambda }}}\left\{\mathrm{\Delta }ph_0+p_0p_0\mathrm{}+[\stackrel{~}{\lambda }+(\lambda \stackrel{~}{\lambda })h_0]U_{ss}^{(0)}\right\}.`$ (A41)
The advantage of this equation is that it provides us with the full $`\rho `$ dependence of $`U_{\rho \rho }^{(0)}`$, allowing the explicit evaluation of integrals on $`\rho `$ containing the strains. An analogous procedure for the second strain equation determines $`U_{s\rho }^{(0)}`$ to be equal to zero.
Now we proceed to the next-order equation for $`\mathrm{\Phi }`$. Written out explicitly, it reads
$`v_\rho \mathrm{\Phi }_0`$ $`=`$ $`{\displaystyle \frac{\gamma }{\stackrel{~}{k}}}\{_\rho ^2\mathrm{\Phi }_1+\kappa _\rho \mathrm{\Phi }_02g_0^{\prime \prime }\mathrm{\Phi }_1`$ (A44)
$`{\displaystyle \frac{h_0^{}}{3\gamma }}(\mu [U_{\rho \rho }^{(0)}{}_{}{}^{2}+U_{ss}^{(0)}{}_{}{}^{2}]+{\displaystyle \frac{\lambda \stackrel{~}{\lambda }}{2}}(U_{\rho \rho }^{(0)}+U_{ss}^{(0)})^2`$
$`+\mathrm{\Delta }p(U_{\rho \rho }^{(0)}+U_{ss}^{(0)})+\mathrm{\Delta }W+\mathrm{\Delta }\rho gz(s))\}`$
where we have used $`U_{s\rho }^{(0)}=0`$. With the help of (A27) and (A35), we can arrange this as
$`L\mathrm{\Phi }_1`$ $`=`$ $`{\displaystyle \frac{h_0^{}}{3\gamma }}\{\stackrel{~}{k}v+\gamma \kappa +\mu [U_{\rho \rho }^{(0)}{}_{}{}^{2}+U_{ss}^{(0)}{}_{}{}^{2}]+{\displaystyle \frac{\lambda \stackrel{~}{\lambda }}{2}}(U_{\rho \rho }^{(0)}+U_{ss}^{(0)})^2`$ (A46)
$`+\mathrm{\Delta }p(U_{\rho \rho }^{(0)}+U_{ss}^{(0)})+\mathrm{\Delta }W+\mathrm{\Delta }\rho gz(s)\},`$
with $`L_\rho ^22g_0^{\prime \prime }`$ being a self-adjoint linear operator. The solvability condition for this inhomogeneous linear equation is that the right-hand side be orthogonal to the left-sided null eigenspace of $`L`$. Since $`L`$ is hermitean, we know that the translational mode $`_\rho \mathrm{\Phi }_0`$ is an appropriate eigenvector:
$$_\rho \mathrm{\Phi }_0L=L_\rho \mathrm{\Phi }_0=0.$$
(A47)
Multiplying (A46) by $`3\gamma _\rho \mathrm{\Phi }_0`$ from the left and integrating on $`\rho `$, we find
$`0`$ $`=`$ $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}d\rho \{\stackrel{~}{k}v+\gamma \kappa +\mu [U_{\rho \rho }^{(0)}{}_{}{}^{2}+U_{ss}^{(0)}{}_{}{}^{2}]+{\displaystyle \frac{\lambda \stackrel{~}{\lambda }}{2}}(U_{\rho \rho }^{(0)}+U_{ss}^{(0)})^2`$ (A49)
$`+\mathrm{\Delta }p(U_{\rho \rho }^{(0)}+U_{ss}^{(0)})+\mathrm{\Delta }W+\mathrm{\Delta }\rho gz(s)\}h^{}_0_\rho \mathrm{\Phi }_0.`$
Now we can exploit (A41), telling us that the $`\rho `$ dependence of the braces in (A49) is fully contained in their dependence on $`h_0`$. All the integrals can be done analytically, using
$$I_{\mathrm{}}^{\mathrm{}}𝑑\rho f(h_0)h_0^{}_\rho \mathrm{\Phi }_0=_1^0𝑑\mathrm{\Phi }_0f(h(\mathrm{\Phi }_0))h^{}(\mathrm{\Phi }_0)=_0^1𝑑hf(h).$$
(A50)
Integrals that appear in (A49) are
$`I_1`$ $`=`$ $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\rho h_0^{}_\rho \mathrm{\Phi }_0=1,`$ (A51)
$`I_2`$ $`=`$ $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\rho U_{\rho \rho }^{(0)}h_0^{}_\rho \mathrm{\Phi }_0,`$ (A52)
$`I_3`$ $`=`$ $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\rho U_{\rho \rho }^{(0)}{}_{}{}^{2}h_0^{}_\rho \mathrm{\Phi }_0.`$ (A53)
The evaluation of the latter two integrals is as straightforward as that of the first, although a little more tedious. We just give the final result for the solvability condition, taking $`p_0\mathrm{}=p_0`$ for simplicity:
$$\stackrel{~}{k}v=\gamma \kappa +\mathrm{\Delta }\rho gz+\frac{\mu }{2(\mu +\lambda )(2\mu +\lambda )}\left[2(\mu +\lambda )U_{ss}^{(0)}+\mathrm{\Delta }p\right]^2\frac{2\mu +\lambda }{8\mu (\mu +\lambda )}\sigma _{00}^2,$$
(A54)
At this point, we may specify our choice of reference state for the solid. First, let us assume that the unstrained state corresponds to a state of equal hydrodynamic pressure in the two phases, i.e. $`p_{0s}=p_0\mathrm{}`$, or $`\mathrm{\Delta }p=0`$. This is the KM choice. Then taking the limit $`\rho \mathrm{}`$ of (A41) we get ($`U_{ss}^{(0)}=u_{ss}^{(0)}`$)
$$u_{rr}^{(0)}=\frac{\lambda u_{ss}^{(0)}}{2\mu +\lambda },$$
(A55)
implying $`\sigma _{ss}^{(0)}\sigma _{rr}^{(0)}=2\mu (u_{ss}^{(0)}u_{rr}^{(0)})=4\mu (\mu +\lambda )u_{ss}^{(0)}/(2\mu +\lambda )`$, from which we obtain
$$u_{ss}^{(0)}=\frac{2\mu +\lambda }{4\mu (\mu +\lambda )}(\sigma _{tt}^{(0)}\sigma _{nn}^{(0)}),$$
(A56)
where we have now switched to the conventional notation for the principal components of the stress tensor in the normal and tangential directions ($`\sigma _{rr}=\sigma _{nn}`$, $`\sigma _{ss}=\sigma _{tt}`$). Finally, expressing the Lamé constants by Young’s modulus and the Poisson ratio, we arrive at
$$v=\frac{1}{\stackrel{~}{k}}\left\{\frac{1\nu ^2}{2E}\left[(\sigma _{tt}^{(0)}\sigma _{nn}^{(0)})^2\sigma _{00}^2\right]+\gamma \kappa +\mathrm{\Delta }\rho gz\right\},$$
(A57)
which is the desired sharp-interface limit. \[In eqs. (1, 2), $`\sigma _{00}=0`$.\]
A remark is in order here. The phase-field equations imply the continuity of $`u_{ss}^{(0)}`$ across the interface. As this quantity is obviously nonzero whenever the solid is strained, this means that we will not have $`u_{ss}^{(0)}=0`$ in the liquid. However, we will still have $`u_{ss}^{(0)}+u_{rr}^{(0)}=0`$, i.e., the divergence of the displacement vector vanishes in the liquid. But this is all that matters, because it is only this quantity that enters the description of the liquid.
The reason for $`u_{ss}^{(0)}0`$ in the liquid is that the phase-field description imposes coherence of the strain across the interface, which ultimately goes back to our viewing the liquid as a (shear-free, but nonetheless) solid. For a true liquid in contact with a solid there is no such coherence condition as it is free to slip on the solid surface. Therefore, it could always keep its strain tensor isotropic (if such a notion made much sense at all for a liquid). The reason why we can nevertheless model the liquid as a solid is the additional degree of freedom that arises in the description of a liquid by having two fields $`u_x`$ and $`u_z`$ at our disposal even though only their combination $`𝐮`$ enters the free-energy expression. Therefore, we can compensate, so to speak, for imposing (nonphysical) coherence by allowing (equally nonphysical) anisotropic strain in the liquid.
At this point we may also note that if we were to model the elastic properties of two real solids by the current phase-field approach, we would necessarily impose coherence at the interface. The treatment of noncoherent solid-solid interfaces via phase fields would require some rethinking of the method.
Concerning computational purposes, one disadvantage of the chosen reference state is that it is not very well-suited for the use of periodic boundary conditions, meaning periodically varying stresses and strains. The field equations are set up in terms of the displacements which acquire linearly increasing or decreasing components in directions where the strain has a nonzero average. This observation motivates the consideration of a different reference state in the solid, in which the average strain due to the external stress is subtracted. This is the MG choice (see App. B). If we impose a constant stress $`\sigma _0`$ in the $`x`$ direction, the stress tensor in the solid is $`\sigma _{ij}=p_0\delta _{ij}+\sigma _0\delta _{i1}\delta _{1j}`$ and requiring $`u_{xx}=0`$, we find that this is achieved by setting
$$\mathrm{\Delta }p=\frac{2\mu +\lambda }{2\mu }\sigma _0.$$
(A58)
The corresponding homogeneous strain tensor is given by $`u_{xx}=0`$, $`u_{xz}=0`$, and $`u_{zz}=\mathrm{\Delta }p/(2\mu +\lambda )`$. We then obtain taking the limit $`\rho \mathrm{}`$ of (A41)
$$u_{rr}^{(0)}=\frac{\mathrm{\Delta }p\lambda u_{ss}^{(0)}}{2\mu +\lambda }.$$
(A59)
This can be used to express
$$\sigma _{ss}^{(0)}\sigma _{rr}^{(0)}=2\mu (u_{ss}^{(0)}u_{rr}^{(0)})=\frac{2\mu }{2\mu +\lambda }\left[2(\mu +\lambda )u_{ss}^{(0)}+\mathrm{\Delta }p\right],$$
(A60)
wherefrom we obtain $`2(\mu +\lambda )u_{ss}^{(0)}+\mathrm{\Delta }p=(2\mu +\lambda )(\sigma _{tt}^{(0)}\sigma _{nn}^{(0)})/2\mu `$, which on insertion in (A54) leads back to (A57).
Note that even here we cannot require $`u_{ss}^{(0)}=u_{rr}^{(0)}=0`$ in the liquid, which would imply $`u_{ss}^{(0)}=0`$ at the interface and thus, according to (A41), $`u_{rr}^{(0)}=\mathrm{\Delta }p/(2(\mu +\lambda ))`$, i.e., $`u_{rr}^{(0)}`$ would be constant along the interface. Then also $`\sigma _{tt}^{(0)}\sigma _{nn}^{(0)}`$ would have to be constant, which would lead to a dynamics entirely different from that of the Grinfeld instability \[where $`(\sigma _{tt}^{(0)}\sigma _{nn}^{(0)})^2`$ increases in the grooves and diminishes on the peaks\]. Hence, once again we are obliged to make use of the additional degree of freedom of the fields inside the liquid, even though now we can impose $`u_{xx}^{(0)}=u_{zz}^{(0)}=0`$ in the liquid as an initial condition (and as a far-field boundary condition), because this satisfies the periodicity requirement.
## B Mapping of the Müller-Grant model to the present formulation
The main difference between the form of the MG model given in and the one given here is a different choice of the functions $`g(\varphi )`$ and $`h(\varphi )`$. To make this conspicuous, we will rename their original functions to $`\stackrel{~}{g}(\varphi )`$ and $`\stackrel{~}{h}(\varphi )`$ (since the second function was called $`g(\varphi )`$ in , our renaming is also useful to avoid unnecessary confusion here). We shall leave the gravity and shift terms, $`f_{\mathrm{grav}}`$ and $`f_\mathrm{c}`$ out of the consideration, since they were not used by MG.
Their double well potential is defined as
$$f_{\mathrm{dw}}(\varphi )=\frac{1}{a}\stackrel{~}{g}(\varphi ),$$
(B1)
with $`\stackrel{~}{g}(\varphi )=\varphi ^2(1\varphi ^2)^2`$, which is a sixth-order polynomial and actually has a third minimum at $`\varphi =1`$. The latter does not, however, play any role in the dynamics, provided no negative $`\varphi `$ values are given in the initial condition. $`a`$ is a constant to be identified via the sharp-interface limit.
Second, there is an elastic contribution to the free energy which they give as
$$f_{\mathrm{el}}(\varphi ,\{u_{ij}\})=\frac{1}{2}K(𝐮)^2+\stackrel{~}{\mu }\underset{ij}{}\left(u_{ij}\frac{\delta _{ij}}{d}𝐮\right)^2,$$
(B2)
where $`K`$ is the bulk modulus and $`\stackrel{~}{\mu }`$ the shear modulus which is $`\varphi `$ dependent:
$$\stackrel{~}{\mu }=\mu _1\stackrel{~}{h}(\varphi ).$$
(B3)
The convenient choice
$$\stackrel{~}{h}(\varphi )=\frac{1}{2}\varphi ^2\frac{1}{4}\varphi ^4,$$
(B4)
guarantees that both bulk phases keep their equilibrium values at $`\varphi =0`$ (liquid) and $`\varphi =1`$ (solid). This is due to the fact that $`\stackrel{~}{h}^{}(0)=\stackrel{~}{h}^{}(1)=0`$, a property $`\stackrel{~}{h}(\varphi )`$ shares with $`h(\varphi )`$ from the KM model (see Sec. II C). Obviously, the true shear modulus of the solid is $`\mu =\mu _1\stackrel{~}{h}(1)=\mu _1/4`$.
For simplicity and since it does not change the behavior qualitatively, the bulk modulus is assumed to be the same in both phases. However, this restriction can be easily dropped by replacing $`K`$ with
$$K=K_0+K_1\stackrel{~}{h}(\varphi ).$$
(B5)
As reference state they chose a prestressed state of the solid with $`\sigma _{xx}=\sigma _0`$, in which the strains $`u_{xx}`$ and $`u_{xz}`$ for a flat surface vanish. This entails that the state in which all strains vanish is a hydrostatic state with a different stress value. It is described by MG using a parameter $`\eta _0`$ and \[as may be verified easily from Eq. (B10) below\], in this state we have $`\sigma _{xx}=\sigma _{zz}=\eta _0\stackrel{~}{h}(1)=\eta _0/4`$. As a result, there is a relation between the new parameter $`\eta _0`$ and the external stress $`\sigma _0`$ in the uniaxially stressed reference state
$$\eta _0=\frac{8(K+\mu _1/4)}{\mu _1}\sigma _0.$$
(B6)
The free energy density is then given as the sum of $`f_{\mathrm{dw}}(\varphi )`$ and $`f_{\mathrm{el}}(\varphi ,\{u_{ij}\})`$ and two additional terms $`(\eta _0^2/2K)\stackrel{~}{h}(\varphi )^2`$ and $`\eta _0\stackrel{~}{h}(\varphi )𝐮`$. The first of these terms describes an energy shift, the second an additional coupling between the phase field and the elastic field (beyond that already implied by the $`\varphi `$ dependence of the shear modulus and, possibly, the bulk modulus). A nice feature of the approach given in the present paper is that these terms are automatically generated by accounting for the fact that the equilibrium state does not have vanishing strain when the MG reference state is used: these are the terms containing $`p_0p_{0s}`$ in (17) and the corresponding terms $`\mathrm{\Delta }pu_{ii}`$ and $`\mathrm{\Delta }W`$ in (39).
The free energy density is then given by:
$$f(\varphi ,u_{ij})=\frac{1}{a}\stackrel{~}{g}(\varphi )+\frac{\eta _0^2}{2K}\stackrel{~}{h}(\varphi )^2+\eta _0\stackrel{~}{h}(\varphi )𝐮+\frac{1}{2}K(𝐮)^2+\stackrel{~}{\mu }\underset{ij}{}\left(u_{ij}\frac{\delta _{ij}}{d}𝐮\right)^2.$$
(B7)
The first term is the double well potential. The second and third terms are due to the particular choice of reference frame, and $`\eta _0`$ is related to the externally applied stress as described by Eq. (B6). Applying the same line of reasoning as in Sec. II C, they obtain a system of coupled partial differential equations:
$`{\displaystyle \frac{\varphi }{t}}`$ $`=`$ $`\stackrel{~}{\mathrm{\Gamma }}[{\displaystyle \frac{1}{a}}\stackrel{~}{g}^{}(\varphi )l^2^2\varphi `$ (B8)
$`+`$ $`{\displaystyle \frac{\eta _0^2}{K}}\stackrel{~}{h}(\varphi )\stackrel{~}{h}^{}(\varphi )+\eta _0\stackrel{~}{h}^{}(\varphi )𝐮+\mu _1\stackrel{~}{h}^{}(\varphi ){\displaystyle \underset{ij}{}}(u_{ij}{\displaystyle \frac{\delta _{ij}}{d}}𝐮)^2],`$ (B9)
and
$$\frac{\sigma _{ij}}{x_j}=\frac{}{x_i}[\eta _0\stackrel{~}{h}(\varphi )+K𝐮]+2\mu _1\frac{}{x_j}\left[\stackrel{~}{h}(\varphi )\left(u_{ij}\frac{\delta _{ij}}{d}𝐮\right)\right]=0.$$
(B10)
They show in that the phase field equations of this model also converge to the sharp interface equations. By expanding the solution of the mechanical equilibrium condition to first order in the shear modulus they were able to integrate out the elastic fields, so that they were left with an equation for $`\varphi `$ only. That allowed them to use a pseudospectral method with which they could study wide periodic systems and three dimensional systems .
Whether this expansion is entirely consistent is an open question, since they consider $`\eta _0`$ an independent parameter that is $`O(1)`$, whereas for fixed $`\sigma _0`$ we actually have $`\eta _0=O(1/\mu _1)`$, and $`\mu _1`$ is the small quantity in their expansion. Probably this does not really matter in the absence of gravity where the precise value of $`\eta _0`$ is immaterial as it only sets the time scale. The problem may be more awkward in the presence of gravity.
Equations (B9) and (B10) are in a form that allows direct comparison with Eqs. (39) and (42), respectively. First note that in two dimensions $`K=\lambda +\mu =\lambda +\mu _1/4`$ and that we must set $`\stackrel{~}{\lambda }=K`$ to have the same elastic constants in the two sets of equations. In 2D, we have $`[u_{ij}(\delta _{ij}/d)𝐮]^2=u_{ij}u_{ij}\frac{1}{2}u_{ii}^2`$ (summations over $`i`$ and $`j`$ are implied) and because of $`\lambda \stackrel{~}{\lambda }=\mu `$, we obtain
$$\mu u_{ij}u_{ij}+\frac{\lambda \stackrel{~}{\lambda }}{2}u_{ii}^2=\frac{1}{4}\mu _1\left(u_{ij}\frac{\delta _{ij}}{d}𝐮\right)^2,$$
(B11)
which shows that on replacing $`\stackrel{~}{h}(\varphi )`$ with $`\frac{1}{4}h(\varphi )`$ the term of (B9) that is quadratic in the strains becomes equal to the corresponding term of (39). Common prefactors will be discussed below. We then see immediately, that the choice $`\mathrm{\Delta }p=\eta _0/4`$ will make the linear terms equal. This choice is the right one as is revealed by a quick comparison of Eq. (B6) with Eq. (A58), derived in appendix A. Next we have to compare the constant terms which are $`\eta _0^2/K`$ and $`\mathrm{\Delta }W`$. Here we notice that in , it was assumed that $`p_0\mathrm{}=p_0=0`$. If we furthermore set $`\sigma _{00}=0`$, then we infer from (41) that $`\mathrm{\Delta }W=\mathrm{\Delta }p^2/2K=\eta _0^2/32K`$. Now the $`\varphi `$-dependent factor of $`\eta _0^2/K`$ in (B9) is $`\stackrel{~}{h}(\varphi )\stackrel{~}{h}^{}(\varphi )=\frac{1}{2}[\stackrel{~}{h}(\varphi )^2]/\varphi `$. We have checked by directly performing the sharp-interface limit of the original MG model, that this limit does not change when $`\stackrel{~}{h}(\varphi )^2`$ is replaced with $`\frac{1}{4}\stackrel{~}{h}(\varphi )`$ (the solid and liquid phase limits are obviously unchanged). Doing this replacement first and then substituting $`\frac{1}{4}h(\varphi )`$ for $`\stackrel{~}{h}(\varphi )`$, we get identity of the constant terms, too.
Finally, the prefactors should be discussed. In order to make the prefactor of the elastic expressions the same in both equations, we must set $`\stackrel{~}{\mathrm{\Gamma }}=1/(3\stackrel{~}{k}ϵ)`$, which shows that $`l^2=3\gamma ϵ`$. To determine the factor $`1/a`$ of the double well potential in (B9), one must actually perform the sharp-interface limit \[because the potential is not the same as in (39)\], which yields $`a=3ϵ/8\gamma `$. With these choices, Eqs. (B9) and (39) are asymptotically equivalent.
The comparison of the equations describing mechanical equilibrium is even more straightforward. Inserting the expressions (37) and (38) with $`p_0\mathrm{}=0`$ and $`p_{0s}=\mathrm{\Delta }p=\eta _0/4`$ into (42) and using $`\stackrel{~}{\lambda }=K`$, we get
$`0`$ $`=`$ $`{\displaystyle \frac{}{x_j}}\left\{h(\varphi )\left[{\displaystyle \frac{\eta _0}{4}}\delta _{ij}+Ku_{kk}\delta _{ij}+2\mu \left(u_{ij}{\displaystyle \frac{1}{2}}u_{kk}\delta _{ij}\right)\right]+[1h(\varphi )]Ku_{kk}\delta _{ij}\right\}`$ (B12)
$`=`$ $`{\displaystyle \frac{}{x_i}}\left[{\displaystyle \frac{\eta _0}{4}}h(\varphi )+Ku_{kk}\right]+2\mu {\displaystyle \frac{}{x_j}}\left[h(\varphi )\left(u_{ij}{\displaystyle \frac{1}{2}}u_{kk}\delta _{ij}\right)\right],`$ (B13)
which is obviously identical to (B10), once we replace $`h(\varphi )/4`$ by $`\stackrel{~}{h}(\varphi )`$.
## C Analytic solution of the elastic problem for the double cycloid
In this appendix (only!), we switch back from our notation for geometric relations in the plane which was based on a coordinate system spanned by the $`x`$ and $`z`$ axes (thus reminding ourselves that in reality we have a three-dimensional system and we simply suppress deformations in the $`y`$ direction) to the more conventional use of $`x`$ and $`y`$ for the planar coordinates. This way we “liberate” the symbol $`z`$ for use as a complex variable: $`z=x+iy`$.
We wish to solve the elastic problem in a half-infinite geometry, the top of which is bounded by the double cycloid given by Eq. (52) with $`z`$ replaced by $`y`$. In the complex plane, this curve is described by
$$z=x+iy=\xi iAe^{ik\xi }iBe^{2ik\xi }$$
(C1)
and because the parameter $`\xi `$ is real, this equation also defines a conformal mapping frome the $`z`$ plane to the $`\zeta `$ plane, where $`\zeta =\xi +i\eta `$, mapping the curve $`z(x)`$ to the $`\xi `$ axis.
To rephrase the elastic problem in the complex plane, we use the Goursat function formalism. Its basic statements can be easily inferred from the representation of two-dimensional elasticity in terms of a single scalar function, the Airy function $`\chi (x,y)`$. Setting $`\sigma _{xx}=_y^2\chi `$, $`\sigma _{yy}=_x^2\chi `$, and $`\sigma _{xy}=_x_y\chi `$, the mechanical equilibrium equations $`_j\sigma _{ij}=0`$ are automatically satisfied. Hooke’s law for isotropic bodies then implies that $`\chi `$ obeys the biharmonic equation:
$$^4\chi =0.$$
(C2)
In terms of the complex variables $`z`$ and $`\overline{z}=xiy`$, the Laplacian becomes $`^2=4_{\overline{z}}_z`$ \[because $`_x=_z+_{\overline{z}}`$, $`_y=i(_z_{\overline{z}})`$\], hence the most general form of the solution to (C2) is given by
$$\chi (x,y)\stackrel{~}{\chi }(z,\overline{z})=\overline{z}f_1(z)+g_1(z)+zf_2(\overline{z})+g_2(\overline{z}),$$
(C3)
where $`f_j`$, $`g_j`$ ($`j=1,2`$) are analytic functions of their arguments. Since we are looking for real solutions, we can restrict ourselves to two independent complex functions instead of four, i.e., we can write
$$\stackrel{~}{\chi }(z,\overline{z})=\frac{1}{2}\left\{\overline{z}\varphi (z)+\psi (z)𝑑z+z\overline{\varphi (z)}+\overline{\psi (z)}𝑑\overline{z}\right\},$$
(C4)
where an overbar denotes complex conjugation. In this formula, $`\varphi `$ and $`\psi `$ are the Goursat functions.
From (C4), we obtain by direct differentiation
$`\sigma _{xx}+\sigma _{yy}=^2\stackrel{~}{\chi }`$ $`=`$ $`2\left[\varphi ^{}(z)+\overline{\varphi ^{}(z)}\right],`$ (C5)
$`\sigma _{yy}\sigma _{xx}+2i\sigma _{xy}`$ $`=`$ $`2\left[\overline{z}\varphi ^{\prime \prime }(z)+\psi ^{}(z)\right].`$ (C6)
The displacements can also be expressed by the Goursat functions \[using Hooke’s law and (C5,C6)\], but since we do not need the corresponding relation, we omit it here.
We have boundary conditions for the stresses on the curve given by Eq. (C1), which one could attempt to use directly in (C5,C6) to obtain equations for the Goursat functions. However, since we are going to employ conformal mapping, it is useful to keep the order of derivatives small – the analytic evaluation of higher-order derivatives can be quite cumbersome. Thus it is desirable to reformulate the boundary conditions in partially integrated form. The force $`(f_x,f_y)`$ on the boundary can be condensed into a single complex number
$$f_x+if_y=\sigma _{xj}n_j+i\sigma _{yj}n_j=(\sigma _{xx}+i\sigma _{xy})n_x+(\sigma _{yy}i\sigma _{xy})in_y,$$
(C7)
where $`(n_x,n_y)`$ is the normal vector to the boundary. Introducing the arclength $`s`$, we have (directing $`s`$ such that $`s\mathrm{}`$ corresponds to $`x\mathrm{}`$)
$$n_x+in_y=\frac{dy}{ds}+i\frac{dx}{ds}=i\frac{dz}{ds}$$
(C8)
and thus
$`f_x+if_y`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(\sigma _{xx}+\sigma _{yy}\right)i{\displaystyle \frac{dz}{ds}}+{\displaystyle \frac{1}{2}}\left(\sigma _{yy}\sigma _{xx}2i\sigma _{xy}\right)i{\displaystyle \frac{d\overline{z}}{ds}}`$ (C9)
$`=`$ $`\left[\varphi ^{}(z)+\overline{\varphi ^{}(z)}\right]i{\displaystyle \frac{dz}{ds}}+\left[z\overline{\varphi ^{\prime \prime }(z)}+\overline{\psi ^{}(z)}\right]i{\displaystyle \frac{d\overline{z}}{ds}}`$ (C10)
$`=`$ $`i{\displaystyle \frac{d}{ds}}\left\{\varphi (z)+z\overline{\varphi ^{}(z)}+\overline{\psi (z)}\right\}.`$ (C11)
Integrating this local relation along the boundary, we obtain
$$ifi\left(f_x+if_y\right)𝑑s=\varphi (z)+z\overline{\varphi ^{}(z)}+\overline{\psi (z)},$$
(C12)
which allows to apply the boundary condition in expressions involving first-order derivatives only.
We substract the stress at infinity from our (linear) elastic equations to be able to work with analytic functions that are bounded at infinity. Hence, we set
$$\sigma _{ij}=\sigma _{ij}^{(0)}+\sigma _0\delta _{ix}\delta _{jx}$$
(C13)
and replace $`\sigma _{ij}`$ with $`\sigma _{ij}^{(0)}`$ in Eqs. (C5,C6) above. Taking the equilibrium pressure in the liquid equal to zero (which we can do without loss of generality), the boundary conditions at the liquid-solid interface ($`\sigma _{ij}n_j=0`$) become
$$\sigma _{ij}^{(0)}n_j=\sigma _0\delta _{ix}n_x,$$
(C14)
translating into
$$f_x+if_y=\sigma _0\frac{dy}{ds}if=\frac{1}{2}(\overline{z}z)\sigma _0,$$
(C15)
where we have dropped an arbitrary constant of integration.
Before embarking on the actual calculation, we ought to ponder one more point. We would like to solve a problem with periodic boundary conditions for strains and stresses in the $`x`$ direction. But periodicity of the strains does not imply periodicity of the displacements nor does it imply periodicity of the Goursat functions. On the other hand, the use of periodic functions greatly facilitates the derivation. As noted by Spencer and Meiron , we can express the Goursat functions by periodic functions $`\varphi _0`$ and $`\psi _0`$ via
$`\varphi (z)`$ $`=`$ $`\varphi _0(z),`$ (C16)
$`\psi (z)`$ $`=`$ $`\psi _0(z)z\varphi _0^{}(z).`$ (C17)
The mathematical problem is then to find two periodic functions $`\varphi _0`$ and $`\psi _0`$, analytic in the domain occupied by the solid, satisfying
$$\varphi _0(z)+(z\overline{z})\overline{\varphi _0^{}(z)}+\overline{\psi _0(z)}=\frac{1}{2}(z\overline{z})\sigma _0$$
(C18)
at the interface and remaining bounded for $`y\mathrm{}`$. The solution to this problem must be unique apart from possible additive constants to the functions $`\varphi _0(z)`$ and $`\psi _0(z)`$.
We transform to the $`\zeta `$ plane, using the analytic continuation of (C1)
$$z=\zeta iAe^{ik\zeta }iBe^{2ik\zeta }\omega (\zeta ).$$
(C19)
This maps the interface to the real axis and the solid to the half plane $`\eta <0`$. To designate functions in the $`\zeta `$ plane, we put a tilde on the letter they have in the $`z`$ plane:
$`\varphi _0(z)`$ $`=`$ $`\varphi _0(\omega (\zeta ))\stackrel{~}{\varphi }_0(\zeta ),`$ (C20)
$`\psi _0(z)`$ $`=`$ $`\stackrel{~}{\psi }_0(\zeta ).`$ (C21)
The derivative of $`\varphi _0`$ transforms as follows:
$$\varphi _0^{}(z)=\stackrel{~}{\varphi }_0^{}(\zeta )\frac{d\zeta }{dz}=\frac{\stackrel{~}{\varphi }_0^{}(\zeta )}{\omega ^{}(\zeta )}.$$
(C22)
Our task therefore is to construct two analytic functions satisfying
$$\stackrel{~}{\varphi }_0(\xi )+[\omega (\xi )\overline{\omega }(\xi )]\frac{\overline{\stackrel{~}{\varphi }_0^{}}(\xi )}{\overline{\omega ^{}}(\xi )}+\overline{\stackrel{~}{\psi }}_0(\xi )=\frac{\sigma _0}{2}[\omega (\xi )\overline{\omega }(\xi )]$$
(C23)
on the real axis ($`\eta =0`$) and remaining bounded as $`\eta \mathrm{}`$. $`\omega (\xi )`$ is given by Eq. (C19), hence
$$\omega ^{}(\xi )=1Ake^{ik\xi }2Bke^{2ik\xi }$$
(C24)
and Eq. (C23) becomes ($`A`$ and $`B`$ are real)
$`\stackrel{~}{\varphi }_0(\xi )i\left(Ae^{ik\xi }+Be^{2ik\xi }+Ae^{ik\xi }+Be^{2ik\xi }\right){\displaystyle \frac{\overline{\stackrel{~}{\varphi }_0^{}}(\xi )}{1Ake^{ik\xi }2Bke^{2ik\xi }}}+\overline{\stackrel{~}{\psi }_0}(\xi )`$ (C25)
$`=i{\displaystyle \frac{\sigma _0}{2}}\left(Ae^{ik\xi }+Be^{2ik\xi }+Ae^{ik\xi }+Be^{2ik\xi }\right).`$ (C26)
Note that $`\overline{\stackrel{~}{\varphi }_0^{}}(\xi )`$ and $`\overline{\stackrel{~}{\psi }_0}(\xi )`$ are functions that should be analytically continued to the upper half plane, for if $`\overline{\stackrel{~}{\varphi }_0^{}}(\overline{\zeta })`$ \[$`\overline{\stackrel{~}{\psi }_0}(\overline{\zeta })`$\] is analytic in the upper half plane, $`\stackrel{~}{\varphi }_0^{}(\zeta )`$ \[$`\stackrel{~}{\psi }_0(\zeta )`$\] is analytic in the lower half plane which is what we need. The basic idea in constructing the solution is to divide the terms in (C26) into two groups, one of which corresponds to functions analytic in the upper half plane, the other to those analytic in the lower half plane. The equality between these two groups implies that each of them is equal to a constant, which gives us two equations for the two functions sought. It is clear that $`\stackrel{~}{\varphi }_0(\xi )`$ belongs to the terms of (C26) that are analytic in the lower half plane (on replacement of $`\xi `$ by $`\zeta `$) whereas $`\overline{\stackrel{~}{\psi }_0}(\xi )`$ has to be analytic in the upper half plane. The difficult term is the middle one on the left hand side as it contains some expressions that are analytic and bounded in the upper half plane (e.g., $`e^{ik\xi }`$) but also some for which this is the case in the lower half plane (e.g., $`e^{ik\xi }`$). One way to proceed is to expand both $`\stackrel{~}{\varphi }_0`$ and $`\stackrel{~}{\psi }_0`$ in a series in powers of $`e^{ik\xi }`$ (a one-sided Fourier series, so to speak), to multiply Eq. (C26) by the denominator of the middle expression, and to separate terms with plus signs and minus signs in the exponents. This gives a two-termed recursion for $`\stackrel{~}{\varphi }_0`$ containing two constants that have to be determined from the analyticity properties. As it turns out, the series for $`\stackrel{~}{\varphi }_0`$ is finite, only the two first terms are nonzero. This suggests that a close look at Eq. (C26) would have revealed this property, allowing to avoid the tedious expansion procedure. Indeed, there is a more elegant way leading to this result. Its discovery is left as an exercise to the astute reader. Here, we simply take
$$\stackrel{~}{\varphi }_0(\xi )=a_1e^{ik\xi }+a_2e^{2ik\xi }$$
(C27)
as an ansatz. Inserting this into (C26), we get
$`a_1e^{ik\xi }+a_2e^{2ik\xi }+\left(Ae^{ik\xi }+Be^{2ik\xi }+Ae^{ik\xi }+Be^{2ik\xi }\right){\displaystyle \frac{\overline{a}_1ke^{ik\xi }+2\overline{a}_2ke^{2ik\xi }}{1Ake^{ik\xi }2Bke^{2ik\xi }}}`$ (C28)
$`i{\displaystyle \frac{\sigma _0}{2}}\left(Ae^{ik\xi }+Be^{2ik\xi }+Ae^{ik\xi }+Be^{2ik\xi }\right)=\overline{\stackrel{~}{\psi }_0}(\xi ).`$ (C29)
Since we have just $`\stackrel{~}{\psi }_0(\xi )`$ on the right-hand side, which must be analytically continued to the upper half plane and remain bounded there, the left-hand side must not contain, after substitution of $`\zeta `$ for $`\xi `$, any “dangerous” terms diverging as $`\eta \mathrm{}`$. Dangerous terms would obviously be the terms $`e^{ik\zeta }`$ and $`e^{2ik\zeta }`$ as well as the zeros of the denominator. Now, we have the condition $`Ak+2Bk<1`$ and we have $`|e^{ik\zeta }|1`$ for $`\eta 0`$. Therefore, the denominator is always different from zero in the upper half plane. All we have to do then is to choose $`a_1`$ and $`a_2`$ such that the dangerous terms cancel. This is straightforward for $`a_2`$, since there are only two terms containing $`e^{2ik\zeta }`$, after the prefactor of the second term has been multiplied with the numerator. There remains a term proportional to $`e^{ik\zeta }`$, however, in the numerator, which must also be canceled by the choice of $`a_1`$. The result is
$`a_1`$ $`=`$ $`{\displaystyle \frac{\sigma _0}{2}}i{\displaystyle \frac{A}{1Bk}},`$ (C30)
$`a_2`$ $`=`$ $`{\displaystyle \frac{\sigma _0}{2}}iB.`$ (C31)
From this, we immediately get $`\stackrel{~}{\varphi }_0`$ as
$$\stackrel{~}{\varphi }_0(\zeta )=\frac{\sigma _0}{2}i\left\{\frac{A}{1Bk}e^{ik\zeta }+Be^{2ik\zeta }\right\}.$$
(C32)
Once $`\stackrel{~}{\varphi }_0`$ is known, $`\stackrel{~}{\psi }_0`$ is obtained from (C26) or, using (C31), from (C29). We give the result for reference purposes, even though it will not be needed in the following
$`\stackrel{~}{\psi }_0(\zeta )`$ $`=`$ $`{\displaystyle \frac{\sigma _0}{2}}i{\displaystyle \frac{1}{1Ake^{ik\zeta }2Bke^{2ik\zeta }}}\{k(2B^2+A^2{\displaystyle \frac{1+Bk}{1Bk}})+{\displaystyle \frac{2ABk}{1Bk}}e^{ik\zeta }`$ (C34)
$`\text{+}(Ae^{ik\zeta }+Be^{2ik\zeta })(1+{\displaystyle \frac{ABk^2}{1Bk}}e^{ik\zeta })\}.`$
Note that only exponentials of negative multiples of $`ik\zeta `$ appear; the numerators are thus analytic and bounded in the lower half-plane and the denominator remains nonzero because of $`Ak+2Bk<1`$, hence $`\stackrel{~}{\psi }_0(\zeta )`$ satisfies all the required analyticity and boundedness conditions. In the case $`A=0`$ or $`B=0`$, these results reduce to those of Gao et al. for the simple cycloid (after transforming back to the nonperiodic Goursat functions).
We need only $`\stackrel{~}{\varphi }_0`$ to compute the tangential stresses on the boundary. Since we know the normal stresses to be equal to zero from the boundary condition, we can write $`\sigma _{tt}=\text{tr}\sigma =\sigma _{xx}^{(0)}+\sigma _{yy}^{(0)}+\sigma _0`$. We have
$$\sigma _{xx}^{(0)}+\sigma _{yy}^{(0)}=2\left[\frac{\stackrel{~}{\varphi }_0^{}(\zeta )}{\omega ^{}(\zeta )}+\frac{\overline{\stackrel{~}{\varphi }_0^{}(\zeta )}}{\overline{\omega ^{}(\zeta )}}\right],$$
(C35)
which, when specialized to the boundary, gives us
$$\sigma _{tt}=\frac{\sigma _0}{2}\{\frac{1+\frac{1+Bk}{1Bk}Ake^{ik\xi }+2Bke^{2ik\xi }}{1Ake^{ik\xi }2Bke^{2ik\xi }}+c.c.\}$$
(C36)
In the grooves of the pattern, we have $`k\xi =m\pi `$ hence $`e^{ik\xi }=(1)^m`$, $`e^{2ik\xi }=1`$, from which we obtain the tangential stress as
$$\sigma _{tt}=\sigma _0\frac{1+\frac{1+Bk}{1Bk}Ak(1)^m+2Bk}{1Ak(1)^m2Bk}.$$
(C37)
To obtain the normal velocity in the grooves, we need the curvature there. The curvature is easily calculated from the parametric representation of the double cycloid
$`\kappa `$ $`=`$ $`{\displaystyle \frac{x^{}(\xi )y^{\prime \prime }(\xi )y^{}(\xi )x^{\prime \prime }(\xi )}{\left(x^{}(\xi )^2+y^{}(\xi )^2\right)^{3/2}}}`$ (C38)
$`=`$ $`{\displaystyle \frac{k^3\left[A^2+8B^2+6AB\mathrm{cos}k\xi \right]k^2\left[A\mathrm{cos}k\xi +4B\mathrm{cos}2k\xi \right]}{\left(A^2k^2+4B^2k^2+4ABk^2\mathrm{cos}k\xi +12Ak\mathrm{cos}k\xi 4Bk\mathrm{cos}2k\xi \right)^{3/2}}}`$ (C39)
and its value in the bottom of the grooves is ($`\mathrm{cos}k\xi =\pm 1`$)
$$\kappa =\frac{k^2\left[A(1)^m+4B\right]}{\left[1Ak(1)^m2Bk\right]^2}.$$
(C40)
Note that as a cusp is approached ($`Ak+2Bk1`$), $`\sigma _{tt}^2`$ and the curvature diverge with the same denominator, for even $`m`$.
Inserting (C37) and (C40) into Eq. (44) for the nondimensional normal velocity and neglecting the gravity term we obtain
$$\stackrel{~}{v}_n=\frac{1}{2\left[1Ak(1)^m2Bk\right]^2}\left\{\left(1+2Bk+\frac{1+Bk}{1Bk}Ak(1)^m\right)^22\mathrm{}_1k^2\left[A(1)^m+4B\right]\right\},$$
(C41)
where we have nondimensionalized the curvature via multiplication by $`\mathrm{}_1`$. Setting $`\alpha =2\mathrm{}_1k`$, we generate Eq. (53). Note that if we consider the amplitude $`A`$ to be a small perturbation of a cycloid determined by $`B`$, the basic wave number is $`2k`$, not $`k`$. Then $`\alpha `$ is just the ratio of the basic wave number and the wave number of the fastest-growing mode.
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# 1 Introduction
## 1 Introduction
Vector meson production in lepton-proton collisions is a powerful probe to investigate the nature of diffraction. At HERA, because of the wide kinematic ranges in the photon virtuality, $`Q^2`$, and in the hadronic centre of mass energy, $`W`$, the details of the production mechanism can be studied. It is also possible to select different vector mesons, allowing the cross section for different quark types to be studied. Recent measurements of $`\rho `$ meson electroproduction for high $`Q^2`$ values ($`Q^2`$ $`\text{ }>`$ 10 $`\mathrm{GeV}^2`$) and of $`J/\psi `$ meson photo– and electroproduction show a strong energy dependence of the $`\gamma ^{}pVp`$ cross sections. This behaviour indicates that the mass of the $`c`$ quark or a high $`Q^2`$ value provides a hard scale in the interaction, and we study the elastic cross sections as a function of the scale ($`Q^2+M_V^2`$), where $`M_V`$ is the mass of the vector meson.
This paper presents a measurement of elastic $`\varphi `$ meson electroproduction
$$e^++pe^++\varphi +p;\varphi K^++K^{},$$
(1)
in the $`Q^2`$ range from 1 to 15 $`\mathrm{GeV}^2`$, and in the $`W`$ range from 40 to 130 GeV. The data were obtained with the H1 detector in two running periods when the HERA collider was operated with 820 GeV protons and 27.5 GeV positrons. A low $`Q^2`$ data set ($`1<Q^2<5`$ $`\mathrm{GeV}^2`$) with integrated luminosity of 125 $`\mathrm{nb}^1`$ was obtained from a special run in 1995, with the $`ep`$ interaction vertex shifted by 70 cm in the outgoing proton beam direction. This results in a higher acceptance for low $`Q^2`$ production. A larger sample of integrated luminosity of 3$`\mathrm{pb}^1`$ with $`2.5<Q^2<15`$ $`\mathrm{GeV}^2`$ was obtained in 1996 under normal running conditions. The present measurements provide detailed new information in the region $`1\text{ }<Q^2\text{ }<6`$ $`\mathrm{GeV}^2`$ and they increase the precision of the H1 measurement of $`\varphi `$ electroproduction with $`Q^2>6`$ $`\mathrm{GeV}^2`$, which was first performed using data collected in 1994 . They are compared to results of the ZEUS experiment in photoproduction and at $`Q^2>7`$ $`\mathrm{GeV}^2`$ . The elastic $`\varphi `$ meson cross section is also compared to elastic $`\rho `$ , $`\omega `$ , $`J/\psi `$ and $`\mathrm{{\rm Y}}`$ meson production results from H1 and ZEUS.
The event selection and the $`K^+K^{}`$ mass distribution is presented in section 2. In section 3, the elastic $`\varphi `$ cross section is presented as a function of $`Q^2`$ and $`W`$. In order to minimise the uncertainties, the cross section is measured as a ratio to elastic $`\rho `$ production, and the absolute elastic $`\varphi `$ cross section is then extracted using the results for $`\rho `$ production from . A compilation of the $`\rho `$, $`\omega `$, $`\varphi `$, $`J/\psi `$, and $`\mathrm{{\rm Y}}`$ cross sections is presented as a function of ($`Q^2`$ \+ $`M_V^2`$). The $`t`$ dependence of the elastic $`\varphi `$ cross section is analysed in section 4. A detailed analysis of the photon and $`\varphi `$ meson polarisations is performed in section 5 and the 15 spin density matrix elements are extracted. The ratio $`R`$ of the longitudinal to transverse $`\varphi `$ cross sections is obtained as a function of $`Q^2`$. A compilation of the measurements of $`R`$ for elastic $`\rho `$, $`\varphi `$ and $`J/\psi `$ meson production is presented as a function of $`Q^2`$/$`M_V^2`$.
The present analysis uses to a large extent the techniques described in the H1 publication on elastic $`\rho `$ production .
## 2 Data selection
Elastic $`\varphi `$ electroproduction events are selected on the basis of their topology in the H1 detector<sup>1</sup><sup>1</sup>1 A detailed description of the H1 detector can be found in .. They must have a positron candidate and two oppositely charged hadron candidates, originating from a vertex situated in the nominal $`e^+p`$ interaction region, with K<sup>+</sup>K<sup>-</sup> invariant mass in the range from 1.00 to 1.04 GeV. The scattered positron is identified as an electromagnetic cluster of energy larger than 15 GeV detected in the H1 backward electromagnetic calorimeter SPACAL <sup>2</sup><sup>2</sup>2 H1 uses a right-handed coordinate system with the $`z`$ axis taken along the beam direction, the $`+z`$ direction being that of the outgoing proton beam. The $`x`$ axis points towards the centre of the HERA ring.. The two hadron candidates are recognised as tracks of opposite signs, with a momentum transverse to the beam direction larger than 100 MeV, reconstructed in the H1 central tracking detector with a polar angle in the range from 20 to 160. The vertex must lie within 30 cm along the beam axis from the nominal interaction point. The nature of the hadrons is not explicitly identified. Their charge and momentum are measured in the central part of the detector by means of a uniform 1.15 T magnetic field. No other activity must be observed in the detector since the scattered proton remains in the beam pipe and is not detected because of the small momentum transfer to the target in diffractive interactions. Events were therefore rejected if there were signals in the forward part of the detector (forward muon and forward proton tagger detectors) and if there were clusters in the liquid argon calorimeter with an energy above 0.5 GeV not associated with the hadron candidates. To reduce effects of QED radiative corrections, the selected events have to satisfy $`_{e,h}(Ep_z)`$ $`>`$ 45 GeV.
The $`Q^2`$ variable is reconstructed from the incident electron beam energy and the polar angles of the positron and of the $`\varphi `$ meson candidates (double angle method ). The $`W`$ variable is reconstructed using in addition the energy and the longitudinal momentum of the $`\varphi `$ meson candidate.
The variable $`t`$ is the square of the four-momentum transfer to the target proton. At HERA energies, to a very good precision, the absolute value of $`t`$ is equal to the square of the transverse momentum of the outgoing proton. The latter is computed, under the assumption that the selected event corresponds to reaction (1), as the square of the vector sum of the transverse momenta of the $`\varphi `$ meson candidate and of the scattered positron. Events with $`|t|`$ $`<`$ 0.5 $`\mathrm{GeV}^2`$ are selected in order to reduce the remaining production of proton dissociation events which have a flatter $`t`$ distribution, and to suppress the production of hadron systems of which the $`\varphi `$ is only part and in which the remaining particles were not detected.
The distribution of $`m_{KK}`$, the two particle invariant mass computed under the assumption that the hadron candidates are kaons, is presented in Fig. 1a and Fig. 1b for $`m_{KK}`$ $`<`$ 1.12 GeV and for $`m_{KK}`$ $`<`$ 2.00 GeV, respectively. A clear $`\varphi `$ signal is observed in the data, with 424 events in the range 1.00 $`<`$ $`m_{KK}`$ $`<`$ 1.04 GeV.
The main backgrounds to reaction (1) are due to diffractive $`\varphi `$ events in which the proton is excited into a system of higher mass which subsequently dissociates, and to the elastic production of $`\rho `$ and $`\omega `$ vector mesons. The other backgrounds (other $`\varphi `$ decay channels, higher mass resonances or non resonant production) are estimated to be less than a few percent. The fraction of proton dissociation background is assumed to be the same for $`\varphi `$ as for $`\rho `$ meson production and is taken to be 11$`\pm `$5 % as in . The background due to $`\rho `$ and $`\omega `$ production is estimated using the DIFFVM simulation . The DIFFVM Monte Carlo simulation program is based on Regge theory and on the vector meson dominance model. The $`\rho `$ and $`\omega `$ backgrounds are normalised to the $`m_{\pi \pi }`$ distribution observed in the data, where $`m_{\pi \pi }`$ is the invariant mass computed under the pion hypothesis for the hadron candidates. This is shown in Fig. 1c after the $`\varphi `$ signal has been removed by selecting $`m_{KK}`$ $`>`$ 1.04 GeV. The background under the $`\varphi `$ peak from $`\rho `$ and $`\omega `$ meson production is $`Q^2`$ dependent and varies from 15 % to 4 %. For the full sample (2.5 $`<`$ $`Q^2`$ $`<`$ 15 $`\mathrm{GeV}^2`$) this background is found to be 9$`\pm `$5 %.
The data are corrected for acceptances, efficiencies and detector resolution effects using the DIFFVM Monte Carlo simulation. The response of the H1 detector is fully simulated.
## 3 Elastic cross section
The elastic $`\mathit{\varphi }`$ meson production cross section is obtained by first measuring the ratio of the $`\mathit{\varphi }`$ to $`𝝆`$ cross sections and then using the $`𝝆`$ cross sections which were precisely measured as described in . In the ratio, several uncertainties cancel, most notably the luminosity uncertainty, the contribution of the proton dissociation background and the trigger efficiency, which is very similar for both data samples since it is mostly based on the positron detection in the SPACAL. The remaining corrections account for the mass selection range and the differences in acceptances and $`𝒕`$ distribution. The correction for the accepted mass range (0.6 $`\mathbf{<}`$ $`𝒎_{𝝅𝝅}`$ $`\mathbf{<}`$ 1.1 GeV) in the $`𝝆`$ sample is 1.16 $`\mathbf{\pm }`$ 0.02 $`{}_{\mathbf{}\mathbf{0.00}}{}^{}{}_{}{}^{\mathbf{+}\mathbf{0.05}}`$ ; the correction for the mass range (1.00 $`\mathbf{<}`$ $`𝒎_{𝑲𝑲}`$ $`\mathbf{<}`$ 1.04 GeV) in the $`\mathit{\varphi }`$ sample is estimated to be 1.03 $`\mathbf{\pm }`$ 0.01 using the DIFFVM simulation. In both samples, hadron tracks must be detected in the central tracker with $`\mathrm{𝟐𝟎}^{\mathbf{}}\mathbf{<}𝜽\mathbf{<}\mathrm{𝟏𝟔𝟎}^{\mathbf{}}`$. Differences in the acceptances for the two samples, due to the different decay hadron and vector meson masses, are estimated as a function of $`𝑸^\mathrm{𝟐}`$, $`𝑾`$ and $`𝒕`$ using the DIFFVM simulations for $`𝝆`$ and $`\mathit{\varphi }`$ production. Differences in the detector efficiency for pions and kaons are taken into account in the detector simulation. Finally the correction for events with $`\mathbf{|}𝒕\mathbf{|}`$ $`\mathbf{>}`$ 0.5 $`\mathrm{𝐆𝐞𝐕}^\mathrm{𝟐}`$ in the $`𝝆`$ and $`\mathit{\varphi }`$ samples is estimated by assuming an exponentially falling $`\mathbf{|}𝒕\mathbf{|}`$ distribution, using recent measurements for the slope parameter of the exponential . The correction factor on the $`\mathit{\varphi }`$/$`𝝆`$ cross section ratio is 1.03 $`\mathbf{\pm }`$ 0.02, independent of $`𝑸^\mathrm{𝟐}`$. The branching ratios of 0.49 and 1.0 were used for the decays $`\mathit{\varphi }`$ $`\mathbf{}`$ $`𝑲^\mathbf{+}𝑲^{\mathbf{}}`$ and $`𝝆`$ $`\mathbf{}`$ $`𝝅^\mathbf{+}𝝅^{\mathbf{}}`$ respectively.
Systematic errors on the measurement of the cross section ratio are estimated by varying all corrections within the errors. In addition, in both the $`𝝆`$ and $`\mathit{\varphi }`$ simulations the cross section dependence on $`𝑸^\mathrm{𝟐}`$, $`𝑾`$, $`𝒕`$ and the vector meson angular decay distributions were varied by amounts allowed by the present and most recent measurements .
The $`𝑸^\mathrm{𝟐}`$ dependence of the $`\mathit{\varphi }`$ to $`𝝆`$ elastic cross section ratio is presented in Fig. 2 together with previous H1 and ZEUS results. The values of the ratio are given in Table 1. The present measurements confirm the significant rise of the cross section ratio with $`𝑸^\mathrm{𝟐}`$. As $`𝑸^\mathrm{𝟐}`$ increases, the HERA cross section ratios approach the value $`\mathrm{𝟐}\mathbf{/}\mathrm{𝟗}`$ expected from quark charge counting and SU(5). It should be noted that calculations based on perturbative QCD predict that the cross section ratio should exceed this value at very large $`𝑸^\mathrm{𝟐}`$ . The $`𝑾`$ dependence of the $`\mathit{\varphi }`$ to $`𝝆`$ elastic cross section ratio is measured in the range 40 $`\mathbf{<}`$ W $`\mathbf{<}`$ 130 GeV and is observed to be constant, within the experimental uncertainties.
To extract the $`𝜸^{\mathbf{}}𝒑\mathbf{}\mathit{\varphi }𝒑`$ cross section, the measurement of the $`\mathit{\varphi }\mathbf{/}𝝆`$ cross section ratio is multiplied by the $`𝜸^{\mathbf{}}𝒑\mathbf{}𝝆𝒑`$ cross section calculated from the fit in . The values are given in Table 1. The systematic errors on the $`\mathit{\varphi }`$ cross section measurement include the systematic errors on the ratio of $`\mathit{\varphi }`$ to $`𝝆`$ cross sections, as well as an 8.4 % contribution coming from the parametrisation error in the fit of the $`𝝆`$ cross section (see ), added in quadrature.
In Fig. 3, the cross section for the elastic production of $`\mathit{\varphi }`$ mesons (full squares) is presented together with other vector mesons $`𝑽`$ and for various values of $`𝑸^\mathrm{𝟐}`$, as a function of the variable ($`𝑸^\mathrm{𝟐}`$ \+ $`𝑴_𝑽^\mathrm{𝟐}`$). The data in Fig. 3 compile the HERA measurements of the $`𝜸^{\mathbf{}}𝒑\mathbf{}𝑽𝒑`$ cross sections (see also ). The cross sections were scaled by SU(5) factors, according to the quark charge content of the vector meson, which amount to 1 for the $`𝝆`$, 9 for the $`𝝎`$, 9/2 for the $`\mathit{\varphi }`$, 9/8 for the $`𝑱\mathbf{/}𝝍`$ and 9/2 for the $`𝚼`$ meson. The cross sections are measured at $`𝑾`$ = 75 GeV, or are moved to that value according to the parametrisation $`𝝈\mathbf{}𝑾^𝜹`$, using the $`𝜹`$ value measured by the corresponding experiment. The ZEUS $`𝝆`$ and $`\mathit{\varphi }`$ cross sections were corrected ($`\text{ }\mathbf{}\mathbf{<}`$ 7 %) for the unmeasured signal with $`\mathbf{|}𝒕\mathbf{|}`$ $`\mathbf{>}`$ 0.5 (or 0.6) $`\mathrm{𝐆𝐞𝐕}^\mathrm{𝟐}`$ by assuming a simple exponential fall of d$`𝝈\mathbf{/}𝒅𝒕\mathbf{}𝒆^{𝒃𝒕}`$. In this procedure the observed $`𝑸^\mathrm{𝟐}`$ dependence of the $`𝒃`$ slope was taken into account.
Within the experimental errors, the total cross sections for vector meson production, including the SU(5) normalisation factors, appear to lie on a universal curve when plotted as a function of the scale $`\mathbf{(}𝑸^\mathrm{𝟐}\mathbf{+}𝑴_𝑽^\mathrm{𝟐}\mathbf{)}`$, except possibly for the $`𝚼`$ photoproduction<sup>3</sup><sup>3</sup>3The cross sections $`\sigma (\gamma p\mathrm{{\rm Y}}(1\mathrm{S})p)`$ measured by H1 and ZEUS at $`W`$ = 143 and 120 GeV respectively , were moved to the value $`W`$ = 75 GeV using the parametrisation $`\sigma W^\delta `$, with $`\delta `$ = 1.7. This high value of the parameter $`\delta `$ comes from the prediction of . Note that if the value $`\delta `$ = 0.8 is used (a value measured in case of $`J/\psi `$ photoproduction), the cross sections increase by a factor 1.5 for ZEUS and 1.8 for H1.. A fit performed on the H1 and ZEUS $`𝝆`$ data using the parametrisation $`𝝈\mathbf{=}𝒂_\mathrm{𝟏}\mathbf{(}𝑸^\mathrm{𝟐}\mathbf{+}𝑴_𝑽^\mathrm{𝟐}\mathbf{+}𝒂_\mathrm{𝟐}\mathbf{)}^{𝒂_\mathrm{𝟑}}`$, with $`𝒂_\mathrm{𝟏}`$ = 10689 $`\mathbf{\pm }`$ 165 nb, $`𝒂_\mathrm{𝟐}`$ = 0.42 $`\mathbf{\pm }`$ 0.09 $`\mathrm{𝐆𝐞𝐕}^\mathrm{𝟐}`$ and $`𝒂_\mathrm{𝟑}`$ = – 2.37 $`\mathbf{\pm }`$ 0.10 ($`𝝌^\mathrm{𝟐}\mathbf{/}𝒏𝒅𝒇`$ = 0.67) is shown as the curve in Fig. 3. The ratio of the $`𝝎`$, $`\mathit{\varphi }`$ and $`𝑱\mathbf{/}𝝍`$ cross sections to this parametrisation is presented in the insert of Fig. 3. Note that the universal $`\mathbf{(}𝑸^\mathrm{𝟐}\mathbf{+}𝑴_𝑽^\mathrm{𝟐}\mathbf{)}`$ dependence is for the total cross section measurements only. The separate behaviour of the longitudinal and transverse cross sections is described in ref. .
## 4 Dependence on $`𝒕`$
In this section and the following one, the elastic $`\mathit{\varphi }`$ meson production is studied using the $`\mathit{\varphi }`$ sample defined above, with the additional selection: the centre of gravity of the scattered positron cluster was required to lie outside the innermost part of the SPACAL calorimeter $`\mathbf{}\mathrm{𝟏𝟔}\mathbf{<}𝒙\mathbf{<}\mathrm{𝟖}`$ cm and $`\mathbf{}\mathrm{𝟖}\mathbf{<}𝒚\mathbf{<}\mathrm{𝟏𝟔}`$ cm in order to obtain good ($`\mathbf{>}`$ 95 %) and uniform trigger efficiency. The number of elastic $`\mathit{\varphi }`$ candidates is then reduced to 221 events for 2.5 $`\mathbf{<}`$ $`𝑸^\mathrm{𝟐}`$ $`\mathbf{<}`$ 15 $`\mathrm{𝐆𝐞𝐕}^\mathrm{𝟐}`$.
The measured $`\mathbf{|}𝒕\mathbf{|}`$ dependence is shown in Fig. 4 and the characteristic falling exponential distribution is observed. To take into account the contribution of different backgrounds, the $`\mathbf{|}𝒕\mathbf{|}`$ distribution is fitted by the sum of three exponentials corresponding to the elastic $`\mathit{\varphi }`$ component, the diffractive $`\mathit{\varphi }`$ component with proton dissociation and the $`𝝎`$ and $`𝝆`$ production. The elastic $`\mathit{\varphi }`$ component is fitted with a free normalisation and a free slope parameter $`𝒃`$, whereas the other contributions are fixed to their calculated values. The contribution of diffractive $`\mathit{\varphi }`$ events with proton dissociation of $`\mathrm{𝟏𝟏}\mathbf{\pm }\mathrm{𝟓}\mathbf{\%}`$ of the elastic signal and a slope parameter of $`\mathbf{2.5}\mathbf{\pm }\mathbf{1.0}`$ $`\mathrm{𝐆𝐞𝐕}^\mathbf{}\mathrm{𝟐}`$, was taken from . The $`𝝎`$ and $`𝝆`$ background contributions, amounting to $`\mathrm{𝟗}\mathbf{\pm }\mathrm{𝟓}\mathbf{\%}`$ of the signal (see section 2), have an effective slope parameter $`𝒃`$ = 2.9 $`\mathbf{\pm }`$ 0.6 $`\mathrm{𝐆𝐞𝐕}^\mathbf{}\mathrm{𝟐}`$, computed using the DIFFVM simulation.
The fitted exponential slope parameter for elastic $`\mathit{\varphi }`$ events is found to be $`𝒃\mathbf{=}`$5.8$`\mathbf{\pm }`$0.5 (stat.) $`\mathbf{\pm }`$0.6 (syst.) $`\mathrm{𝐆𝐞𝐕}^\mathbf{}\mathrm{𝟐}`$, for an average $`𝑸^\mathrm{𝟐}`$ value of 4.5 $`\mathrm{𝐆𝐞𝐕}^\mathrm{𝟐}`$ and $`\mathbf{}𝑾\mathbf{}`$=75 GeV. The systematic error is computed by varying the amounts of the background contributions and their slopes within the quoted errors, and by varying the binning and the limits of the fit. The effect of the QED radiative corrections on the $`𝒃`$ measurement is estimated using the simulation DIFFVM including a HERACLES interface, and is found to decrease the value of the $`𝒃`$ measurement by 0.13 $`\mathrm{𝐆𝐞𝐕}^\mathbf{}\mathrm{𝟐}`$ (the $`𝒃`$ value given above is not corrected for this effect). This result can be compared with other measurements, $`𝒃\mathbf{=}\mathbf{7.3}\mathbf{\pm }\mathbf{1.0}\mathbf{\pm }\mathbf{0.8}`$ $`\mathrm{𝐆𝐞𝐕}^\mathbf{}\mathrm{𝟐}`$ in photoproduction and $`𝒃\mathbf{=}\mathbf{5.2}\mathbf{\pm }\mathbf{1.6}\mathbf{\pm }\mathbf{1.0}`$ $`\mathrm{𝐆𝐞𝐕}^\mathbf{}\mathrm{𝟐}`$ for $`\mathbf{}𝑸^\mathrm{𝟐}\mathbf{}`$ = 10 $`\mathrm{𝐆𝐞𝐕}^\mathrm{𝟐}`$ . The data are consistent with a decrease of the slope parameter as $`𝑸^\mathrm{𝟐}`$ increases; this would be expected from the decrease of the transverse size of the virtual photon.
The value of the $`𝒃`$ slope parameter is in agreement within the errors with the one obtained in elastic $`𝝆`$ meson production: $`𝒃`$ = 5.5 $`\mathbf{\pm }`$ 0.5 (stat.) $`{}_{\mathbf{}\mathbf{0.2}}{}^{}{}_{}{}^{\mathbf{+}\mathbf{0.5}}`$ (syst.) $`\mathrm{𝐆𝐞𝐕}^\mathbf{}\mathrm{𝟐}`$, at $`𝑸^\mathrm{𝟐}`$ = 4.8 $`\mathrm{𝐆𝐞𝐕}^\mathrm{𝟐}`$ .
## 5 Polarisation studies
The study of the angular distributions of the production and decay of the $`\mathit{\varphi }`$ meson provides information on the photon and $`\mathit{\varphi }`$ meson polarisation states. In the helicity system , three angles are defined as follows. The angle $`𝚽`$, defined in the hadronic centre of mass system (cms), is the azimuthal angle between the electron scattering plane and the plane containing the $`\mathit{\varphi }`$ meson and the scattered proton. The $`\mathit{\varphi }`$ meson decay is described by the polar angle $`𝜽`$ and the azimuthal angle $`𝝋`$ of the positive kaon in the $`𝑲^\mathbf{+}𝑲^{\mathbf{}}`$ rest frame, with the quantisation axis taken as the direction opposite to that of the outgoing proton in the hadronic cms (the so called helicity frame). Details of the kinematics and the mathematical formalism can be found in and . The normalised angular decay distribution $`𝑭`$($`\mathrm{𝐜𝐨𝐬}𝜽`$, $`𝝋`$, $`𝚽`$) is expressed as a function of 15 spin density matrix elements corresponding to different bilinear combinations of the helicity amplitudes $`𝑻_{𝝀_\mathit{\varphi }\mathbf{,}𝝀_𝜸}`$, where $`𝝀_\mathit{\varphi }`$ and $`𝝀_𝜸`$ are the helicities of the $`\mathit{\varphi }`$ meson and of the photon, respectively. In the case of $`𝒔`$-channel helicity conservation (SCHC), the helicities of the $`\mathit{\varphi }`$ meson and the photon are equal, only the amplitudes $`𝑻_{\mathrm{𝟎𝟎}}`$, $`𝑻_{\mathrm{𝟏𝟏}}`$, and $`𝑻_{\mathbf{}\mathrm{𝟏}\mathbf{}\mathrm{𝟏}}`$ are different from zero and 10 of the 15 matrix elements are zero.
The matrix elements are measured using projections of the decay angular distribution onto orthogonal trigonometric functions of the angles $`𝜽`$, $`𝝋`$ and $`𝚽`$ . The results are presented in Fig. 5 in two $`𝑸^\mathrm{𝟐}`$ bins: $`\mathbf{2.5}\mathbf{<}𝑸^\mathrm{𝟐}\mathbf{<}\mathbf{4.5}`$ $`\mathrm{𝐆𝐞𝐕}^\mathrm{𝟐}`$ and $`\mathbf{4.5}\mathbf{<}𝑸^\mathrm{𝟐}\mathbf{<}\mathrm{𝟏𝟓}`$ $`\mathrm{𝐆𝐞𝐕}^\mathrm{𝟐}`$. In Fig. 5, the results are not corrected for the small effects due to proton dissociation, $`𝝎`$ and $`𝝆`$ production backgrounds and radiative effects.
The matrix elements generally follow the SCHC predictions, except for the elements $`𝒓_{\mathrm{𝟎𝟎}}^\mathrm{𝟏}`$ and $`𝒓_{\mathrm{𝟎𝟎}}^\mathrm{𝟓}`$, which may indicate a small violation of SCHC. The matrix element $`𝒓_{\mathrm{𝟎𝟎}}^\mathrm{𝟓}`$ is proportional to the product $`𝑻_{\mathrm{𝟎𝟎}}^{\mathbf{}}𝑻_{\mathrm{𝟎𝟏}}`$ of helicity amplitudes, the dominant SCHC violating amplitude being $`𝑻_{\mathrm{𝟎𝟏}}`$ ($`𝝀_\mathit{\varphi }`$ = 0 and $`𝝀_𝜸`$ = 1).
Predictions from recent models based on perturbative QCD are compared to the measurement of the 15 matrix elements. The models are expected to be valid at high $`𝑸^\mathrm{𝟐}`$ (providing a scale for the perturbative expansion) and at high energy: $`𝑾^\mathrm{𝟐}\mathbf{}𝑸^\mathrm{𝟐}\mathbf{}𝚲_{\mathrm{𝐐𝐂𝐃}}^\mathrm{𝟐}`$. The $`\mathit{\varphi }`$ meson production is factorised, in the proton rest frame, into three parts involving different time scales: the fluctuation of the photon into a $`𝒒\overline{𝒒}`$ state, at a large distance from the target, the hard scattering of the $`𝒒\overline{𝒒}`$ pair with the proton, modelled as two-gluon exchange, and the $`𝒒\overline{𝒒}`$ pair recombination into a $`\mathit{\varphi }`$ meson wave. The amplitudes are computed separately for the different helicities of the photon and the $`\mathit{\varphi }`$ meson. In models , the gluon density in the proton is used for the computation of the hard scattering amplitude. Differences between the models are related to the way of introducing quark off-shellness and Fermi motion. All models describe the data relatively well, predicting in particular a non-zero value for the $`𝒓_{\mathrm{𝟎𝟎}}^\mathrm{𝟓}`$ matrix element (see Fig. 5). The model gives a poorer description of the $`𝑸^\mathrm{𝟐}`$ dependence of the $`𝒓_{\mathrm{𝟎𝟎}}^{\mathrm{𝟎𝟒}}`$, $`𝒓_{\mathrm{𝟏}\mathbf{}\mathrm{𝟏}}^\mathrm{𝟏}`$ and Im $`𝒓_{\mathrm{𝟏}\mathbf{}\mathrm{𝟏}}^\mathrm{𝟐}`$ matrix elements, which are correlated, than the models of .
Another way to study the violation of SCHC is to measure the $`𝚽`$ angular distribution: $`𝑭\mathbf{(}𝚽\mathbf{)}\mathbf{}\mathrm{𝟏}\mathbf{}𝜺\mathrm{𝐜𝐨𝐬}\mathrm{𝟐}𝚽\mathbf{(}\mathrm{𝟐}𝒓_{\mathrm{𝟏𝟏}}^\mathrm{𝟏}\mathbf{+}𝒓_{\mathrm{𝟎𝟎}}^\mathrm{𝟏}\mathbf{)}\mathbf{+}\sqrt{\mathrm{𝟐}𝜺\mathbf{(}\mathrm{𝟏}\mathbf{+}𝜺\mathbf{)}}\mathrm{𝐜𝐨𝐬}𝚽\mathbf{(}\mathrm{𝟐}𝒓_{\mathrm{𝟏𝟏}}^\mathrm{𝟓}\mathbf{+}𝒓_{\mathrm{𝟎𝟎}}^\mathrm{𝟓}\mathbf{)}`$, where $`𝜺`$ is the polarisation parameter of the virtual photon. In the case of SCHC, this distribution is predicted to be uniform, the matrix elements $`𝒓_{\mathrm{𝟏𝟏}}^\mathrm{𝟏}`$, $`𝒓_{\mathrm{𝟎𝟎}}^\mathrm{𝟏}`$, $`𝒓_{\mathrm{𝟏𝟏}}^\mathrm{𝟓}`$ and $`𝒓_{\mathrm{𝟎𝟎}}^\mathrm{𝟓}`$ being zero.
The $`𝚽`$ distribution for the elastic $`\mathit{\varphi }`$ meson production is presented in Fig. 6a. The distribution is corrected for the presence of $`𝝆`$ and $`𝝎`$ backgrounds (hashed area). The result of the fit to the function $`𝑭\mathbf{(}𝚽\mathbf{)}`$ is given as the full line and shows a clear $`\mathrm{𝐜𝐨𝐬}𝚽`$ dependence with a small $`\mathrm{𝐜𝐨𝐬}\mathrm{𝟐}𝚽`$ modulation. The extracted values for the combination $`\mathbf{(}\mathrm{𝟐}𝒓_{\mathrm{𝟏𝟏}}^\mathrm{𝟓}\mathbf{+}𝒓_{\mathrm{𝟎𝟎}}^\mathrm{𝟓}\mathbf{)}`$ are presented in Fig. 6b for three bins in $`𝑸^\mathrm{𝟐}`$. The $`𝝆`$ and $`𝝎`$ background subtraction in the $`𝚽`$ distribution reduces the value of the combination $`\mathbf{(}\mathrm{𝟐}𝒓_{\mathrm{𝟏𝟏}}^\mathrm{𝟓}\mathbf{+}𝒓_{\mathrm{𝟎𝟎}}^\mathrm{𝟓}\mathbf{)}`$ by 13 % (around half of the statistical error). In Fig. 6b, the effect of QED radiative corrections on the measurement of the combination $`\mathbf{(}\mathrm{𝟐}𝒓_{\mathrm{𝟏𝟏}}^\mathrm{𝟓}\mathbf{+}𝒓_{\mathrm{𝟎𝟎}}^\mathrm{𝟓}\mathbf{)}`$ was taken into account. This effect was estimated using the DIFFVM simulation including a HERACLES interface, and reduces the observed value of the combination $`\mathbf{(}\mathrm{𝟐}𝒓_{\mathrm{𝟏𝟏}}^\mathrm{𝟓}\mathbf{+}𝒓_{\mathrm{𝟎𝟎}}^\mathrm{𝟓}\mathbf{)}`$ by 17 %. The combination $`\mathbf{(}\mathrm{𝟐}𝒓_{\mathrm{𝟏𝟏}}^\mathrm{𝟓}\mathbf{+}𝒓_{\mathrm{𝟎𝟎}}^\mathrm{𝟓}\mathbf{)}`$ obtained from the fit deviates significantly from the zero prediction of SCHC (a 5 $`𝝈`$ effect). The values of the combination $`\mathbf{(}\mathrm{𝟐}𝒓_{\mathrm{𝟏𝟏}}^\mathrm{𝟓}\mathbf{+}𝒓_{\mathrm{𝟎𝟎}}^\mathrm{𝟓}\mathbf{)}`$ are similar to the ones obtained in case of elastic $`𝝆`$ meson production .
From the measurement of the spin density matrix elements, the ratio $`𝑹`$ of cross sections for $`\mathit{\varphi }`$ meson production by longitudinal and transverse virtual photons can be extracted. As the SCHC violating amplitudes are small compared to the helicity conserving amplitudes, one can make<sup>4</sup><sup>4</sup>4 The effect of SCHC violation on the measurement of $`R`$ is of the order of 3 %. the SCHC approximation in order to estimate $`𝑹`$, which is then obtained directly from the measurement of the matrix element $`𝒓_{\mathrm{𝟎𝟎}}^{\mathrm{𝟎𝟒}}`$ .
The $`𝑸^\mathrm{𝟐}`$ dependence of $`𝑹`$ is presented in Fig. 7a, together with other measurements performed under the SCHC approximation , see also table 2. It is observed that $`𝑹`$ rises steeply with $`𝑸^\mathrm{𝟐}`$, and that the longitudinal cross section dominates over the transverse cross section for $`𝑸^\mathrm{𝟐}\text{ }\mathbf{}\mathbf{>}`$$`\mathrm{𝐆𝐞𝐕}^\mathrm{𝟐}`$. The rise of $`𝑹`$ with $`𝑸^\mathrm{𝟐}`$ for $`\mathit{\varphi }`$ meson production is slower than for the $`𝝆`$ meson . However, when plotted as a function of $`𝑸^\mathrm{𝟐}\mathbf{/}𝑴_𝑽^\mathrm{𝟐}`$, the ratio $`𝑹`$ appears to show a common dependence for different vector mesons , see Fig. 7b (for further details see ref. ).
## 6 Summary
The elastic electroproduction of $`\mathit{\varphi }`$ mesons has been studied with the H1 detector in the kinematic range 1$`\mathbf{<}`$ $`𝑸^\mathrm{𝟐}`$ $`\mathbf{<}`$ 15 $`\mathrm{𝐆𝐞𝐕}^\mathrm{𝟐}`$ and 40 $`\mathbf{<}`$ $`𝑾`$ $`\mathbf{<}`$ 130 GeV. The $`𝑸^\mathrm{𝟐}`$ dependence of the cross section is presented in the form of the ratio to the elastic $`𝝆`$ meson cross section. A significant rise of the ratio with $`𝑸^\mathrm{𝟐}`$ is observed. The elastic $`\mathit{\varphi }`$ meson cross section is extracted using recent H1 results of elastic $`𝝆`$ meson production. A compilation of the elastic $`𝝆`$, $`𝝎`$, $`\mathit{\varphi }`$, $`𝑱\mathbf{/}𝝍`$ and $`𝚼`$ meson cross sections, scaled by SU(5) factors, is presented as a function of ($`𝑸^\mathrm{𝟐}`$ \+ $`𝑴_𝑽^\mathrm{𝟐}`$). A common dependence is observed within experimental errors. The $`\mathbf{|}𝒕\mathbf{|}`$ dependence of the elastic $`\mathit{\varphi }`$ meson cross section is well described by an exponentially falling distribution. The full set of spin density matrix elements is measured in two $`𝑸^\mathrm{𝟐}`$ bins. Predictions based on perturbative QCD are compared to the measurements. The combination $`\mathbf{(}\mathrm{𝟐}𝒓_{\mathrm{𝟏𝟏}}^\mathrm{𝟓}\mathbf{+}𝒓_{\mathrm{𝟎𝟎}}^\mathrm{𝟓}\mathbf{)}`$ is extracted from the $`𝚽`$ angle distribution and is observed to deviate from zero, which indicates a small but significant violation of the $`𝒔`$-channel helicity conservation (SCHC) approximation. The ratio $`𝑹`$ of longitudinal to transverse $`\mathit{\varphi }`$ meson production cross sections is observed to increase with $`𝑸^\mathrm{𝟐}`$. A common dependence for $`𝑹`$ as a function of $`𝑸^\mathrm{𝟐}\mathbf{/}𝑴_𝑽^\mathrm{𝟐}`$ is observed for elastic $`𝝆`$, $`\mathit{\varphi }`$ and $`𝑱\mathbf{/}𝝍`$ meson production.
## Acknowledgements
We are grateful to the HERA machine group whose outstanding efforts have made and continue to make this experiment possible. We thank the engineers and technicians for their work in constructing and now maintaining the H1 detector, our funding agencies for financial support, the DESY technical staff for continual assistance, and the DESY directorate for the hospitality which they extend to the non DESY members of the collaboration. We thank further I. Akushevich, J.-R. Cudell, D.Yu. Ivanov, N. Nikolaev and I. Royen for useful discussions and for providing us with their model predictions.
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# On the Leibniz cohomology of vector fields
## 1 Introduction
Let $`g`$ be a Lie algebra. The Leibniz cohomology $`HL^{}(g)`$ of $`g`$ with trivial coefficients is the cohomology of the complex with cochain modules $`CL^n(g)=\mathrm{Hom}(g^n,)`$ and coboundary operators $`d:CL^n(g)CL^{n+1}(g)`$, defined by
$$df(x_1,\mathrm{},x_{n+1})=\underset{i<j}{}(1)^{j+1}f(x_1,\mathrm{},x_{i1},[x_i,x_j],x_{i+1},\mathrm{},\widehat{x}_j,\mathrm{},x_{n+1}),$$
where $`(x_1,\mathrm{},x_k)`$ denotes the element $`x_1\mathrm{}x_k`$ in the tensor product $`g^k`$. The Leibniz cohomology is a generalization of the Chevalley-Eilenberg cohomology of Lie algebras, in the sense that it constitutes the natural cohomology theory for Leibniz algebras, which are non-antisymmetric generalizations of Lie algebras. It was introduced by J.-L. Loday and T. Pirashvili in . When computed on Lie algebras, the Leibniz cohomology may detect new invariants in dimensions higher then $`1`$, since $`HL^0(g)=H^0(g)`$ and $`HL^1(g)=H^1(g)`$. The method to compare the Lie and the Leibniz cohomology of a Lie algebra in higher dimensions is given by the Pirashvili long exact sequence
$$\begin{array}{c}0H^2(g)HL^2(g)H_{rel}^0(g)H^3(g)\mathrm{}\hfill \\ \mathrm{}H^n(g)HL^n(g)H_{rel}^{n2}(g)H^{n+1}(g)\mathrm{}\hfill \end{array}$$
induced by the projections $`g^n\mathrm{\Lambda }^ng`$, where $`H_{rel}^{}(g)`$ is a relative cohomology described in .
Loday and J.-M. Oudom showed in , that the Leibniz cohomology $`HL^{}(g)`$ is endowed with the structure of a graded dual Leibniz algebra (in the sense of Koszul duality for operads). In fact, there is a cup product $`:HL^p(g)HL^q(g)HL^{p+q}(g)`$ defined on the cohomology classes of two cocycles $`\alpha CL^p(g)`$, $`\beta CL^q(g)`$ as the cohomology class of the cocycle
$$\alpha \beta (x_1,\mathrm{},x_{p+q})=\underset{\sigma Sh_{p1,q}}{}\mathrm{sgn}(\sigma )\alpha (x_1,x_{\sigma (2)},\mathrm{},x_{\sigma (p)})\beta (x_{\sigma (p+1)},\mathrm{},x_{\sigma (p+q)}),$$
where $`Sh_{p1,q}`$ denotes the set of the $`(p1,q)`$-shuffles among the permutations on $`\{2,3,\mathrm{},p+q\}`$. The cup product satisfies the relation of dual Leibniz algebras
$$(\alpha \beta )\gamma \alpha (\beta \gamma )=(1)^{|\beta ||\gamma |}\alpha (\gamma \beta ).$$
Therefore, in order to describe the Leibniz cohomology, it suffices to give its generators as a dual Leibniz algebra.
When $`g`$ is a topological Lie algebra, we can consider the submodules of continuous cochains $`CL_{cont}^n(g)=\mathrm{Hom}_{cont}(g^n,)\mathrm{Hom}(g^n,)=CL^n(g)`$. Since the Leibniz differential $`df`$ of any continuous cochain $`f`$ is still continuous, the sequence of continuous cochains $`d:\mathrm{Hom}_{cont}(g^n,)\mathrm{Hom}_{cont}(g^{n+1},)`$ forms a subcomplex of the Leibniz complex, whose cohomology is called continuous Leibniz cohomology of $`g`$. It is still denoted by $`HL^{}(g)`$ when no confusion can arise.
If $`g`$ is the topological Lie algebra $`\mathrm{Vect}M`$ of vector fields on a differentiable manifold $`M`$, the Leibniz cohomology $`HL^{}(\mathrm{Vect}M)`$ is a generalization of the Lie cohomology $`H^{}(\mathrm{Vect}M)`$ studied by I. M. Gelfand, D. B. Fuks, R. Bott, G. Segal and A. Haefliger in , , , and references therein.
In order to study the Leibniz cohomology of $`\mathrm{Vect}^n`$, J. Lodder starts with the Lie algebra of formal vector fields $`W_n=[[x_1,\mathrm{},x_n]]\{\frac{}{x_1},\mathrm{},\frac{}{x_n}\}`$, where $`[[x_1,\mathrm{},x_n]]`$ denotes the algebra of formal power series in $`n`$ variables, and $`\frac{}{x_i}`$ the partial derivation in respect of the variable $`x_i`$. He showed in that the first non-zero class in $`HL^{}(W_n)`$ is the Godbillon-Vey class in dimension $`2n+1`$. Using Pirashvili’s sequence, he describes $`HL^{}(W_n)`$ as a dual Leibniz algebra with generators and relations given in . For $`n=1`$, this dual Leibniz algebra has one generator in degree $`3`$ (the Godbillon-Vey class $`\theta `$) and a new generator in degree $`4`$ (called $`\alpha `$) with relations $`\theta ^2=0`$ and $`\alpha \theta =0`$.
The case of $`HL^{}(\mathrm{Vect}M)`$, for general smooth manifolds, is much more complicated. Even in the simplest case $`M=S^1`$ there are too many non-zero terms appearing in the sequence, which prevent us to determine the dimension of the Leibniz cohomology groups. Hence Pirashvili’s method cannot be reproduced in general.
In this paper, instead, we generalize the classical method of Gelfand and Fuks, regarding continuous Leibniz cochains on $`\mathrm{Vect}M`$ as generalized sections of the tensor powers of the tangent bundle $`TM`$, over the manifold $`M\times \mathrm{}\times M`$. This guarantees for Leibniz cohomology the same general results previously determined by Gelfand and Fuks for Lie cohomology, and some explicit computations for manifolds with few non-trivial de Rham cohomology classes.
In particular, in section 2 we describe the spectral sequence associated to the diagonal filtration (filtration of a complex of distributions by their support). The first term of the sequence is given by some quotient complexes.
In section 3 we describe a spectral sequence converging to the cohomology of these quotient complexes. It is associated to the order filtration, i.e. the filtration of a complex of distributions which are supported on a given submanifold by their order.
In the last section, we show that the order spectral sequence for the Leibniz cohomology of $`\mathrm{Vect}S^1`$ collapses. The resulting diagonal cohomology is $`4`$-periodic, with dimension respectively $`0,1,2,1`$. Finally, we describe the diagonal cocycles which determine the new generators.
For the collapsing of the diagonal spectral sequence, and thus for an explicit desciption of all continuous Leibniz cohomology of $`\mathrm{Vect}S^1`$, we would need the dual Leibniz algebra structure which is still a mystery for us. The knowledge of this structure leads to study the relationship with Leibniz minimal models (as introduced by M. Livernet in ) and ‘real’ homotopy types in the spirit of Bott, Haefliger and Segal .
## 2 Generalized sections and diagonal filtration
Let $`M`$ be a smooth oriented manifold of dimension $`n`$. Let $`\mathrm{Vect}M`$ be the Lie algebra of vector fields on $`M`$, equipped with the $`C^{\mathrm{}}`$ topology which makes it a topological Fréchet nuclear Lie algebra. Let $`CL_{cont}^{}(\mathrm{Vect}M)`$ be the complex of continuous Leibniz cochains of $`\mathrm{Vect}M`$. Then, for any $`m0`$, the module of Leibniz $`m`$-cochains is the space of generalized sections of the bundle of tensor powers $`TM^m`$ over the cartesian product manifold $`M^m`$, meaning $`TM^m=_{i=1}^m\pi _i^{}(TM)`$ where $`\pi _i:M^mM`$ is the projection on the $`i`$th factor. In other words, the Leibniz $`m`$-cochains of $`\mathrm{Vect}M`$ are the distributions on $`M^m`$ with values in the bundle $`TM^m`$. Remark that the Chevalley-Eilenberg $`m`$-cochains of $`\mathrm{Vect}M`$, that is, the elements of $`C_{cont}^m(\mathrm{Vect}M)`$, can be seen as conveniently anti-symmetrized generalized sections of the same bundle, cf , . Throughout this section we follow the notations of .
For any $`k1`$, let $`M_k^m`$ be the subset of $`M^m`$ of $`m`$-uples $`(x_1,\mathrm{},x_m)`$ which do not have more then $`k`$ different entries $`x_i`$, that is,
$$M_k^m=\{(x_1,\mathrm{},x_m)M^m|(i_1,\mathrm{},i_{k+1})(1,\mathrm{},m)(i_r,i_s):x_{i_r}=x_{i_s}\}.$$
Then, the submanifold $`M_1^m=\{(x,\mathrm{},x)M^m\}=\mathrm{}`$ coincides with the diagonal of $`M^m`$, all the subsets $`M_k^m`$ for $`km`$ coincide with the whole $`M^m`$, and there is a sequence of inclusions
$$0=M_0^mM_1^m=\mathrm{}M_2^m\mathrm{}M_{m1}^mM_m^m=M^m.$$
Denote by $`CL_k^m(M)`$ the subspace of $`CL_{cont}^m(\mathrm{Vect}M)`$ of generalized sections of the bundle $`TM^m`$ with support on the subset $`M_k^m`$. Then there is an induced sequence of inclusions of cochains modules
$$\{0\}=CL_0^m(M)CL_1^m(M)\mathrm{}CL_{m1}^m(M)CL_m^m(M)=CL_{cont}^m(\mathrm{Vect}M).$$
In other words, as a distribution, a Leibniz $`m`$-cochain is concentrated on some subspace $`M_k^m`$ (perhaps on the whole $`M^m`$).
The above inclusions of Leibniz cochain modules define an increasing multiplicative filtration of the complex $`CL_{cont}^{}(\mathrm{Vect}M)`$ called diagonal filtration, that is,
$$d(CL_k^m(M))CL_k^{m+1}(M),$$
for all $`k`$ and $`m`$, and
$$CL_k^m(M)CL_h^l(M)CL_{k+h}^{m+l}(M),$$
for all $`k,h`$ and $`m,l`$. To see this, it suffices to understand that the subspace $`CL_k^m(M)`$ consists exactely of Leibniz cochains which vanish on any family of $`m`$ vector fields having the property $`(\mathrm{}_k)`$. (Recall that a family of $`m`$-vector fields $`\mathrm{\Gamma }\mathrm{Vect}M`$ has the property $`(\mathrm{}_k)`$ if for any set of $`k`$ points $`SM`$ at least one vector field of $`\mathrm{\Gamma }`$ vanishes in a neighbourhood of $`S`$.)
Hence, there exists a spectral sequence abutting to the Leibniz cohomology of $`\mathrm{Vect}M`$. To avoid confusion, we denote this spectral sequence by $`B_{}(M)`$. The $`0`$-th term is the bicomplex of the quotients
$$B_0^{m,k}(M):=CL_k^m(M)/CL_{k1}^m(M),$$
with differentials $`d_v:B_0^{m,k}(M)B_0^{m+1,k}(M)`$ and $`d_t:B_0^{m,k}(M)B_0^{m+1,k1}(M)`$ (of bidegree $`(1,0)`$ and $`(1,1)`$) induced by the Leibniz coboundary. For $`k=1`$, the differential $`d_v`$ coincides with the Leibniz differential, hence the family $`B_0^{m,1}(M)=CL_{\mathrm{}}^m(M)`$ is a subcomplex called diagonal complex. For $`k=0`$ we have $`B_0^{m,0}(M)=0`$, except for $`m=0`$ where $`B_0^{0,0}(M)=`$.
Then the first term of the spectral sequence is the cohomology
$$\begin{array}{cc}B_1^{m,k}(M)=H^m(CL_k^{}(M)/CL_{k1}^{}(M),d_v),\hfill & \text{for }k>1\hfill \\ B_1^{m,1}(M)=H^m(CL_{\mathrm{}}^{}(M),d_v),\hfill & \text{for }k=1\hfill \\ B_1^{m,0}(M)=\{\begin{array}{cc},\hfill & m=0\hfill \\ 0,\hfill & m>0\hfill \end{array},\hfill & \text{for }k=0.\hfill \end{array}$$
The cohomology of the complex $`(CL_{\mathrm{}}^{}(M),d_v)`$ is also called the diagonal Leibniz cohomology of $`\mathrm{Vect}M`$, and denoted by $`HL_{\mathrm{}}^{}(\mathrm{Vect}M)`$. By abuse of notation, we call $`k`$-diagonal cochains the elements of $`B_0^{,k}(M)`$, and $`k`$-diagonal cohomology the term $`B_1^{,k}(M)`$, also denoted $`HL_{(k)}^{}(\mathrm{Vect}M)`$. Of course, if $`\stackrel{~}{d_t}`$ is the differential induced by $`d_t`$ on the first term, the second term of the spectral sequence is the cohomology $`B_2^{m,k}(M)=H^k(B_1^{m+,}(M),\stackrel{~}{d_t})`$. The Leibniz cohomology of $`\mathrm{Vect}M`$ is then
$$HL^m(\mathrm{Vect}M)=\underset{km}{}B_{\mathrm{}}^{m,k}(M).$$
In this generality, as for the Gelfand-Fuks computations, the spectral sequence is rather intractable. In order to compute its first term, we introduce following Gelfand and Fuks, for each quotient complex $`B_0^{,k}(M)=CL_k^{}(M)/C_{k1}^{}(M)`$ a spectral sequence abutting to its cohomology $`B_1^{m,k}(M)`$. These spectral sequences are defined by the order filtration for Leibniz cochains, which is discussed in the next section.
## 3 The order filtration for diagonal cohomology
The set $`B_0^{m,k}(M)`$ of $`k`$-diagonal Leibniz cochains contains the distributions concentrated on the subset $`M_k^m`$ modulo those concentrated on the subset $`M_{k1}^m`$. The set $`M_{(k)}^m=M_k^mM_{k1}^m`$ is now a submanifold without singularities, and a $`k`$-diagonal Leibniz cochain can then be described as a distributions defined on the jets of sections of the bundle $`TM^m`$, where the jet expansion is taken in the normal direction to the submanifold $`M_{(k)}^m`$ in $`M^m`$.
Recall that a generalized section of a vector bundle $`E`$ over $`M`$, concentrated on a subset $`SM`$, is of order $`l`$ on $`S`$ if it vanishes on any section having trivial $`l`$-jet at every point of $`S`$.
For any integer $`p`$, denote by $`F^pB_0^{m,k}(M)`$ the subspace of Leibniz $`k`$-diagonal $`m`$-cochains of $`\mathrm{Vect}M`$ which are of order $`mp`$. The family of such spaces defines a decreasing filtration of the complex $`CL_k^{}(M)/CL_{k1}^{}(M)`$ of $`k`$-diagonal Leibniz cochains. In particular, for $`k=1`$, the family $`F^pB_0^{m,1}(M)=F^pCL_{\mathrm{}}^m(M)`$ gives a filtration of the diagonal complex.
###### Theorem 3.1
For each $`k>0`$ there is a spectral sequence abutting to the $`k`$-diagonal Leibniz cohomology of $`\mathrm{Vect}M`$, with the following $`E_2`$-term:
E2p,q(k)=Hp(Mk,Mk1k)
=+q1qkq
>q1,,qk0
HLq1(Wn)HLqk(Wn),superscriptsuperscriptsubscript𝐸2𝑝𝑞𝑘tensor-productsubscript𝐻𝑝superscript𝑀𝑘subscriptsuperscript𝑀𝑘𝑘1subscriptdirect-sum
=+q1qkq
>q1,,qk0
tensor-product𝐻superscript𝐿subscript𝑞1subscript𝑊𝑛𝐻superscript𝐿subscript𝑞𝑘subscript𝑊𝑛{}^{(k)}E_{2}^{p,q}=H_{-p}(M^{k},M^{k}_{k-1})\otimes\bigoplus_{\mbox{\scriptsize\vbox{\hbox{$q_{1}+\ldots+q_{k}=q$}\hbox{$q_{1},\ldots,q_{k}>0$}}}}HL^{q_{1}}(W_{n})\otimes\ldots\otimes HL^{q_{k}}(W_{n}),
where $`H_p(M^k,M_{k1}^k)`$ denotes the relative cohomology of the manifold $`M^k`$ with respect to the subspace $`M_{k1}^k`$.
In particular, for $`k=1`$, there is a spectral sequence
$$E_2^{p,q}=H_p(M)HL^q(W_n)HL_{\mathrm{}}^{p+q}(\mathrm{Vect}M)$$
abutting to the diagonal cohomology of $`\mathrm{Vect}M`$.
Proof: Since the proof for Leibniz cohomology does not substantially differs from that for Lie cohomology, we only stress the main differences in the simple case of the diagonal complex.
By definition, the $`0`$-th term of the spectral sequence defined by the order filtration on the diagonal complex $`CL_{\mathrm{}}^{}(M)`$ is the bicomplex
$$E_0^{p,q}(M)=F^pCL_{\mathrm{}}^{p+q}(M)/F^{p+1}CL_{\mathrm{}}^{p+q}(M)$$
that is, $`E_0^{p,q}`$ is the quotient space of cochains from $`CL_{\mathrm{}}^{p+q}(\mathrm{Vect}M)`$ which are of order $`q`$ with respect to the diagonal $`\mathrm{}`$, modulo those of order $`<q`$. In other words, an element of the space $`E_0^{p,q}`$ is a generalized section of the bundle over $`M`$
$$_0^{p,q}=\mathrm{Hom}(S^q\mathrm{norm}_{M^{p+q}}\mathrm{},TM^{p+q}|_{\mathrm{}}).$$
Here, $`\mathrm{norm}_{M^{p+q}}\mathrm{}`$ is the normal bundle of the diagonal (closed) submanifold $`\mathrm{}M^{p+q}`$, and is related to the fact that the jet expansion is in the normal direction with respect to the diagonal. The $`q`$th symmetric power $`S^q`$ arises from the fact that the generalized sections are non-zero exactly on those elements with non-trivial $`q`$-jet on the diagonal modulo those with trivial $`q+1`$-jet. The fiber of the bundle $`_0^{p,q}`$ in a point $`xM`$ is
$$\mathrm{Hom}(S^q(\stackrel{p+q}{}V/V_{\mathrm{}}),V^{p+q}),$$
where $`V=T_xM`$ is the fiber of $`TM`$ and $`V_{\mathrm{}}`$ denotes the image of the diagonal inclusion $`VV\mathrm{}V`$. The vector space $`S^q(^{p+q}V/V_{\mathrm{}})`$ admits the standard Koszul resolution
$$\begin{array}{c}\hfill 0S^q((^{p+q}V)/V_{\mathrm{}})\stackrel{pr}{}S^q(^{p+q}V)S^{q1}(^{p+q}V)V\mathrm{}\\ \hfill \mathrm{}S^{qi+1}(^{p+q}V)\mathrm{\Lambda }^{i1}V\mathrm{}\end{array}$$
Using the isomorphisms
$$\begin{array}{cc}\mathrm{Hom}(S^{qi}(\stackrel{p+q}{}V)\mathrm{\Lambda }^iV,V^{p+q})\hfill & \mathrm{\Lambda }^iV^{}\mathrm{Hom}(S^{qi}(\stackrel{p+q}{}V),V^{p+q})\hfill \\ & \mathrm{\Lambda }^iV^{}\mathrm{Hom}(S^{qi}(\stackrel{p+q}{}V)(V^{})^{p+q},),\hfill \end{array}$$
and
$$S^{}(\stackrel{p+q}{}V)(V^{})^{p+q}(S^{}(V)(V^{}))^{p+q}W_n^{p+q},$$
where $`W_n`$ is the Lie algebra of formal vector fields on $`^n`$, we finally have an isomorphism
$$\mathrm{Hom}(S^{qi}(\stackrel{p+q}{}V)\mathrm{\Lambda }^iV,V^{p+q})\mathrm{\Lambda }^iV^{}CL_{(pi)}^{p+q}(W_n)$$
where the index $`l`$ in $`CL_{(l)}^{p+q}(W_n)`$ means the weight with respect to the adjoint action of the Euler field $`e_0=_{i=1}^nx_i\frac{d}{dx_i}`$. Then the Koszul resolution gives rise to an exact sequence of fibres. ¿From this, we get the corresponding exact sequence of bundles, and of sections, because the global sections functor is exact on $`𝒞^{\mathrm{}}`$ fibre bundles. Finally, we pass to the exact sequence of generalized sections by a lemma on duals of Fréchet nuclear spaces.
In conclusion, we have an exact sequence of complexes (see p.146 for details)
$$0E_0^{p,}Sec^{}\xi _0^{p,}Sec^{}\xi _1^{p,}\mathrm{}$$
where $`\xi _i^{p,q}`$ the bundle associated with $`TM`$ with typical fiber $`\mathrm{\Lambda }^i(TM)^{}\gamma _i^{p,q}`$, and $`\gamma _i^{p,q}`$ is the bundle associated with $`TM`$ with typical fiber $`CL_{(pi)}^{p+q}(W_n)^{}`$. We use lemma (A.1) to get the acyclicity of all complexes in the above sequence except $`E_0^{p,}`$ and $`Sec^{}\xi _p^{p,}`$. The identification of the differential $`d_1^{p,q}`$ is clearly the same as in , so we get the stated result. $`\mathrm{}`$
This proof can easily be adapted to obtain the higher diagonal $`E_2`$-terms.
## 4 Leibniz cohomology of the Lie algebra of vector fields on $`S^1`$
Now we consider the circle $`S^1`$. We choose coordinates $`z=\mathrm{exp}(i\phi )`$, for $`\phi [0,2\pi ]`$. As a topological Lie algebra, $`\mathrm{Vect}S^1`$ is spanned by the vector fields $`e_k:=i\mathrm{exp}(ik\phi )\frac{d}{d\phi }`$, for all integers $`k`$, with Lie bracket $`[e_k,e_l]=(kl)e_{k+l}`$.
The Lie cohomology of $`\mathrm{Vect}S^1`$ was computed by Gelfand and Fuks. It is a free graded-commutative algebra with one even generator in degree $`2`$, represented by the Gelfand-Fuks cocycle $`\omega `$, and one odd generator in degree $`3`$, represented by the Godbillon-Vey cocycle $`\theta `$. Hence
$$H^{}(\mathrm{Vect}S^1)=\underset{k0}{}[\omega ^k]\underset{k0}{}[\theta \omega ^k].$$
For each $`\phi _0S^1`$, the Taylor expansion of vector fields defines a map
$$\begin{array}{ccc}\hfill \pi _{\phi _0}:\mathrm{Vect}S^1& & W_1\hfill \\ \hfill f(\phi )\frac{d}{d\phi }& & _{n0}\frac{d^nf(\phi _0)}{d\phi ^n}x^n\frac{d}{dx},\hfill \end{array}$$
where $`W_1=[[x]]\{\frac{d}{dx}\}`$ is the complex Lie algebra of formal vector fields in the origin of $``$ and $`x=\phi \phi _0`$. The pull back $`\pi ^{}`$ gives a map from the Lie cocycles of $`W_1`$ to the Lie cocycles of $`\mathrm{Vect}S^1`$, and similarly a map between the Leibniz cocycles. A cocycle $`\gamma _{\phi _0}=\pi _{\phi _0}^{}\gamma `$ on $`\mathrm{Vect}S^1`$ which is pulled back from a cocycle $`\gamma `$ on $`W_1`$ is local on $`S^1`$, in the sense that as a distribution it has support on the single point $`\phi _0`$. Changing the point $`\phi _0`$ of evaluation produces cocycles which are cohomologous, hence a local cocycle admits an $`S^1`$-invariant form given by the integration
$$\gamma :=_0^{2\pi }\gamma _\phi 𝑑\phi ,$$
where $`S^1`$ is meant to act on $`\mathrm{Vect}S^1`$ by rotations.
The Godbillon-Vey cocycle $`\theta `$ on $`\mathrm{Vect}S^1`$ is precisely pulled back from a cocycle on $`W_1`$, hence it is a local. Instead, the Gelfand-Fuks cocycle $`\omega `$ is a diagonal cocycle, that is, it is a distribution with support on the diagonal of $`S^1\times S^1`$.
The Leibniz cohomology of $`W_1`$, computed by Lodder, is the dual Leibniz algebra generated by the Godbillon-Vey cocycle $`\theta `$ and a new generator $`\alpha `$ in degree $`4`$. The cup products $`\theta ^2`$ and $`\alpha \theta `$ are zero, because of the definition of the cup product. The cup product $`\alpha ^2\alpha `$ is equal to $`2\alpha \alpha ^2`$, because of the relation of dual Leibniz algebra and the fact that $`\alpha `$ has even degree. Hence the only higher order Leibniz cocycles for $`W_1`$ are $`\alpha ^k:=\alpha \alpha ^{k1}`$ and $`\theta \alpha ^k`$.
The pull back on $`\mathrm{Vect}S^1`$ of these cocycles gives all the local Leibniz cocycles which appear in $`HL^{}(\mathrm{Vect}S^1)`$. We now determine the diagonal ones.
### The diagonal Leibniz cohomology
The diagonal cohomology $`HL_\mathrm{\Delta }^{}(\mathrm{Vect}S^1)`$ is given by the cohomology classes which are represented by diagonal cocycles. In particular, it also contains local classes. The subset of local classes forms a dual Leibniz algebra with the cup product. In fact, the cup product of two local cocycles is still a local cocycle. Instead, the cup product of a local cocycle by a diagonal one, or the product of two diagonal cocycles, is not diagonal anymore, but $`2`$-diagonal. In general, the cup product of $`k`$ diagonal cocycle (with at least a non-local one) is $`k`$-diagonal. Therefore, the diagonal cohomology $`HL_\mathrm{\Delta }^{}(\mathrm{Vect}S^1)`$ is not a dual Leibniz algebra with the cup product, and it can only be described as a vector space.
###### Theorem 4.1
The diagonal cohomology of $`\mathrm{Vect}S^1`$ is the graded vector space spanned by the classes of the local cocycles
$$\begin{array}{cc}\theta \alpha ^r& \text{in degree }3+4r\text{, for }r0\hfill \\ \alpha ^s& \text{in degree }4s\text{, for }s1\hfill \end{array}$$
and by the classes of the diagonal cocycles
$$\begin{array}{cc}\omega _r& \text{in degree }2+4r\text{, for }r0\hfill \\ \beta _s& \text{in degree }4s1\text{, for }s1\hfill \end{array}$$
where $`\omega _0=\omega `$ is the Gelfand-Fuks cocycle in degree $`2`$ and $`\omega _r`$, $`\beta _s`$ determine new invariants for $`r,s1`$.
Proof: The spectral sequence of theorem (3.1), which abuts to $`HL_{\mathrm{}}^{}(\mathrm{Vect}S^1)`$, degenerates at the second term for $`M=S^1`$, because
$$E_2^{p,q}(S^1)=H_p(S^1)HL^q(W_1)$$
differs from $`0`$ only for $`p=0,1`$, hence all the induced differentials in the higher terms are zero. So, we have
HLm(VectS1)=
=+pqm
>q0
Hp(S1)HLq(W1),𝐻subscriptsuperscript𝐿𝑚Vectsuperscript𝑆1subscriptdirect-sum
=+pqm
>q0
tensor-productsubscript𝐻𝑝superscript𝑆1𝐻superscript𝐿𝑞subscript𝑊1HL^{m}_{\triangle}({\rm Vect}\,S^{1})=\bigoplus_{\mbox{\scriptsize\vbox{\hbox{$p+q=m$}\hbox{$q>0$}}}}H_{-p}(S^{1})\otimes HL^{q}(W_{1}),
for all $`m>0`$, where the Leibniz cohomology of the formal vector fields $`W_1`$ was computed by Lodder, in , as
$$HL^q(W_1)=\underset{r+s=q}{}\mathrm{\Lambda }^r[\theta ]T^s[\alpha ].$$
Thus, we obtain
HLm(VectS1)=
=+prsm
>+rs0
Λp[η]Λr[θ]Ts[α],𝐻subscriptsuperscript𝐿𝑚Vectsuperscript𝑆1subscriptdirect-sum
=+prsm
>+rs0
tensor-producttensor-productsuperscriptΛ𝑝delimited-[]𝜂superscriptΛ𝑟delimited-[]𝜃superscript𝑇𝑠delimited-[]𝛼HL^{m}_{\triangle}({\rm Vect}\,S^{1})=\bigoplus_{\mbox{\scriptsize\vbox{\hbox{$p+r+s=m$}\hbox{$r+s>0$}}}}\Lambda^{p}[\eta]\otimes\Lambda^{r}[\theta]\otimes T^{s}[\alpha],
where $`\eta `$ is a generator in degree $`1`$, $`\theta `$ is the Godbillon-Vey generator in degree $`3`$, and $`\alpha `$ is the Lodder generator in degree $`4`$. Since $`\mathrm{\Lambda }^p[\eta ]`$ differs from $`0`$ only for $`p=0,1`$, and $`\mathrm{\Lambda }^r[\theta ]`$ differs from $`0`$ only for $`r=0,3`$, we have
$$\begin{array}{c}HL_{\mathrm{}}^m(\mathrm{Vect}S^1)=\underset{r+s=m+1}{}[\eta ]\mathrm{\Lambda }^r[\theta ]T^s[\alpha ]\underset{r+s=m}{}\mathrm{\Lambda }^r[\theta ]T^s[\alpha ]\hfill \\ =[\eta \theta ]T^{m2}[\alpha ][\theta ]T^{m3}[\alpha ][\eta ]T^{m+1}[\alpha ]T^m[\alpha ].\hfill \end{array}$$
If we call the new diagonal cocycles
$$\begin{array}{cc}\omega _r:=\eta \theta \alpha ^r,\hfill & \text{for }r0\hfill \\ \beta _s:=\eta \alpha ^s,\hfill & \text{for }s1,\hfill \end{array}$$
and we remark that the higher-degree local cocycles are represented by the cup products, we get the final result
$$HL_{\mathrm{}}^{}(\mathrm{Vect}S^1)=\mathrm{\Lambda }^{}[\theta ]T^{}[\alpha ]\underset{r0}{}[\omega _r]\underset{s1}{}[\beta _s].$$
$`\mathrm{}`$
###### Corollary 4.2
The diagonal Leibniz cohomology $`HL_{\mathrm{}}^{}(\mathrm{Vect}S^1)`$ is periodic of period $`4`$. The dimensions of $`HL_{\mathrm{}}^{n+4k}(\mathrm{Vect}S^1)`$ for $`n=1,2,3,4`$ are respectively $`0,1,2,1`$.
In particular, the second cohomology group $`HL^2(\mathrm{Vect}S^1)`$ has dimension $`1`$, generated by the Gelfand-Fuks class. This result was obtained by Loday and Pirashvili in .
### Representative cocycles for the cohomology classes
Recall that the Godbillon-Vey class in $`H^3(W_1)`$ has a representative cocycle
$$\theta (F,G,H)=\left|\begin{array}{ccc}f\hfill & g\hfill & h\hfill \\ f^{}\hfill & g^{}\hfill & h^{}\hfill \\ f^{\prime \prime }\hfill & g^{\prime \prime }\hfill & h^{\prime \prime }\hfill \end{array}\right|_{x=0},$$
where $`F=f(x)\frac{d}{dx}`$, $`G=g(x)\frac{d}{dx}`$ and $`H=h(x)\frac{d}{dx}`$ denote some formal vector fields, and all the functions are evaluated at $`x=0`$. Now, there are two ways to derive from $`\theta `$ the Gelfand-Fuks cocycle.
Under the pull back $`\pi _{\phi _0}^{}:C^{}(W_1)C^{}(\mathrm{Vect}S^1)`$, $`\theta `$ determines a local cohomology class in $`H^3(\mathrm{Vect}S^1)`$. Its representative cocycle $`\theta _{\phi _0}(F,G,H)`$ can be expressed by the same local formula, where $`F=f(\phi )\frac{d}{d\phi }`$, $`G`$ and $`H`$ are now vector fields on $`S^1`$ and the functions are evaluated at $`\phi =\phi _0`$. In fact, the integral form of local cocycles has an advantage, since it may allow to determine automatically diagonal cocycles of one degree less, by contracting any of the variables with the Euler field $`e_0=i\frac{d}{d\phi }`$. We explain this in detail.
The only cochains which may contribute to the cohomology are those of zero weight, i.e. those $`\gamma `$ for which $`ad_{e_0}(\gamma )=0`$, because the subcomplex of non-zero weight cochains is contractible, as shown in appendix A. Among these, those of the form $`\gamma =\gamma ^{}ϵ^0`$, where $`ϵ^0=e_0^{}`$ is the dual cochain of the Euler vector field and $`\gamma ^{}`$ has zero weight, are such that the contraction gives $`i_{e_0}(\gamma )=\gamma ^{}`$. Since the Lie coboundary operator satisfies Cartan’s formula, we have
$$d(\gamma )=d(\gamma ^{})ϵ^0.$$
Hence, $`\gamma `$ is a cocycle if and only if $`\gamma ^{}`$ is.
However, if the local cocycle $`\gamma `$ is not in its integral form, the dependence $`\gamma =\gamma _{\phi _0}`$ from the evaluation point $`\phi _0`$ is reflected onto $`\gamma ^{}=\gamma _{\phi _0}^{}`$, while the new cocycle $`\gamma ^{}`$ is surely not local. Indeed, the term which leads to locality is precisely the $`ϵ^0`$ that $`\gamma ^{}`$, being antisymmetric, cannot contain. Since the contraction and the differential commute with the integration, we can avoid the dependence on the point of evaluation, by applying the procedure to the integral mean of $`\gamma ϵ^0`$.
The Godbillon-Vey cocycle is an example of a cocycle of the form $`\gamma ϵ^0`$, with $`\gamma =i/2ϵ^1ϵ^1`$ and $`ϵ^p=e_p^{}`$. For the Gelfand-Fuks class in $`H^2(\mathrm{Vect}S^1)`$, the well known representative cocycle
$$\omega (F,G)=_{S^1}\left|\begin{array}{cc}f^{}(\phi )\hfill & g^{}(\phi )\hfill \\ f^{\prime \prime }(\phi )\hfill & g^{\prime \prime }(\phi )\hfill \end{array}\right|𝑑\phi $$
can be obtained from the Godbillon-Vey cocycle as follows,
$$\omega =i_{e_0}_{S^1}\theta _\phi 𝑑\phi .$$
We apply the same method to get diagonal $`k`$-cocycles on $`S^1`$ from local Leibniz $`k+1`$-cocycles.
The Godbillon-Vey class in $`HL^3(W_1)`$ has the same representative cocycle $`\theta `$, and by pull back we obtain the local class $`\theta _{\phi _0}`$ in $`HL^3(\mathrm{Vect}S^1)`$. By contracting its integral form we recover the Leibniz cocycle $`\omega `$ which represents the Gelfand-Fuks class.
In fact, as a Leibniz cocycle, $`\theta `$ is also cohomologous to the cocycle
$$\stackrel{~}{\theta }(F,G,H)=(fgh^{\prime \prime \prime }fg^{\prime \prime \prime }h)|_{x=0},$$
which, by integration and contraction, produces again the Gelfand-Fuks cocycle. However, we prefer to use the antisymmetric form of $`\theta `$ since it is advantageous in the computation of the cup products.
The Lodder class in $`HL^4(W_1)`$ is the image of the non trivial Lie cohomology class in $`H^3(W_1;W_1^{})`$ with values in the adjoint representation, under the map $`H^3(W_1;W_1^{})HL^3(W_1;W_1^{})H^4(W_1)`$. It is represented by the cocycle
$$\alpha (L,F,G,H)=l^{}(0)\left|\begin{array}{ccc}f\hfill & g\hfill & h\hfill \\ f^{}\hfill & g^{}\hfill & h^{}\hfill \\ f^{\prime \prime }\hfill & g^{\prime \prime }\hfill & h^{\prime \prime }\hfill \end{array}\right|_{x=0},$$
where $`L=l(x)\frac{d}{dx}`$ denote another formal vector field. Then, the first new class for the Leibniz cohomology of $`\mathrm{Vect}S^1`$ is represented by the cocycle
$$\beta _1(L,F,G)=_{S^1}l^{}(\phi )\left|\begin{array}{cc}f^{}(\phi )\hfill & g^{}(\phi )\hfill \\ f^{\prime \prime }(\phi )\hfill & g^{\prime \prime }(\phi )\hfill \end{array}\right|𝑑\phi ,$$
where, as before, $`L,F,G`$ now denote vector fields on $`S^1`$.
The second method to obtain the Gelfand-Fuks cocycle from the Godbillon-Vey cocycle does not differ too much from the first one: one can show that
$$\frac{d}{d\varphi }\theta _\varphi =d\left(\left|\begin{array}{cc}^{}& ^{}\\ ^{\prime \prime }& ^{\prime \prime }\end{array}\right|\right).$$
This permits to define $`\omega `$ as integral over $`S^1`$ of the expression in paranthesis on the right hand side, because it gives automatically a cocycle by the above equation. It is easily seen that $`\beta _1`$ can also be obtained via this method, i.e. that in the Leibniz setting
$$\frac{d}{d\varphi }\alpha _\varphi =d\left((^{})\left|\begin{array}{cc}^{}& ^{}\\ ^{\prime \prime }& ^{\prime \prime }\end{array}\right|\right).$$
The representative cocycles for the other new diagonal cohomology classes can be obtained in the same manner, starting from the cup products of the two local cocycles $`\theta `$ and $`\alpha `$. The formulas are though complicated by the presence of the shuffles, and we omit them.
To end this section, let us just remark that the higher diagonal Leibniz cohomology spectral sequences also collapse at the second term. This is due to the fact that $`H^{}((S^1)^k,(S^1)_{k1}^k)`$ is non-zero only in two dimensions, see .
## Aknowledgements
The authors would like to thank Olivier Mathieu for many useful discussions and suggestions on the subject of the paper and Muriel Livernet for kindly remarking some oversights.
## Appendix A Reduction to Euler-invariant cochains
Suppose a Lie algebra $`g`$ possesses an element $`e_0`$, called the Euler field, and a basis of eigenvectors for the adjoint operator $`ad_{e_0}:X[X,e_0]`$. For instance, let $`\{e_k,k,k0\}`$ be the basis of elements such that $`ad_{e_0}(e_k)=ke_k`$.
The adjoint operator $`ad_{e_0}`$ is a derivation of the Lie algebra $`g`$. It can be extended to the tensor powers of $`g`$ as a graded derivation (with respect to the tensor product), and consequently to the Leibniz cochains $`CL^{}(g)`$.
For any $`\lambda `$, let $`CL_{(\lambda )}^{}(g)`$ be the subset of Leibniz cochains $`\gamma `$ such that $`ad_{e_0}(c)=\lambda c`$. Formula $`(iv)`$ of proposition (3.1) in means that the adjoint operator $`ad_{e_0}`$ commutes with the Leibniz differential. Hence, $`CL_{(\lambda )}^{}(g)`$ is a subcomplex for any $`\lambda `$, and the Leibniz complex $`CL^{}(g)=_\lambda CL_{(\lambda )}^{}(g)`$ splits up into a direct sum of such subcomplexes. In particular, the elements $`c`$ of $`CL_{(0)}^{}(g)`$ are called Euler-invariant cochains, because $`ad_{e_0}(c)=0`$.
The same splitting, as a completed direct sum, occurs if the Lie algebra has a topology and the basis $`\{e_k,k,k0\}`$ is a topological one.
As for Lie cohomology, cf. theorem 1.5.2 , we then have:
###### Lemma A.1
Suppose $`g`$ is a Lie algebra. Under the previous assumptions the Euler-invariant cohomology determines the Leibniz cohomology of $`g`$, that is,
$$HL_{cont}^{}(g)H^{}(CL_{(0)}^{}(g)).$$
$`\mathrm{𝐏𝐫𝐨𝐨𝐟}:`$We construct a contracting homotopy for the cochain complexes $`CL_{(\lambda )}^{}(g)`$ with $`\lambda 0`$. For $`\lambda 0`$, let $`D_{(\lambda )}^p:CL_{(\lambda )}^{p+1}(g)CL_{(\lambda )}^p(g)`$ be defined by
$$(D_{(\lambda )}^pc)(x_1,\mathrm{},x_{p1})=c(x_1,\mathrm{},x_{p1},e_0)$$
for any $`p`$-cochain $`c`$. Cartan’s formula (formula $`(i)`$ in proposition (3.1) ) shows that
$$(d\stackrel{~}{D}_{(\lambda )}^pc)(x_1,\mathrm{},x_p)=(\lambda \stackrel{~}{D}_{(\lambda )}^{p+1}d)(c)(x_1,\mathrm{},x_p).$$
where we have set $`\lambda =_{i=1}^p\lambda _i`$ with $`ad_{e_0}(x_i)=:\lambda _ix_i`$, and $`\stackrel{~}{D}_{(\lambda )}^p:=(1)^pD_{(\lambda )}^p`$. This means that $`\stackrel{~}{D}_{(\lambda )}^p`$ is a contracting homotopy for $`CL_{(\lambda )}^{}(g)`$ for $`\lambda 0`$. $`\mathrm{}`$
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# The Connection between Globular Cluster Systems and the Host Galaxies
## 1 Background Motivation
Globular clusters are relatively homogeneous entities containing stars of a single age and metallicity. Although we do not fully understand the relative efficiency of forming globular cluster (GC) stars over galactic field stars, there is good evidence that when a starburst occurs in a galaxy, GCs will also form (e.g. Ho 1997). Thus GCs can provide a unique probe of the star formation and chemical enrichment history of galaxies. To fully exploit this ‘probe’ one needs to have an idea for how GCs form in galaxies. Current scenarios include mergers (Schweizer 1987; Ashman & Zepf 1992), two phase in situ (Forbes, Brodie & Grillmair 1997) and tidal stripping/accretion (Forte, Martinez & Muzzio 1982; Cote, Marzke & West 1998).
A direct comparison of GC and host galaxy colours was first attempted in the late 1970s and early 1980s (e.g. Hanes 1977; Forte, Strom & Strom 1981). A difference in colour was observed, in the sense that GCs were $``$0.5 dex more metal–poor than galaxy stars. This lead Forte et al. (1981) to conclude that “…the chemical enrichment history of the globular cluster system and of the spheroidal component must have differed, and that the globular clusters are likely to form a component dynamically as well as chemically distinct from the spheroid.”
However another line of argument has developed indicating that there is a link between GCs and their host galaxy. This involves the correlation between the mean metallicity of the GC system and galaxy luminosity. Such a trend was first suggested by van den Bergh (1975) and shown by Brodie & Huchra (1991), but for only 10 galaxies. It was later disputed by Ashman & Bird (1993) claiming that there was no correlation, and by Perelmuter (1995) claiming that the relation only existed for spirals. In 1996, Forbes et al. showed a strong trend over 10 magnitudes for 35 early type galaxies. Subsequent larger samples have confirmed a correlation (e.g. Durrell et al. 1996). Although the scatter is large, a linear fit has a slope very similar to that for the galaxy stellar metallicity versus galaxy luminosity relation, i.e. Z $``$ L<sup>0.4</sup> (see Brodie & Huchra 1991).
The galaxy relations, like the Mg<sub>2</sub>$`\sigma `$ and the colour – magnitude relations, are generally believed to be correlations of metallicity and mass (Kodama & Arimoto 1997; Forbes, Ponman & Brown 1998). Thus by analogy, the mean metallicity of the GC system ‘knows about’ the mass of the galaxy in which it resides; the chemical enrichment of GCs is directly linked to the galaxy formation process.
It wasn’t until the 1990s that we could easily understand the apparent metallicity difference between the GC system and the host galaxy (as noted by Forte et al. 1981) – it is due to the presence of two distinct GC subpopulations (e.g. Lee & Geisler 1993). The metallicity difference being due to the relative mix of metal–rich and metal–poor GC subpopulations. It appears that the metal–rich GCs have a mean metallicity that is similar to the host galaxy stars with the metal–poor ones $``$1 dex lower (e.g. Geisler et al. 1996; Forbes, Brodie & Huchra 1997). In the Forte et al. (1981) study, their sampling was dominated by the blue GCs leading them to suggest that the GC system and the host galaxy were ‘distinct’.
Revisiting the metallicity – luminosity relation, Forbes et al. (1997) showed that for 11 GC systems the metallicity of the metal–rich population correlated with galaxy luminosity, but the metal–poor one did not (see also Burgarella, Kissler–Patig & Buat 2000). Forbes et al. (1997) concluded that only the metal–rich GCs were closely linked to the formation process of the galaxy.
However one caveat about the interpretation of the GC metallicity – galaxy luminosity relation is that metallicities are usually based on colours that have been converted into \[Fe/H\] assuming a Galactic relation. This has two limitations:
i) The Galactic relation is strictly only valid for very old GCs. The limited current evidence suggests that GCs around ellipticals are indeed very old (Kissler–Patig et al. 1998; Cohen, Blakeslee & Rhyzov 1998; Puzia et al. 1999), but this may not always be the case.
ii) Few Galactic GCs have solar or supersolar metallicities, whereas this is often the case for elliptical galaxy GCs. Recently a conversion from V–I to supersolar metallicities has been derived by Kissler–Patig et al. (1998), but it is still somewhat uncertain for other colour systems.
So although it seems fairly probable that the metal–rich GCs formed from gas that has been processed within the potential well of the galaxy, previously arguments relied on the assumption that the Galactic relation is applicable. In this paper, we use direct measurements of GC colours and compare these to galaxy internal velocity dispersions for a sample of 28 galaxies. Our approach alleviates some of the problems associated with metallicity – luminosity relations mentioned above. In addition, we avoid the distance dependence present in the metallicity – luminosity relation. We also examine the GC and galaxy halo colours in detail for a few individual galaxies. We find that the mean colour of the red GCs correlates with galaxy velocity dispersion. Furthermore, in well–studied galaxies the mean colour of the red GCs is almost identical to the galaxy halo (bulge) stars. We briefly discuss the implications of our findings for GC formation scenarios.
## 2 Conversion to a Single Colour System
The bulk of high quality GC colours, available in the literature, come from Hubble Space Telescope studies, and most of these use the F555W (V) and F814W (I) filters. Although this choice of filters does not provide as much metallicity ‘leverage’ as say B–I or C–T1 (on the Washington system), a large number of early type galaxies are now known to possess GC systems with bimodal V–I colour distributions. Several more are known from observations in other colour systems. Before compiling a homogeneous dataset of GC colours we need to convert the bimodality seen using other colour systems into V–I. We will use a V–I vs colour relation based on Galactic GCs, however we first need to confirm that it is applicable for the GC systems of early type galaxies.
In Fig. 1 we show the extinction corrected (V–I)<sub>o</sub> versus (C–T1)<sub>o</sub> and (B–I)<sub>o</sub> colour for Galactic GCs with low reddenings (i.e. E(B–V) $`<`$ 0.1) taken from Reed, Hesser & Shaw (1988) and Harris & Canterna (1977). The best bisector fits are shown by a solid lines. We have also made a fit to transform B–R to V–I (not shown in Fig. 1). The fits are the following:
(V–I)<sub>o</sub> = 0.49(C–T1)<sub>o</sub> \+ 0.32
(V–I)<sub>o</sub> = 0.51(B–I)<sub>o</sub> \+ 0.11
(V–I)<sub>o</sub> = 0.68(B–R)<sub>o</sub> \+ 0.15
The typical rms about the fits is $`\pm `$ 0.03 mag. Thus one could expect to estimate V–I for a Galactic type GC with an accuracy of $`0.03`$ mag. from a different colour.
Fig. 1 also shows the mean blue and red colours for the GCs in the well–studied galaxies NGC 1399, NGC 4472 (M49) and NGC 4486 (M87). Here we have plotted the (C–T1)<sub>o</sub> measurements for NGC 1399 (Ostrov, Forte & Geisler 1998), NGC 4472 (Geisler, Lee & Kim 1996), and NGC 4486 (Lee & Geisler 1993) and the (B–I)<sub>o</sub> values for NGC 1399 (Forbes et al. 1998). The corresponding (V–I)<sub>o</sub> values are NGC 1399 (Kissler–Patig et al. 1997), NGC 4472 (Puzia et al. 1999) and NGC 4486 (Kundu et al. 1999). The GC systems of these galaxies are consistent with the Galactic relation, including an extrapolation to redder (more metal–rich) colours. Thus we feel confident in our transformation from C–T1 and B–I into V–I. We do not have any independent confirmation of the B–R transformation but have no reason to doubt that it too is applicable to GCs in ellipticals.
In Table 1 we list the mean (V–I)<sub>o</sub> colours for the GC subpopulations in 28 galaxies. Of these, only four have been transformed from other colour systems (as noted in Table 1). The estimated photometric errors for the mean colours are typically $`\pm `$ 0.05 mag. (this includes the uncertainty of transformation from one colour system to V–I). The confidence of bimodality is given by ‘Yes’ for definite and ‘Likely’ for probable. In most cases the statistical significance of the bimodality can be found in the original reference.
## 3 Results and Discussion
### 3.1 Trends with Velocity Dispersion
In Fig. 2 we show the mean (V–I)<sub>o</sub> colour of the GC subpopulations versus the galaxy velocity dispersion for all of the galaxies listed in Table 1. The velocity dispersion data come from Prugniel & Simien (1996).
The blue GCs show no correlation with galaxy velocity dispersion. Indeed they reveal a fairly constant colour of (V–I)<sub>o</sub> = 0.954 $`\pm `$ 0.008. This is very similar to the overall mean colours of the Milky Way (0.94 $`\pm `$ 0.01) and M31 (0.96 $`\pm `$ 0.01) GC systems (see Barmby et al. 2000). If we exclude the disk/bulge populations of GCs in these spirals, then the halo GCs in both galaxies have a mean (V–I)<sub>o</sub> $``$ 0.92 (Barmby et al. 2000). Thus the blue GCs in spirals and early type galaxies have a very similar age and metallicity, but there are hints that those in spirals may be slightly younger and/or more metal–poor.
The red GCs, on the other hand, reveal a strong correlation (at 3$`\sigma `$ significance) of mean colour with galaxy velocity dispersion. The best fit line for the red GCs (V–I = 0.23 log $`\sigma `$ \+ 0.61) is shown in Fig. 2. Galaxy velocity dispersion is a tracer of galaxy mass. Thus the colour (and presumably metallicity) of the red GCs is directly linked to the depth of the galaxy’s potential well.
We note that if we had converted the C–T1 measurements of Secker et al. (1995) for NGC 3311 then the colours of the blue (V–I = 1.15) and red (V–I = 1.28) GCs would deviate strongly from the trends seen in Fig. 2. However we used the more recent results of Brodie, Larsen & Kissler–Patig (2000) which indicate normal GC colours for NGC 3311. We suspect zero point errors in the Secker et al. (1995) analysis. This may also affect the C–T1 colours (Zepf et al. 1995) for the NGC 3923 GC system which was observed on the same run (we have assigned it a larger error). Indeed NGC 3923 is one of the more deviant points in Fig. 2.
### 3.2 Trends with Galaxy Halo Colour
Having shown that the red GC subpopulation has a direct connection with the host galaxy, while the blue subpopulation is unconnected we now explore the trend between mean GC colour and galaxy halo colour.
In order to best compare GC and galaxy colours, the colours should come from the same observations, and sample a similar galactocentric annulus around any given galaxy as radial gradients may exist in the galaxy and GC subpopulations (e.g. Ostrov et al. 1998; Forte et al. 2000). We have also decided to restrict our analysis to the wide–area studies conducted using C–T1 (which has the best metallicity sensitivity of GC studies). This excludes HST studies which typically probe only the galaxy inner regions in V–I.
With these criteria we found four galaxies (the Secker et al. (1995) C–T1 data on NGC 3311 has been excluded for reasons mentioned above). The four are NGC 1399 (Ostrov et al. 1998), NGC 1427 (Forte et al. 2000), NGC 3923 (Zepf et al. 1995), NGC 4472 (Geisler et al. 1996; Lee & Kim 2000) and NGC 4486 (Geisler et al. 2000). We note that the zero points and hence colours of NGC 3923 may be in error (as discussed above) but as the galaxy and GCs colours come from the same observation this will not effect the relative colour difference. The two papers referenced for NGC 4472 use the same data collected in 1993. Where possible, colours have been taken from similar galactocentric radii. In the case of NGC 1399, Ostrov et al. (1998) quote galaxy and GC colours in three radial bins for the red subpopulation and two bins for the blue ones.
The mean GC colours versus galaxy halo colour is shown in Fig. 3. The mean blue GC colours do not reveal a strong trend with halo colour, and in particular lie well away from a one-to-one relation. However the red GCs do lie close to a one-to-one relation. Within the errors, the red GCs have identical C–T1 colours to those of their host galaxy. So not only do red GCs have a metallicity that is linked to the depth of the galaxy’s potential well (Fig. 2), for at least some galaxies they also have the same metallicity as the galaxy halo stars. This suggests a very well coordinated star formation and chemical enrichment history.
For the NGC 1399 red GCs, the agreement is not only restricted to the mean colours discussed above, but also in the sense that they share the same colour gradient as the underlying galaxy halo (see table 5 in Ostrov et al. 1998). If the galaxy halo was in fact related to the overall GC system, then halo colour gradients would reflect the varying ratio of the blue to red GCs with galactocentric radius. However if this were the case, we would expect the halo colour gradient to get significantly bluer with increasing radius (as the blue GCs dominate). For NGC 1399 (Ostrov et al. 1998) this is not the case indicating that the halo colours do not reflect the blue GCs, but rather the red ones.
Finally, we note that in recent ground–breaking work, Harris et al. (1999) used HST to resolve the stars in the halo of NGC 5128. They showed that the red halo stars and the metal–rich GCs have the same metallicity, and concluded that they formed in the same star formation event.
### 3.3 How did the Red Globular Clusters Form ?
The simplest explanation for the fact that the red GCs and the galaxy stars may have similar colours, and that red GCs colours correlate with galaxy velocity dispersion, is that they formed at the same time from the same chemically enriched gas. This would naturally favour the two phase collapse scenario (Forbes et al. 1997). In this scenario the red GCs form almost contemporaneously with the galaxy field stars, and hence share their properties. The trends described here are a natural consequence of this idea.
Are these trends expected in the merger picture (Ashman & Zepf 1992) ? After a gaseous merger the galaxy will consist of at least two stellar and two GC populations (the first associated with the progenitor galaxies and the other created in the merger). The relative ratios of these two components depend on how much gas is available to create the new stars and GCs. The gas fraction would be high at early epochs. To explain Fig. 2 one could argue that the more metal–rich (redder) GCs formed in the larger (higher velocity dispersion) galaxies. However to explain the nearly identical colours of the galaxy halo and red GCs (Fig. 3), one must then argue that the halo is dominated by new stars formed at the same time, from the same gas, as the new (red) GCs. This would imply that the progenitors were largely gaseous and hence the merger occurred at high redshift. The distinction between the merger and collapse pictures becomes blurred. Various other difficulties for the merger model are discussed in detail by Forbes et al. (1997).
In the accretion model of Cote et al. (1998), the red GCs form in situ but the blue ones are accreted, along with starlight, from stripped dwarf galaxies. This accumulation of metal–poor material into the halo of the primary galaxy would tend to make the halo light bluer than the red GCs – in contradiction to what is seen in Fig. 3.
## 4 Concluding Remarks
Here we have presented evidence that the mean colour of the blue GCs in early type galaxies are unrelated to their host galaxy, supporting the work of Forbes et al. (1997) and Burgarella et al. (2000). In terms of a collapse scenario, where the blue GCs form in a proto–galactic dark matter halos, they do not appear to have ‘memory’ of the collapse phase suggesting they formed pre–collapse or possibly quite independently of the eventual host galaxy (Burgarella et al. 2000). We note however, that this is probably not true for the Milky Way halo globular clusters (the ‘blue’ subpopulation). These metal–poor globular clusters do show evidence that the majority formed within our Galactic halo (van den Bergh 1996).
The red GCs, on the other hand, do ‘know’ about the galaxy they occupy in terms of its potential well and halo colour. This suggests that the red GCs and the host galaxy share a common formation event. This naturally favours the two phase collapse picture, in which the red GCs and the halo stars are formed together (Forbes et al. 1997). The merger picture (Ashman & Zepf 1992) is only compatiable if the merger took place at early epochs involving gaseous progenitors.
Acknowledgements
We thank A. Georgakakis, S. Larsen and S. van den Bergh for useful discussions. S. Larsen also kindly provided results on NGC 3311 before publication. DAF thanks La Plata Observatory for their hospitality during August 1999 when much of work was carried out. Finally we thank the referee J. Brodie for her many useful comments and suggestions.
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# Rigidity properties of diagram groups
## 1 Introduction
The class of diagram groups was introduced by Meakin and Sapir in 1993. Kilibarda obtained the first results about diagram groups in , . The theory was further developed in , . It turned out that many important groups (including the R. Thompson group $`F`$) are diagram groups. On the other hand, diagram groups satisfy some interesting properties, and there exists a deep similarity between combinatorics on diagrams and combinatorics on words. Recent results by D. Farley show that diagram groups act by isometries on CAT(0)-spaces. This allowed him to prove that the R. Thompson group satisfies the rational Novikov conjecture.
The first (and still very useful) definition of diagram groups (see ,) was algebraic. From this point of view, every diagram group $`𝒟(𝒫,w)`$ is determined by a semigroup presentation $`𝒫`$ and a distinguished word $`w`$. One can give an equivalent topological definition of diagram groups . From the topological point of view, diagram groups are determined by a directed 2-complex $`K`$ (all edges have directions, every cell is bounded by two positive paths, the top and the bottom), and a distinguished positive path $`p`$. Diagram groups are similar to second relative homotopy groups of 2-complexes, only one needs to consider directed 2-complexes and homotopies consisting of positive paths only (we call them directed homotopies).
Here is an informal definition of diagram groups (see for details).
Let $`𝒫=\mathrm{\Sigma }`$ be a semigroup presentation where $`\mathrm{\Sigma }`$ is an alphabet and $``$ is the set of defining relations.
Any diagram over $`𝒫`$ is obtained as follows. Start with a positive (horizontal) path $`p`$ on the plane labeled by some word $`w`$ over $`\mathrm{\Sigma }`$ (that is, a linear oriented labeled graph with $`|w|`$ edges which form a path, whose label is $`w`$). This is a trivial $`(w,w)`$-diagram, and $`p`$ is the top and the bottom path of this diagram.
Next find a subword in $`w`$ which is equal to $`u`$ (or $`v`$) for some relation $`u=v`$ in $``$: $`p=p^{}qp^{\prime \prime }`$ where the label of $`q`$ is $`u`$ (resp. $`v`$). Below $`p`$, draw a path $`q^{}`$ labeled by $`v`$ (resp. $`u`$) whose initial and terminal vertices coincide with the initial and terminal vertices of $`q`$. The path $`q(q^{})^1`$ must bound a region on the plane (called a cell). The result of this operation is a one-cell diagram whose top path is labeled by $`w`$ and the bottom path is labeled by the word obtained from $`w`$ by replacing $`u`$ by $`v`$ (resp. $`v`$ by $`u`$). Attaching a new cell to the bottom path of the diagram, we get a diagram with two cells, etc. Every diagram $`\mathrm{\Delta }`$ is a plane labeled oriented graph which tesselates a region of the plane between two positive paths $`\text{top}(\mathrm{\Delta })`$ and $`\text{bot}(\mathrm{\Delta })`$. If $`w`$ is the label of $`\text{top}(\mathrm{\Delta })`$ and $`w^{}`$ is the label of $`\text{bot}(\mathrm{\Delta })`$ then $`\mathrm{\Delta }`$ is called a $`(w,w^{})`$-diagram.
Two diagrams are called equal if there exists an isotopy of the plane which takes one of the diagrams to the other one.
A diagram is called reduced if it does not contain dipoles. A dipole is a pair of cells such that the bottom path of one of them coincides with the top path of the other one and these cells are mirror images of each other. If a diagram contains a dipole, the two cells forming the dipole can be removed. So every diagram can be reduced. By the theorem of Kilibarda the reduced form of every diagram is unique.
Fix a word $`w`$ and consider the set $`𝒟(𝒫,w)`$ of all $`(w,w)`$-diagrams over $`𝒫`$. One can multiply two diagrams $`\mathrm{\Delta }_1`$ and $`\mathrm{\Delta }_2`$ in $`𝒟(𝒫,w)`$ by gluing together $`\text{bot}(\mathrm{\Delta }_1)`$ and $`\text{top}(\mathrm{\Delta }_2)`$ and reducing the resulting diagram. This operation is associative, the trivial $`(w,w)`$-diagram plays the role of the identity element, and every diagram $`\mathrm{\Delta }`$ has an inverse, the mirror image of $`\mathrm{\Delta }`$. Thus $`𝒟(𝒫,w)`$ is a group which is called the diagram group over the presentation $`𝒫`$ with base word $`w`$.
Since we are going to use only the algebraic definition of diagram groups, we do not give here a precise topological definition. Let us only mention that the directed complex corresponding to a semigroup presentation is similar to the standard 2-complex of a group presentation. It has one vertex, one oriented edge for each generator and one oriented cell for each relation $`u=v`$ with bottom path $`u`$ and the top path $`v`$. Then the word $`w`$ in the algebraic definition of a diagram group turns into a positive path $`w`$ in the directed 2-complex, and every $`(w,w)`$-diagram is a planar representative of a directed homotopy from $`w`$ to $`w`$. Conversely, every directed $`(w,w)`$-homotopy is represented by a $`(w,w)`$-diagram. The product of homotopies corresponds to the product of diagrams. Equivalent diagrams correspond to equivalent (isotopic) homotopies. This allows one to translate every statement about diagram groups from the algebraic language to the topological language and back.
The relation between diagram groups and semigroup presentations (directed complexes) is not rigid. For example, if the presentation $`𝒫`$ is aspherical, then the diagram groups are trivial (regardless of the base). On the other hand, presentations of finite semigroups may correspond to “large” diagram groups. In particular, the diagram group corresponding to the presentation $`xx^2=x`$ of the trivial semigroup is the well known R. Thompson group $`F`$ (for every base). The directed complex corresponding to this presentation is the well known Dunce hat which can be obtained from the triangle
by gluing all three sides according to their direction. It is easy to construct other semigroup presentations (directed complexes) with diagram groups isomorphic to $`F`$. Nevertheless in this paper we show that if $`F`$ appears in a diagram group of a directed complex (resp. presentation of a semigroup) then Dunce hat maps into the complex (the semigroup contains an idempotent). Thus there is a rigid relationship between $`F`$ and the Dunce hat (the presentation $`xx^2=x`$).
Recall that the group $`F`$ can be given by the following infinite presentation:
$$x_0,x_1,\mathrm{}x_j^{x_i}=x_{j+1}(j>i).$$
It has also a finite presentation
$$x_0,x_1x_2^{x_1}=x_3,x_3^{x_1}=x_4,$$
(1)
where $`x_2=x_1^{x_0}`$, $`x_3=x_2^{x_0}`$, $`x_4=x_3^{x_0}`$ by definition.
###### Theorem 1.
The following conditions are equivalent.
1. For some word $`w`$, the diagram group $`𝒟(𝒫,w)`$ contains an isomorphic copy of the R. Thompson group $`F`$.
2. The semigroup given by $`\mathrm{\Sigma }`$ contains an idempotent.
A part of this theorem, namely the implication $`21`$, has been proved in \[6, Theorem 25\]. We asked \[6, Problem 2\] whether the converse is true. Theorem 1 gives an affirmative answer to this question.
The topological formulation of Theorem 1 is the following
###### Theorem 2.
Let $`K`$ be a directed complex. Then the following conditions are equivalent.
1. A diagram group corresponding to $`K`$ contains an isomorphic copy of the R. Thompson group $`F`$.
2. The complex $`K`$ contains a positive non-empty path $`t`$ which is directly homotopic to its square.
3. There exists a directed morphism from the Dunce hat to $`K`$.
Clearly Theorem 1 and 2 are equivalent We shall prove the theorem in the first formulation.
Recall also that in \[6, Theorem 24\], we have proved a similar rigidity theorem for the restricted wreath product $`𝐙wr𝐙`$. Similar rigidity theorems might be true for other diagram groups as well.
## 2 Proof of the rigidity theorem
We need one auxiliary geometric fact. Let $`𝒫`$ be a semigroup presentation and let $`\mathrm{\Delta }`$ be a diagram over $`𝒫`$. For any two vertices $`o^{}`$, $`o^{\prime \prime }`$ in $`\mathrm{\Delta }`$ we put $`o^{}o^{\prime \prime }`$ whenever there exists a positive path in $`\mathrm{\Delta }`$ from $`o^{}`$ to $`o^{\prime \prime }`$. It is easy to see that the labels of any two positive paths from $`o^{}`$ to $`o^{\prime \prime }`$ are equal modulo $`𝒫`$ (see ). So one can define the element $`\mu (o^{},o^{\prime \prime })`$ in the monoid $`M`$ presented by $`𝒫`$. This element is represented in $`M`$ by the label of any positive path from $`o^{}`$ to $`o^{\prime \prime }`$.
Recall also that for every $`(u,v)`$-diagram $`\mathrm{\Delta }`$ and $`(u^{},v^{})`$-diagram $`\mathrm{\Delta }^{}`$ one can define the sum $`\mathrm{\Delta }+\mathrm{\Delta }^{}`$ by gluing the terminal vertex of $`\text{top}(\mathrm{\Delta })`$ with the initial vertex of $`\text{top}(\mathrm{\Delta }^{})`$. The result is a $`(uu^{},vv^{})`$-diagram.
###### Lemma 3.
Let $`𝒫`$ be a semigroup presentation. Let $`M`$ denote the monoid presented by $`𝒫`$. Suppose that $`\mathrm{\Delta }`$ is a $`(uv,uv)`$-diagram over $`𝒫`$. Let $`o_1`$ $`(`$resp. $`o_2)`$ be the vertex in the top (bottom) path of $`\mathrm{\Delta }`$ that subdivides it into a product of two paths labeled by $`u`$ and $`v`$. Suppose that $`o`$ is a vertex in $`\mathrm{\Delta }`$, where $`oo_1`$, $`oo_2`$. If $`\mathrm{\Delta }`$ is equivalent to a sum of a $`(u,u)`$-diagram and a $`(v,v)`$-diagram, then $`\mu (o,o_1)=\mu (o,o_2)`$. $`(`$It would be more precise to write $`\mu _\mathrm{\Delta }`$ but it will always be clear what diagram we refer to$`)`$.
Proof. Obviously, the reduced form of $`\mathrm{\Delta }`$ is a sum of a $`(u,u)`$-diagram and a $`(v,v)`$-diagram. Suppose that we need to cancel $`m0`$ pairs of dipoles in order to reduce $`\mathrm{\Delta }`$. We prove the claim by induction on $`m`$. If $`m=0`$ then the conclusion is obvious since in this case $`o_1=o_2`$. Let $`m>0`$. Cancel a dipole that consists of two cells $`\pi _1`$ and $`\pi _2`$, where the bottom path of $`\pi _1`$ coincides with the top path of $`\pi _2`$. As a result, we get a diagram $`\mathrm{\Delta }^{}`$ that can be reduced in $`m1`$ step. Let $`p_1`$ be the top path of $`\pi _1`$, $`p_2`$ be the bottom path of $`\pi _2`$, and let $`p`$ be the common boundary of $`\pi _1`$ and $`\pi _2`$.
Suppose first that $`o`$ is a vertex that does not disappear in $`\mathrm{\Delta }^{}`$, that is, $`o`$ does not belong to $`p`$ as an inner point. In this case, for any positive path from $`o`$ to $`o_1`$ in $`\mathrm{\Delta }`$, we can find a positive path from $`o`$ to $`o_1`$, which does not contain $`p`$ as a subpath (just replace $`p`$ by $`p_1`$). The same is true for positive paths from $`o`$ to $`o_2`$. The vertices $`o`$, $`o_1`$, $`o_2`$ still exist in $`\mathrm{\Delta }^{}`$, and the elements $`\mu (o,o_1)`$, $`\mu (o,o_2)`$ do not change when we replace $`\mathrm{\Delta }`$ by $`\mathrm{\Delta }^{}`$. Applying our inductive assumption to $`\mathrm{\Delta }^{}`$, we see that these elements are equal there. Thus they are equal for $`\mathrm{\Delta }`$, too.
Now suppose that $`o`$ disappears in $`\mathrm{\Delta }^{}`$. Thus $`p`$ is subdivided by $`o`$ into two paths, say, $`q`$ and $`r`$. Let $`\overline{o}`$ be the terminal point of $`r`$. Obviously, any positive path from $`o`$ to $`o_1`$ or $`o_2`$ begins with $`r`$. By the previous paragraph, $`\mu (\overline{o},o_1)=\mu (\overline{o},o_2)`$. Now it remains to notice that $`\mu (o,o_i)=\nu \mu (\overline{o},o_i)`$, where $`\nu M`$ is the element represented by the label of $`r`$ ($`i=1,2`$). This completes the proof.
Let $`H`$ be a group, $`y_0,y_1H`$. Suppose that $`y_0`$, $`y_1`$ do not commute in $`H`$ and satisfy relations (1), that is, $`y_1^{y_0y_1}=y_1^{y_0^2}`$, $`y_1^{y_0^2y_1}=y_1^{y_0^3}`$. Since all proper homomorphic images of $`F`$ are abelian , it is clear that $`y_0`$, $`y_1`$ generate $`F`$ as a subgroup of $`H`$. In this case, we say that an ordered pair $`y_0`$, $`y_1`$ generates $`F`$ canonically. We also introduce elements $`y_i`$ for $`i2`$ by $`y_i=y_1^{y_0^{i1}}`$.
Let us recall some definitions. We refer to \[5, Section 15\] for details. Let $`𝒫=\mathrm{\Sigma }`$ be any semigroup presentation. A $`(w,w)`$-diagram $`\mathrm{\Delta }`$ over $`𝒫`$ is called absolutely reduced provided $`\mathrm{\Delta }^n`$ is reduced for every $`n1`$. For any $`(w,w)`$-diagram $`\mathrm{\Delta }`$ over $`𝒫`$, where $`w\mathrm{\Sigma }^+`$, there exists a word $`v\mathrm{\Sigma }^+`$, a $`(w,v)`$-diagram $`\mathrm{\Psi }`$ and an absolutely reduced $`(v,v)`$-diagram $`\overline{\mathrm{\Delta }}`$ such that $`\mathrm{\Delta }=\mathrm{\Psi }\overline{\mathrm{\Delta }}\mathrm{\Psi }^1`$. One can decompose $`\overline{\mathrm{\Delta }}`$ into a sum $`A_1+\mathrm{}+A_m`$ of spherical diagrams. Here each nontrivial summand cannot be decomposed into a sum of spherical diagrams. We also assume that for any $`i`$ ($`1i<m`$) at least one of the diagrams $`A_i`$, $`A_{i+1}`$ is nontrivial. The summands $`A_i`$ ($`1im`$) are called components of $`\overline{\mathrm{\Delta }}`$. The number of nontrivial components does not depend on the choice of $`\overline{\mathrm{\Delta }}`$. So it can be denoted by $`𝐜𝐨𝐦𝐩(\mathrm{\Delta })`$.
Let $`G=𝒟(𝒫,w)`$ be a diagram group. For any $`(w,v)`$-diagram $`\mathrm{\Psi }`$ over $`𝒫`$, where $`v\mathrm{\Sigma }^+`$, we have an isomorphism $`\psi :GH=𝒟(𝒫,v)`$ that takes any diagram $`\mathrm{\Delta }G`$ to $`\mathrm{\Psi }^1\mathrm{\Delta }\mathrm{\Psi }`$. For any $`(w,w)`$-diagram $`\mathrm{\Delta }`$ over $`𝒫`$ we can construct an isomorphism $`\psi `$ defined above such that the diagram $`\psi (\mathrm{\Delta })=\mathrm{\Psi }^1\mathrm{\Delta }\mathrm{\Psi }`$ will be absolutely reduced. In this case we will often assume without loss of generality that $`\mathrm{\Delta }`$ is absolutely reduced up to changing the base of our diagram group.
Suppose that $`A=A_1+\mathrm{}+A_m`$ and $`B=B_1+\mathrm{}+B_n`$ are absolutely reduced diagrams each decomposed into a sum of components. Let $`A_i`$ ($`1im`$) and $`B_j`$ ($`1jn`$) be $`(v_i,v_i)`$\- and $`(w_j,w_j)`$-diagrams, respectively. If there exists a $`(w,v)`$-diagram $`\mathrm{\Gamma }`$ such that $`A=\mathrm{\Gamma }^1B\mathrm{\Gamma }`$, where $`v=v_1\mathrm{}v_m`$, $`w=w_1\mathrm{}w_n`$, then $`m=n`$ and $`\mathrm{\Gamma }`$ can be decomposed into a sum $`\mathrm{\Gamma }_1+\mathrm{}+\mathrm{\Gamma }_m`$ of $`(w_i,v_i)`$-diagrams $`\mathrm{\Gamma }_i`$ such that $`A_i=\mathrm{\Gamma }_i^1B_i\mathrm{\Gamma }_i`$ ($`1im`$). Any element $`C`$ in the centralizer of $`A`$ in $`𝒟(𝒫,v)`$ can be decomposed into a sum $`C=C_1+\mathrm{}+C_m`$, where $`C_i`$ is a $`(v_i,v_i)`$-diagram that commutes with $`A_i`$ ($`1im`$). If $`A_i`$ is nontrivial then its centralizer is cyclic so $`A_i`$ and $`C_i`$ belong to the same cyclic subgroup. It is easy to see that one can change the base in such a way that both diagrams $`A`$ and $`C`$ become cyclically reduced (see \[6, Theorem 17\]).
The following theorem is stronger than the implication $`12`$ in Theorem 1.
###### Theorem 4.
Let $`𝒫=\mathrm{\Sigma }`$ be a semigroup presentation, $`w\mathrm{\Sigma }^+`$. If the diagram group $`G=𝒟(𝒫,w)`$ contains an isomorphic copy of R. Thompson’s group $`F`$, then the semigroup $`S`$ presented by $`𝒫`$ contains an idempotent. Moreover, $`G`$ contains a copy of $`F`$ if and only if there exist words $`w_1,w_2\mathrm{\Sigma }^{}`$, $`e\mathrm{\Sigma }^+`$ such that equalities $`w=w_1ew_2`$, $`e^2=e`$ hold modulo $`𝒫`$.
Proof. Suppose that $`G=𝒟(𝒫,w)`$ contains an isomorphic copy of $`F`$. Then there exist $`(w,w)`$-diagrams $`Y_0`$, $`Y_1`$ over $`𝒫`$ that generate $`F`$ canonically. We assume that the total number of their components, that is, $`𝐜𝐨𝐦𝐩(Y_0)+𝐜𝐨𝐦𝐩(Y_1)`$, is minimal possible. Note that this number does not change if we replace $`Y_0`$, $`Y_1`$ by their conjugates $`Y_0^D`$, $`Y_1^D`$ for any $`(w,v)`$-diagram $`D`$ over $`𝒫`$, where $`v`$ is a nonempty word over $`\mathrm{\Sigma }`$.
It is easy to see that the element $`x_2x_3x_2^2F`$ commutes with $`x_i`$ for all $`i3`$. So it also commutes with $`x_3x_4x_3^2`$. Changing the base $`w`$, we can assume without loss of generality that $`D_2=Y_2Y_3Y_2^2`$ is a cyclically reduced diagram over $`𝒫`$ decomposed into the sum of components $`A_1+\mathrm{}+A_m`$, where $`A_i`$ is a $`(v_i,v_i`$)-diagram $`(1im)`$. Obviously, $`D_2`$ is nontrivial. (Otherwise $`Y_2=Y_3`$ and $`Y_0`$ commutes with $`Y_1`$ .) Since $`D_3=Y_3Y_4Y_3^2`$ is in the centralizer of $`D_2`$, we can assume that both diagrams $`D_2`$, $`D_3`$ are absolutely reduced and $`D_3=B_1+\mathrm{}+B_m`$, where $`B_i`$ commutes with $`A_i`$ for all $`1im`$. (Note that the summands $`B_i`$ are not necessarily components of $`D_3`$.)
Suppose that $`A_i`$ is nontrivial for some $`i`$. Let it be the $`j`$th nontrivial component of $`D_2`$ counting from left to right. It is clear that $`D_3=D_2^{Y_0}=D_2^{Y_1}`$. Thus $`D_2`$ and $`D_3`$ conjugate and so they have the same structure of components. Let $`B^{}`$ be the $`j`$th nontrivial component of $`D_3`$ counting from left to right. There are three possible cases: $`B^{}`$ is contained in either 1) $`B_i`$, or 2) $`B_1+\mathrm{}+B_{i1}`$, or 3) $`B_{i+1}+\mathrm{}+B_m`$. Clearly, the third case is symmetric to the second one. So we consider only the first two cases.
Case 1. It is obvious that $`B^{}=B_i`$. Hence the conjugation of $`D_2`$ by each of $`Y_0`$, $`Y_1`$ takes $`A_i`$ to $`B_i`$. This implies that each of the diagrams $`Y_0`$, $`Y_1`$ can be decomposed into a sum of three spherical diagrams with bases $`u_1=v_1\mathrm{}v_{i1}`$, $`u_2=v_i`$, $`u_3=v_{i+1}\mathrm{}v_m`$, respectively. So we have an injective homomorphism $`\varphi `$ from the Thompson group $`F`$ (generated by $`Y_0`$, $`Y_1`$) to the direct product $`𝒟(𝒫,u_1)\times 𝒟(𝒫,u_2)\times 𝒟(𝒫,u_3)`$. Denote by $`H_k`$ the projection of $`F`$ onto $`k`$th factor and let $`\psi _k`$ be the homomorphism from $`F`$ onto $`H_k`$ ($`k=1,2,3`$). The group $`F`$ embeds into $`H_1\times H_2\times H_3`$. Therefore, at least one of the three groups $`H_k`$ is not abelian. Then it must be isomorphic to $`F`$ because all proper homomorphic images of $`F`$ are abelian. So let $`H_k`$ be non-abelian. Let us show that $`k=1`$ or $`k=3`$.
We know that the diagrams $`A_i`$, $`B_i`$ belong to the same cyclic subgroup. By \[5, Theorem 15.30\], we may assume that they belong to the maximal cyclic subgroup $`K`$ of the diagram group $`𝒟(𝒫,v_i)`$. Let us establish that any $`(v_i,v_i)`$-diagram $`D`$ over $`𝒫`$ such that $`A_i^D=B_i`$, also belongs to $`K`$. Let $`C`$ be the generator of $`K`$. By definition, $`A_i`$ is nontrivial. So $`B_i`$ is also nontrivial and so we have $`A_i=C^r`$, $`B_i=C^s`$, where $`r`$, $`s`$ are non-zero integers. We now have $`(C^D)^r=C^s`$. So we can apply \[5, Corollary 15.28\] to conclude that there is a diagram $`C_0`$ and some integers $`p`$, $`q`$ such that $`C^D=C_0^p`$, $`C=C_0^q`$ and $`pr=qs`$. Since $`C`$ generates maximal cyclic subgroup, we have $`|p|=|q|=1`$. Thus $`C^D=C^{\pm 1}`$. If $`C^D=C^1`$, then $`(CD)^2=D^2`$. Using the fact that diagram groups have the unique extraction of roots property (\[5, Section 15\]), we deduce that $`C`$ is trivial. This is a contradiction. So $`C^D=C`$. Hence $`D`$ belongs to $`K`$ because $`K`$ coincides with its centralizer. Now we can conclude that the images of $`Y_0`$, $`Y_1`$ under $`\psi _2`$ belong to the same cyclic subgroup. So $`H_2=\psi _2(F)`$ is abelian.
We have proved that either $`H_1`$ or $`H_3`$ is isomorphic to $`F`$. It is obvious that for any diagram $`\mathrm{\Delta }`$ from the subgroup generated by $`Y_0`$, $`Y_1`$, one has $`_{k=1}^3𝐜𝐨𝐦𝐩(\psi _k(\mathrm{\Delta }))=𝐜𝐨𝐦𝐩(\mathrm{\Delta })`$. Since $`\psi _2(F)`$ is nontrivial, we see that $`𝐜𝐨𝐦𝐩(\psi _2(Y_0))+𝐜𝐨𝐦𝐩(\psi _2(Y_1))>0`$. So for any $`k=1,3`$ we have $`𝐜𝐨𝐦𝐩(\psi _k(Y_0))+𝐜𝐨𝐦𝐩(\psi _k(Y_1))<𝐜𝐨𝐦𝐩(Y_0)+𝐜𝐨𝐦𝐩(Y_1)`$. Now we can take the value of $`k`$ such that $`\psi _k(F)F`$ and replace the elements of our canonical generating pair $`Y_0`$, $`Y_1`$ by their images under $`\psi _k`$. We get another canonical generating pair with smaller total number of components. This is a contradiction, so Case 1 is impossible.
Case 2. Let $`B^{}`$ be contained in $`B_1+\mathrm{}+B_{i1}`$ as a subdiagram. We have $`B_1+\mathrm{}+B_{i1}=\mathrm{\Xi }_1+B^{}+\mathrm{\Xi }_2`$ for some spherical diagrams $`\mathrm{\Xi }_1`$, $`\mathrm{\Xi }_2`$. Let $`z`$ be the base of the diagram $`B_{i+1}+\mathrm{}+B_m`$ and let $`t`$ be the base of $`\mathrm{\Xi }_2+B_i`$. Obviously, $`t`$ is nonempty because it has a terminal segment $`v_i`$. We will show that $`t^2=t^3`$ modulo $`𝒫`$ so $`e=t^2`$ represents an idempotent in $`S`$. It will be also clear that $`w`$ belongs to the two-sided ideal in $`M`$ generated by $`e`$, where $`M=S^1`$ is the monoid presented by $`𝒫`$.
Let $`D`$ be $`Y_0`$ or $`Y_1`$. We use the fact that $`D_2^D=D_3`$. Each of the diagrams $`D_2`$, $`D_3`$ is a sum of components. According to the above description, $`D`$ can be naturally decomposed into a sum of $`m`$ diagrams (not necessarily spherical) such that the conjugation by the $`k`$th summand $`(1km`$) takes $`A_k`$ (the $`k`$th component of $`D_2`$) to the $`k`$th component of $`D_3`$ (recall that this component may not coincide with $`B_k`$). Then $`A_i`$, the $`j`$th nontrivial component of $`D_2`$, is taken to $`B^{}`$, the $`j`$th nontrivial component of $`D_3`$. The bases of diagrams to the right of $`A_i`$, $`B^{}`$ in $`D_2`$ and $`D_3`$, respectively, are $`z`$ and $`tz`$. This means that $`D`$ is a sum of an $`(xt,x)`$-diagram and a $`(z,tz)`$-diagram, where $`x`$ is the base of $`\mathrm{\Xi }_1`$.
Note that $`x_2x_3x_2^2`$ commutes with $`x_3`$. So $`Y_3`$ belongs to the centralizer of $`D_2`$. Hence $`Y_3`$ is a sum of an $`(xt,xt)`$-diagram and a $`(z,z)`$-diagram. The diagram
$$\mathrm{\Delta }Y_0^1Y_0^1Y_1Y_0Y_0,$$
equivalent to $`Y_3`$, has the following structure:
Here $`o`$ is the vertex in $`Y_1`$ that subdivides it into the sum of an $`(xt,x)`$\- and a $`(z,zt)`$-diagrams. By $`o_1`$ ($`o_2`$) we denote the vertex on the top (bottom) path of $`\mathrm{\Delta }`$ that subdivides this path into a product of paths with labels $`xt`$ and $`z`$. Clearly, there is a path in $`\mathrm{\Delta }`$ from $`o`$ to $`o_1`$ labeled by $`t^2`$ and there is a path in $`\mathrm{\Delta }`$ from $`o`$ to $`o_2`$ labeled by $`t^3`$. Applying Lemma 3, we conclude that $`t^2=t^3`$ modulo $`𝒫`$. (It is obvious that $`w`$ belongs to $`Mt^2M`$ as an element in $`S`$.)
The converse is proved in \[6, Theorem 25\].
The proof is complete.
###### Remark 5.
Given a finite semigroup presentation $`𝒫`$ and a word $`w\mathrm{\Sigma }^+`$, we cannot decide algorithmically whether the diagram group $`𝒟(𝒫,w)`$ contains $`F`$ as a subgroup. Indeed, the property of a finitely presented semigroup not to have an idempotent, is a Markov property. Let $`a`$, $`b`$ be new letters that do not belong to $`\mathrm{\Sigma }`$. Adding them to $`\mathrm{\Sigma }`$ and adding relations of the form $`ax=a`$, $`xb=b`$ ($`x\mathrm{\Sigma }`$), we get a new semigroup presentation $`𝒬`$. The diagram group $`𝒟(𝒬,ab)`$ contains $`F`$ as a subgroup if and only if $`S`$ has an idempotent, where $`S`$ is the semigroup presented by $`𝒫`$. This is clear because all idempotents in the semigroup presented by $`𝒬`$ are represented by words over $`\mathrm{\Sigma }`$ and $`ab`$ belongs to the two-sided ideal generated by any word over $`\mathrm{\Sigma }`$.
Victor Guba
Vologda State Pedagogical University,
S. Orlov Street 6,
Vologda, 160600
Russia
guba@uni-vologda.ac.ru
Mark Sapir
Vanderbilt University,
Nashville, TN 37240, U.S.A.
msapir@math.vanderbilt.edu
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# Error-correcting code on a cactus: a solvable model
## Abstract
An exact solution to a family of parity check error-correcting codes is provided by mapping the problem onto a Husimi cactus. The solution obtained in the thermodynamic limit recovers the replica symmetric theory results and provides a very good approximation to finite systems of moderate size. The probability propagation decoding algorithm emerges naturally from the analysis. A phase transition between decoding success and failure phases is found to coincide with an information-theoretic upper bound. The method is employed to compare Gallager and MN codes.
The theory of error-correcting codes concentrates on the efficient introduction of redundancy to given messages for protecting the information content against corruption. The theoretical foundations of this area were laid by Shannon’s seminal work and have been developing ever since (see and references therein). One of the main results obtained in this field is the celebrated channel coding theorem stating that there exists a code such that the average message error probability $`P_E`$, when maximum likelihood decoding is used, is upper bounded by $`P_E<e^{ME(R)}`$, where $`M`$ is the length of the encoded transmission and $`R=(`$ message information content $`)/M`$ is the code rate. The exponent $`E(R)`$ is positive for code rates below the channel capacity, corresponding to the maximal mutual information between the received and the transmitted signals, and vanishes above it. For rates $`R`$ below the channel capacity, commonly termed Shannon’s bound, the error probability can be made arbitrarily small.
The channel coding theorem is based on unstructured random codes and impractical decoders as maximum likelihood or typical sets . In the last fifty years several practical methods have been proposed and implemented, but none has been able to saturate Shannon’s bound. In 1963 Gallager proposed a coding scheme which involves sparse linear transformations of binary messages that was forgotten soon after in part due to the success of convolutional codes and the computational limitations of the time. Gallager codes have been recently rediscovered by MacKay and Neal (MN) that independently proposed a closely related code . This almost coincided with the breakthrough discovery of the high performance turbo codes . Variations of Gallager codes have displayed performance comparable (and sometimes superior) to turbo codes , qualifying them as state-of-the-art codes.
Statistical physics has been applied to the analysis of error-correcting codes as an alternative to information theory methods yielding some new interesting directions and suggesting new high-performance codes . Sourlas was the first to relate error-correcting codes to spin glass models , showing that the Random Energy Model (REM) can be thought of as an ideal code capable of saturating Shannon’s bound at vanishing code rates. This work was extended recently to the case of finite code rates and has been further developed for analysing MN codes of various structures . All of the analyses mentioned above as well as the recent turbo codes analysis relied on the replica approach under the assumption of replica symmetry. It is also worthwhile mentioning a different approach, used in the analysis of convolutional codes , of employing the transfer-matrix formalism and power series expansions. However, to date, the only model that can be analysed exactly is the REM that corresponds to an impractical coding scheme of a vanishing code rate.
In this letter we present an exact analysis to the performance of Gallager error-correcting codes on a generalisation of Bethe lattices known as Husimi cactus . We solve the model recovering results obtained by the replica symmetric theory and finding the noise level that corresponds to the phase transition between perfect decoding and a decoding failure phase, this appears to coincide with existing information-theoretic upper bounds. We experimentally show that the solution accurately approximates Gallager codes of moderate size. We also show that the probability propagation (PP) decoding algorithm emerges naturally from this framework allowing for the analysis of the practical decoding performance. Finally, we summarise the differences between Gallager and MN codes, which are somewhat obscure in the information theory literature but become explicit in this framework.
We will concentrate here on a simple communication model whereby messages are represented by binary vectors and are communicated through a Binary Symmetric Channel (BSC) where uncorrelated bit flips appear with probability $`f`$. A Gallager code is defined by a binary matrix $`𝑨=[𝑪_\mathrm{𝟏}𝑪_\mathrm{𝟐}]`$, concatenating two very sparse matrices known to both sender and receiver, with $`𝑪_\mathrm{𝟐}`$ (of dimensionality $`(MN)\times (MN)`$) being invertible - the matrix $`𝑪_\mathrm{𝟏}`$ is of dimensionality $`(MN)\times N`$.
Encoding refers to the production of a $`M`$ dimensional binary code word $`𝒕\{0,1\}^M`$ ($`M>N`$) from the original message $`𝝃\{0,1\}^N`$ by $`𝒕=𝑮^𝑻𝝃\text{(mod 2)}`$, where all operations are performed in the field $`\{0,1\}`$ and are indicated by (mod 2). The generator matrix is $`𝑮=[𝑰𝑪_\mathrm{𝟐}^\mathbf{}\mathrm{𝟏}𝑪_\mathrm{𝟏}]\text{ (mod 2)}`$, where $`𝑰`$ is the $`N\times N`$ identity matrix, implying that $`𝑨𝑮^𝑻\text{ (mod 2)}=0`$ and that the first $`N`$ bits of $`𝒕`$ are set to the message $`𝝃`$. In regular Gallager codes the number of non-zero elements in each row of $`𝑨`$ is chosen to be exactly $`K`$. The number of elements per column is then $`C=(1R)K`$, where the code rate is $`R=N/M`$ (for unbiased messages). The encoded vector $`𝒕`$ is then corrupted by noise represented by the vector $`𝜻\{0,1\}^M`$ with components independently drawn from $`P(\zeta )=(1f)\delta (\zeta )+f\delta (\zeta 1)`$. The received vector takes the form $`𝒓=𝑮^𝑻𝝃+𝜻\text{ (mod 2)}`$.
Decoding is carried out by multiplying the received message by the matrix $`𝑨`$ to produce the syndrome vector $`𝒛=𝑨𝒓=𝑨𝜻\text{ (mod 2)}`$ from which an estimate $`\widehat{𝝉}`$ for the noise vector can be produced. An estimate for the original message is then obtained as the first $`N`$ bits of $`𝒓+\widehat{𝝉}\text{ (mod 2)}`$. The Bayes optimal estimator (also known as marginal posterior maximiser, MPM) for the noise is defined as $`\widehat{\tau }_j=\text{argmax}_{\tau _j}P(\tau _jz)`$. The performance of this estimator can be measured by the probability of bit error $`p_b=11/M_{j=1}^M\delta [\widehat{\tau }_j;\zeta _j]`$, where $`\delta [;]`$ is Kronecker’s delta. Knowing the matrices $`𝑪_\mathrm{𝟐}`$ and $`𝑪_\mathrm{𝟏}`$, the syndrome vector $`𝒛`$ and the noise level $`f`$ it is possible to apply Bayes’ theorem and compute the posterior probability
$$P(𝝉𝒛)=\frac{1}{Z}\chi \left[𝒛=𝑨𝝉\text{(mod 2)}\right]P(𝝉),$$
(1)
where $`\chi [X]`$ is an indicator function providing $`1`$ if $`X`$ is true and $`0`$ otherwise. To compute the MPM one has to compute the marginal posterior $`P(\tau _j𝒛)=_{ij}P(𝝉𝒛)`$, which in general requires $`𝒪(2^M)`$ operations, thus becoming impractical for long messages. To solve this problem one can use the sparseness of $`𝑨`$ to design algorithms that require $`𝒪(M)`$ operations to perform the same task. One of these methods is the probability propagation algorithm (PP), also known as belief propagation, sum-product algorithm (see ) or generalised distributive law .
The connection to statistical physics becomes clear when the field $`\{0,1\}`$ is replaced by Ising spins $`\{\pm 1\}`$ and mod $`2`$ sums by products . The syndrome vector acquires the form of a multi-spin coupling $`𝒥_\mu =_{j(\mu )}\zeta _j`$ where $`j=1,\mathrm{},M`$ and $`\mu =1,\mathrm{},(MN)`$. The $`K`$ indices of nonzero elements in the row $`\mu `$ of $`𝑨`$ are given by $`(\mu )=\{j_1,\mathrm{},j_K\}`$, and in a column $`l`$ are given by $`(l)=\{\mu _1,\mathrm{},\mu _C\}`$.
The posterior (1) can be written as the Gibbs distribution :
$`P(𝝉𝒥)`$ $`=`$ $`{\displaystyle \frac{1}{Z}}\underset{\beta \mathrm{}}{lim}\text{exp}\left[\beta _\beta (𝝉;𝒥)\right]`$ (2)
$`_\beta (𝝉;𝒥)`$ $`=`$ $`{\displaystyle \underset{\mu =1}{\overset{MN}{}}}\left(𝒥_\mu {\displaystyle \underset{j(\mu )}{}}\tau _j1\right){\displaystyle \frac{F}{\beta }}{\displaystyle \underset{j=1}{\overset{M}{}}}\tau _j.`$
The external field corresponds to the prior probability over the noise and has the form $`F=\text{atanh}(12f)`$. Note that the Hamiltonian itself depends on the inverse temperature $`\beta `$. The disorder is trivial and can be gauged as $`𝒥_\mu 1`$ by using $`\tau _j\tau _j\zeta _j`$. The resulting Hamiltonian is a multi-spin ferromagnet with finite connectivity in a random field $`h_j=\beta ^1F\zeta _j`$. The decoding process corresponds to finding zero temperature local magnetisations $`m_j=lim_\beta \mathrm{}\tau _j_\beta `$ and calculating estimates as $`\widehat{\tau }_j=\text{sgn}(m_j)`$.
In the $`\{\pm 1\}`$ representation the probability of bit error, acquires the form
$$p_b=\frac{1}{2}\frac{1}{2M}\underset{j=1}{\overset{M}{}}\zeta _j\text{ sgn}(m_j),$$
(3)
connecting the code performance with the computation of local magnetisations.
A Husimi cactus with connectivity $`C`$ is generated starting with a polygon of $`K`$ vertices with one Ising spin in each vertex (generation $`0`$). All spins in a polygon interact through a single coupling $`𝒥_\mu `$ and one of them is called the base spin. In figure 1 we show the first step in the construction of a Husimi cactus, in a generic step the base spins of the $`n1`$ generation polygons, numbering $`(C1)(K1)`$, are attached to $`K1`$ vertices of a generation $`n`$ polygon. This process is iterated until a maximum generation $`n_{\text{max}}`$ is reached, the graph is then completed by attaching $`C`$ uncorrelated branches of $`n_{\text{max}}`$ generations at their base spins. In that way each spin inside the graph is connected to exactly $`C`$ polygons. The local magnetisation at the centre $`m_j`$ can be obtained by fixing boundary (initial) conditions in the $`0`$-th generation and iterating recursion equations until generation $`n_{\text{max}}`$ is reached. Carrying out the calculation in the thermodynamic limit corresponds to having $`n_{\text{max}}\mathrm{ln}M`$ generations and $`M\mathrm{}`$.
The Hamiltonian of the model has the form (2) where $`(\mu )`$ denotes the polygon $`\mu `$ of the lattice. Due to the tree-like structure, local quantities far from the boundary can be calculated recursively by specifying boundary conditions. The typical decoding performance can therefore be computed exactly without resorting to replica calculations .
We adopt the approach presented in where recursion relations for the probability distribution $`P_{\mu k}(\tau _k)`$ for the base spin of the polygon $`\mu `$ is connected to $`(C1)(K1)`$ distributions $`P_{\nu j}(\tau _j)`$, with $`\nu (j)\mu `$ (all polygons linked to $`j`$ but $`\mu `$) of polygons in the previous generation:
$$P_{\mu k}(\tau _k)=\frac{1}{𝒩}\text{ Tr}_{\{\tau _j\}}\mathrm{exp}\left[\beta \left(𝒥_\mu \tau _k\underset{j(\mu )k}{}\tau _j1\right)+F\tau _k\right]\underset{\nu (j)\mu }{}\underset{j(\mu )k}{}P_{\nu j}(\tau _j),$$
(4)
where the trace is over the spins $`\tau _j`$ such that $`j(\mu )k`$.
The effective field $`\widehat{x}_{\nu j}`$ on a base spin $`j`$ due to neighbours in polygon $`\nu `$ can be written as :
$$\mathrm{exp}\left(2\widehat{x}_{\nu j}\right)=\text{e}^{2F}\frac{P_{\nu j}()}{P_{\nu j}(+)},$$
(5)
Combining (4) and (5) one finds the recursion relation:
$$\mathrm{exp}\left(2\widehat{x}_{\mu k}\right)=\frac{\text{ Tr}_{\{\tau _j\}}\mathrm{exp}\left[\beta 𝒥_\mu _{j(\mu )k}\tau _j+_{j(\mu )k}(F+_{\nu (j)\mu }\widehat{x}_{\nu j})\tau _j\right]}{\text{ Tr}_{\{\tau _j\}}\mathrm{exp}\left[+\beta 𝒥_\mu _{j(\mu )k}\tau _j+_{j(\mu )k}(F+_{\nu (j)\mu }\widehat{x}_{\nu j})\tau _j\right]}.$$
(6)
By computing the traces and taking $`\beta \mathrm{}`$ one obtains:
$$\widehat{x}_{\mu k}=\text{atanh}\left[𝒥_\mu \underset{j(\mu )k}{}\text{tanh}(F+\underset{\nu (j)\mu }{}\widehat{x}_{\nu j})\right]$$
(7)
The effective local magnetisation due to interactions with the nearest neighbours in one branch is given by $`\widehat{m}_{\mu j}=\text{tanh}(\widehat{x}_{\mu j})`$. The effective local field on a base spin $`j`$ of a polygon $`\mu `$ due to $`C1`$ branches in the previous generation and due to the external field is $`x_{\mu j}=F+_{\nu (j)\mu }\widehat{x}_{\nu j}`$; the effective local magnetisation is, therefore, $`m_{\mu j}=\text{tanh}(x_{\mu j})`$. Equation (7) can then be rewritten in terms of $`\widehat{m}_{\mu j}`$ and $`m_{\mu j}`$ and the PP equations can be recovered:
$$m_{\mu k}=\text{tanh}\left(F+\underset{\nu (j)\mu }{}\text{atanh }(\widehat{m}_{\nu k})\right)\widehat{m}_{\mu k}=𝒥_\mu \underset{j(\mu )k}{}m_{\mu j}$$
(8)
Once the magnetisations on the boundary ($`0`$-th generation) are assigned, the local magnetisation $`m_j`$ in the central site is determined by iterating (8) and computing :
$$m_j=\text{tanh}\left(F+\underset{\nu (j)}{}\text{atanh }(\widehat{m}_{\nu j})\right)$$
(9)
The free energy can be obtained by integration as (8) represents extrema of a free energy .
By applying the gauge transformation $`𝒥_\mu 1`$ and $`\tau _j\tau _j\zeta _j`$, assigning the probability distributions $`P_0(x)`$ to boundary fields and averaging over random local fields $`F\zeta `$ one obtains from (7) the recursion relation in the space of probability distributions $`P(x)`$ :
$`P_n(x)`$ $`=`$ $`{\displaystyle \underset{l=1}{\overset{C1}{}}d\widehat{x}_l\widehat{P}_{n1}(\widehat{x}_l)\delta \left[xF\zeta \underset{l=1}{\overset{C1}{}}\widehat{x}_l\right]_\zeta }`$
$`\widehat{P}_{n1}(\widehat{x})`$ $`=`$ $`{\displaystyle \underset{j=1}{\overset{K1}{}}dx_jP_{n1}(x_j)\delta \left[\widehat{x}\text{atanh}\left(\underset{j=1}{\overset{K1}{}}\text{tanh}(x_j)\right)\right]},`$ (10)
where $`P_n(x)`$ is the distribution of effective fields at the $`n`$-th generation due to the previous generations and external fields, in the thermodynamic limit the distribution far from the boundary will be $`P_{\mathrm{}}(x)`$ (generation $`n\mathrm{}`$). The local field distribution at the central site is computed by replacing $`C1`$ by $`C`$ in (10), taking into account $`C`$ polygons in the generation just before the central site, and inserting the distribution $`P_{\mathrm{}}(x)`$. Equations (10) are identical to those obtained by the replica symmetric theory as in .
By setting initial (boundary) conditions $`P_0(x)`$ and numerically iterating (10), for $`C3`$ one can find, up to some noise level $`f_s`$, a single stable fixed point at infinite fields, corresponding to a totally aligned state (successful decoding). At $`f_s`$ a bifurcation occurs and two other fixed points appear, stable and unstable, the former corresponding to a misaligned state (decoding failure). This situation is identical to that one observed in . In terms of the local fields distribution $`P_n(x)`$, the aligned state corresponds to a runaway wave travelling to $`x(n)\mathrm{}`$ with $`n`$ being the time variable. The misaligned state corresponds to a stable wave located at $`x(n)𝒪(1)`$. Representing the distributions (10) by the first cummulants only, one can obtain a rough approximation in terms of one dimensional maps showing a bifurcation at some noise level $`\stackrel{~}{f}_s`$, this approach will be further exploited elsewhere.
The practical PP decoding is performed by setting initial conditions as $`m_{\mu j}=12f`$ to correspond to the prior probabilities and iterating (8) until stationarity or a maximum number of iterations is attained . The estimate for the noise vector is then produced by computing $`\widehat{\tau }_j=\text{sign}(m_j)`$. At each decoding step the system can be described by histograms of the variables (8), this is equivalent to iterating (10) (a similar idea was presented in ). Below $`f_s`$ the process always converges to the successful decoding state, above $`f_s`$ it converges to the successful decoding only if the initial conditions are fine tuned, in general the process converges to the failure state. In Fig.2a we show the theoretical mean overlap between actual noise $`𝜻`$ and the estimate $`\widehat{𝝉}`$ as a function of the noise level $`f`$ as well as results obtained with PP decoding.
Information theory provides an upper bound for the maximum attainable code rate by equalising the maximal information contents of the syndrome vector $`𝒛`$ and of the noise estimate $`\widehat{𝝉}`$ . The thermodynamic phase transition obtained by finding the stable fixed points of (10) and their free energies interestingly coincides with this upper bound within the precision of the numerical calculation. Note that the performance predicted by thermodynamics is not practical as it requires $`𝒪(2^M)`$ operations for an exhaustive search for the global minimum of the free energy. In Fig.2b we show the thermodynamic transition for $`K=6`$ and compare with the upper bound, Shannon’s bound and $`f_s`$ values.
The geometrical structure of a Gallager code defined by the matrix $`𝑨`$ can be represented by a bipartite graph (Tanner graph) with bit and check nodes. Each column $`j`$ of $`𝑨`$ represents a bit node and each row $`\mu `$ represents a check node, $`A_{\mu j}=1`$ means that there is an edge linking bit $`j`$ to check $`\mu `$. It is possible to show that for a random ensemble of regular codes, the probability of completing a cycle after walking $`l`$ edges starting from an arbitrary node is upper bounded by $`𝒫[l;K,C,M]l^2K^l/M`$. It implies that for very large $`M`$ only cycles of at least order $`\mathrm{ln}M`$ survive. In the thermodynamic limit $`M\mathrm{}`$ the probability $`𝒫[l;K,C,M]0`$ for any finite $`l`$ and the bulk of the system is effectively tree-like. By mapping each check node to a polygon with $`K`$ bit nodes as vertices, one can map a Tanner graph into a Husimi lattice that is effectively a tree for any number of generations of order less than $`\mathrm{ln}M`$. It is experimentally observed that the number of iterations of (8) required for convergence does not scale with the system size, therefore, it is expected that the interior of a tree-like lattice approximates a Gallager code with increasing accuracy as the system size increases. Fig.2a shows that the approximation is fairly good even for sizes as small as $`M=100`$. Note that although the local magnetisations $`m_j`$ for a loopy graph are not generally expected to converge to the values computed in a tree, $`\text{sgn}(m_j)`$ seems to do so. A thorough discussion on this respect for some specific graphical models can be found in .
In MacKay and Neal introduced a variation on Gallager codes termed MN codes. The main difference between these codes is that for MN codes the syndrome vector contains also information on the original message in the form $`𝒛=𝑪_𝒔𝝃+𝑪_𝒏𝜻`$. The message itself is directly estimated and there is no need for recovering the noise vector. MacKay has formulated and proved a number of theorems simultaneously for both codes using the fact that if both message and noise are sampled from the same distribution, these codes can be formulated as the same estimation problem, to say, finding the most probable vector $`𝒙`$ that satisfies $`𝒛=𝑨𝒙`$, given the matrix $`𝑨`$ and a prior distribution $`P(𝒙)`$. Using statistical physics, we previously analysed MN codes . It is interesting to note that in spite of the similarity between the two codes, there are some important differences in their dependence on the parameters $`K`$ and $`C`$. In particular, Shannon’s bound is only attainable by Gallager codes if $`K\mathrm{}`$, in contrast to results obtained for MN codes. Decoding of unbiased messages is generally possible with Gallager codes, but successful convergence is only guaranteed (in the thermodynamic limit) for $`K=1,2`$ in the MN codes. We outlined those differences in table I.
To summarise, we solved exactly, without resorting to the replica method, a system representing a Gallager code on a Husimi cactus. The results obtained are in agreement with the replica symmetric calculation and with numerical experiments carried out in systems of moderate size. The framework can be easily extended to MN and similar codes. We believe that methods of statistical physics are complimentary to those used in the statistical inference community and can enhance our understanding of general graphical models beyond error-correcting codes.
\***
We acknowledge support from EPSRC (GR/N00562), The Royal Society (RV,DS) and from the JSPS RFTF program (YK).
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# Gauge Field Theory Coherent States (GCS) : I. General Properties
## 1 Introduction
Quantum General Relativity (QGR) has matured over the past decade to a mathematically well-defined theory of quantum gravity. In contrast to string theory, by definition QGR is a manifestly background independent, diffeomorphism invariant and non-perturbative theory. The obvious advantage is that one will never have to postulate the existence of a non-perturbative extension of the theory, which in string theory has been called the still unknown M(ystery)-Theory.
The disadvantage of a non-perturbative and background independent formulation is, of course, that one is faced with new and interesting mathematical problems so that one cannot just go ahead and “start calculating scattering amplitudes”: As there is no background around which one could perturb, rather the full metric is fluctuating, one is not doing quantum field theory on a spacetime but only on a differential manifold. Once there is no (Minkowski) metric at our disposal, one loses familiar notions such as causality, locality, Poincaré group and so forth, in other words, the theory is not a theory to which the Wightman axioms apply. Therefore, one must build an entirely new mathematical apparatus to treat the resulting quantum field theory which is drastically different from the Fock space picture to which particle physicists are used to.
As a consequence, the mathematical formulation of the theory was the main focus of research in the field over the past decade. The main achievements to date are the following (more or less in chronological order) :
* Kinematical Framework
The starting point was the introduction of new field variables for the gravitational field which are better suited to a background independent formulation of the quantum theory than the ones employed until that time. In its original version these variables were complex valued, however, currently their real valued version, considered first in for classical Euclidean gravity and later in for classical Lorentzian gravity, is preferred because to date it seems that it is only with these variables that one can rigorously define the kinematics and dynamics of Euclidean or Lorentzian quantum gravity .
These variables are coordinates for the infinite dimensional phase space of an $`SU(2)`$ gauge theory subject to further constraints besides the Gauss law, that is, a connection and a canonically conjugate electric field. As such, it is very natural to introduce smeared functions of these variables, specifically Wilson loop and electric flux functions. (Notice that one does not need a metric to define these functions, that is, they are background independent). This had been done for ordinary gauge fields already before in and was then reconsidered for gravity (see e.g. ).
The next step was the choice of a representation of the canonical commutation relations between the electric and magnetic degrees of freedom. This involves the choice of a suitable space of distributional connections and a faithful measure thereon which, as one can show , is $`\sigma `$-additive. The proof that the resulting Hilbert space indeed solves the adjointness relations induced by the reality structure of the classical theory as well as the canonical commutation relations induced by the symplectic structure of the classical theory can be found in . Independently, a second representation of the canonical commutation relations, called the loop representation, had been advocated (see e.g. and especially and references therein) but both representations were shown to be unitarily equivalent in (see also for a different method of proof).
This is then the first major achievement : The theory is based on a rigorously defined kinematical framework.
* Geometrical Operators
The second major achievement concerns the spectra of positive semi-definite, self-adjoint geometrical operators measuring lengths , areas and volumes of curves, surfaces and regions in spacetime. These spectra are pure point (discete) and imply a discrete Planck scale structure. It should be pointed out that the discreteness is, in contrast to other approaches to quantum gravity, not put in by hand but it is a prediction !
* Regularization- and Renormalization Techniques
The third major achievement is that there is a new regularization and renormalization technique for diffeomorphism covariant, density-one-valued operators at our disposal which was successfully tested in model theories . This technique can be applied, in particular, to the standard model coupled to gravity and to the Poincaré generators at spatial infinity . In particular, it works for Lorentzian gravity while all earlier proposals could at best work in the Euclidean context only (see, e.g. and references therein). The algebra of important operators of the resulting quantum field theories was shown to be consistent . Most surprisingly, these operators are UV and IR finite ! Notice that this result, at least as far as these operators are concerned, is stronger than the believed but unproved finiteness of scattering amplitudes order by order in perturbation theory of the five critical string theories, in a sense we claim that the perturbation series converges. The absence of the divergences that usually plague interacting quantum fields propagating on a Minkowski background can be understood intuitively from the diffeomorphism invariance of the theory : “short and long distances are gauge equivalent”. We will elaborate more on this point in future publications.
* Spin Foam Models
After the construction of the densely defined Hamiltonian constraint operator of , a formal, Euclidean functional integral was constructed in and gave rise to the so-called spin foam models (a spin foam is a history of a graph with faces as the history of edges) . Spin foam models are in close connection with causal spin-network evolutions , state sum models and topological quantum field theory, in particular BF theory . To date most results are at a formal level and for the Euclidean version of the theory only but the programme is exciting since it may restore manifest four-dimensional diffeomorphism invariance which in the Hamiltonian formulation is somewhat hidden.
* Finally, the fifth major achievement is the existence of a rigorous and satisfactory framework for the quantum statistical description of black holes which reproduces the Bekenstein-Hawking Entropy-Area relation and applies, in particular, to physical Schwarzschild black holes while stringy black holes so far are under control only for extremal charged black holes.
Summarizing, the work of the past decade has now culminated in a promising starting point for a quantum theory of the gravitational field plus matter and the stage is set to pose and answer physical questions.
The most basic and most important question that one should ask is : Does the theory have classical general relativity as its classical limit ? Notice that even if the answer is negative, the existence of a consistent, interacting, diffeomorphism invariant quantum field theory in four dimensions is already a quite non-trivial result. However, we can claim to have a satisfactory quantum theory of Einstein’s theory only if the answer is positive.
To settle this issue we have launched an attack based on coherent states which has culminated in a series of papers called “Gauge Field Theory Coherent States” and this paper is the first one in this collection which is going to be extended further.
The organization of this series is the following :
* General Properties
In this paper we describe a fairly general method to generate families of diffeomorphism covariant coherent states with the usual desired properties such as annihilation operator eigenstate nature, expectation value reproduction for annihilation and creation operators and saturation of the Heisenberg uncertainty bound. If certain analytical conditions are met, overcompleteness can be established as well. The construction is based on the so-called configuration space complexifier method described in detail in . The latter work arose as an abstraction of the results of Hall who chose a very special, but very convenient configuration space complexification for the case that the configuration space is a compact, connected Lie group. Hall’s results were later generalized to diffemorphism invariant gauge theories in . In this paper we focus on general properties of such states for a general complexification such as gauge invariance and diffeomorphism covariance. Besides such physical features also analytical properties are addressed and it is a mixture of the two that will determine one’s choice of the complexification. In fact, in the remainder of this series we will mostly deal with a generalization of the complexification chosen by Hall. Our main reason for this choice is simply mathematical convenience : The spectrum of the operator that generates the configuration space complexification is explicitly known and sufficiently simple. This allows us to get started, but it should be kept in mind that other choices are available that may prove physically more interesting later on in our programme.
* Peakedness Properties
Associated with the configuration space complexification is a so-called coherent state transform and both of focussed on the unitarity of that transform while the properties of the coherent states themselves remained untouched. Moreover, it remained unclear how the complexified connection $`A^{\text{ }\mathrm{C}}`$ looks like in terms of the coordinates $`(A,E)`$ of the real phase space and without this an interpretation of the label $`A^{\text{ }\mathrm{C}}`$ of the coherent state and thus expectation values, fluctuations and so forth remain veiled. Here, $`A`$ is a connection for a compact gauge group and $`E`$ is a canonically conjugate electric field. To fill both of these gaps is the purpose of the second paper in this series. First of all, we find the expected result, namely that roughly speaking $`A^{\text{ }\mathrm{C}}=AiE`$ in a suitably smeared sense. Secondly, we analyze in detail the peakedness properties of the coherent states for diffeomorphism invariant gauge theories in the configuration –, momentum – and the Segal-Bargmann representation. We find that these states are very sharply peaked at the point $`A`$, $`E`$ or $`(A,E)`$ respectively of the configuration –, momentum – and phase space respectively. That paper also contains extensive graphics to demonstrate these peakedness properties pictorially and while there are important differences, the states display the essential Gaussian decay of the harmonic oscillator coherent states.
* Ehrenfest Theorems
In the third paper of this series we prove Ehrenfest theorems for our coherent states. That is, we show that the expectation value not only of normal ordered polynomials of creation an annihilation operators but of all polynomials of the elementary operators associated with $`\widehat{A},\widehat{E}`$ equals, to leading order in $`\mathrm{}`$, precisely the labels $`A,E`$ of the coherent state. This result can be extended to certain operators that are non-polynomial in the basic ones and that appear in the Hamiltonian constraint of quantum general relativity coupled to matter . Moreover, we show that commutators between these operators divided by $`i\mathrm{}`$ have an expectation value which equals to leading order in $`\mathrm{}`$ the correspending Poisson bracket evaluated at the label $`(A,E)`$ of the coherent state. Together, these results imply that the classical limit of the Hamiltonian constraint operator and its infinitesimal quantum dynamics correspond to its classical counterparts.
Both of mainly deal with $`G=U(1),SU(2)`$ but we sketch how all the results can be extended to groups of higher rank, an issue which we will examine in detail in .
* Infinite Tensor Product and Thermodynamical Limit
The states that one considered in Quantum General Relativity until now are labelled by piecewise analytic, finite graphs (an extension to finite collections of smooth curves with controlled intersection properties is possible, see later on). However, finite graphs are suitable to describe semiclassical physics on physically interesting spacetimes only if the underlying manifold is spatially compact. The most interesting applications, flat space or an entire black hole spacetime (and not only the horizon region) cannot be treated with finite graphs. To extend the framework it turns out that piecewise analytical, countably infinite graphs together with the framework of the Infinite Tensor Product (ITP) construction introduced by von Neumann more than sixty years ago are appropriate. To the best of the knowledge of the author, the first time that truly infinite graphs and infinite tensor product states were considered in QGR in the context of a Hilbert space structure, was in section 3.2 of which dealt with the asymptotic Poincaré group of asymptotically flat spacetimes, however, the overall mathematical framework of such constructions was not described there. In we deliver this structure and embed it into our coherent states framework. In particular, we are able to connect mathematical notions with physical ones, an example being the following :
A state $`f`$ in the infinite tensor product Hilbert space over an infinite graph which is a direct product of normalized states, one for each edge of the graph, generates so-called strong and weak equivalence classes of so-called $`C_0`$-sequences. It turns out that the corresponding $`C_0`$-vector plays the role of a cyclic vector (vacuum state) for a Fock-like tiny closed subspace of the complete ITP Hilbert space, called an $`f`$-adic incomplete ITP. Fock-like spaces corresponding to different strong and weak equivalence classes are mutually orthogonal. Those Fock-like spaces that correspond to the same weak class but different strong classes are unitarily equivalent while those that correspond to different strong and weak classes are unitarily inequivalent. This way the ITP gives rise to an uncountably infinite number of mutually unitarily inequivalent representations of the canonical commutation relations. The representation theory of operator algebras becomes especially interesting, the enveloping framework being that of factors of von Neumann algebras.
Generically, incomplete ITP’s generated by different weak equivalence classes correspond to physical situations which differ drastically with respect to certain physical quantities such as energy, volume or topology. For instance, the Ashtekar-Isham-Lewandowski Hilbert space based on finite graphs describes finite volume and/or compact topology while a $`C_0`$ vector of infinite volume can be constructed by using our coherent states, appropriate to approximate a flat Minkowski space geometry. The two Hilbert spaces are mutually orthogonal closed subspaces within our complete ITP Hilbert space corresponding to different weak classes. The vacuum underlying the Ashtekar-Isham-Lewandowski Hilbert space via the GNS construction is based on a $`C_0`$ vector which equals unity for every edge of any possible graph. It can be shown that such a state, in the context of non-compact topologies, is a pure quantum vacuum in the sense that it describes metrics of almost everywhere zero spatial volume.
It should be clear from these considerations that the ITP is possibly able to describe all phyically different situations at once and might enable us to describe topology change within canonical quantum general relativity and therefore to get rid off the embedding spacetime manifold that one started with classically !
The infinite tensor product opens the gate to a plethora of other physical and mathematical disciplines, such as thermodynamics and statistical field theory, Tomita-Takesaki (or modular) theory necessary to classify the appearing types of type III factors of von Neumann algebras etc.
* Higgs Fields and Fermions
The framework described so far is sufficient for pure quantum gauge theories coupled to quantum general relativity only. By combining the framework of with the infinite tensor product construction and existing results for coherent states for fermions (e.g. and references therein) we can extend the framework to all matter of the standard model including possible supersymmetric extensions. The details are described in .
* Photons and Gravitons
Most of the criticism directed towards quantum general relativity coming from the particle physics community is that the programme, being manifestly non-perturbative by construction, seems to be infinitely far away from any established perturbative results such as (free) quantum field theory on curved backgrounds (widely believed to be the first approximation to full quantum gravity), perturbative quantum (super)gravity (non-renormalizable) and perturbative quantum superstring theory. In we make a first contact with these programmes. Namely, we try to construct a map between the perturbative Photon or Graviton Hilbert spaces and a fully non-perturbative incomplete $`f`$-adic ITP subspace where the $`C_0`$-vector corresponding to $`f`$ is a best approximation state to the Minkowski space solution of the Einstein-Maxwell equations. This work is aimed at demonstrating how perturbative notions such as particles can be absorbed into our fully non-perturbative programme.
* The Non-Perturbative $`\gamma `$-Ray Burst Effect
Many serious theorists and experimentalists nowadays discuss the possibility to actually measure quantum gravity effects, a prominent example being the so-called $`\gamma `$-ray burst experiment (see, e.g. ). In all these types of experiments one exploits the fact that the incredibly tiny quantum gravity effects may accumulate over vast periods of time of the order of the age of the universe to a measurable size. In particular, the theoretical mechanism of the $`\gamma `$-ray burst effect can be roughly described as follows : the quantum metric depends on canonically conjugate magnetic and electric degrees of freedom and thus the Heisenberg uncertainty obstruction tells us that there is no state that can describe the Minkowski vacuum exactly. In other words, there is no Poincaré invariant state in the theory, the best one can do is to construct a coherent state peaked on Minkowski space. The expectation value of the Einstein-Maxwell-Hamiltonian with respect to the gravitational field will therefore include corrections to the classical Minkowski metric which give rise to Poincaré invariance violating dispersion relations. Thus, if one could measure the arrival times of $`\gamma `$-ray photons of different energies they should differ by an amount proportional to the travelling time from the source.
The challenge is now to precisely compute these corrections from our fully non-perturbative framework, in particular, what is the precise power of the Planck mass that the effect is proportional to. This is the subject of which will improve the pioneering work in two respects : First, the latter was based on so-called weave states which, however, approximate only half of the number of degrees of freedom and, secondly, in contrast to our coherent states the existence of weave itself with the assumed semi-classical properties was not proved to exist.
To compute the effect exactly turns out to be a hard piece of analysis due to the non-linear, even non-analytic (interacting) nature of the theory, a property which carries over to our coherent states. In particular, the complicated spectrum of the volume operator makes the enterprise not an easy one. On the other hand, it is absolutely crucial to know the precise spectrum and not only of, say, its main series : If one would do the same with the area operator then, as has been beautifully demonstrated in , one would reach the conclusion that the black hole Hawking radiation spectrum is discrete rather than the quasi-continuous one of a black body, in other words, the spectrum has direct bearing on observation !
It is at this point that super-computers may enter the stage as analytic computations start becoming too hard and lengthy. Notice, however, that in contrast to usual perturbation series in perturbative quantum field theory the computational error is always under good control. The series that we are dealing with are manifestly absolutely converging and there are precise estimates on the error that one creates when keeping only the dominant terms. We will display such error controlled estimates in the next two issues of this series.
* The Classical Limit
As an immediate application of coherent states and the ITP framework one can now precisely prove in detail how it happens that the Hamiltonian constraint constructed in obeys the correct quantum algebra.
More work is in progress. The following list of projects associated with our coherent states represents just the tip of the iceberg, in principle it would would be interesting to repeat all perturbative calculations that have been performed so far with our non-perturbative tools and to provide the error bars.
A)
To relate standard perturbative quantum field theory on curved backgrounds with non-perturbative quantum general relativity one would like to understand why the UV singularities of the former have disappeared in the latter. The naive answer is that the renormalization group has been absorbed into the diffeomorphism group (large and small momenta are gauge related) but one would like to understand this and related notions like bare and renormalized charges, effective actions, renormalization transformations, Epstein-Glaser formalism and the importance of Hadamard states for quantum field theory on curved backgrounds etc. in more detail from the non-perturbative point of view. In particular, it would be nice to map the usual Feynman rules into our framework. This research project will be started in .
B)
An ever fascinating research object has been the black hole. The coherent states provide a natural new setting in which to study quantum black holes and Hawking radiation, in principle one “just” has to take the coherent state that approximates a Kruskal spacetime together with its excitations in order to provide the Kruskal – spacetime – adic incomplete closed ITP Hilbert space structure (that is, a vacuum and excitations). Notice that while the Bekenstein – Hawking entropy has been successfully computed in both canonical quantum gravity and string theory as mentioned above, what would be new here is that one can treat the full spacetime in a Hilbert space context and not only its near horizon structure (charges). Also, there are a priori no constraints such as (near-) stationarity or extremality of the black hole. Finally, one would like to understand what happens to the classical singularity theorems, the information paradoxon, cosmic censorship etc. in the quantum theory. These and rlated issues will be the topic of .
C)
As already mentioned, von Neumann algebras and their representation theory appear quite naturally in the Infinite Tensor Product construction. For the latter, the decomposition of a von Neumann algebra into factors is of particular importance and the basic tool to characterize factors of type III, which typically appear in quantum field theory, is provided by modular theory. This brings us into close contact with algebraic quantum field theory, although presumably in a generalized setting, since the notion of locality plays, almost by definition, a less dominant role in a diffeomorphism invariant quantum field theory. These and related issues will be examined in .
D)
The most effective way to derive a path integral formulation for kinematically linear field theories from the Hamiltonian formulation of the theory is via coherent states, see e.g. and references therein. Thus, it is natural to expect this to be the case also for our coherent states. This may bring us into contact with the formely mentioned spin foam models which have recently attracted quite some attention after the appearence of and will be studied in .
E)
Finally, our coherent states are pure states. The semi-classical behaviour of such states may yet be improved by superimposing them to a so-called mixed state which makes use of random lattices. For weaves, such a framework already exists and has been studied in . We intend to combine both frameworks in .
This article is assembled as follows :
In section two we recall the classical and quantum kinematics of diffeomorphism invariant gauge field theories.
In section three we recall the complexifier method to generate Bargmann-Segal representations for general theories and gauge theories in particular. We comment on the physical and mathematical requirements to be imposed on the complexifier, that is, the canonical generator of the transform that complexifies the configuration space and identifies it with the phase space. In three related subsections we propose three candidate families of coherent states for gauge theories. The first one leads to an actual complex connection, the second only to a complexified holonomy without underlying complex connection and the third one maps the problem at hand in principle to coherent states for an (in)finite collection of uncoupled harmonic oscillators. We describe the advantages and disadvantages of these states as compared to each other. All of this will be done mostly for gauge – and diffeomorphism variant coherent states.
In sections four and five respectively we will deal with the issue of how to construct gauge – and diffeomorphism invariant coherent states respectively. Some of these can even be chosen to be annihilated by the Hamiltonian constraint.
Finally, in section six we display a simple example for gauge invariant coherent states with an actual complex connection in 2+1 gravity and study some of their peakedness properties.
## 2 Kinematical Structure of Diffeomorphism Invariant Quantum Gauge Theories
In this section we will recall the main ingredients of the mathematical formulation of (Lorentzian) diffeomorphism invariant classical and quantum field theories of connections with local degrees of freedom in any dimension and for any compact gauge group. See and references therein for more details.
### 2.1 Classical Theory
Let $`G`$ be a compact gauge group, $`\mathrm{\Sigma }`$ a $`D`$dimensional manifold admitting a principal $`G`$bundle with connection over $`\mathrm{\Sigma }`$. Let us denote the pull-back to $`\mathrm{\Sigma }`$ of the connection by local sections by $`A_a^i`$ where $`a,b,c,..=1,..,D`$ denote tensorial indices and $`i,j,k,..=1,..,dim(G)`$ denote indices for the Lie algebra of $`G`$. Likewise, consider a density-one vector bundle of electric fields, whose pull-back to $`\mathrm{\Sigma }`$ by local sections (their Hodge dual is a $`D1`$ form) is a Lie algebra valued vector density of weight one. We will denote the set of generators of the rank $`N1`$ Lie algebra of $`G`$ by $`\tau _i`$ which are normalized according to $`\text{tr}(\tau _i\tau _j)=N\delta _{ij}`$ and $`[\tau _i,\tau _j]=2f_{ij}^k\tau _k`$ defines the structure constants of $`Lie(G)`$.
Let $`F_i^a`$ be a Lie algebra valued vector density test field of weight one and let $`f_a^i`$ be a Lie algebra valued covector test field. We consider the smeared quantities
$$F(A):=_\mathrm{\Sigma }d^DxF_i^aA_a^i\text{ and }E(f):=_\mathrm{\Sigma }d^DxE_i^af_a^i$$
(2.1)
While both of them are diffeomorphism covariant it is only the latter which is gauge covariant and this is one motivation to consider the singular smearings discussed below. The choice of the space of pairs of test fields $`(F,f)𝒮`$ depends on the boundary conditions on the space of connections and electric fields which in turn depends on the topology of $`\mathrm{\Sigma }`$ and will not be specified in what follows.
Let the set of all pairs of smooth functions $`(A,E)`$ on $`\mathrm{\Sigma }`$ such that (2.1) is well defined for any $`(F,f)𝒮`$ be denoted by $`M`$. We define a topology on $`M`$ through the following globally defined metric :
$`d_{\rho ,\sigma }[(A,E),(A^{},E^{})]`$
$`:=`$ $`\sqrt{{\displaystyle \frac{1}{N}}{\displaystyle _\mathrm{\Sigma }}d^Dx[\sqrt{det(\rho )}\rho ^{ab}\text{tr}([A_aA_a^{}][A_bA_b^{}])+{\displaystyle \frac{[\sigma _{ab}\text{tr}([E^aE^a][E^bE^b])}{\sqrt{det(\sigma )}}}]}`$
where $`\rho _{ab},\sigma _{ab}`$ are fiducial metrics on $`\mathrm{\Sigma }`$ of everywhere Euclidean signature. Their fall-off behaviour has to be suited to the boundary conditions of the fields $`A,E`$ at spatial infinity. Notice that the metric (2.1) on $`M`$ is gauge invariant. It can be used in the usual way to equip $`M`$ with the structure of a smooth, infinite dimensional differential manifold modelled on a Banach (in fact Hilbert) space $``$ where $`𝒮\times 𝒮`$. (It is the weighted Sobolev space $`H_{0,\rho }^2\times H_{0,\sigma ^1}^2`$ in the notation of ).
Finally, we equip $`M`$ with the structure of an infinite dimensional symplectic manifold through the following strong (in the sense of ) symplectic structure
$$\mathrm{\Omega }((f,F),(f^{},F^{}))_m:=_\mathrm{\Sigma }d^Dx[F_i^af_a^iF_i^af_a^i](x)$$
(2.3)
for any $`(f,F),(f^{},F^{})`$. We have abused the notation by identifying the tangent space to $`M`$ at $`m`$ with $``$. To prove that $`\mathrm{\Omega }`$ is a strong symplectic structure one uses standard Banach space techniques. Computing the Hamiltonian vector fields (with respect to $`\mathrm{\Omega }`$) of the functions $`E(f),F(A)`$ we obtain the following elementary Poisson brackets
$$\{E(f),E(f^{})\}=\{F(A),F^{}(A)\}=0,\{E(f),A(F)\}=F(f)$$
(2.4)
As a first step towards quantization of the symplectic manifold $`(M,\mathrm{\Omega })`$ one must choose a polarization. As usual in gauge theories, we will use a particular real polarization, specifically connections as the configuration variables and electric fields as canonically conjugate momenta. As a second step one must decide on a complete set of coordinates of $`M`$ which are to become the elementary quantum operators. The analysis just outlined suggests to use the coordinates $`E(f),F(A)`$. However, the well-known immediate problem is that these coordinates are not gauge covariant. Thus, we proceed as follows :
Let $`\mathrm{\Gamma }_0^\omega `$ be the set of all piecewise analytic, finite, oriented graphs $`\gamma `$ embedded into $`\mathrm{\Sigma }`$ and denote by $`E(\gamma )`$ and $`V(\gamma )`$ respectively its sets of oriented edges $`e`$ and vertices $`v`$ respectively. Here finite means that $`E(\gamma )`$ is a finite set. (One can extend the framework to $`\mathrm{\Gamma }_0^{\mathrm{}}`$, the restriction to webs of the set of piecewise smooth graphs but the description becomes more complicated and we refrain from doing this here). It is possible to consider the set $`\mathrm{\Gamma }_\sigma ^\omega `$ of piecewise analytic, infinite graphs with an additional regularity property but for the purpose of this paper it will be sufficient to stick to $`\mathrm{\Gamma }_0^\omega `$. The subscript <sub>0</sub> as usual denotes “of compact support” while <sub>σ</sub> denotes “$`\sigma `$-finite”.
We denote by $`h_e(A)`$ the holonomy of $`A`$ along $`e`$ and say that a function $`f`$ on $`𝒜`$ is cylindrical with respect to $`\gamma `$ if there exists a function $`f_\gamma `$ on $`G^{|E(\gamma )|}`$ such that $`f=p_\gamma ^{}f_\gamma =f_\gamma p_\gamma `$ where $`p_\gamma (A)=\{h_e(A)\}_{eE(\gamma )}`$. Holonomies are invariant under reparameterizations of the edge and in this article we assume that the edges are always analyticity preserving diffeomorphic images from $`[0,1]`$ to a one-dimensional submanifold of $`\mathrm{\Sigma }`$. Gauge transformations are functions $`g:\mathrm{\Sigma }G;xg(x)`$ and they act on holonomies as $`h_eg(e(0))h_eg(e(1))^1`$.
Next, given a graph $`\gamma `$ we choose a polyhedronal decomposition $`P_\gamma `$ of $`\mathrm{\Sigma }`$ dual to $`\gamma `$. The precise definition of a dual polyhedronal decomposition can be found in but for the purposes of the present paper it is sufficient to know that $`P_\gamma `$ assigns to each edge $`e`$ of $`\gamma `$ an open “face” $`S_e`$ (a polyhedron of codimension one embedded into $`\mathrm{\Sigma }`$) with the following properties :
(1) the surfaces $`S_e`$ are mutually non-intersecting,
(2) only the edge $`e`$ intersects $`S_e`$, the intersection is transversal and consists only of one point which is an interior point of both $`e`$ and $`S_e`$,
(3) $`S_e`$ carries the orientation which agrees with the orientation of $`e`$.
Furthermore, we choose a system $`\mathrm{\Pi }_\gamma `$ of paths $`\rho _e(x)S_e,xS_e,eE(\gamma )`$ connecting the intersection point $`p_e=eS_e`$ with $`x`$. The paths vary smoothly with $`x`$ and the triples $`\gamma ,P_\gamma ,\mathrm{\Pi }_\gamma `$ are such that if $`\gamma ,\gamma ^{}`$ are diffeomorphic, so are $`P_\gamma ,P_\gamma ^{}`$ and $`\mathrm{\Pi }_\gamma ,\mathrm{\Pi }_\gamma ^{}`$, see for details.
With these structures we define the following function on $`(M,\mathrm{\Omega })`$
$$P_i^e(A,E):=\frac{1}{N}\text{tr}(\tau _ih_e(0,1/2)[_{S_e}h_{\rho _e(x)}E(x)h_{\rho _e(x)}^1]h_e(0,1/2)^1)$$
(2.5)
where $`h_e(s,t)`$ denotes the holonomy of $`A`$ along $`e`$ between the parameter values $`s<t`$, $``$ denotes the Hodge dual, that is, $`E`$ is a $`(D1)`$form on $`\mathrm{\Sigma }`$, $`E^a:=E_i^a\tau _i`$ and we have chosen a parameterization of $`e`$ such that $`p_e=e(1/2)`$.
Notice that in contrast to similar variables used earlier in the literature the function $`P_i^e`$ is gauge covariant. Namely, under gauge transformations it transforms as $`P^eg(e(0))P^eg(e(0))^1`$, the price to pay being that $`P^e`$ depends on both $`A`$ and $`E`$ and not only on $`E`$. The idea is therefore to use the variables $`h_e,P_i^e`$ for all possible graphs $`\gamma `$ as the coordinates of $`M`$.
The problem with the functions $`h_e(A)`$ and $`P_i^e(A,E)`$ on $`M`$ is that they are not differentiable on $`M`$, that is, $`Dh_e,DP_i^e`$ are nowhere bounded operators on $``$ as one can easily see. The reason for this is, of course, that these are functions on $`M`$ which are not properly smeared with functions from $`𝒮`$, rather they are smeared with distributional test functions with support on $`e`$ or $`S_e`$ respectively. Nevertheless one would like to base the quantization of the theory on these functions as basic variables because of their gauge and diffeomorphism covariance. Indeed, under diffeomorphisms $`h_eh_{\phi ^1(e)},P_i^eP_i^{\phi ^1(e)}`$ where the latter notation means that $`P_e^{\phi ^1(e)}`$ is labelled by $`\phi ^1(S_e),\phi ^1(\mathrm{\Pi }_\gamma )`$. We proceed as follows.
###### Definition 2.1
By $`\overline{M}_\gamma `$ we denote the direct product $`[G\times Lie(G)]^{|E(\gamma )|}`$. The subset of $`\overline{M}_\gamma `$ of pairs $`(h_e(A),P_i^e(A,E))_{eE(\gamma )}`$ as $`(A,E)`$ varies over $`M`$ will be denoted by $`(\overline{M}_\gamma )_{|M}`$. We have a corresponding map $`p_\gamma :M\overline{M}_\gamma `$ which maps $`M`$ onto $`(\overline{M}_\gamma )_{|M}`$.
Notice that the set $`(\overline{M}_\gamma )_{|M}`$ is in general a proper subset of $`\overline{M}_\gamma `$, depending on the boundary conditions on $`(A,E)`$, the topology of $`\mathrm{\Sigma }`$ and the “size” of $`e,S_e`$. For instance, in the limit of $`e,S_eeS_e`$ but holding the number of edges fixed, $`(\overline{M}_\gamma )_{|M}`$ will consist of only one point in $`M_\gamma `$. This follows from the smoothness of the $`(A,E)`$.
We equip a subset $`M_\gamma `$ of $`\overline{M}_\gamma `$ with the structure of a differentiable manifold modelled on the Banach space $`_\gamma =\text{ }\mathrm{R}^{2dim(G)|E(\gamma )|}`$ by using the natural direct product manifold structure of $`[G\times Lie(G)]^{|E(\gamma )|}`$. While $`\overline{M}_\gamma `$ is a kind of distributional phase space, $`M_\gamma `$ satisfies appropriate regularity properties induced by (2.1).
In order to proceed and to give $`M_\gamma `$ a symplectic structure derived from $`(M,\mathrm{\Omega })`$ one must regularize the elementary functions $`h_e,P_i^e`$ by writing them as limits (in which the regulator vanishes) of functions which can be expressed in terms of the $`F(A),E(f)`$. Then one can compute their Poisson brackets with respect to the symplectic structure $`\mathrm{\Omega }`$ at finite regulator and then take the limit pointwise on $`M`$. The result is the following well-defined strong symplectic structure $`\mathrm{\Omega }_\gamma `$ on $`M_\gamma `$.
$`\{h_e,h_e^{}\}_\gamma `$ $`=`$ $`0`$
$`\{P_i^e,h_e^{}\}_\gamma `$ $`=`$ $`\delta _e^{}^e{\displaystyle \frac{\tau _i}{2}}h_e`$
$`\{P_i^e,P_j^e^{}\}_\gamma `$ $`=`$ $`\delta ^{ee^{}}f_{ij}^kP_k^e`$ (2.6)
Since $`\mathrm{\Omega }_\gamma `$ is obviously block diagonal, each block standing for one copy of $`G\times Lie(G)`$, to check that $`\mathrm{\Omega }_\gamma `$ is non-degenerate and closed reduces to doing it for each factor together with an appeal to well-known Hilbert space techniques to establish that $`\mathrm{\Omega }_\gamma `$ is a surjection of $`_\gamma `$. This is done in where it is shown that each copy is isomorphic with the cotangent bundle $`T^{}G`$ equipped with the symplectic structure (2.1) (choose $`e=e^{}`$ and delete the label $`e`$).
Now that we have managed to assign to each graph $`\gamma `$ a symplectic manifold $`(M_\gamma ,\mathrm{\Omega }_\gamma )`$ we can quantize it by using geometric quantization. This can be done in a well-defined way because the relations (2.1) show that the corresponding operators are non-distributional. This is therefore a clean starting point for the regularization of any operator of quantum gauge field theory which can always be written in terms of the $`\widehat{h}_e,\widehat{P}^e,eE(\gamma )`$ if we apply this operator to a function which depends only on the $`h_e,eE(\gamma )`$.
The question is what $`(M_\gamma ,\mathrm{\Omega }_\gamma )`$ has to do with $`(M,\mathrm{\Omega })`$. In it is shown that there exists a partial order $``$ on the set $``$ of triples $`l=(\gamma ,P_\gamma ,\mathrm{\Pi }_\gamma )`$. In particular, $`\gamma \gamma ^{}`$ means $`\gamma \gamma ^{}`$ and $``$ is a directed set so that one can form a generalized projective limit $`M_{\mathrm{}}`$ of the $`M_\gamma `$ (we abuse notation in displaying the dependence of $`M_\gamma `$ on $`\gamma `$ only rather than on $`l`$). For this one verifies that the family of symplectic structures $`\mathrm{\Omega }_\gamma `$ is self-consistent in the sense that if $`(\gamma ,P_\gamma ,\mathrm{\Pi }_\gamma )(\gamma ^{},P_\gamma ^{},\mathrm{\Pi }_\gamma ^{})`$ then $`p_{\gamma ^{}\gamma }^{}\{f,g\}_\gamma =\{p_{\gamma ^{}\gamma }^{}f,p_{\gamma ^{}\gamma }^{}g\}_\gamma ^{}`$ for any $`f,gC^{\mathrm{}}(M_\gamma )`$ and $`p_{\gamma ^{}\gamma }:M_\gamma ^{}M_\gamma `$ is a system of natural projections, more precisely, of (non-invertible) symplectomorphisms.
Now, via the maps $`p_\gamma `$ of definition 2.1 we can identify $`M`$ with a subset of $`M_{\mathrm{}}`$. Moreover, in it is shown that there is a generalized projective sequence $`(\gamma _n,P_{\gamma _n},\mathrm{\Pi }_{\gamma _n})`$ such that $`lim_n\mathrm{}p_{\gamma _n}^{}\mathrm{\Omega }_{\gamma _n}=\mathrm{\Omega }`$ pointwise in $`M`$. This displays $`(M,\mathrm{\Omega })`$ as embedded into a generalized projective limit of the $`(M_\gamma ,\mathrm{\Omega }_\gamma )`$, intuitively speaking, as $`\gamma `$ fills all of $`\mathrm{\Sigma }`$, we recover $`(M,\mathrm{\Omega })`$ from the $`(M_\gamma ,\mathrm{\Omega }_\gamma )`$. Of course, this works with $`\mathrm{\Gamma }_0^\omega `$ only if $`\mathrm{\Sigma }`$ is compact, otherwise we need the extension to $`\mathrm{\Gamma }_\sigma ^\omega `$.
It follows that quantization of $`(M,\mathrm{\Omega })`$, and conversely taking the classical limit, can be studied purely in terms of $`(M_\gamma ,\mathrm{\Omega }_\gamma )`$ for all $`\gamma `$. The quantum kinematical framework is given in the next subsection.
### 2.2 Quantum Theory
Let us denote the set of all smooth connections by $`𝒜`$. This is our classical configuration space and we will choose for its coordinates the holonomies $`h_e(A),e\gamma ,\gamma \mathrm{\Gamma }_0^\omega `$. $`𝒜`$ is naturally equipped with a metric topology induced by (2.1).
Recall the notion of a function cylindrical over a graph from the previous subsection. A particularly useful set of cylindrical functions are the so-called spin-netwok functions . A spin-network function is labelled by a graph $`\gamma `$, a set of non-trivial irreducible representations $`\stackrel{}{\pi }=\{\pi _e\}_{eE(\gamma )}`$ (choose from each equivalence class of equivalent representations once and for all a fixed representant), one for each edge of $`\gamma `$, and a set $`\stackrel{}{c}=\{c_v\}_{vV(\gamma )}`$ of contraction matrices, one for each vertex of $`\gamma `$, which contract the indices of the tensor product $`_{eE(\gamma )}\pi _e(h_e)`$ in such a way that the resulting function is gauge invariant. We denote spin-network functions as $`T_I`$ where $`I=\{\gamma ,\stackrel{}{\pi },\stackrel{}{c}\}`$ is a compound label. One can show that these functions are linearly independent. From now on we denote by $`\stackrel{~}{\mathrm{\Phi }}_\gamma `$ finite linear combinations of spin-network functions over $`\gamma `$, by $`\mathrm{\Phi }_\gamma `$ the finite linear combinations of elements from any possible $`\stackrel{~}{\mathrm{\Phi }}_\gamma ^{},\gamma ^{}\gamma `$ a subgraph of $`\gamma `$ and by $`\mathrm{\Phi }`$ the finite linear combinations of spin-network functions over an arbitrary finite collection of graphs. Clearly $`\stackrel{~}{\mathrm{\Phi }}_\gamma `$ is a subspace of $`\mathrm{\Phi }_\gamma `$. To express this distinction we will say that functions in $`\stackrel{~}{\mathrm{\Phi }}_\gamma `$ are labelled by the “coloured graphs” $`\gamma `$ while functions in $`\mathrm{\Phi }_\gamma `$ are labelled simply by graphs $`\gamma `$ where we abuse notation by using the same symbol $`\gamma `$.
The set $`\mathrm{\Phi }`$ of finite linear combinations of spin-network functions forms an Abelian algebra of functions on $`𝒜`$. By completing it with respect to the sup-norm topology it becomes an Abelian C algebra $``$ (here the compactness of $`G`$ is crucial). The spectrum $`\overline{𝒜}`$ of this algebra, that is, the set of all algebraic homomorphisms $`\text{ }\mathrm{C}`$ is called the quantum configuration space. This space is equipped with the Gel’fand topology, that is, the space of continuous functions $`C^0(\overline{𝒜})`$ on $`\overline{𝒜}`$ is given by the Gel’fand transforms of elements of $``$. Recall that the Gel’fand transform is given by $`\stackrel{~}{f}(\overline{A}):=\overline{A}(f)\overline{A}\overline{𝒜}`$. It is a general result that $`\overline{𝒜}`$ with this topology is a compact Hausdorff space. Obviously, the elements of $`𝒜`$ are contained in $`\overline{𝒜}`$ and one can show that $`𝒜`$ is even dense . Generic elements of $`\overline{𝒜}`$ are, however, distributional.
The idea is now to construct a Hilbert space consisting of square integrable functions on $`\overline{𝒜}`$ with respect to some measure $`\mu `$. Recall that one can define a measure on a locally compact Hausdorff space by prescribing a positive linear functional $`\chi _\mu `$ on the space of continuous functions thereon. The particular measure we choose is given by $`\chi _{\mu _0}(\stackrel{~}{T}_I)=1`$ if $`I=\{\{p\},\stackrel{}{0},\stackrel{}{1}\}`$ and $`\chi _{\mu _0}(\stackrel{~}{T}_I)=0`$ otherwise. Here $`p`$ is any point in $`\mathrm{\Sigma }`$, $`0`$ denotes the trivial representation and $`1`$ the trivial contraction matrix. In other words, (Gel’fand transforms of) spin-network functions play the same role for $`\mu _0`$ as Wick-polynomials do for Gaussian measures and like those they form an orthonormal basis in the Hilbert space $`:=L_2(\overline{𝒜},d\mu _0)`$ obtained by completing their finite linear span $`\mathrm{\Phi }`$.
An equivalent definition of $`\overline{𝒜},\mu _0`$ is as follows :
$`\overline{𝒜}`$ is in one to one correspondence, via the surjective map $`H`$ defined below, with the set $`\overline{𝒜}^{}:=\text{Hom}(𝒳,G)`$ of homomorphisms from the groupoid $`𝒳`$ of composable, holonomically independent, analytical paths into the gauge group. The correspondence is explicitly given by $`\overline{𝒜}\overline{A}H_{\overline{A}}\text{Hom}(𝒳,G)`$ where $`𝒳eH_{\overline{A}}(e):=\overline{A}(h_e)=\stackrel{~}{h}_e(\overline{A})G`$ and $`\stackrel{~}{h}_e`$ is the Gel’fand transform of the function $`𝒜Ah_e(A)G`$. Consider now the restriction of $`𝒳`$ to $`𝒳_\gamma `$, the groupoid of composable edges of the graph $`\gamma `$. One can then show that the projective limit of the corresponding cylindrical sets $`\overline{𝒜}_\gamma ^{}:=\text{Hom}(𝒳_\gamma ,G)`$ coincides with $`\overline{𝒜}^{}`$. Moreover, we have $`\{\{H(e)\}_{eE(\gamma )};H\overline{𝒜}_\gamma ^{}\}=\{\{H_{\overline{A}}(e)\}_{eE(\gamma )};\overline{A}\overline{𝒜}\}=G^{|E(\gamma )|}`$. Let now $`f`$ be a function cylindrical over $`\gamma `$ then
$$\chi _{\mu _0}(\stackrel{~}{f})=_{\overline{𝒜}}𝑑\mu _0(\overline{A})\stackrel{~}{f}(\overline{A})=_{G^{|E(\gamma )|}}_{eE(\gamma )}d\mu _H(h_e)f_\gamma (\{h_e\}_{eE(\gamma )})$$
where $`\mu _H`$ is the Haar measure on $`G`$. As usual, $`𝒜`$ turns out to be contained in a measurable subset of $`\overline{𝒜}`$ which has measure zero with respect to $`\mu _0`$.
Let $`\mathrm{\Phi }_\gamma `$, as before, be the finite linear span of spin-network functions over $`\gamma `$ and $`_\gamma `$ its completion with respect to $`\mu _0`$. Clearly, $``$ itself is the completion of the finite linear span $`\mathrm{\Phi }`$ of vectors from the mutually orthogonal $`\stackrel{~}{\mathrm{\Phi }}_\gamma `$. Our basic coordinates of $`M_\gamma `$ are promoted to operators on $``$ with dense domain $`\mathrm{\Phi }`$. As $`h_e`$ is group-valued and $`P^e`$ is real-valued we must check that the adjointness relations coming from these reality conditions as well as the Poisson brackets (2.1) are implemented on our $``$. This turns out to be precisely the case if we choose $`\widehat{h}_e`$ to be a multiplication operator and $`\widehat{P}_j^e=i\mathrm{}\kappa X_j^e/2`$ where $`X_j^e=X_j(h_e)`$ and $`X^j(h),hG`$ is the vector field on $`G`$ generating left translations into the $`jth`$ coordinate direction of $`Lie(G)T_h(G)`$ (the tangent space of $`G`$ at $`h`$ can be identified with the Lie algebra of $`G`$) and $`\kappa `$ is the coupling constant of the theory. For details see .
## 3 Coherent States from a Coherent State Transform
In the first subsection of this section we will recall the state of the art of families of coherent state transforms which have been defined in the literature already. We point out advantages and disadvantages of one transform as compared to another as well as general properties of every transform and draw attention to some gaps that were left over. In the subsequent subsection we show how some of these gaps can be filled.
### 3.1 Review of Known Results
The first construction of coherent states that are relevant for the quantization of cotangent bundles over connected compact Lie groups $`G`$ is due to Hall who showed how to construct a unitary map between the Hilbert space $`L_2(G,d\mu _H)`$ and a Hilbert space consisting of square integrable holomorphic functions of the complexification $`G^{\text{ }\mathrm{C}}`$ of $`G`$ with respect to some measure $`\nu `$ that he explicitly constructed. In these results were applied to our graph theoretic framework, namely one needs to repeat Hall’s construction, roughly speaking, for every holonomy associated with the various edges of a graph and to glue them together in a cylindrically consistent way. In finally, Hall’s construction was generalized suitably and made applicable to very general phase spaces taking into account also some dynamical aspects. We will now outline the main idea, following :
Central to the subject is the choice of a complex polarization of the classical phase space. In other words, we must choose the analogue of $`z=xip`$ of the harmonic oscillator. This is equivalent to choosing a certain generator $`C`$ (called complexifier in ) of the associated complex symplectomorphism which in the case of the harmonic oscillator consists of the the map $`(x,p)(z,p)`$ and is easily seen to be $`C=p^2/2`$ if, as usual, the symplectic structure is defined by $`\{p,x\}=1`$. Namely we have
$$z=x+i\{x,C\}=\underset{n=0}{\overset{\mathrm{}}{}}\frac{i^n}{n!}\{x,C\}_n$$
(3.1)
where the multiple Poisson bracket is inductively defined by $`\{f,g\}_0=f,\{f,g\}_{n+1}=\{\{f,g\}_n,g\}`$. It is important for the existence of the coherent state transform that the polarization is a positive Kähler polarization, in other words, that the generator $`C`$ is a positive function on the phase space. We will see this in a moment.
The next step consists in the quantization. Following the rule that Poisson brackets be replaced by commutators times $`1/(i\mathrm{})`$ and phase space functions by operators in a suitable ordering we obtain for the harmonic ocillator
$$\widehat{z}=\underset{n=0}{\overset{\mathrm{}}{}}\frac{\mathrm{}^n}{n!}[\widehat{x},\widehat{C}]_n=\widehat{W}_t\widehat{x}\widehat{W}_t^1$$
(3.2)
where we have defined
$$\widehat{W}_t:=e^{\frac{\widehat{C}}{t}}$$
(3.3)
where in this case $`t=\mathrm{}`$ so that $`\widehat{W}_t=\mathrm{exp}(\mathrm{}\mathrm{\Delta }/2)`$ where $`\mathrm{\Delta }=(/x)^2`$ is the Laplacian on $`\mathrm{R}`$. Notice that with our conventions $`\widehat{p}=i\mathrm{}/x`$. One can check that in the case of the harmonic oscillator this gives correctly the annihilation operator $`\widehat{z}=\widehat{x}i\widehat{p}`$. We see that the generator $`C`$ naturally gives rise to the map $`\widehat{W}_t`$ which due to the positiveness of the operator $`\widehat{C}`$ defines a self-adjoint contraction semi-group of bounded operators.
It is for this reason that the following map, called the kernel of the coherent state transform, is well-defined
$$\rho ^t(y,x):=(\widehat{W}_t\delta _y)(x)$$
(3.4)
which for the harmonic oscillator is easily seen to be the standard heat kernel on $`\mathrm{R}`$, $`\delta _y`$ being the $`\delta `$ distribution with respect to $`dx`$ supported at $`x`$.
The coherent states themselves arise as the analytic continuation of the kernel, that is
$$\psi _z^t(x):=\rho ^t(y,x)_{yz}$$
(3.5)
which exists, again, because the operator $`\widehat{C}`$ is positive. It can be shown for the harmonic oscillator that the naturally arising map
$$\widehat{U}_t:=\widehat{K}\widehat{W}_t,$$
(3.6)
where $`\widehat{K}`$ denotes analytic continuation, is a unitary map between $`=L_2(\text{ }\mathrm{R},dx)`$ and $`^{\text{ }\mathrm{C}}=L_2(\text{ }\mathrm{C},d\nu _t)\text{Hol}(\text{ }\mathrm{C})`$ where the latter denotes the space of square integrable holomorphic functions on $`\mathrm{C}`$ with respect to a measure $`\nu _t`$ which is constructed from $`\rho _t`$. For the case of the harmonic oscillator this latter Hilbert space is the familiar Bargmann-Segal-Fock space.
In Hall observed that the case of the harmonic oscillator can be naturally extended to the case of a cotangent bundle over a connected compact Lie group $`G`$, once the following substitutions are made :
$`\text{ }\mathrm{R}G,dxd\mu _H(h),\text{ }\mathrm{C}G^{\text{ }\mathrm{C}},\mathrm{\Delta }\mathrm{\Delta }_G`$ where $`G^{\text{ }\mathrm{C}}`$ is the complexification of $`G`$ and $`\mathrm{\Delta }_G`$ denotes the Lalace-Beltrami operator on $`G`$. In particular, he constructed the map $`\widehat{U}_t`$ and the measure $`\nu _t`$. What he did not analyze, except for phase space bounds, are the anlytical properties of the states $`\psi _g(h)`$ of (3.5), that is, peakedness and Ehrenfest properties. Here and in what follows we will always take $`hG,gG^{\text{ }\mathrm{C}}`$.
In , Hall’s results were applied to the case of a quantum gauge field theory. That is, one applies the coherent state transform as generated by the Laplace Beltrami operator to each copy of the group $`G`$ associated with the edges $`e`$ of a graph of a cylindrical function and obtains a function cylindrical over the same graph but with the holonomies taking values in the complexified gauge group. Thus, coherent states become functions of $`g_eG^{\text{ }\mathrm{C}},eE(\gamma )`$. While this gives a satisfactory mathematical framework for the construction of measures on $`G^{\text{ }\mathrm{C}}`$, the physics of this map was not understood : namely, not only do we need square integrable functions on $`G^{\text{ }\mathrm{C}}`$ but we also need to know what the complex connection is which gives rise to the complexified holonomies, that is, we need to know the map $`(A,E)A^{\text{ }\mathrm{C}}`$ that expresses the complex connection as a function of the real phase space. Otherwise, for instance expectation values which will be functions of the $`g_e`$ cannot be interpreted in terms of the $`(A,E)`$ and thus semi-classical analysis cannot be developed because, say solutions to the Einstein equations, are formulated in terms of the latter.
In order to determine $`A^{\text{ }\mathrm{C}}`$ one must determine the classical limit of the operator which on cylindrical functions reduces to $`\mathrm{\Delta }_\gamma :=_{eE(\gamma )}\mathrm{\Delta }(h_e)`$ where $`h_e`$ is the holonomy of the real connection of the $`G`$bundle along the edge $`e`$ of $`\gamma `$. The problem is, that such a classical limit does not exist !
To see this, notice that roughly $`\mathrm{\Delta }(h_e)(\widehat{E}(S_e)_i/\mathrm{})^2`$ where $`S_e`$ are mutually disjoint analytic surfaces each of which intersects the graph only in one point which is an interior point of both $`e`$ and $`S_e`$ (for definiteness, that intersection can be chosen transversal). However, it is not possible to write down a single operator which reproduces $`\mathrm{\Delta }_\gamma `$ for every $`\gamma `$ and has a classical limit as a well-defined function on the classical phase space $`M`$. Namely, suppose first that $`\mathrm{\Sigma }`$ is compact. Since the graph $`\gamma `$ is arbitrary we may consider a net of finer and finer graphs $`\gamma _ϵ`$ which in the limit $`ϵ0`$ fill all of $`\mathrm{\Sigma }`$. Let us choose the $`\gamma _ϵ`$ to be such that $`\gamma _ϵ\gamma _ϵ^{}`$ for $`ϵ<ϵ^{}`$ and to be (subsets of) cubic lattices of spacing $`ϵ`$ with respect to some spatial background metric. If V is the volume of $`\mathrm{\Sigma }`$ as measured by that metric, then in $`D`$ spatial dimension one will have an order of $`V/ϵ^D`$ vertices in $`\gamma _ϵ`$ each of which accounts for $`D`$ surfaces of area of order $`ϵ^{D1}`$. We see that in the classical limit for sufficiently small $`ϵ`$, using the smoothness of the classical fields
$$\mathrm{\Delta }_{\gamma _ϵ}[ϵ^{2(D1)}\underset{vV(\gamma _ϵ)}{}\underset{I=1}{\overset{D}{}}[E_i^a(v)n_a^I(v)]^2][1+O(ϵ)]$$
(3.7)
where the sum runs over the vertices of $`\gamma _ϵ`$, $`n_a^I(v)`$ is the normal of the surface $`S_{e^I(v)}`$ and $`e^I(v)`$ is an edge of $`\gamma _ϵ`$ that starts at $`v`$ and runs into the $`I`$’th coordinate direction. This object has a chance to converge in the limit $`ϵ0`$ to a well-defined classical quantity only if $`2(D1)=D`$, i.e. $`D=2`$, so that in fact an integral results. For $`D<2`$ this object diverges and for $`D>2`$ it approaches zero for generic field configurations. One could replace $`\mathrm{\Delta }_\gamma `$ by $`_{eE(\gamma )}(\mathrm{\Delta }_e)^{D/(2(D1))}`$ in order to fix this (the eigenvalues would still behave as $`j^{D/(D1)}>j`$), however, while this operator now does have a suitable classical limit at least for the net $`\gamma _ϵ`$, it is no longer diffeomorphism covariant because it carries the sign of the background metric in the definition of the normals $`n_a^I(x)`$. If $`\mathrm{\Sigma }`$ is not compact then (3.7) diverges even in $`D=2`$ (or its just described replacement in any $`D`$) because in gravity the field $`E`$ does not decay at spatial infinity.
In conclusion, there seems to be no classical limit of the cylindrically defined operator $`\mathrm{\Delta }_\gamma `$ as a well-defined, diffeomorphism covariant function on $`M`$ and therefore the interpretation of the $`g_e`$ remains obscure. This state of affairs is clearly unsatisfactory and there are basically two ways out :
Option 1) : One has to choose a different generator of the transform which actually comes from a well-defined function on $`M`$.
Option 2) : One gives up the requirement to have a complex continuum connection $`A^{\text{ }\mathrm{C}}`$ altogether and is satisfied with an interpretation of $`g_e`$ in terms of $`h_e`$ and certain other functions of $`A,E`$ smeared over some surfaces $`S_e`$. Since the latter functions can be interpreted in terms of $`(A,E)`$ one also arrives at an interpretation of $`g_e`$ and this is sufficient in order to do semi-classical physics.
In the next two sections we will describe both options in detail.
Remark :
Before closing this section we would like to point out that a great deal of properties of the coherent states can be obtained already at this point, even if the interpretational issue raised above is not yet answered. Namely, let $`\widehat{C}_\gamma `$ be the cylindrical projections of any complexifier and
$$\psi _{\gamma ,\stackrel{}{g}}^t:=(e^{t\widehat{C}_\gamma }\delta _{\gamma ,\stackrel{}{h}})_{|\stackrel{}{h}\stackrel{}{g}}$$
(3.8)
where $`\stackrel{}{g}=\{g^e\}_{eE(\gamma )}`$ and similar for $`\widehat{h}`$. Moreover, define the annihilation and creation operators respectively ($`A,B,C,..`$ are group indices)
$$\widehat{g}_{AB}^e:=e^{t\widehat{C}_\gamma }\widehat{h}_{AB}^ee^{t\widehat{C}_\gamma }\text{ and }(\widehat{g}^e)_{AB}^{}$$
(3.9)
respectively. Then, without specifying $`\widehat{C}_\gamma `$ at all, the following properties are automatically satisfied (obviously all of this is also theory independent, in the relations below, with the obvious changes, $`\stackrel{}{h}`$ could be any configuration coordinates for its cotangent bundle and $`\stackrel{}{g}`$ their analytical continuations) :
* Eigenvalue Property
The coherent states (3.8) are eigenstates of any of the annihilation operators (3.9)
$`[\widehat{g}_{AB}^e\psi _{\gamma ,\stackrel{}{g}}^t](\stackrel{}{h})=[e^{t\widehat{C}_\gamma }\widehat{h}_{AB}^e\delta _{\gamma ,\stackrel{}{h}^{}}](\stackrel{}{h})_{|\stackrel{}{h}^{}\stackrel{}{g}}`$ (3.10)
$`=`$ $`[e^{t\widehat{C}_\gamma }h_{AB}^e\delta _{\gamma ,\stackrel{}{h}^{}}](\stackrel{}{h})_{|\stackrel{}{h}^{}\stackrel{}{g}}=g_{AB}^e\psi _\stackrel{}{g}^t(\stackrel{}{h})`$
simply because the $`\delta `$-distribution is a generalized eigenfunction of the multiplication operator in the configuration representation.
* Expectation Values for Normal Ordered Operators
From a) it is trivial to see that
$$\frac{<\psi _{\gamma ,\stackrel{}{g}}^t,P(\{\stackrel{}{\widehat{g}}^{},\stackrel{}{\widehat{g}}\})\psi _{\gamma ,\stackrel{}{g}}^t>}{\psi _{\gamma ,\stackrel{}{g}}^t^2}=P(\{\overline{\stackrel{}{g}},\stackrel{}{g}\})$$
(3.11)
where $`P`$ is any normal ordered polynomial of the creation and annihilation operators (annihilation operators to the right).
* Saturation of the Unquenched Heisenberg Uncertainty Relation
Define the symmetric operators
$$\widehat{x}_{AB}^e:=\frac{1}{2}(\widehat{g}_{AB}^e+(\widehat{g}_{AB}^e)^{}),\widehat{y}_{AB}^e:=\frac{1}{2i}(\widehat{g}_{AB}^e(\widehat{g}_{AB}^e)^{})$$
(3.12)
then again with a) it is trivial to see that for the fluctuations we find
$`{\displaystyle \frac{<\psi _{\gamma ,\stackrel{}{g}}^t,(\widehat{x}_{AB}^ex_{AB}^e)^2\psi _{\gamma ,\stackrel{}{g}}^t>}{\psi _{\gamma ,\stackrel{}{g}}^t^2}}={\displaystyle \frac{<\psi _{\gamma ,\stackrel{}{g}}^t,(\widehat{y}_{AB}^ey_{AB}^e)^2\psi _{\gamma ,\stackrel{}{g}}^t>}{\psi _{\gamma ,\stackrel{}{g}}^t^2}}`$ (3.13)
$`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \frac{|<\psi _{\gamma ,\stackrel{}{g}}^t,[\widehat{x}_{AB}^e,\widehat{y}_{AB}^e]\psi _{\gamma ,\stackrel{}{g}}^t>|}{\psi _{\gamma ,\stackrel{}{g}}^t^2}}`$
* Reproducing Property
The connection between the coherent state transform $`\widehat{U}_\gamma ^t:_\gamma _\gamma ^{\text{ }\mathrm{C}}`$, defined analogously to (3.6), and the coherent states is summarized by the following reproducing property, valid for any $`\psi _\gamma `$ :
$$(\widehat{U}_t\psi )(\stackrel{}{g})=<\psi _{\gamma ,\stackrel{}{g}^{}}^t,\psi >$$
(3.14)
where $`gg^{}`$ is the unique involution on $`G^{\text{ }\mathrm{C}}`$ that preserves $`G`$ (this formula can be proved by using the expansion of the group $`\delta `$ distribution in terms of characters according to the Peter&Weyl theorem, see e.g. ).
The additional properties that one would like the coherent states to possess and which do not directly follow from the general form (3.8) are the following, for which we need now the expression for $`g_e`$ in terms of $`A,E`$ :
* Peakedness Properties
We want the coherent states (3.8) to be peaked in the configuration representation at $`A`$, in the momentum representation at $`E`$ and in the Bargmann-Segal representation $`^{\text{ }\mathrm{C}}`$ (the image of $``$ under $`\widehat{U}_t`$ to be defined for general $`\widehat{C}`$ along the lines outlined in ) at $`(A,E)`$. For instance, if with respect to $`\gamma `$ we take as elementary configuration coordinates the holonomies $`h_e`$ and as elementary momentum coordinates the $`E_i(S_e)`$ considered above and if we know the explicit formula $`g_e(\{h_e^{},E(S_e^{})\})`$ which is supposed to be invertible, then we want the probability amplitudes for the coherent state with label $`\stackrel{}{g}`$ in the configuration –, momentum – and Segal-Bargmann Hilbert spaces respectively to be peaked at $`h_e(\stackrel{}{g}),[E(S_e)](\stackrel{}{g}),\stackrel{}{g}`$ respectively. Notice that if we take $`\stackrel{}{g}G^{|E(\gamma )|}`$ then as $`t0`$ $`\psi _{\gamma ,\stackrel{}{g}}^t(\stackrel{}{h})`$ on $`G^{|E(\gamma )|}`$ is supported at $`\stackrel{}{g}=\stackrel{}{h}`$ for any choice of complexifier $`\widehat{C}_\gamma `$ since by its very definition $`\psi _{\gamma ,\stackrel{}{g}}^t`$ approaches $`\delta (\stackrel{}{g},\stackrel{}{h})`$ as $`t0`$.
* Ehrenfest Property
While expectation values of normal ordered polynomials of alternation operators already have the correct expectation values without quantum corrections, we want that to leading order in $`t`$ or $`\mathrm{}`$ also the elementary operators associated with $`h_e,E(S_e)`$ as well as their various commutators divided by $`i\mathrm{}`$ have the correct expectation value guaranteeing the correct infinitesimal quantum dynamics. The fact that the alternation operators do have the correct expectation values makes it plausible that also this property can be verified for any $`\widehat{C}_\gamma `$.
* Overcompleteness
The coherent states should be overcomplete in order to be able to approximate any possible physical situation. Overcompleteness follows automatically if the coherent state transform $`\widehat{U}_t:^{\text{ }\mathrm{C}}`$ is unitary since then that map is onto. More precisely, due to the reproducing property (see e.g. ) :
$$1__\gamma =_{(G^{\text{ }\mathrm{C}})^{|E(\gamma )|}}𝑑\nu _t(\stackrel{}{g})|\psi _{\gamma ,\stackrel{}{g}^{}}^t><\psi _{\gamma ,\stackrel{}{g}^{}}^t|$$
(3.15)
A method for a constructive proof for general $`\widehat{C}`$, up to analytical details, is given in . Namely, the measure $`\nu _t`$ can be uniquely determined if the operator $`\widehat{W}_t`$ is well-defined and if the cylindrical family of measures constructed in can be extended to a $`\sigma `$additive measure on the projective limit of the cylindrical projections of spaces of complex quantum connections that one can define in analogy to .
Overcompleteness is actually also rather plausible for general $`\widehat{C}`$ by inspection because these states arise as the “evolution” under $`\widehat{W}_t`$ of the $`\delta `$ distributions. Now the latter provide a complete basis of generalized functions and $`\widehat{W}_t`$ is invertible on a dense domain of $`\widehat{W}_t^1`$ (the inverse is certainly not bounded).
* Diffeomorphism Covariance
The coherent states should, as all the other states of the Hilbert space, transform covariantly under the diffeomorphism group. This will be the case provided that the operator $`\widehat{C}_\gamma `$ is itself diffeomorphism covariant (does not make use of any background structure), specifically, $`\widehat{U}(\phi )\widehat{C}_\gamma \widehat{U}(\phi )^1=\widehat{C}_{\phi ^1(\gamma )}`$ where $`\text{Diff}(\mathrm{\Sigma })\phi \widehat{U}(\phi )`$ is the unitary representation of the diffeomorphism group described in .
### 3.2 Option 1) : The Volume Operator as the Complexifier
In this section we modify the coherent state transform by choosing a different complexifier. We will argue now that (a suitable power of) the “volume” of a region $`R\mathrm{\Sigma }`$
$$V(R):=_Rd^Dx\sqrt{det(q)(x)}$$
(3.16)
is the most natural candidate. In case that $`\mathrm{\Sigma }`$ is compact or that classically the fields vanish sufficiently fast at spatial infinity as in Yang-Mills theory, we will take $`R=\mathrm{\Sigma }`$ in the sequel. Otherwise, we will take $`R`$ to be a bounded region to begin with and send $`R\mathrm{\Sigma }`$ only after all calculations have been performed. Here,
$$det(q):=\sqrt[D1]{det(E_i^aE_i^b)}$$
(3.17)
and (3.16) is called the volume functional because in the case of general relativity $`E_i^a=\sqrt{det(q)}e_i^a`$ where $`e_i^a`$ is the $`D`$-bein field and $`q_{ab}`$ is the $`D`$-metric intrinsic to $`\mathrm{\Sigma }`$.
The reasons are as follows :
(i)
As it is clear from the discussion in the previous section, it is important that $`C`$ is a positive semi-definite function on the phase space as this translates into a positive definite operator upon quantization. The volume has this property.
(ii)
Notice that even in the case of gauge theories on a background metric the electric field is a Lie algebra valued vector density of weight one. Therefore, $`E_i^aE_i^b=det(q)q^{ab}`$ is in general a gauge invariant tensor density of weight two. Hence, (3.17) is a scalar density of weight two which can be constructed without any background structure and therefore the volume functional is diffeomorphism invariant if $`R=\mathrm{\Sigma }`$ and diffeomorphism covariant if $`R\mathrm{\Sigma }`$ ! This is important in order to obtain diffeomorphism covariant coherent states in the case of diffeomorphim invariant quantum field theories of connections.
(iii)
As we want to start with a Hilbert space which consists of square integrable functions of connections for which the connection operator is a multiplication operator, it is natural to consider an operator which is entirely constructed from the electric field operator so that the analogue of $`z`$ is given by $`Z_a^j=A_a^j+if_a^j(E)`$. The volume density is the simplest scalar density of weight one entirely constructed from electric fields.
(iv)
Using the symplectic structure $`\{A_a^i(x),E_j^b(y)\}=\kappa \delta _a^b\delta _j^i\delta (x,y)`$ where $`\kappa `$ is the coupling constant and
$$C(R):=\frac{1}{\lambda \kappa }V(R)^n$$
(3.18)
where $`n1`$ is a positive real number and $`\lambda `$ is a positive, possibly dimensionful, parameter so chosen that for $`xR`$
$$f_a^j(x):=\{C(R),A_a^j(x)\}=\frac{nV(R)^{n1}}{\lambda }\frac{\sqrt[2(D1)]{det(E_i^bE_i^c)}}{E_j^a}=:nV(R)^{n1}\frac{e_a^i}{\lambda (D1)}$$
(3.19)
has dimension of inverse length we easily see that $`E_j^a`$ can be reconstructed from $`f_a^j`$ and therefore the complex connection $`Z_a^j`$ together with its complex conjuate contains full phase space information. The field $`e_a^i`$ is the co-$`D`$-bein in general relativity.
Notice that $`Z_a^j=A_a^jif_a^j`$ really transforms as a $`G`$connection under gauge transformations since $`\delta Z=d\mathrm{\Lambda }+[\mathrm{\Lambda },A]+i[\mathrm{\Lambda },e]=d\mathrm{\Lambda }+i[\mathrm{\Lambda },Z]`$ so that the coherent state transform is both diffeomorphism covariant and gauge covariant.
(v)
Finally, to be useful, it is necessary that one can quantize the generator. But this is the case for the volume functional in any dimension along the lines of . Moreover, on the Hilbert space that we have chosen in section 2 the spectrum of that operator is entirely discrete and, although very complicated, explicitly known at least in terms of matrix elements .
Upon quantization $`\widehat{E}_i^a=i\mathrm{}\kappa \delta /\delta A_a^i`$ and the generator takes the following form on cylindrical functions
$$\widehat{C}(R)=\frac{(\mathrm{}\kappa )^{n\frac{D}{D1}}}{\lambda \kappa }\widehat{v}$$
(3.20)
where $`\widehat{v}`$ is a dimensionless operator constructed from invariant vector fields corresponding to the copies of the group associated with the edges of graphs. The coherent state transform is then generated by
$$\widehat{W}_t=e^{t\widehat{v}}\text{ where }t=\frac{(\mathrm{}\kappa )^{n\frac{D}{D1}1}}{\lambda }$$
(3.21)
is a dimensionless parameter which vanishes as $`\mathrm{}0`$. For instance, for general relativity in $`3+1`$ dimensions, $`\mathrm{}\kappa `$ is the Planck area.
Next, we define coherent states in analogy to (3.5). The idea is to define coherent states graphwise, which means that the state approximates a certain point in the classical phase space on that graph only. We can do this for every graph which is contained in the region $`R`$. In case that $`R\mathrm{\Sigma }`$ this does not exclude the possibility to have graphs which run to spatial infinity : We can use the asymptotic structure available and allow only such graphs which run to spatial infinity inside fixed “thin tubes” of $`R`$ which have finite Lebesgue measure. Notice that these complications would not be necessary if we would choose $`n=1`$. In general we cannot choose $`n=1`$ for reasons explained below, see also the model described in section 6.
The fundamental definition is
$$\psi _{\gamma ,Z}^t(A):=(\widehat{W}_t\delta _{\gamma ,A^{}})(A)_{|A^{}Z}$$
(3.22)
where the $`\delta `$ distribution in (3.22) is defined by
$$\delta _{\gamma ,A^{}}(A):=\underset{\stackrel{}{j},\stackrel{}{J}}{}\overline{T_{\gamma ,\stackrel{}{j},\stackrel{}{J}}(A)}T_{\gamma ,\stackrel{}{j},\stackrel{}{J}}(A^{})$$
(3.23)
and where the sum is over all possible not necessarily gauge invariant spin-network functions on that graph $`\gamma `$ if we work at the non-gauge invariant level while it is over all possible gauge invariant functions only if we work at the gauge invariant level. It is important to stress that in (3.23) we include only spin-network states whose vector of representations $`\stackrel{}{j}`$ does not contain a zero entry.
We thus obtain coherent states $`\psi _{\gamma ,Z}^t`$ which have the property to be orthogonal,
$`<\psi _{\gamma ,Z}^t,\psi _{\gamma ^{},Z^{}}^t>=0`$, if their underlying graphs are different, $`\gamma \gamma ^{}`$. We also define coherent states of a different type,
$$\mathrm{\Psi }_{\gamma ,Z}^t:=\underset{\gamma ^{}\gamma }{}\psi _{\gamma ^{},Z}^t$$
(3.24)
where the sum extends over all subgraphs of $`\gamma `$ which can be obtained from $`\gamma `$ by deleting edges of $`\gamma `$ in all possible ways (if the state is to be gauge invariant then the sum extends over closed subgraphs only). The idea is not to take inner products of states (3.24) with different $`\gamma `$ but only between those with the same $`\gamma `$ but different $`Z,Z^{}`$. In other words, one first restricts the Hilbert space $``$ to the completion $`_\gamma `$ of the span of spin-network states over closed subgraphs of $`\gamma `$ and then one lets $`\gamma `$ grow. Recall that given two piecewise analytic graphs, their union is still a piecewise analytic graph. (We cannot immediately transfer our definitions to the smooth category of webs because there we do not have the notion of an orthogonal basis unless we restrict ourselves to non-degenerate webs ; we will, however, not go into this subject in the present paper). Therefore, there is a generalized projective structure on the set of piecewise analytical graphs and the final coherent state $`\mathrm{\Psi }_Z^t`$ is a generalized projective limit of the states $`\mathrm{\Psi }_{\gamma ,Z}^t`$.
Let us illustrate the situation by drawing an analogy with the coherent states for, say, an (in)finite number $`N\mathrm{}`$ of harmonic oscillators : The role of the graph label $`\gamma =(e_1,..,e_n),n<\mathrm{}`$, given as a finite collection of edges, is played by the mode label $`\stackrel{}{k}=(k_1,..,k_n),n<N`$, given by a finite collection of non-negative integers. The analogues of the states $`\mathrm{\Psi }_{\gamma ,Z}^t=\mathrm{\Psi }_{\gamma ,\stackrel{}{g}}^t`$ with $`\stackrel{}{g}=(h_{e_1}(Z),..,h_{e_n}(Z))`$ is given by the coherent state for $`n`$ uncoupled harmonic oscillators $`\mathrm{\Psi }_{\stackrel{}{k},\stackrel{}{z}}^t`$ with an array of complex numbers $`\stackrel{}{z}=(z_{k_1},..,z_{k_n})`$ corresponding to the mode vector $`\stackrel{}{k}`$. The projective limit of taking the “biggest possible graph” corresponds to taking the (in)finite direct product limit $`\stackrel{}{k}(1,2,..,N)`$ and one obtains the full coherent state $`\mathrm{\Psi }_Z^t,Z=(z_1,z_2,..,z_N)`$. One does not compute inner products between states with different $`\stackrel{}{k}`$ but only with different $`\stackrel{}{z}`$ for the same $`\stackrel{}{k}`$ which models the properties of $`\mathrm{\Psi }_Z^t`$ on its cylindrical projections $`\mathrm{\Psi }_{\stackrel{}{k},\stackrel{}{z}}^t`$. The analogues of the states $`\psi _{\gamma ,Z}^t`$ are the states $`\psi _{\stackrel{}{k},\stackrel{}{z}}^t=\mathrm{\Psi }_{\stackrel{}{k},\stackrel{}{z}}^t\mathrm{\Omega }<\mathrm{\Omega },\mathrm{\Psi }_{\stackrel{}{k},\stackrel{}{z}}^t`$ where we have taken out the vacuum mode so that $`<\psi _{\stackrel{}{k},\stackrel{}{z}}^t,\psi _{\stackrel{}{k}^{},\stackrel{}{z}^{}}^t>=0`$ for $`\stackrel{}{k}\stackrel{}{k}^{}`$.
Notice that the restriction of $`\mathrm{\Psi }_{\gamma ,Z}^t(A)`$ to real valued $`Z=A^{}`$ is the “heat kernel” $`\rho _{\gamma ,t}(A,A^{})`$ for the “heat equation”
$$[/t+\widehat{v}]\rho _{\gamma ,t}(A,A^{})=0\text{ such that }\rho _{\gamma ,0}(A,A^{})=\delta _\gamma (A,A^{}).$$
(3.25)
As the volume operator is an essentially self-adjoint, positive semi-definite operator with discrete spectrum which leaves the subspace of $``$ spanned by spin-network states of given $`\gamma ,\stackrel{}{j}`$ invariant, we can diagonalize it and define another orthonormal basis of eigenstates $`T_{\gamma ,\lambda ,n}`$ of $`\widehat{v}`$ where $`\lambda `$ labels the eigenvalue and the integer $`n`$ its degeneracy. We can then write (3.23) alternatively as
$$\delta _{\gamma ,A}(A^{}):=\underset{\lambda ,n}{}\overline{T_{\gamma ,\lambda ,n}(A^{})}T_{\gamma ,\lambda ,n}(A)$$
(3.26)
which allows us to explictly compute the coherent states as
$$\psi _{Z,\gamma ,t}(A)=\underset{\lambda ,n}{}e^{t\lambda }T_{\gamma ,\lambda ,n}(Z)\overline{T_{\gamma ,\lambda ,n}(A)}.$$
(3.27)
The function (3.27) is to be understood in the following sense : Given a point in the classical phase space $`A,E`$, compute the $`G^{\text{ }\mathrm{C}}`$ connection $`Z=Aif(E)`$ and from this its holonomies $`h_e^{\text{ }\mathrm{C}}:=h_e(Z)`$ for each edge $`e`$ of $`\gamma `$. Then insert these elments of $`G^{\text{ }\mathrm{C}}`$ into the eigenfunctions appearing in the series in (3.27).
Several points of worry arise when looking at (3.27) :
(i)
Does the series in (3.27) converge, in the sup-norm topology with respect to $`G^n`$, where $`n`$ denotes the number of edges of $`\gamma `$ ? This will, in particular, not be the case if one of the $`\lambda `$ has infinite multiplicity. The volume operator as defined in , however, has presumably precisely this property for the zero eigenvalue, at least in the case of general relativity in 3+1 dimensions which requires $`G=SU(2)`$! Thus, in this case, in order to make sense of (3.27) we must discard the zero volume eigenstates even from the kinematical Hilbert space. (In particular this has to be done at the gauge non-invariant level). This is quite satisfactory because the classical phase space can be viewed as a cotangent bundle over smooth, signature $`(+,..,+)`$ $`D`$-metrics for which vanishing volume, that is, vanishing determinant of the three-metric, is not allowed. That the signature is $`(+,..,+)`$ is guaranteed if we restrict to states with non-vanishing expectation value for the area operator for every surface that intersects the graph.
But even if all eigenvalues have finite multiplicity, the series does not necessarily converge : while $`T_{\gamma ,\lambda ,n}`$ is a bounded function of $`G^n`$, it is not any longer so of $`(G^{\text{ }\mathrm{C}})^n`$ because that group is not compact. What is needed, roughly speaking, is the following :
we can decompose the $`T_{\gamma ,\lambda ,n}`$ in terms of spin-network functions which turns the above series into a series over $`\stackrel{}{j},\stackrel{}{J}`$. The coefficient of $`T_{\gamma ,\stackrel{}{j},\stackrel{}{J}}(Z)`$ is of the form $`e^{t\lambda (\stackrel{}{j},\stackrel{}{J})}`$ times something that grows at most linearly with $`\stackrel{}{j},\stackrel{}{J}`$ while $`T_{\gamma ,\stackrel{}{j},\stackrel{}{J}}(Z)`$ grows exponentially with $`\stackrel{}{j},\stackrel{}{J}`$ for $`Z`$ in the non-compact directions of $`G^{\text{ }\mathrm{C}}`$. Thus, for the series to converge it would be sufficient if
$$\lambda (\stackrel{}{j},\stackrel{}{J})c(\underset{eE(\gamma )}{}j_e+\underset{vV(\gamma }{}J_v)^{1+ϵ}$$
(3.28)
where $`c`$ is a positive number independent of $`\stackrel{}{j},\stackrel{}{J}`$ and $`ϵ`$ can be any positive number. The criterion (3.28) is a condition on the spectrum of $`\widehat{v}`$ which needs to be checked to hold. This is the reason why we have allowed for a power $`n`$ different from $`n=1`$ in (3.18) : by taking $`n`$ sufficiently large we can guarantee that criteron (3.28) is satisfied.
More precisely we have the following :
Looking at (3.16), (3.18) and the explicit expression for the volume operator as derived in we infer that the electric fields get, roughly speaking, replaced by right invariant vector fields $`X_e^i:=X^i(h_e)`$ on the various copies of $`G`$ corresponding to the edges of $`\gamma `$. As those act on spin-network functions roughly by multiplication by $`j_e`$, we find the eigenvalues of the volume opertor to be of the form
$$\lambda (\stackrel{}{j},\stackrel{}{J})=(P_{2D}(\stackrel{}{j},\stackrel{}{J})^{n/(2(D1))}$$
(3.29)
where $`P_{2D}`$ is a homogenous, positive polynomial of degree $`2D`$ which depends non-trivially on all the variables $`\stackrel{}{j},\stackrel{}{J}`$. Obviously, taking $`n>2(D1)`$ we have good chances to satisfy (3.28). Presumably, $`n>D1`$ will be sufficient since because of gauge invariance the $`\stackrel{}{j},\stackrel{}{J}`$ do not have independent ranges. For instance, for $`SU(2)`$, due to gauge invariance the sum of all but one, say $`j_{e_0}`$, of those $`j_e`$ that correspond to $`e`$’s which meet at a common vertex must always exceed the value of $`j_{e_0}`$. Moreover, the value of $`J_v`$ is bounded by the sum of all those $`j_e`$. These relations hold for all of the vertices and thus there is a good chance that we can estimate (3.29) as
$$\lambda (\stackrel{}{j},\stackrel{}{J})c(\mathrm{max}(\stackrel{}{j},\stackrel{}{J}))^{\frac{n}{D1}}$$
(3.30)
which would be sufficient for any $`D`$. However, this must be checked in the case at hand. In particular, it could happen that the function $`\lambda (\stackrel{}{j},\stackrel{}{J})`$ has “degenerate directions” in which case even large $`n`$ would not help to make the series converge.
(ii)
Even if the series converges, are these coherent states square integrable ? We easily see that the convergence of the series is sufficient for this to be the case. Namely, if the series converges, we compute the norm as
$$\psi _Z^2=\underset{\lambda ,n}{}e^{2t\lambda }|T_{\gamma ,\lambda ,n}(Z)|^2$$
(3.31)
which then certainly converges as well.
(iii)
Finally, is the generalized projective limit $`\mathrm{\Psi }_Z^t`$ of the states $`\mathrm{\Psi }_{\gamma ,Z}^t`$ square integrable ? There is no, not even partial aswer to this question available at the moment, however, notice that even the uncountably infinite direct product limit of an uncountably infinite number of harmonic oscillators is square integrable. This follows immediately from the Kolmogorov theorem for an uncountably infinite tensor product of probability measure (here : Gaussian measures) Hilbert spaces. Thus, the normalizability of $`\mathrm{\Psi }_Z^t`$ is indeed conceivable.
(iv)
If we can then verify the properties (a)-(h) mentioned above, what we will have achieved is that we have states that are peaked on a classical configuration $`Z`$ in the sense that the operator $`\widehat{g}^e`$ corresponding to $`g^e=h^e(Z)`$ has expectation value $`g^e`$, saturates the Heisenberg uncertainty bound etc. However, since all these properties (a)-(h) are verified for $`g^e`$ only, we must ask whether we can reconstruct $`Z`$ from all the $`h^e(Z)`$, that is, whether the holonomies separate the points on the space of smooth complexified connections.
This is a non-trivial question due to the presence of so-called null-rotations for non-compact gauge groups and amounts to proving a Giles’ theorem for non-compact gauge groups. At least for $`SU(2)^{\text{ }\mathrm{C}}=SL(2,\text{ }\mathrm{C})`$, this has been answered affirmatively in an appropriate sense in and we believe the proof to be valid generally for complexifications of compact connected gauge groups. If we work at the gauge non-invariant level, the proof is obvious since we just have to consider the limit of infinitesimal open paths.
We now argue that the coherent states (3.24) so constructed have very good chances to satisfy all the properties (d)-(h) mentioned above, assuming that there are no convergence problems even at the gauge non-invariant level. We will indicate the necessary modifications of the analysis when we restrict to the gauge invariant sector. The analysis is in fact quite general and can be generalized to the quantization of any field theory with a generalized projective structure, once a choice of the complexifier $`C`$ and a choice of polarization of the classical phase space has been made.
(d)
The way in which these states are localized is obscure at the moment. In this paper we will just outline how one might prove this property. First of all, notice that the coherent states become, for real connections $`Z`$, just $`\delta `$ distributions on $`\gamma `$ in the semi-classical limit as $`\mathrm{}0`$ (that is, $`t0`$, see (3.21), (3.25)). Thus, in the connection representation, the state $`|Z,\gamma ,t>`$ is certainly peaked at $`A=\mathrm{}(Z)`$ as $`t0`$ for $`\mathrm{}(Z)=0`$ for any $`\gamma `$. What happens if $`\mathrm{}(Z)0`$ is unclear at the moment. Next, we want to study the state $`|Z,\gamma ,t>`$ in the momentum or electric field representation which is nothing else than the spin-network representation (see ). Now, the representation (3.27), with the $`T_{\gamma ,\lambda ,n}`$ written in the spin-network basis, is not immediately useful in order to study the behaviour of the state in the limit $`t0`$ because the exponential terms become unity in the limit $`t0`$, that is, the convergence of the series worsens in the limit $`t0`$. The idea is to use a Poisson summation formula which exists for all compact gauge groups and which should transform the series into a series with coefficients of the form $`\mathrm{exp}(\lambda (\stackrel{}{j},\stackrel{}{J})/t^\alpha )`$ where $`\alpha `$ is a positive number. In the limit $`t0`$ then the leading term would be the one with $`\lambda `$ closest to zero and this would be the peakedness property in the electric field representation. We will actually use this method in the next paper of this series for the original heat kernel complexifier.
(e)
To prove the Ehrenfest property is very much like proving the peakedness property in the Bargmann Segal representation and also should be based on the Poisson summation formula. The reason is that expectation values of polynomials in the basic operators can be expanded, using overcompleteness of the coherent states, as a polynomial in the matrix elements between normalized coherent states $`\xi _{\gamma ,\stackrel{}{g}}^t`$ where the extra variables $`\stackrel{}{g}^{}`$ as in (3.15) are integrated over with respect to $`d\nu _t(\stackrel{}{g}^{})\psi _{\gamma ,\stackrel{}{g}}^t^2`$. But then the Ehrenfest property follows once we find for any elementary operator $`\widehat{O}`$ that
$$<\xi _{\gamma ,\stackrel{}{g}}^t,\widehat{O}\xi _{\gamma ,\stackrel{}{g}^{}}^t>=O(\stackrel{}{g})<\xi _{\gamma ,\stackrel{}{g}}^t,\xi _{\gamma ,\stackrel{}{g}^{}}^t>(1+O(t))$$
which in turn should be easy to establish if the overlap function on the right hand side of this equality is peaked at $`\stackrel{}{g}=\stackrel{}{g}^{}`$. But the latter property is just the same as the peakedness property in the Segal-Bargmann representation which can be seen generally from the reproducing property.
(f)
As already said, the (over)completeness of the coherent states in the kinematical Hilbert space $`=L_2(\overline{𝒜},d\mu _0)`$ would follow trivially if one could establish that the map (3.6), generalized to our context, is a unitary map between $``$ and a suitable $`L_2`$ space of holomorphic functions of complex connections with respect to a measure $`\nu _t`$ because then the map $`\widehat{W}_t`$ would be onto, in particular. In addition, the general comments from the previous section apply.
(h)
The coherent states of this section are diffeomorphism covariant by their very construction.
This concludes the general outline of how one might construct coherent states for quantum gauge theories from a coherent state transform which can also be interpreted in terms of a complex connection $`A^{\text{ }\mathrm{C}}`$. In we will, however, not use the volume operator as the complexifier for the following reasons :
i)
The spectrum of the volume operator is not explicitly known. This lack of knowledge makes analytical proofs very hard although a numerical method is of course possible.
ii)
More serious is the following observation : Unless $`V(R)`$ itself is a polynomial function of the $`E_i(S)`$, then even classically the $`g_{AB}^e,\overline{g}_{AB}^e`$ do not a form a Poisson algebra for $`D>2`$. This becomes obvious from the fact that while $`\{A_a^{j\text{ }\mathrm{C}}(x),A_b^{k\text{ }\mathrm{C}}(y)\}=\{E_j^a(x),E_k^b(y)\}=0,\{A_a^{j\text{ }\mathrm{C}}(x),E_k^b(y)\}=\delta _a^b\delta _k^j\delta (x,y)`$ (the complexifier induces a canonical transformation) we have
$$\{A_a^{j\text{ }\mathrm{C}}(x),\overline{A_b^{k\text{ }\mathrm{C}}(y)}\}=\delta (x,y)\frac{^2V(R)}{E_j^a(x)E_k^b(x)}+\text{more}$$
where “more” is non-distributional. Thus, since the connections are only smeared in one spatial direction inside a holonomy functional, it follows that for $`D>2`$ the Poisson bracket $`\{g_{AB}(A^{\text{ }\mathrm{C}}),\overline{g_{AB}(A^{\text{ }\mathrm{C}})}\}`$ is necessarily distributional or even ill-defined and does not lie in the original Poisson algebra any longer. This means that the fluctutions of the $`\widehat{x}_{AB},\widehat{y}_{AB}`$ are ill-defined if the Ehrenfest property holds because the right hand side of (3.13) will then be proportional to $`\{g_{AB}(A^{\text{ }\mathrm{C}}),\overline{g_{AB}(A^{\text{ }\mathrm{C}})}\}`$ to first order in $`t`$. Whether or not this is bad is unclear, after all it is unnecessary to work with $`\widehat{g}_{AB}`$ itself. On the other hand, due to the eigenvalue property and the similarity with the creation and annihilation operator algebra it would be very convenient to have the $`\widehat{g}_{AB}`$ at one’s disposal.
Due to these difficulties we will turn to option ii) in the remainder of this paper and the subsequent issues of this series. It should be kept in mind, however, that option 1) exists. Its obvious advantage is that one has an actual complex connection which implies that one can work entirely with graphs and never needs the additional dual polyhedronal decompositions which are a source of ambiguity.
### 3.3 Option 2) : The Heat Kernel Complexifier
In this section we will be satisfied with obtaining $`g_e`$ as a definite function of the functions $`h_e,P_e`$ described in section 2.1. We do not require that $`g_e`$ is itself the holonomy along $`e`$ for some complex connection $`A^{\text{ }\mathrm{C}}`$.
The results of this section hold for arbitrary compact, semisimple connected gauge groups and direct products of such with Abelian ones.
As we want to bring in Planck’s constant $`\mathrm{}`$ as a measure of closeness to classical physics, we need to spend a few moments on dimensionalities as in the previous section for the volume functional. The dimension of the time coordinate $`x^0`$ is taken to be the same as that of the spatial coordinates $`x^a`$, namely $`[x^0]=[x^a]=`$cm<sup>1</sup> which can always be achieved by absorbing an appropriate power of the speed of light into the coupling constant $`\kappa `$ of the theory.
We will take our connection one-form to be of dimension $`[A]=`$cm<sup>-1</sup> so that its holonomy is dimensionless. In $`D+1`$ spacetime dimensions the kinetic term of the classical action is given by
$$A_{kin}=\frac{1}{\kappa }_{\text{ }\mathrm{R}}𝑑t_\mathrm{\Sigma }d^DxE_i^a(x)\dot{A}_a^i(x)$$
and its dimension is that of an action, that is, $`[A_{kin}]=[\mathrm{}]`$. In Yang-Mills theories the electric field is a first derivative of $`A_a^i`$ and thus has dimension $`[E_i^a]=`$cm<sup>-2</sup>. In general relativity the metric components, the D-beins and also $`[E_i^a]=`$cm<sup>0</sup> are dimensionfree. It follows that in Yang-Mills (YM) theory the Feinstruktur constant
$$\alpha :=\mathrm{}\kappa $$
(3.32)
has dimension $`[\alpha ]:=`$cm<sup>D-3</sup> and in general relativity (GR) $`[\alpha ]=`$cm<sup>D-1</sup>.
Let now $`\gamma `$ be a graph and consider the symplectic manifold $`(M_\gamma ,\mathrm{\Omega }_\gamma )`$ introduced in section 2.1 with its canonical coordinates $`h_e,P_i^e:eE(\gamma )`$. The electric flux variable (2.5) then has dimension $`[P_i^e]=`$cm<sup>D-3</sup> in YM and cm<sup>D-1</sup> in GR respectively and in general let $`[P_i^e]=`$cm$`^{n_D^{}}`$. Let now $`a`$ be an arbitrary but fixed constant with the dimension of a length, $`[a]=`$cm<sup>1</sup>, say $`a=1`$cm if $`n_D0`$ and let $`a`$ be dimensionfree otherwise. Then we introduce the dimensionfree quantity
$$p_i^e:=\frac{P_i^e}{a^{n_D}}$$
(3.33)
where $`n_D=n_D^{}`$ if $`n_D^{}0`$ and $`n_D=1`$ otherwise. Notice that a natural choice for a dimensionful constant in general relativity in any $`D>1`$ would be $`a=1/\sqrt{|\mathrm{\Lambda }|}`$ where $`\mathrm{\Lambda }`$ is the (supposed to be non-vanishing) cosmological constant.
On the other hand, it is $`E_i^a/\kappa `$ which is canonically conjugate to $`A_a^i`$ rather than $`E_i^a`$ itself, therefore the brackets (2.1) get modified into
$`\{h_e,h_e^{}\}_\gamma `$ $`=`$ $`0`$
$`\{{\displaystyle \frac{P_i^e}{\kappa }},h_e^{}\}_\gamma `$ $`=`$ $`\delta _e^{}^e{\displaystyle \frac{\tau _i}{2}}h_e`$
$`\{{\displaystyle \frac{P_i^e}{\kappa }},{\displaystyle \frac{P_j^e^{}}{\kappa }}\}_\gamma `$ $`=`$ $`\delta ^{ee^{}}f_{ij}^k{\displaystyle \frac{P_k^e}{\kappa }}`$ (3.34)
We are now ready to define the complexifier for the symplectic manifold $`M_\gamma `$, it is given by
$$C_\gamma :=\frac{1}{2\kappa a^{n_D}}\underset{eE(\gamma )}{}\delta ^{ij}P_i^eP_j^e$$
(3.35)
and since $`C_\gamma `$ is gauge invariant it will pass to the reduced phase space. Using the partial order $``$ of or section 2.1 it is immediately clear that $`C_\gamma `$ defines a self-consistently defined function on the $`M_\gamma `$, that is, for $`\gamma \gamma ^{}`$ we have $`\{p_{\gamma ^{}\gamma }^{}C_\gamma ,p_{\gamma ^{}\gamma }^{}f_\gamma \}_\gamma ^{}=p_{\gamma ^{}\gamma }^{}\{C_\gamma ,f_\gamma \}_\gamma `$ for any $`f_\gamma C^{\mathrm{}}(M_\gamma )`$.
We can explicitly compute the complexified holonomy and complexified momenta for any compact, semi-simple gauge group $`G`$. Since $`\{P_i^e,C_\gamma \}=0`$ (gauge invariance of $`C_\gamma `$) we have
$`\{h_e,C_\gamma \}_\gamma `$ $`=`$ $`P_i^e{\displaystyle \frac{\tau _i}{2a^{n_D}}}h_e=p_i^e{\displaystyle \frac{\tau _i}{2}}h_e`$
$`\{h_e,C_\gamma \}_{\gamma (2)}`$ $`=`$ $`{\displaystyle \frac{1}{a^{2n_D}}}P_i^eP_j^e{\displaystyle \frac{\tau _i\tau _j}{4}}h_e=(p_j^e{\displaystyle \frac{\tau _j}{2}})^2h_e`$ (3.36)
where we define generally $`p^e:=\sqrt{p_j^ep_j^e}`$. In the second line of (3.3) we have made use of the fact that $`G`$ is semi-simple so that the structure constants are completely skew and so $`\{p_j^e,C_\gamma \}=0`$.
We therefore conclude that the complexification of $`h_e`$ is given by
$`h_e^{\text{ }\mathrm{C}}`$ $`:=`$ $`g_e={\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{i^n}{n!}}\{h_e,C\}_{(n)}`$ (3.37)
$`=`$ $`[{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{i^n}{n!}}(p_j^e{\displaystyle \frac{\tau _j}{2}})^n]h_e`$
$`=`$ $`e^{i\tau _jp_j^e/2}h_e=:H_eh_e`$
and similarly $`P_i^{e\text{ }\mathrm{C}}=P_i^e`$. Thus we have established the following.
###### Lemma 3.1
The complexification of the holonomy for compact and semisimple $`G`$ is given directly as a left polar decomposition, where the right unitary factor is the holonomy of the compact gauge group while the left positive definite hermitean factor is just the exponential of $`ip_j^e\tau _j/2`$.
For $`G=U(1)`$ the generator $`\tau _j/2`$ has to be replaced by the imaginary unit $`i`$.
Notice that (3.37) makes sense since $`p_j^e`$ is dimensionless. Moreover, we have naturally stumbled on the diffeomorphism
$$\mathrm{\Phi }:T^{}(G)G^{\text{ }\mathrm{C}};(p^j,h)g:=Hh=e^{ip^j\tau _j/2}h.$$
(3.38)
The diffeomorphism (3.38) has a further consequence : $`(T^{}(G),\omega )`$ is a symplectic manifold while $`G^{\text{ }\mathrm{C}}`$ is a complex manifold. Thus, $`T^{}(G)`$ is a symplectic manifold with a complex structure which, as one can show ( and references therein), is $`\omega `$-compatible. In fact, $`\omega `$ is just given by (3.3) with $`P_i^e`$ replaced by $`p_i`$ and the label $`e=e^{}`$ dropped. Therefore, $`T^{}(G)`$ is in fact a Kähler manifold and a Segal-Bargmann representation (wave functions depending on $`g`$) corresponds to a positive Kähler polarization .
Finally, let us compute the Segal-Bargmann transform corresponding to $`C_\gamma `$ as in . As follows from the previous section, we have in the connection representation (wave functions depending on the $`h_e`$)
$$\widehat{P}_j^e=\frac{i\mathrm{}\kappa }{2}X_j^e\text{ where }X_j^e=X_j(h_e),$$
(3.39)
and $`X_j(h)`$ denotes the right invariant vector fields on $`G`$ at $`h`$, that is $`X_j(h):=\text{tr}((\tau _jh)^T/h)`$. Thus, the coherent state transform is (following the notation of )
$$\widehat{W}_{\gamma t}:=e^{\frac{\widehat{C}_\gamma }{\mathrm{}}}=e^{\frac{t}{2}\mathrm{\Delta }_\gamma }$$
(3.40)
where we have defined the Laplacian on $`\gamma `$ by
$$\mathrm{\Delta }_\gamma =\underset{eE(\gamma )}{}\mathrm{\Delta }_e,\mathrm{\Delta }_e=\frac{1}{4}\delta ^{ij}X_i^eX_j^e$$
(3.41)
and the heat kernel time parameter has the following interpretation in terms of the fundamental constants of the theory
$$t:=\frac{\mathrm{}\kappa }{a^{n_D}}.$$
(3.42)
Notice that $`a`$ is just a parameter that we have put in by hand to make things dimensionless, for instance, it could be $`1`$cm in quantum general relativity in $`D+1=4`$ spacetime dimensions or $`a=10^5`$ for Yang-Mills in $`D+1=4`$ and thus is “large”. The semiclassical limit $`\mathrm{}0`$ thus corresponds to $`t0`$. That $`t`$ is a tiny positive real number will be crucial in all the estimates that we are going to perform in this and the next paper of this series.
The factor of $`1/4`$ in the definition of $`\mathrm{\Delta }_e`$ relative to $`(X_j^e)^2`$ is due to the factor of $`1/2`$ in the second Poisson bracket of (3.3) and it is the same factor which gives $`\mathrm{\Delta }_e`$ the standard spectrum $`j(j+1);j=0,\frac{1}{2},1,\frac{3}{2},..`$ for the case of $`G=SU(2)`$.
We can also explicitly compute the quantum operator corresponding to $`g_e`$ in (3.37) for arbitrary $`G`$. We have
$`\widehat{g}_e`$ $`=`$ $`e^{t\mathrm{\Delta }_\gamma /2}\widehat{h}_e^{t\mathrm{\Delta }_\gamma /2}={\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(t)^n}{2^nn!}}[\widehat{h}_e,\mathrm{\Delta }_e]_{(n)}`$
$`[\widehat{h}_e,\mathrm{\Delta }_e]`$ $`=`$ $`{\displaystyle \frac{1}{4}}(X_e^i\tau _i\widehat{h}_e+\tau _i\widehat{h}_eX_e^i)=X_e^i{\displaystyle \frac{\tau _i}{2}}\widehat{h}_e{\displaystyle \frac{(\tau _i)^2}{4}}\widehat{h}_e`$ (3.43)
Since $`\mathrm{\Delta }_\gamma `$ commutes with $`X_e^i`$ we immediately find
$`\widehat{g}_e`$ $`=`$ $`e^{t\widehat{X}_e^i\frac{\tau _i}{4}t\frac{\tau _i^2}{8}}\widehat{h}_e=e^{i\widehat{p}_e^j\frac{\tau _j}{2}\frac{t\tau _j^2}{8}}\widehat{h}_e=e^{i\widehat{p}_e^j\frac{\tau _j}{2}}e^{t\frac{\tau _j^2}{8}}\widehat{h}_e`$ (3.44)
since $`itX_e^j/2=\widehat{p}_j^e`$ and in the third step we used that the matrix $`\tau _j^2`$ commutes with $`\tau _i`$. Since the $`\widehat{p}_j`$ are not mutually commuting the exponential in (3.44) cannot be defined by the spectral theorem, however, we can define it through Nelson’s analytic vector theorem. Thus, we find precisely the quantization of the polar decomposition (3.37) up to a factor of $`e^{\tau _j^2t/8}`$ which tends to unity linear in $`t0`$ as to be expected. Notice that one obtains the first line of (3.3) from (3.37) if one replaces everywhere $`\{.,.\}`$ by $`[.,.]/(i\mathrm{})`$ and phase space functions by operators which holds, of course, by the very construction of the map $`\widehat{W}_t`$ .
This accomplishes our goal to write $`g_e`$ as a function of the $`h_e,P_e`$ and thus an interpretation of $`g_e`$ is indeed possible. As we will discuss all the properties of the corresponding coherent states in great detail in we will refrain from commenting on them here. As we will see, these states in fact enjoy all the properties (a)–(h) that we wanted them to satisfy. In particular, they are diffeomorphism covariant since, in contrast to , we have simply managed to interprete $`\widehat{C}_\gamma `$ as a function of the diffeomorphism covariant functions $`h_e,P^e`$.
We restrict ourselves here to pointing out that the states constructed there will be mainly discussed at the gauge non-invariant and diffeomorphism non-invariant level only. There are two good reasons for this restriction. First of all, both the gauge group and the diffeomorphism group are represented unitarily on the Hilbert space and thus expectation values of gauge – and diffeomorphism invariant operators are in fact gauge – and diffeomorphism invariant. It follows that no redundant information is produced as far as expectation values are concerned which is enough for semi-classical considerations. Secondly, while the gauge transformations generated by the Hamiltonian constraint are not unitarily represented, what we can do is to investigate whether the infinitesimal dynamics of quantum general relativity as advertized in reduces to that of classical general relatity as $`t0`$. This would give faith into the proposal and as we will see, the answer is indeed affirmative .
More ambitiously, however, one may ask whether it is not possible to work directly at the gauge – and diffeomorphism invariant level. The next two sections outline what can be said about this issue.
Remark :
The reader may wonder what happens with the quantization ambiguity labelled by the Immirzi parameter $`\beta `$ (e.g. ) if one combines the quantum theory with the semi-classical considerations started in this paper. It is easy to see that the ambiguity, expectedly, does not affect the classical limit. To see this, recall that the canonical pair is given by $`A_\beta =\mathrm{\Gamma }+\beta K,E/(\kappa \beta )`$ where $`\mathrm{\Gamma }`$ is the spin connection associated with $`E`$ and $`K`$ is related to the extrinsic curvature. Now, for instance, the area of a surface $`S`$ with normal co-vector $`n_a`$ is given by
$$A(S)=_Sd^2x\sqrt{E_j^aE_j^bn_an_b}=\kappa \beta _Sd^2x\sqrt{\frac{E_j^a}{\kappa \beta }\frac{E_j^b}{\kappa \beta }n_an_b}$$
and the area operator in the theory with label $`\beta `$ will be of the form $`\widehat{A}_\beta (S)=\beta \widehat{A}_1(S)`$ where $`\widehat{A}_1(S)`$ has the standard spectrum of, say . Now the Immirzi parameter also modifies the classicality parameter $`t=\beta \kappa \mathrm{}/a^2`$ and the definition of the momenta $`P_\beta ^e(E)=P_1^e(E)/\beta `$. Consider now a coherent state peaked at $`E`$. In the $`\beta `$-theory the coherent state will then be labelled by $`P_\beta ^e(E)`$ and the expectation value of the area operator, which in terms of $`\widehat{P}_\beta ^e`$ is of the form $`\widehat{A}_\beta (S)=\beta _e\sqrt{\widehat{P}_{\beta j}^e\widehat{P}_{\beta j}^e}`$, will be by construction $`<\widehat{A}_\beta (S)>=\beta _e\sqrt{P_{\beta j}^eP_{\beta j}^e}=_e\sqrt{P_{1j}^eP_{1j}^e}`$, that is, independent of $`\beta `$.
## 4 Coherent States Directly for Gauge Invariant Quantities
There are two possibilities for constructing gauge invariant coherent states. The first possibility consists in group avaraging the gauge-variant coherent states of by means of the group averaging method applied to the gauge group which means quantizing before reducing. Precisely, such states will be constructed as
$$\mathrm{\Psi }_{\gamma ,\stackrel{}{g}}^t(\stackrel{}{h})=_{G^{|V(\gamma )|}}[\underset{vV(\gamma )}{}d\mu _H(u_v)]\psi _{\gamma ,\stackrel{}{g}}^t(\{u_{e(0)}h_eu_{e(1)}^1\}_{eE(\gamma )})$$
(4.1)
where we have assumed that the parameterizations of edges are such that parameter values $`0,1`$ respectively correspond to start and end respectively. An interesting feature of the state (4.1) is that it is separately invariant under gauge transformations of both $`\stackrel{}{h},\stackrel{}{g}`$, a property that is not shared by the diffeomorphism group averaged coherent states of the next section. In order to qualify as a state on the reduced phase space with respect to the Gauss constraint one would have to restrict $`\stackrel{}{g}`$, in addition, to the constraint surface which for the variables $`h_e,P_e`$ is explicitly described in . The properties of the states (4.1) will be studied in some detail in so that we can pass on to the second possibility.
This second approach to gauge invariant coherent states is the following one, consisting in reducing before quantizing :
One directly constructs gauge invariant configuration and momentum operators on the constraint surface of the Gauss constraint which leave the space of cylindrical gauge invariant functions over a given graph invariant. Next, one constructs from those new operators with canonical commutation relations and thus has mapped the problem to that of the construction of coherent states for the quantization of a particle moving in a finite number of dimensions for which a natural answer is given by the usual harmonic oscillator coherent states.
We will now outline this idea in some detail :
First we must determine suitable, independent, gauge invariant configuration and momentum operators on a given graph.
Consider a graph $`\gamma `$ with $`E=|E(\gamma )|`$ edges and $`V=|V(\gamma )|`$ vertices. If the gauge group is $`N`$dimensional then for each vertex we have $`N`$ gauge degrees of freedom which allows us to fix $`NV`$ of the $`NE`$ independent components of the $`E`$ holonomies $`h_e,eE(\gamma )`$. This reveals that the number of physical configuration degrees of freedom associated with a graph $`\gamma `$ is given by $`D(\gamma )=N(EV)`$. (We are considering here generic graphs with only at least four-valent vertices in order to have non-vanishing volume; the formula is not correct for the remaining degenerate graphs, for instance the graph consisting of only a single loop still has $`r`$ degrees of freedom while $`E=V=2`$ with $`r`$ the rank of the group. We also consider only semi-simple Lie groups for definiteness).
Before we construct a suitable set of such $`D`$ configuration observables, let us check that the number of gauge invariant momentum observables also equals $`D=N(EV)`$. A suitable set of gauge invariant quantum operators that can be obtained from electrical field operators alone consists of a maximal set of mutually commuting, gauge invariant operators constructed from the left or right invariant vector fields $`{}_{}{}^{L}X_{e}^{i},^RX_e^i`$ on the various copies of the group associated with the edes $`e`$ of the graph. Such a choice of invariants corresponds to the choice of a “recoupling scheme for the associated angular momenta”. Let us outline this for $`G=SU(2)`$ :
We can construct the $`E`$ Laplacians $`\mathrm{\Delta }_e=(^RX_e^i)^2`$ and for each $`n(v)`$-valent vertex $`v`$ we can construct further mutually commuting $`n(v)3`$ invariants given by the squares of the operators $`(^RX_{e_1}^i)+(^RX_{e_2}^i),(^RX_{e_1}^i)+(^RX_{e_2}^i)+(^RX_{e_3}^i),..,(^RX_{e_1}^i)+..(^RX_{e_{n(v)2}}^i)`$. By gauge invariance $`(^RX_{e_1}^i)+..(^RX_{e_{n(v)}}^i)=0`$ so that $`(^RX_{e_1}^i)+..(^RX_{e_{n(v)1}}^i)=(^RX_{e_{n(v)}}^i)`$ is not another independent quantity. The choice of these recoupling momenta corresponds to the choice of a recoupling scheme. Now notice that each edge is connected to two vertices. Thus the number of recoupling degrees of freedom is given by $`_{vV(\gamma )}(n(v)3)=2E3V`$ which amounts together with the $`E`$ Laplacians to precisely $`D=N(EV)=3(EV)`$ momentum degrees of freedom as well.
In the case of a general group, similar arguments apply.
We now come back to the problem of the construction of quantum observables with canonical commutation relations from the basic holonomy and membrane variables $`h_e(A)`$ and $`P^e(A,E)`$ respectively.
Let us first consider the configuration space operators. Notice that by the Euler relation there are $`L(\gamma )=EV+1`$ generators (based at an arbitrary but fixed vertex $`p`$ of $`\gamma `$) of the homotopy group $`\pi _p(\gamma )`$ of $`\gamma `$. Thus, choosing a set of such generators one can construct $`D`$ independent configuration degrees of freedom by forming $`D`$ traces of holonomies along those loops and their compositions (and products of those if $`r=\text{rank}(G)>1`$). However, one must be careful that the ranges of these traces (of products of holonomies along the various generators) in the set of real numbers do not depend on each other. Let us outline this for $`G=SU(2)`$ :
Choose generators $`\alpha _1,..,\alpha _L`$ of $`\pi (\gamma )`$ and define
$$t_I=\frac{1}{2}tr(h_{\alpha _I}),I=1,..,L$$
(4.2)
and since the $`\alpha _I`$ are independent we have that the $`t_I`$ take independently values in $`[1,1]`$.
Notice that so far we did not capture any information about the unit vectors $`n_I`$ in the representation $`h_{\alpha _I}=t_I1+\tau _jn_I^j\sqrt{1t_I^2}`$. The scalar products $`n_I^in_J^i`$ are certainly gauge invariant but they cannot be all independent. Pick one of the generators, say $`\alpha _1`$, and decompose the $`n_J,J=2,..,L`$ into unit vectors parallel and orthogonal $`b_J,J=2,..,L`$ to $`n_1`$
$$n_J=t_{L+J1}n_1+\sqrt{1t_{L+J1}^2}b_J,J=2,..,L$$
(4.3)
where the parameters $`t_J`$ again take independent values in $`[1,1]`$. We can obtain them in terms of traces as
$$t_{L+J1}=\frac{t_1t_J\frac{1}{2}tr(h_{\alpha _1\alpha _J})}{\sqrt{1t_1^2}\sqrt{1t_J^2}}.$$
(4.4)
Finally, we can also decompose $`b_K,K=3,..,L`$ into unit vectors parallel and orthogonal $`c_K`$ to, say, $`b_2`$
$$b_K=t_{2L+K3}b_2+\sqrt{1t_{2L+K3}^2}c_K$$
(4.5)
where, of course, $`c_K=ϵ_Kc_3,ϵ_K=\pm 1,K=4,..,L`$. Clearly, $`n_1,b_2,c_3`$ form an orthonormal basis in $`\text{ }\mathrm{R}^3`$. In terms of traces again :
$$t_{2L+K3}=\frac{\frac{t_2t_K\frac{1}{2}tr(h_{\alpha _2\alpha _K})}{\sqrt{1t_2^2}\sqrt{1t_K^2}}t_{L+1}t_{L+K1}}{\sqrt{1t_{L+1}^2}\sqrt{1t_{L+K1}^2}}.$$
(4.6)
Similarly, we could also express the $`L3`$ dscrete variables $`ϵ_M,L=4,..,M`$ in terms of traces along the lines given above but we will not display the explicit formulae here. Rather, by means of the following trick we can get rid of them : define $`t_{2L+K3}^{}:=t_{2L+K3}`$ and new parameters $`t_{2L+K3},K=4,..,L`$ by
$$t_{2L+K3}^{}=2t_{2L+K3}^21\text{ and }ϵ_K\sqrt{1(t_{2L+K3}^{})^2}=2\sqrt{1t_{2L+K3}^2}t_{2L+K3}$$
(4.7)
with, again, $`t_{2L+K3}[1,1]`$.
Obviously, the above equations (4.2)-(4.7) define precisely $`3(L1)=3(EV)`$ continuous gauge invariant parameters $`t_I,I=1,..,D`$ with independent range in $`[1,1]`$. The map between the selected traces and these variables is singular but the subset of the space $`[1,1]^D`$ where this map is singular is of Lebesgue measure zero and thus is irrelevant for $`L_2`$ functions. In any case, all traces of loops on $`\gamma `$ can be written as definite functions of the $`D`$ variables $`t_I`$ since any such function is a polynomial in the quantities $`n_I^in_J^i`$ and we just need to substitute (4.3), (4.5).
From now on we will assume that we have constructed precisely $`D(\gamma )`$ independent, gauge invariant configuration variables $`t_I,I=1,,D`$ for every graph $`\gamma `$ with range in $`[1,1]`$ along lines similar as above. This suggests the following strategy :
We would like to map the problem at hand to the problem of $`D`$ uncoupled harmonic oscillators. We achieve this by defining new variables
$$x_I:=\text{arctanh}(t_I)=\frac{1}{2}\mathrm{ln}(\frac{1+t_I}{1+t_I})t_I=\text{tanh}(x_I)$$
(4.8)
which take values in the whole real line. We can now consider the Hilbert space $`_\gamma =L_2(\text{ }\mathrm{R}^D,d^Dx)`$ and construct the usual coherent states associated with the annihilation operators $`\widehat{z}_I=\widehat{x}_I+i\frac{t}{\mathrm{}}\widehat{p}_I`$ where $`\widehat{p}_I=i\mathrm{}/x_I`$ and $`t`$ is a dimensionless parameter.
This is, however, not the end of the story. Namely, in order to interprete these coherent states in terms of the original quantities, we must make the connection with the classical theory. For the configuration variables the interpretation is obvious through the formulae (4.2)-(4.6). For the momentum variables this is less obvious. The way to proceed is to first express the operators $`\widehat{p}_I`$ in terms of right invariant vector fields on functions cylindrical with respect to $`\gamma `$ and then to express the latter in terms of the phase space variables. In order to do that we write
$$_{x_I}=\frac{t_J}{x_I}\frac{\theta _{e,i}}{t_J}_{\theta _{e,i}}$$
(4.9)
where $`\theta _{e,i}=\theta _{e,i}(t_I,t_\mu ),eE(\gamma ),i=1,..,N,\mu =1,..,NED`$ are the $`NE`$ angle parameters which coordinatize the $`E`$ copies of $`G`$ and which we can think of as functions of the $`t_I`$ and remaining gauge degrees of freedom $`t_\mu `$. Now, there exists a map
$$_{\theta _{e,i}}=F_{ij}(\theta _{e,k})(^RX_e^j)$$
(4.10)
which generically (that is, almost everywhere) is also non-singular and which allows us to write (4.9) in the form
$$_{x_I}=F_{I,ei}(\{h_e^{}\}_{e^{}E(\gamma })(^RX_e^j).$$
(4.11)
The final step consists in expressing the right invariant vector fields in terms of electric fields in the form of membrane operators which has been done in section 2.1 where they have been called $`P_j^e`$.
We can then finally think of $`_{x_I}`$ as a definite function of the $`\{\widehat{h}_e,\widehat{P}_i^e\}_{eE(\gamma )}`$ with an obvious classical limit. Of course, the formula (4.11) is far from simple.
Notice that in the course of the construction we have defined a new Hilbert space $`_\gamma =L_2(\text{ }\mathrm{R}^{D(\gamma )},d^{D(\gamma )}x)`$ which, however, is unitarily equivalent to the projection of the kinematical Hilbert space $``$ of section 2 to the space of functions cylindrical over $`\gamma `$ (after integrating out gauge degrees of freedom) which also shows that these Hilbert spaces are cylindrically consistent so that they line up to a big Hilbert space in the projective limit, unitarily equivalent to $``$.
We close this section with a number of comments :
(i)
The advantage of this approach as compared to the one outlined in the previous section is that we are guaranteed to fulfill all the requirements (a)-(h) without going through considerable amount of functional analytic work since we can just copy all the results known from the harmonic oscillator coherent states.
(ii)
A disadvantage is that the coherent states so constructed in terms of the $`x_I`$ are not easily expressed in terms of the gauge invariant spin-network functions in terms of which the spectra of important operators, such as the geometrical ones , are well known.
(iii)
Finally, the reader may ask why we did not work at the gauge non-invariant level to begin with, obtain harmonic oscillator kind of coherent states for the gauge-variant quantities and only then solve the Gauss constraint. While this would simplify the analysis considerably since all the gauge angles $`\theta _e^j`$ could be taken as independent configuration variables and we could relate the conjugate derivative operators much more easily to the right invariant vector fields, unfortunately the gauge invariant subspace of the coherent states constructed from the gauge non-invariant quantities $`\theta _e^j`$ is not explictly known. The only known procedure is to write them in terms of non-gauge invariant spin-network functions and then to keep only the gauge invariant combinations (this can be done alternatively by integrating those states over the gauge degrees of freedom as in (4.1)). However, the coherent states are an infinite superposition of harmonic oscillator eigenstates each of which is an infinite superposition of spin-network states (in the $`L_2`$ sense) because the relation between the $`\theta _{e,I}`$ and the $`h_e`$ are not at all polynomial. Thus, the amount of work to be done to solve the Gauss constraint is considerably larger, if possible at all, than to define gauge invariant coherent states directly.
iv)
Finally, the complications mentioned in ii) of course also apply if one works entirely with gauge variant variables $`\theta _e^j`$ mentioned in iii) without caring about the Gauss constraint, the only simplication as compared to iii) is that the construction of the $`t_I`$ is not necessary.
To summarize, the coherent states defined in sections 3.2, 3.3 may reveal the required properties (a)-(h) less obviously, on the other hand, the operators that appear in applications have a much simpler action on these than on the ones that were constructed in the present section. Thus altogether, at least for analytical purposes the set of states of section 3.3 seems to be preferred.
## 5 Diffeomorphism Invariant Coherent States
Given a coherent state $`\psi _{\gamma ,Z}^t`$ we can group average it with respect to the diffeomorphism constraint and obtain (we discard certain technicalities that come from graph symmetry factors, see , whose notation we follow, for details)
$$\eta _{Diff}\psi _{\gamma ,\stackrel{}{g}}^t=\underset{\lambda ,n}{}e^{t\lambda }T_{\gamma ,\lambda ,n}(\stackrel{}{g})[T_{\gamma ,\lambda ,n}]\text{ and }\eta _{Diff}\mathrm{\Psi }_{\gamma ,\stackrel{}{g}}^t=\underset{\gamma ^{}\gamma }{}\eta _{Diff}\psi _{\gamma ,\stackrel{}{g}}^t$$
(5.1)
where $`[\psi ]`$ denotes the orbit of the state $`\psi `$ under $`Diff(\mathrm{\Sigma })`$, typically
$$[T_{\gamma ,\lambda ,n}]=\underset{\gamma ^{}[\gamma ]}{}T_{\gamma ^{},\lambda ,n}.$$
(5.2)
where $`[\gamma ]`$ denotes the orbit of $`\gamma `$. Here, as in section 3.1 we have written coherent states in terms of eigenfunctions $`T_{\gamma ,\lambda ,n}`$ of a general complexifier with eigenvalue $`\lambda `$ and degeneracy level $`n`$ each of which can be decomposed in terms of spin-network functions with non-trivial dependence on every edge of that graph. This requirement is very important in order for group averaging to be well-defined and thus excludes, in particular, the possibility to average infinite graphs as we will see in , at least not without some kind of renormlization as discussed there, see for a general discussion.
If $`\widehat{C}`$ is diffeomorphism invariant, as it is the case for $`\widehat{v}`$ above then $`\lambda `$ is a diffeomorphism invariant quantity.
Although the state (5.1) is certainly diffeomorphism invariant, being a linear combination of diffeomorphism invariant states, it depends not only on $`[\gamma ]`$ and the equivalence class of complex holonomies under diffeomorphisms $`[\stackrel{}{g}]`$, but explicitly on the representants. In other words, while under diffeomorphisms also $`g_e\phi g_e=g_{\phi ^1(e)}`$, (5.2) is not invariant under mapping $`\stackrel{}{g}`$ by a diffeomorphism in contrast to what happened in (4.1) with respect to the gauge group. This is unsatisfactory because on the diffeomorphism invariant Hilbert space $`_{Diff}`$, which is the Cauchy completion of states of the form $`\eta _{Diff}f,fCyl`$ under the inner product
$$<\eta _{Diff}f,\eta _{Diff}g>_{Diff}:=[\eta _{Diff}f](g)$$
(5.3)
where the latter denotes the application of the distribution $`\eta _{Diff}f`$ to the test function $`g`$, the inner product between diffeomorphism invariant coherent states should depend only on $`[\stackrel{}{g}],[\stackrel{}{g}^{}]`$ and not on the representants. In particular, this leads to the following problem : Suppose that $`(A,E)`$ and $`(A^{},E^{})`$ are diffeomorphic points of the classical phase space and compute from these $`g_e=g_e(A,E)`$ as in section 3.1 or 3.2 and similar for the primed quantities. Then if these quantities differ in the range of $`\gamma `$ then the inner product
$$<\eta _{Diff}\psi _{\gamma ,\stackrel{}{g}}^t,\eta _{Diff}\psi _{\gamma ,\stackrel{}{g}^{}}^t>_{Diff}=\underset{\lambda ,n}{}e^{2t\lambda }\overline{T_{\gamma ,\lambda ,n}(\stackrel{}{g})}T_{\gamma ,\lambda ,n}(\stackrel{}{g}^{})$$
(5.4)
will be small by the very definition of a coherent state, that is, these states are almost orthogonal with respect to $`<.,.>_{Diff}`$. This is certainly not what we want.
The reason for this is, of course, that there are too many of the states $`\eta _{Diff}\mathrm{\Psi }_{\gamma ,Z}^t`$. We should identify all those that are labelled by those $`\stackrel{}{g}^{}`$ which lie in the same equivalence class under diffeomorphisms as $`\stackrel{}{g}`$. This can be done by choosing a representant $`Z_0([Z])`$ in every equivalence class $`[Z]`$ where $`Z`$, as before, stands for phase space points $`(A,E)`$ or an actual complex connection depending on whether we choose coherent states based on otion 2) or 1). Notice that this is not, in general, equivalent to fixing a gauge because choosing a representant is possible also if there does not exist a global gauge fixing condition as it is typically the case in field theories.
One might think that one could alternatively define diffeomorphism invariant coherent states by heat kernel evolution, followed by analytical continuation, of the $`\delta `$ distribution with respect to $`<.,.>_{Diff}`$ given by (notice that $`T_{[\gamma ],\lambda ,n}(A)=T_{[\gamma ],\lambda ,n}([A])`$)
$$\delta _{[\gamma ],[A]}([A^{}]):=\underset{\lambda ,n}{}T_{[\gamma ],\lambda ,n}(A)\overline{T_{[\gamma ],\lambda ,n}(A^{})},$$
(5.5)
however, the resulting state
$$\psi _{[\gamma ],[Z]}^t([A])=\underset{\lambda ,n}{}e^{t\lambda }T_{[\gamma ],\lambda ,n}(Z)\overline{T_{[\gamma ],\lambda ,n}(A)},$$
(5.6)
is no longer normalizable with respect to $`<.,.>_{Diff}`$ so that we are forced to adopt the above strategy.
To summarize, we pick arbitrary but fixed representant functions
$`\gamma _0`$ $`:`$ $`[\mathrm{\Gamma }_0^\omega ]\mathrm{\Gamma }_0^\omega ;[\gamma ]\gamma _0([\gamma ])\text{ and}`$
$`Z_0`$ $`:`$ $`M_{Diff}M;[Z]Z_0([Z])`$ (5.7)
from the sets of equivalence classes under diffeomorphisms of piecewise analytical graphs and from the phase space $`M_{Diff}`$ reduced with respect to the diffeomorphism constraint to the full phase space $`M`$ respectively and we define diffeomorphism invariant coherent states by
$$\mathrm{\Psi }_{[\gamma ],[Z]}^{t,Z_0,\gamma _0}:=\eta _{Diff}\mathrm{\Psi }_{\gamma _0([\gamma ]),Z_0([Z])}^t.$$
(5.8)
The function $`\gamma _0`$ is necessary on top of $`Z_0`$ since a coherent state on $`\gamma `$ depends on $`Z`$ only at $`\gamma `$ and not everywhere. The inner product between these states is given through
$$<\psi _{[\gamma ],[Z]}^{t,Z_0,\gamma _0},\psi _{[\gamma ^{}],[Z^{}]}^{t,Z_0,\gamma _0}>_{Diff}=<\psi _{\gamma _0([\gamma ]),Z_0([Z])}^t,\psi _{\gamma _0([\gamma ^{}]),Z_0[Z^{}]}^t>$$
(5.9)
using the orthogonality of the $`\psi _{\gamma ,Z}^t`$ for different $`\gamma `$. Notice that in the last line we just have the kinematical inner product on $``$. It follows from (5.9) immediately that the diffeomorphism invariant coherent states so defined are localized in the same way as the kinematical ones are. The Ehrenfest properties cannot be verified because we would need a complete set of observables on the Hilbert space $`_{Diff}`$ but it is sufficient to know that these states are peaked on $`[Z]`$ for every $`[\gamma ]`$ in order to make semi-classicl approximations. Moreover, it also follows from (3.18) that the group average of the projective limit $`\mathrm{\Psi }_Z^t`$ coherent state is normalizable with respect to $`<.,.>_{Diff}`$ if and only if $`\mathrm{\Psi }_Z^t`$ is normalizable with respect to $`<.,.>`$.
As we have explicitly indicated in (5.8), the coherent states depend on the representant functions (5). But
$$\eta _{Diff}\mathrm{\Psi }_{\gamma ,Z}^t=\eta _{Diff}\mathrm{\Psi }_{\gamma _0([\gamma ]),\phi _0^{}Z}^t$$
(5.10)
where $`\phi _0(\gamma _0([\gamma ]))=\gamma `$ and $`\phi _0^{}`$ is the action of diffeomorphisms on phase space points $`Z`$. Thus, it is only the relation between $`\gamma _0`$ and $`Z_0`$ which makes a difference (has a dffeomorphism invariant meaning) because in the pair $`\gamma _0^{},Z_0^{}`$ we can always replace $`\gamma _0^{}`$ by $`\gamma _0`$ at the price of changing $`Z_0^{}`$. In other words, if we fix $`\gamma _0`$ once and for all as we can without loss of generality, then our choice of diffeomorphism invariant coherent states is entirely labelled by $`Z_0`$. This choice is to be interpreted as a choice of basis of diffeomorphism invariant coherent states. The inner products between members of different bases have no definite locality properties as we have shown in (5.4). But this is in general true for different sets of coherent states even in systems with only a finite number of degrees of freedom. After all, the requirement of localization does not determine a coherent state uniquely, not even up to unitary equivalence because all that is required is that the inner product between such states is unity if their labels coincide and is “small” otherwise where the notion of smallness depends on the basis.
Thus the dependence of the states on $`Z_0`$ is not a bad but in fact an expected property.
Notice further that some of these diffeomorphism invariant coherent states also lie in the kernel of the Hamiltonian constraint operator defined in : we just have to choose $`[\gamma ]`$ in such a way that the range of the Hamiltonian constraint in the set of linear combinations of spin-network functions cannot contain a spin-network state whose underlying graph lies in the class $`[\gamma ]`$. As shown in , there are an infinite number of such states. This observation may be a starting point for the construction of semiclassical states which lie in the kernel of all three types of constraints : the Gauss-, Diffeomorphism- and Hamiltonian constraint.
## 6 Model for Gauge Invariant Coherent States : Euclidean 2+1 gravity
As we have mentioned in section 3.1, the volume operator qualifies best as a complexifier in $`D=2`$. For Euclidean 2+1 gravity we have $`D=2`$ and $`G=SU(2)`$. The volume operator in two dimensions was derived in . The spectrum of that operator for at most three-valent vertices was also computed analytically there. In this section we focus on a Hilbert space for this theory which has finite linear combinations of spin-network states on at most three-valent graphs as a dense subset because otherwise the spectrum is only known numerically.
There are two cases to consider : either (A) no two of the three edges $`e_1,e_2,e_3`$ meeting at a vertex are co-linear or (B) there is a co-linear pair, say $`e_1,e_2`$ (the third case, that all three edges are co-linear is excluded because the volume would vanish). Let $`\stackrel{}{j}=(j_1,j_2,j_3)`$ be the spins with which the three edges are coloured (for at most three-valent graphs the space of vertex-contractors is one dimensional and thus $`J_v=1`$ is suppressed in what follows; this is also the reason why these spin-network states are eigenstates of the volume operator).
The square of the eigenvalues of the volume operator for a given vertex in the two cases are
$`\lambda _v(\stackrel{}{j})`$ $`=`$ $`{\displaystyle \frac{9}{4}}[2(\mathrm{\Delta }_1\mathrm{\Delta }_2+\mathrm{\Delta }_2\mathrm{\Delta }_3+\mathrm{\Delta }_3\mathrm{\Delta }_1)(\mathrm{\Delta }_1^2+\mathrm{\Delta }_2^2+\mathrm{\Delta }_3^2)]{\displaystyle \frac{1}{2}}[\mathrm{\Delta }_1+\mathrm{\Delta }_2+\mathrm{\Delta }_3]\text{ (A)}`$
$`\lambda _v(\stackrel{}{j})`$ $`=`$ $`[2(\mathrm{\Delta }_1\mathrm{\Delta }_2+\mathrm{\Delta }_2\mathrm{\Delta }_3+\mathrm{\Delta }_3\mathrm{\Delta }_1)(\mathrm{\Delta }_1^2+\mathrm{\Delta }_2^2+\mathrm{\Delta }_3^2)]\mathrm{\Delta }_3\text{ (B)}`$ (6.1)
where $`\mathrm{\Delta }_I=j_I(j_I+1)`$. At first sight it seems that in this case we can eve take $`n=1`$ since the $`\mathrm{\Delta }`$’s appear squared in leading order which would be sufficient to make the series of the coherent state converge. However, this is not the case : For instance we can consider the case that $`j_1=const.`$ and $`j_2\mathrm{}`$. Then, due to gauge invariance $`j_3`$ is of the same order as $`j_2`$ and therefore the leading order of the square bracket in (6) is only $`j_2^2`$. We choose $`n=2`$ in what follows. It is then easy to see that in this case the complex connection is explicitly given by $`Z_a^j=A_a^jif_a^j`$ where
$$f_a^j\frac{V}{\sqrt{det(q)}}ϵ_{jkl}ϵ_{ab}(ϵ_{kmn}ϵ_{cd}E_m^cE_n^d)E_l^b$$
We will for our example analyze only the simplest non-trivial graph, a kink (or double kink) $`\alpha `$ with two edges and one (or two) two-valent vertices. This corresponds to, say $`j_2=0`$, in (6) and $`j:=j_1=j_3`$. Then we obtain the simple eigenvalue $`\lambda _j=\lambda _v^2=\mathrm{\Delta }=j(j+1)`$, in other words, on this graph the volume operator reduces to (two times) the square root of the Laplacian on the copy of the group corresponding to $`h:=h_\alpha ,\alpha =e_1e_3^1`$. A complete orthonormal basis of gauge invariant spin-network functions is given by the characters $`\chi _n(h)=tr(\pi _{n/2}(h)),n=0,1,2,..`$.
On the kink, the coherent state is simply given by
$$\mathrm{\Psi }_{g,\alpha ,t}(h)=\underset{n=0}{\overset{\mathrm{}}{}}e^{\frac{t}{4}\lambda _n}\chi _n(g)\overline{\chi _n(h)}$$
(6.2)
where $`g=h_\alpha (Z),\lambda _n=n(n+2)=4(\lambda _j)_{n=2j}`$. The series (6.2) converges for any $`gSL(2,\text{ }\mathrm{C})`$ as shown in .
We wish to show that the state (6.2) diagonalizes the gauge invariant operator
$$\widehat{T}^{\text{ }\mathrm{C}}:=\widehat{W}_t\widehat{T}(\widehat{W}_T)^1,\widehat{W}_t=e^{t\mathrm{\Delta }},\widehat{T}=tr(\widehat{h}).$$
(6.3)
Denoting $`T_n=tr(h^n),n=0,1,..,T=T_1`$ we notice the identity
$$\chi _n=\{\begin{array}{cc}1+T_2+T_4+..+T_N& \text{ : n even}\\ T_1+T_3+T_5+..+T_N& \text{ : n odd}\end{array}$$
(6.4)
from which follows that $`T\chi _n=\chi _{n+1}+\chi _{n1}`$ by using the $`SU(2)`$ Mandelstam identity $`TT_n=T_{n+1}+T_{n1}`$. It is understood that $`T_1=0`$. Let $`T^{\text{ }\mathrm{C}}=tr(g)`$, then
$`(\widehat{T}^{\text{ }\mathrm{C}}\mathrm{\Psi }_{g,t})(h)=e^{t\mathrm{\Delta }}{\displaystyle \underset{n}{}}\chi _n(g)T\chi _n(h)=e^{t\mathrm{\Delta }}{\displaystyle \underset{n}{}}\chi _n(g)[\chi _{n+1}(h)+\chi _{n1}(h)]`$ (6.5)
$`=`$ $`{\displaystyle \underset{n}{}}\chi _n(g)[e^{t\lambda _{n+1}}\chi _{n+1}(h)+e^{t\lambda _{n1}}\chi _{n1}(h)]`$
$`=`$ $`{\displaystyle \underset{n}{}}[\chi _{n+1}(g)+\chi _{n1}(g)e^{t\lambda _n}\chi _n(h)]=T^{\text{ }\mathrm{C}}\psi _g(h).`$
From this and the general discussion in section 3 it easily follows that $`\widehat{T}^{\text{ }\mathrm{C}}`$ and $`(\widehat{T}^{\text{ }\mathrm{C}})^{}`$ respectively have expectation values $`T^{\text{ }\mathrm{C}}`$ and $`\overline{T^{\text{ }\mathrm{C}}}`$ respectively, moreover, we have the uncertainty relation
$$<(\mathrm{\Delta }\widehat{x})^2><(\mathrm{\Delta }\widehat{y})^2>\frac{|<[\widehat{x},\widehat{y}]>|^2}{4},$$
(6.6)
with $`\widehat{x}=\frac{1}{2}(\widehat{T}^{\text{ }\mathrm{C}}+(\widehat{T}^{\text{ }\mathrm{C}})^{}),\widehat{y}=\frac{1}{2i}(\widehat{T}^{\text{ }\mathrm{C}}(\widehat{T}^{\text{ }\mathrm{C}})^{})`$.
The inner product between two coherent states is given by (we suppress the label $`\alpha `$ in what follows)
$$<\mathrm{\Psi }_g^t,\mathrm{\Psi }_g^{}^t>=\underset{n}{}e^{2t\lambda _n}\overline{\chi _n(g)}\chi _n(g^{})$$
(6.7)
We will now show that the overlap integral
$$I^t(g,g^{}):=\frac{|<\mathrm{\Psi }_g^t,\mathrm{\Psi }_g^{}^t>|^2}{<\mathrm{\Psi }_g^t,\mathrm{\Psi }_g^t><\mathrm{\Psi }_g^{}^t,\mathrm{\Psi }_g^{}^t>}$$
(6.8)
decays exponentially fast with $`|tr(g)tr(g^{\text{ }\mathrm{C}})|`$ as $`t0`$, i.e. in the classical limit $`\mathrm{}0`$. The proof uses the Euler-MacLaurin estimate for the difference between a series and its replacement by an integral which turns out to vanish in our limit $`t0`$.
To begin with, recall that the characters are explicitly given by
$$\chi _n(h)=\frac{\mathrm{sin}((n+1)\varphi )}{\mathrm{sin}(\varphi )}\text{ where }2\mathrm{cos}(\varphi ):=tr(h),\varphi [0,\pi ]$$
(6.9)
for any $`hSU(2)`$. Formula (6.9) is entire analytic in $`\varphi `$ and is readily extended to $`gSL(2,\text{ }\mathrm{C})=SU(2)^{\text{ }\mathrm{C}}`$
$$\chi _n(g)=\frac{\mathrm{sin}((n+1)\theta )}{\mathrm{sin}(\theta )}\text{ where }2\mathrm{cos}(\theta ):=tr(g),\theta =\varphi is,\varphi [0,\pi ],s\text{ }\mathrm{R}.$$
(6.10)
Since the characters are class functions, we can always rotate $`g`$ into a maximal torus of $`SU(2)`$ and so we can think of $`g`$ as given by $`g=\mathrm{exp}(\theta \tau _3)`$. Notice that the Weyl subgroup acts on the torus by $`\theta \theta `$ and indeed (6.9), (6.10) are still invariant under it. We can therefore restrict, without loss of generality to $`s[0,\mathrm{}]`$. In general, $`\theta =\pm \sqrt{(\theta _i)^2},g=\mathrm{exp}(\theta _i\tau _i)`$.
We now compute
$`<\mathrm{\Psi }_{g,t},\mathrm{\Psi }_{g^{},t}>`$ $`=`$ $`{\displaystyle \frac{1}{\mathrm{sin}(\overline{\theta })\mathrm{sin}(\theta ^{})}}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}e^{\frac{t}{2}[(n+1)^21]}\mathrm{sin}((n+1)\overline{\theta })\mathrm{sin}((n+1)\theta ^{})`$ (6.11)
$`=`$ $`{\displaystyle \frac{e^{t/2}}{\mathrm{sin}(\overline{\theta })\mathrm{sin}(\theta ^{})}}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}e^{\frac{t}{2}n^2}\mathrm{sin}(n\overline{\theta })\mathrm{sin}(n\theta ^{})`$
$`=`$ $`{\displaystyle \frac{e^{t/2}}{4\mathrm{sin}(\overline{\theta })\mathrm{sin}(\theta ^{})}}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}e^{\frac{t}{2}n^2}\times `$
$`\times [\mathrm{exp}(in[\overline{\theta }+\theta ^{}])+\mathrm{exp}(in[\overline{\theta }+\theta ^{}])\mathrm{exp}(in[\overline{\theta }\theta ^{}])\mathrm{exp}(in[\overline{\theta }\theta ^{}])]`$
$`=`$ $`{\displaystyle \frac{e^{t/2}}{4\mathrm{sin}(\overline{\theta })\mathrm{sin}(\theta ^{})}}{\displaystyle \underset{nZ\{0\}}{}}e^{\frac{t}{2}n^2}[\mathrm{exp}(in[\overline{\theta }+\theta ^{}])\mathrm{exp}(in[\overline{\theta }\theta ^{}])]`$
$`=`$ $`{\displaystyle \frac{e^{t/2}}{4\mathrm{sin}(\overline{\theta })\mathrm{sin}(\theta ^{})}}{\displaystyle \underset{nZ}{}}e^{\frac{t}{2}n^2}[\mathrm{exp}(in[\overline{\theta }+\theta ^{}])\mathrm{exp}(in[\overline{\theta }\theta ^{}])]`$
where in the last step we have noticed that the term $`n=0`$ vanishes. Let now $`x_n:=\sqrt{t}n,\mathrm{\Delta }x:=x_{n+1}x_n=\sqrt{t}`$, then (6.11) can be written in the form
$$<\mathrm{\Psi }_{g,t},\mathrm{\Psi }_{g^{},t}>=\frac{e^{t/2}}{4\sqrt{t}\mathrm{sin}(\overline{\theta })\mathrm{sin}(\theta ^{})}\underset{nZ}{}\mathrm{\Delta }xe^{\frac{t}{2}x_n^2}[\mathrm{exp}(ix_n\frac{\overline{\theta }+\theta ^{}}{\sqrt{t}})\mathrm{exp}(ix_n\frac{\overline{\theta }\theta ^{}}{\sqrt{t}})]$$
(6.12)
which suggests to replace the sum by a Riemann integral for small $`t`$. It would be literally a Riemann sum if it was not for the explicit $`t`$-dependence of the integrand. Thus, the following expression is only an approximation to (6.12) which becomes exact as $`t0`$
$`i^t(g,g^{})`$ $`=`$ $`{\displaystyle \frac{e^{t/2}}{4\sqrt{t}\mathrm{sin}(\overline{\theta })\mathrm{sin}(\theta ^{})}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑xe^{\frac{t}{2}x^2}[\mathrm{exp}(ix{\displaystyle \frac{\overline{\theta }+\theta ^{}}{\sqrt{t}}})\mathrm{exp}(ix{\displaystyle \frac{\overline{\theta }\theta ^{}}{\sqrt{t}}})]`$ (6.13)
$`=`$ $`{\displaystyle \frac{e^{t/2}}{4\sqrt{2\pi t}\mathrm{sin}(\overline{\theta })\mathrm{sin}(\theta ^{})}}[\mathrm{exp}({\displaystyle \frac{(\overline{\theta }+\theta ^{})^2}{2t}})\mathrm{exp}({\displaystyle \frac{(\overline{\theta }\theta ^{})^2}{2t}})]`$
where we have used a Cauchy integral formula. The overlap integral thus is approximated by
$`\stackrel{~}{I}(g,g^{},t)`$ $`=`$ $`{\displaystyle \frac{|i^t(g,g^{})|^2}{i^t(g,g)i^t(g^{},g^{})}}`$ (6.14)
$`=`$ $`{\displaystyle \frac{|\mathrm{exp}(\frac{(\overline{\theta }+\theta ^{})^2}{2t})\mathrm{exp}(\frac{(\overline{\theta }\theta ^{})^2}{2t})|^2}{[\mathrm{exp}(\frac{(\overline{\theta }+\theta )^2}{2t})\mathrm{exp}(\frac{(\overline{\theta }\theta )^2}{2t})][\mathrm{exp}(\frac{(\overline{\theta }^{}+\theta ^{})^2}{2t})\mathrm{exp}(\frac{(\overline{\theta }^{}\theta ^{})^2}{2t})]}},`$
or when decomposing $`\theta =\varphi +is,\theta ^{}=\varphi ^{}+is^{}`$
$$\stackrel{~}{I}(g,g^{},t)=\frac{|\mathrm{exp}(\frac{([\varphi +\varphi ^{}]i[ss^{}])^2}{2t})\mathrm{exp}(\frac{([\varphi \varphi ^{}]i[s+s^{}])^2}{2t})|^2}{[\mathrm{exp}(2\frac{\varphi ^2}{t})\mathrm{exp}(2\frac{s^2}{t})][\mathrm{exp}(2\frac{(\varphi ^{})^2}{t})\mathrm{exp}(2\frac{(s^{})^2}{t})]}$$
(6.15)
We now multiply numerator and denominator of (6.15) with $`\mathrm{exp}([s^2(s^{})^2+\varphi ^2+(\varphi ^{})^2])/t`$ and obtain
$`\stackrel{~}{I}(g,g^{},t)`$ $`=`$ $`{\displaystyle \frac{|\mathrm{exp}(\frac{[\varphi \varphi ^{}+ss^{}]i[(\varphi +\varphi ^{})(ss^{})]}{t})\mathrm{exp}(\frac{[\varphi \varphi ^{}+ss^{}]+i[(\varphi \varphi ^{})(s+s^{})]}{t})|^2}{4\text{sinh}(\frac{\varphi ^2+s^2}{t})\text{sinh}(\frac{(\varphi ^{})^2+(s^{})^2}{t})}}`$ (6.16)
$`=`$ $`{\displaystyle \frac{\text{cosh}(2\frac{\varphi \varphi ^{}+ss^{}}{t})\text{cos}(2\frac{\varphi s^{}\varphi ^{}s}{t})}{2\text{sinh}(\frac{\varphi ^2+s^2}{t})\text{sinh}(\frac{(\varphi ^{})^2+(s^{})^2}{t})}}.`$
Now, for $`\theta ,\theta ^{}0`$, in the limit $`t0`$ we have (since cos is a bounded function and $`ss^{},\varphi \varphi ^{}0`$)
$$\stackrel{~}{I}(g,g^{},t)\mathrm{exp}(\frac{[\varphi \varphi ^{}]^2+[ss^{}]^2}{t})$$
(6.17)
which is indeed rapidly vanishing as $`t0`$ unless $`\theta =\theta ^{}`$ in which case it equals unity as it should.
If either of $`\theta ,\theta ^{}`$ vanishes, say $`s^{}=\varphi ^{}=\theta ^{}=0`$ then expression (6.16) is of the type $`0/0`$ and we can evaluate it, provided the limit exists, by picking up the leading order terms of numerator and denominator by Cauchy’s formula. It turns out to be sufficient to keep the terms of second order in $`s^{},\varphi ^{}`$. The numerator becomes
$$\frac{1}{2}[(2\frac{\varphi \varphi ^{}+ss^{}}{t})^2+(2\frac{\varphi s^{}\varphi ^{}s}{t})^2]+O((\theta ^{})^3=2\frac{(\varphi ^2+s^2)((\varphi ^{})^2+(s^{})^2)}{t^2}+O((\theta ^{})^3)$$
while the denominator becomes
$$2\text{sinh}(\frac{s^2+\varphi ^2}{t})\frac{(\varphi ^{})^2+(s^{})^2}{t}+..$$
where the dots denote terms of at least fourth order in $`s^{},\varphi ^{}`$. Thus, (6.16) has the well-defined limit
$$\stackrel{~}{I}(g,g^{}=1,t)=\frac{s^2+\varphi ^2}{t\text{sinh}(\frac{s^2+\varphi ^2}{t})}$$
(6.18)
which is again exponentially damped as $`t0`$ unless $`s=\varphi =0`$ in which case it equals unity as it should.
To conclude, the overlap integral $`I^t(g,g^{})_{t0}\stackrel{~}{I}(g,g^{},t)`$ is exponentially damped unless $`T^{\text{ }\mathrm{C}}=(T^{})^{\text{ }\mathrm{C}}`$.
Next, we should also show that the normalized coherent state itself, in both the configuration and the momentum representation, is peaked. These and other issues will be much more systematically analyzed for general graphs in by using the Poisson summation formula.
Acknowledgements
We thank the Institute of Theoretical Physics at Santa Barbara for hospitality where part of this work has been completed. Special thanks to the organizers of the workshop “Classical and Quantum Physics of Strong Gravitational Fields”, Santa Barbara, January – June 1999, for establishing an inspriring research environment at the ITP. We also thank A. Ashtekar, L. Bombelli, R. Gambini, J. Lewandowski, R. Loll and J. Pullin for for general discussions about semi-classical quantum general relativity, A. Ashtekar and L. Bombelli for explaining the idea of random weaves to us and O. Winkler for a careful reading of the manuscript.
This research project was supported in part by the National Science Foundation of the USA under grant PHY94-07194 to the ITP, Santa Barbara.
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# Geometric phases for mixed states in interferometry
## Abstract
We provide a physical prescription based on interferometry for introducing the total phase of a mixed state undergoing unitary evolution, which has been an elusive concept in the past. We define the parallel transport condition that provides a connection-form for obtaining the geometric phase for mixed states. The expression for the geometric phase for mixed state reduces to well known formulas in the pure state case when a system undergoes noncyclic and unitary quantum evolution.
When a pure quantal state undergoes cyclic evolution the system returns to its original state but may acquire a phase factor of purely geometric origin. Though this was realized in the adiabatic context , the nonadiabatic generalization was found in . Based on Pancharatnam’s earlier work, this concept was generalized to noncyclic evolutions of quantum systems . Subsequently, the kinematic approach and gauge potential description of geometric phases for noncyclic and non-Schrödinger evolutions were provided. The adiabatic Berry phase and Hannay angle for open paths were introduced and discussed . The noncyclic geometric phase has been generalized to non-Abelian cases . Applications of geometric phase have been found in molecular dynamics , response function of many-body system , and geometric quantum computation . Noncyclic geometric phase for entangled states has also been studied . In all these developments the geometric phase has been discussed only for pure states. However, in some applications, in particular geometric fault tolerant quantum computation , we are primarily interested in mixed state cases. Uhlmann was probably the first to address the issue of mixed state holonomy, but as a purely mathematical problem . In contrast, here we provide a new formalism of geometric phase for mixed states in the experimental context of quantum interferometry.
The purpose of this Letter is to provide an operationally well defined notion of phase for unitarily evolving mixed quantal states in interferometry, which has been an elusive concept in the past. This phase fulfills two central properties that makes it a natural generalization of the pure case: (i) it gives rise to a linear shift of the interference oscillations produced by a variable $`U(1)`$ phase, and (ii) it reduces to the Pancharatnam connection for pure states. We introduce the notion of parallel transport based on our defintion of total phase. We moreover introduce a concept of geometric phase for unitarily evolving mixed quantal states. This geometric phase reduces to the standard geometric phase for pure states undergoing noncyclic unitary evolution.
Mixed states, phases and interference: Consider a conventional Mach-Zehnder interferometer in which the beam-pair spans a two dimensional Hilbert space $`\stackrel{~}{}=\{|\stackrel{~}{0},|\stackrel{~}{1}\}`$. The state vectors $`|\stackrel{~}{0}`$ and $`|\stackrel{~}{1}`$ can be taken as wave packets that move in two given directions defined by the geometry of the interferometer. In this basis, we may represent mirrors, beam-splitters and relative $`U(1)`$ phase shifts by the unitary operators
$`\stackrel{~}{U}_M`$ $`=`$ $`\left(\begin{array}{cc}\hfill 0& \hfill 1\\ \hfill 1& \hfill 0\end{array}\right),\stackrel{~}{U}_B={\displaystyle \frac{1}{\sqrt{2}}}\left(\begin{array}{cc}\hfill 1& \hfill 1\\ \hfill 1& \hfill 1\end{array}\right),`$ (5)
$`\stackrel{~}{U}(1)`$ $`=`$ $`\left(\begin{array}{cc}e^{i\chi }& \hfill 0\\ 0& \hfill 1\end{array}\right),`$ (8)
respectively. An input pure state $`\stackrel{~}{\rho }_{\text{in}}=|\stackrel{~}{0}\stackrel{~}{0}|`$ of the interferometer transforms into the output state
$`\stackrel{~}{\rho }_{\text{out}}`$ $`=`$ $`\stackrel{~}{U}_B\stackrel{~}{U}_M\stackrel{~}{U}(1)\stackrel{~}{U}_B\stackrel{~}{\rho }_{\text{in}}\stackrel{~}{U}_B^{}\stackrel{~}{U}^{}(1)\stackrel{~}{U}_M^{}\stackrel{~}{U}_B^{}`$ (9)
$`=`$ $`{\displaystyle \frac{1}{2}}\left(\begin{array}{cc}1+\mathrm{cos}\chi & i\mathrm{sin}\chi \\ i\mathrm{sin}\chi & 1\mathrm{cos}\chi \end{array}\right)`$ (12)
that yields the intensity along $`|\stackrel{~}{0}`$ as $`I1+\mathrm{cos}\chi `$. Thus the relative $`U(1)`$ phase $`\chi `$ could be observed in the output signal of the interferometer.
Now assume that the particles carry additional internal degrees of freedom, e.g., spin. This internal spin space $`_\text{i}𝒞^N`$ is spanned by the vectors $`|k`$, $`k=1,2,\mathrm{}N`$, chosen so that the associated density operator is initially diagonal
$$\rho _0=\underset{k}{}w_k|kk|$$
(13)
with $`w_k`$ the classical probability to find a member of the ensemble in the pure state $`|k`$. The density operator could be made to change inside the interferometer
$$\rho _0U_\text{i}\rho _0U_\text{i}^{}$$
(14)
with $`U_\text{i}`$ a unitary transformation acting only on the internal degrees of freedom. Mirrors and beam-splitters are assumed to leave the internal state unchanged so that we may replace $`\stackrel{~}{U}_M`$ and $`\stackrel{~}{U}_B`$ by $`𝐔_M=\stackrel{~}{U}_M1_\text{i}`$ and $`𝐔_B=\stackrel{~}{U}_B1_\text{i}`$, respectively, $`1_\text{i}`$ being the internal unit operator. Furthermore, we introduce the unitary transformation
$$𝐔=\left(\begin{array}{cc}\hfill 0& \hfill 0\\ \hfill 0& \hfill 1\end{array}\right)U_\text{i}+\left(\begin{array}{cc}e^{i\chi }& \hfill 0\\ 0& \hfill 0\end{array}\right)1_\text{i}.$$
(15)
The operators $`𝐔_M`$, $`𝐔_B`$, and $`𝐔`$ act on the full Hilbert space $`\stackrel{~}{}_\text{i}`$. $`𝐔`$ correponds to the application of $`U_\text{i}`$ along the $`|\stackrel{~}{1}`$ path and the $`U(1)`$ phase $`\chi `$ similarly along $`|\stackrel{~}{0}`$. We shall use $`𝐔`$ to generalize the notion of phase to unitarily evolving mixed states.
Let an incoming state given by the density matrix $`\varrho _{\text{in}}=\stackrel{~}{\rho }_{\text{in}}\rho _0=|\stackrel{~}{0}\stackrel{~}{0}|\rho _0`$ be split coherently by a beam-splitter and recombine at a second beam-splitter after being reflected by two mirrors. Suppose that $`𝐔`$ is applied between the first beam-splitter and the mirror pair. The incoming state transforms into the output state
$$\varrho _{\text{out}}=𝐔_B𝐔_M\mathrm{𝐔𝐔}_B\varrho _{\text{in}}𝐔_B^{}𝐔^{}𝐔_M^{}𝐔_B^{}.$$
(16)
Inserting Eqs. (8) and (15) into Eq. (16) yields
$`\varrho _{\text{out}}`$ $`=`$ $`{\displaystyle \frac{1}{4}}[\left(\begin{array}{cc}\hfill 1& \hfill 1\\ \hfill 1& \hfill 1\end{array}\right)U_\text{i}\rho _0U_\text{i}^{}+\left(\begin{array}{cc}\hfill 1& \hfill 1\\ \hfill 1& \hfill 1\end{array}\right)\rho _0`$ (27)
$`+e^{i\chi }\left(\begin{array}{cc}\hfill 1& \hfill 1\\ \hfill 1& \hfill 1\end{array}\right)\rho _0U_\text{i}^{}`$
$`+e^{i\chi }\left(\begin{array}{cc}\hfill 1& \hfill 1\\ \hfill 1& \hfill 1\end{array}\right)U_\text{i}\rho _0].`$
The output intensity along $`|\stackrel{~}{0}`$ is
$`I`$ $``$ $`\text{Tr}\left(U_\text{i}\rho _0U_\text{i}^{}+\rho _0+e^{i\chi }U_\text{i}\rho _0+e^{i\chi }\rho _0U_\text{i}^{}\right)`$ (28)
$``$ $`1+|\text{Tr}\left(U_\text{i}\rho _0\right)|\mathrm{cos}\left[\chi \mathrm{arg}\text{Tr}\left(U_\text{i}\rho _0\right)\right],`$ (29)
where we have used $`\text{Tr}(\rho _0U_\text{i}^{})=\left[\text{Tr}\left(U_\text{i}\rho _0\right)\right]^{}`$.
The important observation from Eq. (29) is that the interference oscillations produced by the variable $`U(1)`$ phase $`\chi `$ is shifted by $`\varphi =\mathrm{arg}\text{Tr}\left(U_\text{i}\rho _0\right)`$ for any internal input state $`\rho _0`$, be it mixed or pure. This phase difference reduces for pure states $`\rho _0=|\psi _0\psi _0|`$ to the Pancharatnam phase difference between $`U_\text{i}|\psi _0`$ and $`|\psi _0`$. These two latter facts are the central properties for $`\varphi `$ being a natural generalization of the pure state phase. Moreover the visibility of the interference pattern is $`\nu =|\text{Tr}\left(U_\text{i}\rho _0\right)|0`$, which reduces to the expected $`\nu =|\psi _0|U_\text{i}|\psi _0|`$ for pure states.
The output intensity in Eq. (29) may be understood as an incoherent weighted average of pure state interference profiles as follows. The state $`k`$ gives rise to the interference profile
$$I_k1+\nu _k\mathrm{cos}\left[\chi \varphi _k\right],$$
(30)
where $`\nu _k=|k|U_\text{i}|k|`$ and $`\varphi _k=\mathrm{arg}k|U_\text{i}k`$. This yields the total output intensity
$$I=\underset{k}{}w_kI_k1+\underset{k}{}w_k\nu _k\mathrm{cos}\left[\chi \varphi _k\right],$$
(31)
which is the incoherent classical average of the above single-state interference profiles weighted by the corresponding probabilities $`w_k`$. Eq. (31) may be written in the desired form $`1+\stackrel{~}{\nu }\mathrm{cos}(\chi \stackrel{~}{\varphi })`$ by making the identifications
$`\stackrel{~}{\varphi }`$ $`=`$ $`\mathrm{arg}\left({\displaystyle \underset{k}{}}w_k\nu _ke^{i\varphi _k}\right)=\mathrm{arg}\text{Tr}\left(U_\text{i}\rho _0\right)=\varphi ,`$ (32)
$`\stackrel{~}{\nu }`$ $`=`$ $`\left|{\displaystyle \underset{k}{}}w_k\nu _ke^{i\varphi _k}\right|=|\text{Tr}\left(U_\text{i}\rho _0\right)|=\nu .`$ (33)
Parallel transport condition and geometric phase: Consider a continuous unitary transformation of the mixed state given by $`\rho (t)=U(t)\rho _0U^{}(t)`$. (From now on, we omit the subscript “i” of $`U`$.) We say that the state of the system $`\rho (t)`$ acquires a phase with respect to $`\rho _0`$ if $`\mathrm{arg}\text{Tr}[U(t)\rho _0]`$ is nonvanishing. Now if we want to parallel transport a mixed state $`\rho (t)`$ along an arbitrary path, then at each instant of time the state must be in-phase with the state at an infinitesimal time. The state at time $`t+dt`$ is related to the state at time $`t`$ as $`\rho (t+dt)=U(t+dt)U^{}(t)\rho (t)U(t)U^{}(t+dt)`$. Therefore, the phase difference between $`\rho (t)`$ and $`\rho (t+dt)`$ is $`\mathrm{arg}\text{Tr}[\rho (t)U(t+dt)U^{}(t)]`$. We can say $`\rho (t)`$ and $`\rho (t+dt)`$ are in phase if $`\text{Tr}[\rho (t)U(t+dt)U^{}(t)]`$ is real and positive. This condition can be regarded as a generalization of Pancharatnam’s connection from pure to mixed states. However, from normalization and Hermiticity of $`\rho (t)`$ it follows that $`\text{Tr}[\rho (t)\dot{U}(t)U^{}(t)]`$ is purely imaginary. Hence the above mixed state generalization of Pancharatnam’s connection can be met only when
$$\text{Tr}[\rho (t)\dot{U}(t)U^{}(t)]=0.$$
(34)
This is the parallel transport condition for mixed states undergoing unitary evolution. On the space of density matrices the above condition can be translated to $`\text{Tr}[\rho dUU^{}]=0`$, where $`d`$ is the exterior derivative on the space of density operators. However, $`\rho (t)`$ determines the $`N\times N`$ matrix $`U(t)`$ ($`N`$ being the dimension of the Hilbert space) up to $`N`$ phase factors, and the single condition Eq. (34) while necessary is not sufficient to determine $`U(t)`$. These $`N`$ phase factors are fixed by the $`N`$ parallel transport conditions
$$k(t)|\dot{U}(t)U^{}(t)|k(t)=0,k=1,2,\mathrm{}N,$$
(35)
where the $`|k(t)`$’s are orthonormal eigenstates of $`\rho (t)`$. These are sufficient to determine the parallel transport operator $`U(t)`$ if we are given a non-degenerate density matrix $`\rho (t)`$.
The parallel transport condition for a mixed state provides us a connection in the space of density operators which can be used to define the geometric phase. Thus a mixed state can acquire pure geometric phase if it undergoes parallel transport along an arbitrary curve. One can check that if we have a pure state density operator $`\rho (t)=|\psi (t)\psi (t)|`$ then the parallel transport condition Eq. (34) reduces to $`\psi (t)|\dot{\psi }(t)=0`$ as has been discussed in which is both necessary and sufficient.
Now we can define a geometric phase for mixed state evolution. Let the state trace out an open unitary curve $`\mathrm{\Gamma }:t[0,\tau ]\rho (t)=U(t)\rho _0U^{}(t)`$ in the space of density operators with “end-points” $`\rho (0)=\rho _0`$ and $`\rho (\tau )`$, where $`U(t)`$ satisfies Eq. (34). The evolution need not be cyclic, i.e. $`\rho (\tau )\rho _0`$. We can naturally assign a geometric phase $`\gamma _g[\mathrm{\Gamma }]`$ to this curve once we notice that the dynamical phase vanishes identically. The dynamical phase is the time integral of the average of Hamiltonian and can be defined as
$`\gamma _d={\displaystyle \frac{1}{\mathrm{}}}{\displaystyle _0^\tau }𝑑t\text{Tr}[\rho (t)H(t)]`$ (36)
$`=i{\displaystyle _0^\tau }𝑑t\text{Tr}[\rho _0U^{}(t)\dot{U}(t)].`$ (37)
Since the density matrix undergoes parallel transport evolution the dynamical phase vanishes identically. Moreover, the parallel transport operator $`U(t)`$ should fulfill the stronger condition Eq. (35). Thus the geometric phase for a mixed state is defined as
$`\gamma _g[\mathrm{\Gamma }]`$ $`=`$ $`\varphi =\mathrm{arg}\text{Tr}[\rho _0U(t)]=\mathrm{arg}\left({\displaystyle \underset{k}{}}w_k\nu _ke^{i\beta _k}\right),`$ (38)
where $`\mathrm{exp}(i\beta _k)`$ are geometric phase factors associated with the individual pure state paths in the given ensemble. The above geometric phase can be given a gauge potential description such that the line integral will give the open path geometric phase for mixed state evolution. Indeed the mixed state holonomy can be expressed as
$`\gamma _g[\mathrm{\Gamma }]`$ $`=`$ $`{\displaystyle 𝑑ti\text{Tr}[\rho _0W^{}(t)\dot{W}(t)]}`$ (39)
$`=`$ $`{\displaystyle _\mathrm{\Gamma }}i\text{Tr}\left[\rho _0W^{}dW\right]={\displaystyle _\mathrm{\Gamma }}𝑑\mathrm{\Omega },`$ (40)
where
$$W(t)=\frac{\text{Tr}[\rho _0U^{}(t)]}{|\text{Tr}[\rho _0U^{}(t)]|}U(t).$$
(41)
and $`U(t)`$ satisfies (35). The quantity $`\mathrm{\Omega }=i\text{Tr}\left[\rho _0W^{}dW\right]`$ can be regarded as a gauge potential on the space of density operators pertaining to the system.
The geometric phase defined above is manifestly gauge invariant, does not depend explicitly on the dynamics but it depends only on the geometry of the open unitary path $`\mathrm{\Gamma }`$ in the space of density operators pertaining to the system. It is also independent of the rate at which the system is transported in the quantum state space. The geometric phase Eq. (40) can also be expressed in terms of an average connection form
$`\gamma _g[\mathrm{\Gamma }]`$ $`=`$ $`{\displaystyle _\mathrm{\Gamma }}{\displaystyle \underset{k}{}}w_ki\chi _k|d\chi _k={\displaystyle _\mathrm{\Gamma }}{\displaystyle \underset{k}{}}w_k\mathrm{\Omega }_k,`$ (42)
where $`\mathrm{\Omega }_k`$ is connection-form and $`|\chi _k(t)=W(t)|k`$ is the “reference-section” for $`k`$th component in the ensemble. To be sure, what we have defined is consistent with known results, we can check that this expression reduces to the standard geometric phase
$`\gamma _g[\mathrm{\Gamma }]`$ $`=`$ $`\mathrm{arg}\psi (0)|\psi (\tau )={\displaystyle _0^\tau }𝑑ti\chi (t)|\dot{\chi }(t)`$ (43)
for a pure state $`\rho (t)=|\psi (t)\psi (t)|`$ when it satisfies parallel transport condition. Here, $`|\chi (t)`$ is a reference state, which gives the generalised connection one-form .
Purification: An alternative approach to the above results is given by lifting the mixed state into a purified state $`|\mathrm{\Psi }`$ by attaching an ancilla. We can imagine that any mixed state can be obtained by tracing out some degrees of freedom of a larger system which was in a pure state
$$|\mathrm{\Psi }=\underset{k}{}\sqrt{w_k}|k_s|k_a,$$
(44)
where $`|k_a`$ is a basis in an auxiliary Hilbert space, describing everything else apart from the spatial and the spin degrees of freedom. The existence of the above purification requires that the dimensionality of the auxiliary Hilbert space is at least as large as that of the internal Hilbert space. If $`|\mathrm{\Psi }`$ is the state at time $`t=0`$ and it is transformed to $`|\mathrm{\Psi }(t)`$ by a local unitary operator $`U(t)=U_s(t)I_a`$ then
$$|\mathrm{\Psi }(t)=\underset{k}{}\sqrt{w_k}U_s(t)|k_s|k_a.$$
(45)
The inner-product of initial and final state
$$\mathrm{\Psi }(0)|\mathrm{\Psi }(t)=\underset{k}{}w_kk|U(t)|k=\text{Tr}(U(t)\rho _0)$$
(46)
gives the full description of the modified interference. Indeed by comparing Eqs. (29) and (46), we see that $`\mathrm{arg}\mathrm{\Psi }(0)|\mathrm{\Psi }(t)`$ is the phase shift and $`|\mathrm{\Psi }(0)|\mathrm{\Psi }(t)|`$ is the visibility of the output intensity obtained in an interferometer.
The parallel transport condition, given by Eq. (34), follows immediately from the pure state case when applied to any purification of $`\rho _0`$, i.e.
$`0`$ $`=`$ $`\mathrm{\Psi }(t)|\dot{\mathrm{\Psi }}(t)={\displaystyle \underset{k}{}}w_kk|U^{}(t)\dot{U}(t)|k`$ (47)
$`=`$ $`\text{Tr}[\rho _0U^{}(t)\dot{U}(t)]=\text{Tr}[\rho (t)\dot{U}(t)U^{}(t)].`$ (48)
Thus a parallel transport of a density operator $`\rho (t)`$ amounts to a parallel transport of any of its purifications.
Example: Consider a qubit (a spin-$`\frac{1}{2}`$ particle) whose density matrix can be written as
$$\rho =\frac{1}{2}(1+r\widehat{𝐫}𝝈),$$
(49)
where $`\widehat{𝐫}`$ is a unit vector and $`r`$ is constant for unitary evolution. The pure states $`r=1`$ define the unit Bloch sphere containing the mixed states $`r<1`$. Suppose that during the time evolution $`\widehat{𝐫}`$ traces out a curve on the Bloch sphere that subtends a geodesically closed solid angle $`\mathrm{\Omega }`$ . The two pure states $`|\pm ;\widehat{𝐫}𝝈`$ acquire noncyclic geometric phase $`\mathrm{\Omega }/2`$ and identical visibility $`\nu _+=\nu _{}\eta `$. Using Eq. (38) we obtain the geometric phase for $`\mathrm{\Gamma }`$
$$\varphi =\gamma _g[\mathrm{\Gamma }]=\mathrm{arctan}\left(r\mathrm{tan}\frac{\mathrm{\Omega }}{2}\right).$$
(50)
The visibility $`\nu =|\text{Tr}\left(U\rho _0\right)|`$ is given by
$`\nu `$ $`=`$ $`\eta \sqrt{\mathrm{cos}^2{\displaystyle \frac{\mathrm{\Omega }}{2}}+r^2\mathrm{sin}^2{\displaystyle \frac{\mathrm{\Omega }}{2}}}.`$ (51)
For cyclic evolution we have $`\eta =1`$ but the mixed state $`\nu 1`$ due to the square root factor on the right-hand side of Eq. (51). Moreover Eqs. (50) and (51) reduce to the usual expressions for pure states $`\varphi =\mathrm{\Omega }/2`$ and $`\nu =\eta `$ by letting $`r=1`$.
In the case of maximally mixed states $`r=0`$ we obtain $`\varphi =\mathrm{arg}\mathrm{cos}(\mathrm{\Omega }/2)`$ and $`\nu =|\mathrm{cos}(\mathrm{\Omega }/2)|`$. Thus the output intensity for such states is
$`I`$ $``$ $`1+|\mathrm{cos}{\displaystyle \frac{\mathrm{\Omega }}{2}}|\mathrm{cos}\left[\chi \mathrm{arg}\mathrm{cos}{\displaystyle \frac{\mathrm{\Omega }}{2}}\right]`$ (52)
$`=`$ $`1+\mathrm{cos}{\displaystyle \frac{\mathrm{\Omega }}{2}}\mathrm{cos}\chi .`$ (53)
Early experiments to test the $`4\pi `$ symmetry of spinors utilized unpolarized neutrons. Eq. (53) show that in these experiments the sign change for $`\mathrm{\Omega }=2\pi `$ is a consequence of the phase shift $`\varphi =\mathrm{arg}\mathrm{cos}\pi =\pi `$.
Note that $`\gamma _g[\mathrm{\Gamma }]`$ in Eq. (50) equals the geodesically closed solid angle on the Poincaré sphere iff $`r=1`$. In the mixed state case the geometric phase factor is weighted average of the solid angles subtended by the two pure state paths on the Bloch sphere.
In conclusion, we have provided a physical prescription based on interferometry for introducing a concept of total phase for mixed states undergoing unitary evolution. We have provided the necessary and sufficient condition for parallel transport of a mixed state and introduced a concept of geometric phase for mixed states when it undergoes parallel transport. This reduces to known formulas for pure state case when the system follows a noncyclic and unitary quantum evolutions. We have also provided a gauge potential for noncyclic evolutions of mixed states whose line integral gives the geometric phase. We hope this will lead to experimental test of geometric phases for mixed states and further generalization of it to nonunitary and nonlinear evolutions.
The work by E.S. was financed by the Swedish Natural Science Research Council (NFR). A.K.P. acknowledges EPSRC for financial support and UK Quantum Computing Network for supporting his visit to Centre for Quantum Computation, Oxford. J.S.A. thanks Y. Aharonov and A. Pines for useful discussions and NSF and ONR grants for financial support.M.E. acknowledges financial support from the European Science Foundation. D.K.L.O. acknowledges financial support from CESG.
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# 1. Introduction
## 1. Introduction
In a recent series of papers Meyer , Kent and Clifton and Kent (to whom we will subsequently refer as MKC) claim to have “nullified” the Kochen-Specker theorem . They infer that “there is no truly compelling argument establishing that non-relativistic quantum mechanics describes classically inexplicable physics” . They suggest that this may have significant consequences for quantum information theory and quantum computing.
The purpose of this paper is to criticize MKC’s conclusions. It is true that MKC have circumvented the *particular kind* of non-classicality which features in the Kochen-Specker theorem. It is also true that in doing so they have significantly deepened our understanding of the conceptual implications of quantum mechanics. However, when it comes to the central question, as to whether non-relativistic quantum mechanics is classically explicable, it appears to us that a closer examination of their models leads to a different conclusion. We will argue that, although MKC have nullified the Kochen-Specker theorem *strictly so-called*, there are other, related propositions which are not nullified. The argument we will give is a development of some of the points made in Mermin’s critique (for other critical comments see Havlicek *et al* , Cabello and Basu *et al* ).
As Mermin points out, MKC’s analysis of finite precision measurements is not entirely adequate. MKC only consider one source of non-ideality: namely, the non-ideality which is due to inaccuracies in the specification of the observables to be measured. In every other respect the measurements they consider<sup>1</sup><sup>1</sup>1In the main part of their argument. The part of their argument which concerns (in their terminology) “positive operator measurements” will be discussed below (see Section 8). are perfectly ideal (using the word “ideal” in the sense defined in Section 2). Such measurements might be rather better described (following Mermin) as ideal measurements which are not precisely specified. Consequently, the analysis of MKC is incomplete. In order to make it complete one needs to extend the analysis to the case of measurements which are not ideal *in any respect*: not ideal in respect of the target observables, which the apparatus is intended to measure; and not ideal in respect of any other observables either. The purpose of this paper is to present such an extended analysis. In the first part of the paper we give a more comprehensive account of approximate quantum mechanical measurements (based on ideas previously presented in Appleby ). In the second part we apply these results to the MKC models.
It is important to distinguish the specific, technical result proved by Kochen and Specker, and the essential point of their argument. By the “essential point” we mean the proposition that quantum mechanics (whether relativistic or not) is inconsistent with classical conceptions of physical reality.
In the theories of classical physics it was tacitly assumed
1. To each observable quantity characterising a system there corresponds an objective physical quantity, which has a determinate value at every instant.
2. An ideal, perfectly precise measurement gives, *with certainty*, a value which *exactly coincides* with the value which the quantity being measured objectively did possess, immediately before the measurement process was initiated.
Of course, real laboratory measurements are not perfectly precise; and this fact was acknowledged in classical physics, just as it is in quantum physics. Consequently, the above propositions ought to be supplemented:
1. A non-ideal, approximate measurement gives, *with high probability*, a value which is *close* to the value which the quantity being measured objectively did possess, immediately before the measurement process was initiated.
We will refer to these three propositions collectively as the principle of accessible objective values, or the AOV principle for short. Of course, if the AOV principle is not true, it does not necessarily follow that objective values do not exist. However, if the postulated objective values are typically quite different from the values obtained by measurement, then it is difficult to see what is achieved by assuming them. Consequently, failure of the AOV principle can be taken (though need not necessarily be taken) to justify a positivistic view: on the grounds that “a wheel that can be turned though nothing else moves with it, is not part of the mechanism” (as Wittgenstein succinctly put it, in a different context). This was (in essence) the perception which motivated the Copenhagen Interpretation.
The significance of the Kochen-Specker theorem is that it seems to provide a rigorous proof that the AOV principle is inconsistent with the predictions of quantum mechanics. Kochen and Specker show that, if it is possible to make joint, ideal measurements of any set of commuting observables then, in a hidden variables theory, the result of making an ideal measurement of one observable must, in general, depend on which other commuting observables are jointly and ideally measured with it. This property is not consistent with clause 2 of the AOV principle.
The weakness in Kochen and Specker’s argument was identified by MKC, who noted that the observables to be measured cannot be specified with perfect precision. Consequently, in an experiment which is intended to measure one set of commuting observables, the possibility cannot be excluded that what is actually measured is another, slightly different set of commuting observables. MKC use this freedom to construct a hidden variables theory which *does* satisfy clause 2 of the AOV principle. It should be noted that the theory they construct is not strictly equivalent to standard quantum mechanics (because they postulate that an observable can only be measured if it belongs to a particular, proper subset of the set of all self-adjoint operators). Consequently, they have not shown that clause 2 of the AOV principle is consistent with all the predictions of quantum mechanics (*i.e.*, they have not *refuted* the Kochen-Specker theorem). On the other hand, they have shown that this clause is consistent with the predictions of quantum mechanics in so far as these are empirically verifiable \[*i.e.*, they have *nullified* the Kochen-Specker theorem (strictly so-called)\].
However, just as MKC have noted a significant weakness in the argument of Kochen and Specker, so in turn Mermin has noted a significant weakness in theirs. Although MKC allow for the fact that the observables which are actually measured may be slightly and uncontrollably different from the target observables, which the experiment is intended to measure, they nevertheless follow Kochen and Specker in assuming that the observables which are actually measured still are strictly commuting. But, as Mermin points out, MKC’s own assumptions suggest that the observables which are actually measured will almost certainly *not* be strictly commuting. Also (and, as it turns out, closely connected with this point) one may ask: do the MKC models satisfy clause 3 of the AOV principle? After all, if one accepts MKC’s starting point, that perfect precision is practically unattainable, then clause 3 of the AOV principle, relating as it does to approximate measurements, must be the one which really matters. Clause 2, by contrast, relating as it does to ideal measurements—which is to say practically unrealizable measurements—must be regarded as being of negligible importance. Yet clause 2 is the only one on which their argument bears.
In order to appreciate the force of Mermin’s point it will be helpful to specialise to the standard example of a spin-1 particle, with angular momentum $`\widehat{𝐋}`$. Kochen and Specker consider joint measurements of the three commuting projectors $`(𝐞_r\widehat{𝐋})^2`$ for $`r=1,2,3`$ where $`𝐞_r`$ is any orthonormal triad in $`^3`$. A schematic arrangement for performing such a measurement using three separate analyzers is illustrated in Fig. 1.
MKC correctly observe that, in such an arrangement, it would not practically be possible to align the analyzers precisely along the three directions $`𝐞_r`$. In practice one would expect there to be some uncontrollable errors, so that what are actually measured are the projections $`(𝐞_r^{}\widehat{𝐋})^2`$, where the triad $`𝐞_r^{}`$ is close, but not exactly coincident with the triad $`𝐞_r`$. Nevetheless, MKC assume that the triad $`𝐞_r^{}`$ is precisely orthonormal. Yet it seems clear that, given that the errors are random and uncontrollable, and given that the analyzers are separate instruments, one would typically expect there to be some slight departures from strict orthogonality. At the least, there are no evident grounds for assuming the contrary.
MKC assume that the triad $`𝐞_r^{}`$ must be exactly orthonormal because they follow Kochen and Specker in relying on the principle that it is only sets of commuting observables which can jointly be measured with perfect accuracy. However, if one relaxes the requirement that the measurements be perfectly accurate, then this principle is no longer valid. There is now an extensive literature on the subject of joint, inexact measurements of non-commuting observables. To date, the topics which have received most attention are joint measurements of position and momentum , and joint measurements of the components of spin . For recent reviews, and additional references, the reader may consult Busch *et al* , and Leonhardt . It should be stressed that recent advances in the field of quantum optics mean that such measurements can now be realized, in the laboratory.
The purpose of this paper is to investigate the consequences of extending the analysis of MKC, so as to include approximate joint measurements of non-commuting observable.
The paper is in two main parts. The first part, comprising Sections 25, is concerned with the theory of approximate measurements. The discussion in these sections is based on ideas previously presented in Appleby (also see Appleby ). In our earlier papers we were concerned with the two special cases, of approximate joint measurements of position and momentum , and approximate joint measurements of the components of spin . In Sections 2 and 3 we show how the same methods can be used to analyse approximate measurements for any system having a finite dimensional state space (the extension to the case of a system having an infinite dimensional state space is straightforward, but unnecessary for present purposes). In Section 4 we discuss Uffink’s criticisms of the description of approximate joint measurement processes which (unlike the approach taken in this paper) is based on the concept of an unsharp observable. In Section 5 we specialise the discussion to the case of approximate joint measurements of the operators $`(𝐞_r^{}\widehat{𝐋})^2`$ introduced above, where $`\widehat{𝐋}`$ is the angular momentum for a spin-1 system, and where the triad $`𝐞_r^{}`$ is not assumed to be orthonormal.
The account in Sections 25 is somewhat lengthy. This is because there are some subtleties, and serious potential confusions, which make it necessary to give a detailed discussion of the underlying concepts. Uffink (also see Fleming ) has identified some obscurities in the theory of joint measurements of non-commuting observables as it is presented by (for example) Busch *et al* . It so happens that these objections have a direct bearing on the questions addressed in this paper. One of the advantages of our approach is that Uffink’s objections do not apply to it (the other advantage of our approach being that it leads to an improved definition of measurement accuracy). The importance of this fact will appear in Section 8, where we consider the argument of Clifton and Kent which (they claim) “rule\[s\] out falsifications of non-contextual models based on generalized observables, represented by POV measures”.
In Section 6 we apply the concepts and methods developed in Sections 25 to joint measurements of the target observables $`(𝐞_r\widehat{𝐋})^2`$, under circumstances where the triad $`𝐞_r`$ is not precisely specified, so that the observables $`(𝐞_r^{}\widehat{𝐋})^2`$ which are actually measured cannot be assumed to be precisely commuting. We show that, if the errors in the alignments of the vectors $`𝐞_r`$ are statistically independent, then it is possible to prove a modified version of the Kochen-Specker theorem (a Kochen-Specker theorem for approximate measurements, as it might be called), from which it follows that clause 3 of the AOV principle is not satisfied. This argument was originally inspired by some of the points made by Mermin (although it appears to us that our formulation is considerably sharper than Mermin’s: moreover, Mermin does not remark on the need to assume that the errors are independent).
The result proved in Section 6 shows that, if the alignment errors are independent, then the outcome of an approximate measurement must, in general, be strongly dependent on the particular manner in which the measurement is carried out. In other words, the theory must exhibit a kind of contextuality. However, it may be that there are theories of the MKC type for which the errors are not independent. This possibility is discussed in Section 7. We begin by remarking that, although it is conceivable that there exist theories of the type proposed by MKC for which the errors are not independent, and which do satisfy all three clauses of the AOV principle, it would not be straightforward actually to prove that this was the case. The distribution of errors is not a feature of the theory which one is free simply to postulate. In a complete theory it should be a consequence of the detailed dynamics of the interaction between the system, the measuring apparatus, and the environment. The models proposed by MKC are, as they stand, incomplete, since they do not include any dynamical postulate. In order to show that they satisfy clause 3 the AOV principle it would be necessary, first to specify the dynamical evolution of the hidden variables which characterise the interacting system+apparatus+environment composite, and then to work out the distribution of errors which this implies. Moreover, one would need to establish that the assumption of independence fails in just the way that is required for clause 3 of the AOV principle to be satisfied; and one would need to show that this is the case for every possible system, and every possible set of measurements. Such a program would constitute a highly non-trivial theoretical undertaking. We may therefore conclude, in the first place, that it remains an open question, whether there exists a hidden variables theory satisfying clause 3 of the AOV principle. Whether or not this clause *can* be nullified, it has not been nullified *yet*.
In the second place it is to be observed that such a theory, if it could be constructed, would entail the existence of a delicately adjusted collaboration between the ostensibly random fluctuations in the different parts of a composite apparatus. In Section 7 we argue that this would itself represent a kind of contextuality: for it would mean that the fluctuations in each component of a complex apparatus were, in general, intricately and inescapably dependent on the overall experimental context in which that component was employed. In other words, a theory of this kind would not so much eliminate the phenomenon of contextuality, as shift the locus of the contextuality, from the system, onto the fluctuations in the measuring apparatus.
Our overall conclusion consequently is that, no matter how the theoretical postulates are adjusted, some kind of contextuality must appear somewhere.
Finally, in Section 8 we discuss Clifton and Kent’s theorem 2 which is intended to “rule out falsifications of non-contextual models based on generalized observables, represented by POV measures” (also see Kent ). Approximate joint measurements of non-commuting observables are most conveniently described using a POVM, and so it may at first sight seem that Clifton and Kent’s theorem 2 contradicts the result proved in Section 6 of this paper. In fact, this is not the case, as we show in Section 8. The reason is connected with the point made in Section 4: namely, that although the concept of an approximate measurement involves the concept of a POVM, it does not involve the concept of a new kind of “generalized observable”, distinct from the ordinary kind of observable which is represented by a self-adjoint operator. We go on to discuss some other difficulties which arise from the way in which Clifton and Kent use the concept of a generalized observable.
## 2. Approximate Measurements
The purpose of this section and the one following is to give a general characterisation of approximate measurement(s) performed on a system having a finite dimensional state space. The observables being measured may be commuting or non-commuting. Our approach is based on ideas previously presented in Appleby , in connection with approximate joint measurements of position and momentum (also see Appleby ). As discussed in Section 4, our approach differs from the one taken by many other authors in that it makes no use of the concept of an unsharp observable. This will prove relevant in Section 8, when we discuss Clifton and Kent’s theorem 2. The basic physical ideas are described in this section. The mathematical elaboration in terms of POVM’s is described in Section 3.
An approximate measurement is a measurement which is less than perfectly accurate. It follows, that in order properly to characterise an approximate measurement it is necessary first to arrive at a satisfactory, quantum mechanical concept of measurement accuracy. This is the problem to which we now turn. We begin by considering the accuracy of an imperfect measurement of a single observable. We then extend the discussion to the case of simultaneous, imperfect measurements of a set of several different observables (commuting or non-commuting).
The ordinary, intuitive concept of accuracy involves a comparison between the result of the measurement, and the original value which the quantity being measured did take, immediately before the measurement was carried out. In a quantum mechanical context this concept becomes problematic. The reason for this is the very feature of quantum mechanics which the Kochen-Specker theorem was intended to establish: namely, the fact that in quantum mechanics the concept of “the original value of the observable being measured” is not always well-defined. Of course, one is free to make it well-defined, by taking a hidden variables approach. However, this way of arriving at a concept of quantum mechanical accuracy is not satisfactory because, quite apart from the fact that it compels one to favour a hidden variables approach over all the many alternatives, it makes the accuracy strongly dependent on which particular hidden variables theory is adopted. It is arbitrary, in other words. What one wants is a concept of accuracy which is (1) a natural generalization of the classical concept and (2) independent of the way in which the theory is interpreted (so that it is a feature of quantum mechanics *as such*, and not simply a feature of this or that particular interpretation). In the following we will present a solution to this problem.
Let us start with the standard, elementary textbook example of a measurement process. Consider a system, with finite dimensional state space $`_{\mathrm{sy}}`$, and an apparatus, with finite dimensional state space $`_{\mathrm{ap}}`$. Let $`\widehat{A}`$ be a system observable acting on $`_{\mathrm{sy}}`$, and let $`\widehat{\alpha }`$ be a pointer observable acting on $`_{\mathrm{ap}}`$. Suppose that $`\widehat{A}`$ and $`\widehat{\alpha }`$ have the same set of eigenvalues $`\{a\}`$, which for simplicity we will assume to be non-degenerate. Let $`|a_{\mathrm{sy}}`$ be the corresponding eigenvectors of the operator $`\widehat{A}`$, and let $`|a_{\mathrm{ap}}`$ be the eigenvectors of $`\widehat{\alpha }`$. Let $`|\varphi _0_{\mathrm{ap}}`$ be the initial “zeroed” or “ready” state of the apparatus, and let $`_ac_a|a_{\mathrm{sy}}`$ be the initial state of the system. We then obtain an idealised measurement process by postulating that the unitary evolution operator $`\widehat{U}`$ describing the interaction between system and apparatus is such that
$$\widehat{U}\left(\left(\underset{a}{}c_a|a_{\mathrm{sy}}\right)|\varphi _0_{\mathrm{ap}}\right)=\underset{a}{}c_a\left(|a_{\mathrm{sy}}|a_{\mathrm{ap}}\right)$$
(1)
What makes this a measurement is the fact that it establishes a correlation between the system and pointer observables. What makes it ideal is the fact that the correlation is, in a certain sense, perfect. Specifically:
1. The measurement is retrodictively ideal in the sense that, if the system was initially in the eigenstate of $`\widehat{A}`$ with eigenvalue $`a`$, then there is probability 1 that the recorded value of the pointer observable will also be $`a`$. Consequently, if the system was prepared in some unknown eigenstate of $`\widehat{A}`$, the result of the measurement can be used to retrodict, with certainty, which particular eigenstate it was.
2. The measurement is predictively ideal in the sense that, if the pointer observable is recorded as having the value $`a`$ immediately after the measurement, then one can predict, with probability 1, that a second, immediately subsequent retrodictively ideal measurement of $`\widehat{A}`$ will give the same value $`a`$.
It is easily seen that these two properties, of retrodictive and predictive ideality, are independent. That is, there exist unitary evolution operators $`\widehat{U}`$ describing processes which are retrodictively but not predictively ideal; and operators $`\widehat{U}`$ describing processes which are predictively but not retrodictively ideal.
Practically speaking perfection is seldom, if ever attainable. Consequently, one does not expect a real measurement process to be either retrodictively or predictively ideal. A more realistic model of a measurement process is obtained if, instead of Eq. (1), we take the evolution to be described by
$$\widehat{U}\left(\left(\underset{a}{}c_a|a_{\mathrm{sy}}\right)|\varphi _0_{\mathrm{ap}}\right)=\underset{a}{}c_a\left(|a_{\mathrm{sy}}|a_{\mathrm{ap}}\right)+\underset{a,b,d}{}c_aϵ_{a,bd}\left(|b_{\mathrm{sy}}|d_{\mathrm{ap}}\right)$$
(2)
where $`|ϵ_{a,bd}|1`$ for all $`a,b,d`$. Of course, this model does not include all the complications which one might expect to find in a real measurement process. A complete account should allow for the existence of other apparatus degrees of freedom, additional to $`\widehat{\alpha }`$. It should also allow for the interaction with the environment , and for the fact that $`\widehat{A}`$ and $`\widehat{\alpha }`$ may not have exactly the same spectrum. However, the model just indicated has the merit of simplicity, and it will serve to illustrate the essential ideas.
If the coefficients $`ϵ_{a,bd}`$ are sufficiently small, then the process described by Eq. (2) may be regarded as an approximate measurement of $`\widehat{A}`$: for, corresponding to the properties 1 and 2 above, we have
1. The measurement is retrodictively good in the sense that, if the system was prepared in some unknown eigenstate of $`\widehat{A}`$, then the result of the measurement can be used to retrodict, with a high degree of confidence, which particular eigenstate it was.
2. The measurement is predictively good in the sense that, if the pointer observable is recorded as having the value $`a`$ immediately after the measurement, then one can predict, with probability close to 1, that a second, immediately subsequent retrodictively ideal measurement of $`\widehat{A}`$ will give the same value $`a`$.
We next show how it is possible to quantify the degree of accuracy of the measurement. Define
$`\widehat{A}_\mathrm{f}`$ $`=\widehat{U}^{}\widehat{A}\widehat{U}`$
$`\widehat{\alpha }_\mathrm{f}`$ $`=\widehat{U}^{}\widehat{\alpha }\widehat{U}`$
$`\widehat{A}_\mathrm{f}`$, $`\widehat{\alpha }_\mathrm{f}`$ are the final Heisenberg picture observables, defined at the moment the measurement interaction is completed. Let $`\widehat{A}_\mathrm{i}=\widehat{A}`$ denote the initial Heisenberg picture system observable, defined at the moment the measurement interaction begins. Define the retrodictive error operator $`\widehat{ϵ}_\mathrm{i}`$ and predictive error operator $`\widehat{ϵ}_\mathrm{f}`$ by
$`\widehat{ϵ}_\mathrm{i}`$ $`=\widehat{\alpha }_\mathrm{f}\widehat{A}_\mathrm{i}`$ (3)
$`\widehat{ϵ}_\mathrm{f}`$ $`=\widehat{\alpha }_\mathrm{f}\widehat{A}_\mathrm{f}`$ (4)
Let $`𝒮_{\mathrm{sy}}`$ denote the unit sphere $`_{\mathrm{sy}}`$. Following the discussion in Appleby we now define the maximal rms error of retrodiction, $`\mathrm{\Delta }_{\mathrm{ei}}A`$ by
$`\mathrm{\Delta }_{\mathrm{ei}}A`$ $`=\left(\underset{\psi 𝒮_{\mathrm{sy}}}{sup}\left(\psi \varphi _0|\widehat{ϵ}_\mathrm{i}^2|\psi \varphi _0\right)\right)^{\frac{1}{2}}`$ (5)
and the maximal rms error of prediction, $`\mathrm{\Delta }_{\mathrm{ef}}A`$ by
$`\mathrm{\Delta }_{\mathrm{ef}}A`$ $`=\left(\underset{\psi 𝒮_{\mathrm{sy}}}{sup}\left(\psi \varphi _0|\widehat{ϵ}_\mathrm{f}^2|\psi \varphi _0\right)\right)^{\frac{1}{2}}`$ (6)
Of these two quantities the predictive error $`\mathrm{\Delta }_{\mathrm{ef}}A`$ is the easier to interpret because $`\widehat{ϵ}_\mathrm{f}`$ (unlike $`\widehat{ϵ}_\mathrm{i}`$) connects Heisenberg picture observables defined at the same instant of time. Let $`|\psi _{\mathrm{sy}}`$ be the (normalised) initial system state. Then, reverting to the Schrödinger picture,
$$\left(\psi \varphi _0|\widehat{U}^{}(\widehat{\alpha }\widehat{A})^2\widehat{U}|\psi \varphi _0\right)^{\frac{1}{2}}\mathrm{\Delta }_{\mathrm{ef}}A$$
from which we see that, the smaller $`\mathrm{\Delta }_{\mathrm{ef}}A`$, the more closely the result of a second, immediately subsequent, retrodictively ideal measurement of $`\widehat{A}`$ may be expected to approximate the result of the (non-ideal) measurement under discussion. In particular, if $`\mathrm{\Delta }_{\mathrm{ef}}A=0`$, then the measurement is predictively ideal. It is not difficult to see that the condition $`\mathrm{\Delta }_{\mathrm{ef}}A=0`$ is in fact, not only sufficient, but also necessary for the measurement to be predictively ideal. This justifies the interpretation of $`\mathrm{\Delta }_{\mathrm{ef}}A`$ as providing a quantitative indication of the degree of predictive accuracy.
Let us now consider the interpretation of the quantity $`\mathrm{\Delta }_{\mathrm{ei}}A`$. Suppose, to begin with, that the initial system state $`|\psi `$ is an eigenstate of $`\widehat{A}`$ with eigenvalue $`a`$. Then
$$\left(\psi \varphi _0|\widehat{U}^{}(\widehat{\alpha }a)^2\widehat{U}|\psi \varphi _0\right)^{\frac{1}{2}}=\left(\psi \varphi _0|(\widehat{\alpha }_\mathrm{f}a)^2|\psi \varphi _0\right)^{\frac{1}{2}}\mathrm{\Delta }_{\mathrm{ei}}A$$
from which we see that, the smaller $`\mathrm{\Delta }_{\mathrm{ei}}A`$, the more closely the recorded value of the pointer observable may be expected to approximate $`a`$, and the more accurate the measurement is retrodictively. In particular, if $`\mathrm{\Delta }_{\mathrm{ei}}A=0`$, then the measurement is retrodictively ideal. It is not difficult to see that the condition $`\mathrm{\Delta }_{\mathrm{ei}}A=0`$ is in fact both necessary and sufficient for the measurement to be retrodictively ideal.
It is also possible to say something about the result of the measurement in the case when the initial system state $`|\psi `$ is not an eigenstate of $`\widehat{A}`$. Let $`\overline{A}`$ and $`\mathrm{\Delta }A`$ denote the initial state mean and uncertainty:
$`\overline{A}`$ $`=\psi |\widehat{A}|\psi `$
$`\mathrm{\Delta }A`$ $`=\left(\psi |(\widehat{A}\overline{A})^2|\psi \right)^{\frac{1}{2}}`$
Then the spread of measured values about the initial state mean satisfies the inequality
$`\left(\psi \varphi _0|(\widehat{\alpha }_\mathrm{f}\overline{A})^2|\psi \varphi _0\right)^{\frac{1}{2}}`$ $`\left(\psi \varphi _0|(\widehat{\alpha }_\mathrm{f}\widehat{A}_\mathrm{i})^2|\psi \varphi _0\right)^{\frac{1}{2}}`$
$`+\left(\psi \varphi _0|(\widehat{A}_\mathrm{i}\overline{A})^2|\psi \varphi _0\right)^{\frac{1}{2}}`$
$`\mathrm{\Delta }_{\mathrm{ei}}A+\mathrm{\Delta }A`$ (7)
We see from this that there are two components to the spread of measured values. $`\mathrm{\Delta }A`$ represents the intrinsic uncertainty of the initial system state. $`\mathrm{\Delta }_{\mathrm{ei}}A`$ represents an upper bound on the extrinsic uncertainty, attributable to the noise introduced by the measuring procedure.
These considerations justify the interpretation of $`\mathrm{\Delta }_{\mathrm{ei}}A`$ as providing a quantitative indication of the degree of retrodictive accuracy.
Finally we note that the necessary and sufficient condition for the coefficients $`ϵ_{a,bd}`$ in Eq. (2) all to be zero (so that the measurement is completely ideal) is that $`\mathrm{\Delta }_{\mathrm{ei}}A=\mathrm{\Delta }_{\mathrm{ef}}A=0`$.
Let us now consider a joint measurement of several different observables. If the observables are mutually commuting then the above discussion generalises in the obvious way. However, the point which is important for the argument of this paper is that it also generalises, in a manner which is only slightly less obvious, to the case when the observables are *not* mutually commuting. It is true that one cannot make completely ideal joint measurements of a set of non-commuting observables. However, there is nothing to preclude one from making joint measurements which are only approximate.
Let $`\widehat{A}_1,\mathrm{},\widehat{A}_n`$ be the non-commuting observables to be measured, acting on the system state space $`_{\mathrm{sy}}`$. Corresponding to these observables we introduce a set of $`n`$ pointer observables $`\widehat{\alpha }_1,\mathrm{},\widehat{\alpha }_n`$ acting on the apparatus state space $`_{\mathrm{ap}}`$. We take it that the observables $`\widehat{\alpha }_1,\mathrm{},\widehat{\alpha }_n`$, unlike the observables $`\widehat{A}_1,\mathrm{},\widehat{A}_n`$, are mutually commuting. Consequently, their joint eigenvectors constitute an orthonormal basis for $`_{\mathrm{ap}}`$. Let $`|a_1,\mathrm{},a_n_{\mathrm{ap}}`$ denote the joint eigenvector with eigenvalues $`a_1,\mathrm{},a_n`$ (for the sake of simplicity we assume that the eigenstates are non-degenerate). As before, let $`|\varphi _0`$ be the initial apparatus “zeroed” or “ready” state, and let $`\widehat{U}:_{\mathrm{sy}}_{\mathrm{ap}}_{\mathrm{sy}}_{\mathrm{ap}}`$ be the unitary evolution operator describing the measurement interaction.
The fact that the observables $`\widehat{A}_1,\mathrm{},\widehat{A}_n`$ are non-commuting means that we cannot choose a basis for $`_{\mathrm{sy}}`$ which consists of their joint eigenvectors. However, we can choose, for each $`r`$ separately, a basis which consists of eigenvectors just of $`\widehat{A}_r`$. Let $`|a,x_r`$ be such a basis (where $`a`$ denotes the eigenvalue, and the additional index $`x`$ is to allow for possible degeneracies). We may then write
$$\widehat{U}\left(\left(\underset{a,x}{}c_{ax}|a,x_r\right)|\varphi _0\right)=\underset{b,y,d_1,\mathrm{},d_n}{}c_{ax}f_{ax;by;d_1\mathrm{}d_n}^{(r)}\left(|b,y_r|d_1,\mathrm{},d_n_{\mathrm{ap}}\right)$$
for suitable coefficients $`f_{ax;by;d_1\mathrm{}d_n}^{(r)}`$. Suppose that, for all $`r`$, these coefficients have the property that $`f_{ax;by;d_1\mathrm{}d_n}^{(r)}`$ is small except when $`a=b=d_r`$. Then, comparing this equation with Eq. (2), we see that, for each $`r`$, the pointer $`\widehat{\alpha }_r`$ provides an approximate measurement of the system observable $`\widehat{A}_r`$. This situation may appropriately be described by saying that the process provides an approximate joint measurement of the set of observables $`\widehat{A}_1,\mathrm{},\widehat{A}_n`$.
As in the case of approximate measurements of a single observable, we may obtain a quantitative indication of the accuracy by making use of the Heisenberg picture observables $`\widehat{A}_{r\mathrm{f}}=\widehat{U}^{}\widehat{A_r}\widehat{U}`$, $`\widehat{\alpha }_{r\mathrm{f}}=\widehat{U}^{}\widehat{\alpha }_r\widehat{U}`$, $`\widehat{A}_{r\mathrm{i}}=\widehat{A}_r`$. By analogy with Eqs. (3) and (4) define
$`\widehat{ϵ}_{r\mathrm{i}}`$ $`=\widehat{\alpha }_{r\mathrm{f}}\widehat{A}_{r\mathrm{i}}`$
$`\widehat{ϵ}_{r\mathrm{f}}`$ $`=\widehat{\alpha }_{r\mathrm{f}}\widehat{A}_{r\mathrm{f}}`$
We then obtain $`n`$ maximal rms errors of retrodiction
$`\mathrm{\Delta }_{\mathrm{ei}}A_r`$ $`=\left(\underset{\psi 𝒮_{\mathrm{sy}}}{sup}\left(\psi \varphi _0|\widehat{ϵ}_{r\mathrm{i}}^2|\psi \varphi _0\right)\right)^{\frac{1}{2}}`$ (8)
and $`n`$ maximal rms errors of prediction
$`\mathrm{\Delta }_{\mathrm{ef}}A_r`$ $`=\left(\underset{\psi 𝒮_{\mathrm{sy}}}{sup}\left(\psi \varphi _0|\widehat{ϵ}_{r\mathrm{f}}^2|\psi \varphi _0\right)\right)^{\frac{1}{2}}`$ (9)
where $`𝒮_{\mathrm{sy}}`$ denotes the unit sphere in the system state space $`_{\mathrm{sy}}`$, as before.
Concerning the interpretation of the quantities $`\mathrm{\Delta }_{\mathrm{ei}}A_r`$, $`\mathrm{\Delta }_{\mathrm{ef}}A_r`$ the same analysis applies to them as was given for the errors characterising an approximate measurement of a single observable, in the paragraphs following Eqs. (5) and (6). In particular, we have, by analogy with Inequality (7),
$$\left(\psi \varphi _0|(\widehat{\alpha }_{\mathrm{rf}}\overline{A}_r)^2|\psi \varphi _0\right)^{\frac{1}{2}}\mathrm{\Delta }_{\mathrm{ei}}A_r+\mathrm{\Delta }A_r$$
for $`r=1,\mathrm{},n`$, where $`\overline{A}_r`$ denotes the initial state mean, and $`\mathrm{\Delta }A_r`$ denotes the initial state uncertainty, as in Inequality (7).
Even though the $`\widehat{A}_r`$ are non-commuting, it may still happen that there exist states for which the intrinsic uncertainties $`\mathrm{\Delta }A_r`$ are all small. If the retrodictive errors $`\mathrm{\Delta }_{\mathrm{ei}}A_r`$ are also small, then the above inequalities show that there is a high probability that, for each $`r`$, the recorded value of $`\widehat{\alpha }_r`$ will be close to the initial state expectation value $`\overline{A}_r`$—which provides a further illustration of the sense in which the processes under discussion may be regarded as approximate joint measurements. For examples of measurement processes to which these comments apply, see Appleby , and Section 5 below.
## 3. Approximate Measurements: POVM
An approximate measurement is most conveniently analysed in terms of the corresponding POVM (positive operator valued measure). We avoided introducing this concept at the outset because we wished to establish that one can give an adequate theoretical description of approximate measurements whilst remaining wholly within the framework of the conventional theory, as it was presented by Dirac and von Neumann . In particular, we wished to establish that one can introduce the concept of an approximate measurement, without being thereby compelled to introduce any unconventional, unsharp or generalized observables. However, it is certainly true the concept of a POVM represents a powerful mathematical tool. Consequently, having established that it is not anything more than a tool (at least in the present context), it is appropriate to indicate how the maximal rms errors defined in Section 2 can be expressed in terms of this construct.
As in the last section, we consider a measurement of $`n`$ non-commuting observables $`\widehat{A}_1,\mathrm{},\widehat{A}_n`$ acting on the system state space $`_{\mathrm{sy}}`$. The system is coupled to $`n`$ commuting pointer observables $`\widehat{\alpha }_1,\mathrm{},\widehat{\alpha }_n`$ acting on the apparatus state space $`_{\mathrm{ap}}`$. Let $`|a_1,\mathrm{},a_n`$ be the joint eigenvector of $`\widehat{\alpha }_1,\mathrm{},\widehat{\alpha }_n`$ with eigenvalues $`a_1,\mathrm{},a_n`$ (which, for simplicity, we assume to be non-degenerate). Let $`|\varphi _0`$ be the initial apparatus state, and let $`\widehat{U}`$ be the unitary evolution operator describing the measurement interaction. Let $`|m`$ be any orthonormal basis for the system space $`_{\mathrm{sy}}`$. Define, for each $`n`$-tuplet $`a_1,\mathrm{},a_n`$,
$$\widehat{T}_{a_1,\mathrm{},a_n}=\underset{m,m^{}}{}\left(m|a_1,\mathrm{},a_n|\right)\widehat{U}\left(|m^{}|\varphi _0\right)|mm^{}|$$
(10)
Unlike $`\widehat{U}`$, which acts on the product space $`_{\mathrm{sy}}_{\mathrm{ap}}`$, the operators $`\widehat{T}_{a_1,\mathrm{},a_n}`$ act just on the system space $`_{\mathrm{sy}}`$.
Let
$$\widehat{E}_{a_1,\mathrm{},a_n}=\widehat{T}_{a_1,\mathrm{},a_n}^{}\widehat{T}_{a_1,\mathrm{},a_n}^{}$$
(11)
It is easily verified that $`\widehat{E}_{a_1,\mathrm{},a_n}`$ is the POVM describing the measurement outcome. In other words, the probability that, immediately after the measurement, the $`n`$ pointers will be recorded as having the values $`a_1,\mathrm{},a_n`$ is
$$p_{a_1,\mathrm{},a_n}=\psi |\widehat{E}_{a_1,\mathrm{},a_n}|\psi $$
where $`|\psi `$ is the initial state of the system, immediately before the measurement.
It is also convenient to define
$$\widehat{E}_{a_r}^{(r)}=\underset{a_1,\mathrm{},a_{r1},a_{r+1},\mathrm{},a_n}{}\widehat{E}_{a_1,\mathrm{},a_n}$$
where the summation is over every index except for $`a_r`$. $`\widehat{E}_{a_r}^{(r)}`$ is the POVM describing the outcome of the measurement just of $`\widehat{A}_r`$, which is obtained by ignoring the other $`n1`$ pointer readings. Thus, the probability that the $`r^{\mathrm{th}}`$ pointer reading will be $`a_r`$ is given by
$$p_{a_r}^{(r)}=\psi |\widehat{E}_{a_r}^{(r)}|\psi $$
Starting from the definition of Eq. (8) it is not difficult to show that the $`r^{\mathrm{th}}`$ retrodictive error $`\mathrm{\Delta }_{\mathrm{ei}}A_r`$ is given by
$$\mathrm{\Delta }_{\mathrm{ei}}A_r=\left(\underset{\psi 𝒮_{\mathrm{sy}}}{sup}\left(\psi |\underset{a_r}{}(\widehat{A}_ra_r)\widehat{E}_{a_r}^{(r)}(\widehat{A}_ra_r)|\psi \right)\right)^{\frac{1}{2}}$$
or, equivalently,
$$\mathrm{\Delta }_{\mathrm{ei}}A_r=\left(\underset{a_r}{}(\widehat{A}_ra_r)\widehat{E}_{a_r}^{(r)}(\widehat{A}_ra_r)\right)^{\frac{1}{2}}$$
(12)
where $``$ denotes the operator norm. Similarly, the $`r^{\mathrm{th}}`$ predictive error may be expressed
$$\mathrm{\Delta }_{\mathrm{ef}}A_r=\left(\underset{a_1,\mathrm{},a_n}{}\widehat{T}_{a_1,\mathrm{},a_n}^{}(\widehat{A}_ra_r)^2\widehat{T}_{a_1,\mathrm{},a_n}^{}\right)^{\frac{1}{2}}$$
(13)
## 4. Approximate Measurements and “Unsharp Observables”
As we stressed earlier, the approach described in Sections 2 and 3 differs from the approach of many other authors in that we make no use of the concept of an “unsharp observable”, or of what Clifton and Kent refer to as a “generalized observable”. In this section we discuss Uffink’s criticisms (also see Fleming ) of this way of describing joint measurements of non-commuting observables. The discussion will prove relevant in Section 8, where we consider Clifton and Kent’s argument “to rule out falsifications of non-contextual models based on generalized observables, represented by POV measures”.
Historically, work on the application of POVM’s to the theory of measurement has been strongly influenced by the fact that, from a mathematical point of view, the concept of a POVM (positive operator valued measure) is a generalization of the concept of a PVM (projection valued measure). There consequently arose the idea that, since observables of the ordinary, orthodox kind are represented by PVMs, therefore a POVM which is not also a PVM must represent an observable of a different, unorthodox kind.
If one takes such a view, then one has to suppose that what would naturally be regarded as an approximate measurement of (for example) position, is in fact a (non-approximate?) measurement of something else—unorthodox, or generalized, or unsharp position as it might be called. This way of thinking is certainly at variance with our ordinary intuitions. It would, for instance, not normally be argued that a ruler cannot be used to measure length properly so-called, but only generalized length. However, this objection is perhaps not crucial, for one does not expect quantum mechanical concepts necessarily to accord with classical intuition. Nevertheless, there are some pertinent questions regarding the interpretation of generalized observables which need to be answered if the concept is to be acceptable. As Uffink puts it: “one would naturally like to know *what* is being measured in a measurement of an unorthodox observable” (his emphasis).
The difficulty becomes particularly acute when it is approximate joint measurements of non-commuting observables which are in question. Proponents of the concept of an unsharp observable argue that, although (for example) the orthodox position and momentum observables cannot jointly be measured, there exists a different pair of unorthodox, “unsharp” observables which are jointly measurable. As Uffink points out, the problem with this approach is that, rather than solving the original problem (the problem of making a joint measurement of a pair of *orthodox* observables), it merely presents us with a solution to a new, ostensibly quite different problem (the problem of making a joint measurement of a pair of *unorthodox* observables). Advocates of the approach attempt to deal with this problem by arguing that the unorthodox observables which one actually measures are related to (are unsharp versions of) the orthodox observables which one would like to measure in such a way that, by making a (non-approximate?) measurement of the former, one acquires approximate information regarding the latter. However, Uffink has identified some problems with this idea .
It appears to us that the source of the difficulty lies in the concept of an unsharp observable which, at least so far as approximate measurements are concerned, adds a wholly unnecessary level of complication to the problem. In classical physics there is no need to introduce the concept of an “unsharp” quantity, and then attempt to show that, by measuring that, one gains approximate information about the ordinary quantity in which one is really interested. It turns out that there is no need to introduce such intermediate quantities in quantum physics either—as we showed in Sections 2 and 3 where (using ideas previously presented in Appleby ) we described approximate quantum mechanical measurements directly, without any unnecessary detours, as measurements of the self-same (orthodox) observables concerning which approximate information is sought.
As was shown in Section 3, the concept of a POVM plays an important role in our analysis. However, its role is simply that of a powerful mathematical construct, which can be used to describe the outcome of an approximate measurement. It is not taken to represent a new kind of observable, distinct from the observable one is trying approximately to measure.
It should be stressed that the above discussion only applies to approximate measurements of (orthodox) observables. In other contexts we would agree that the orthodox identification of “observable” with “self-adjoint operator” is too restrictive—as appears from the fact that, if this identification is correct, then phase and time are not observables (see, for example, Busch *et al* , Pegg and Barnett , Bužek *et al* , Oppenheim *et al* , Egusquiza and Muga , and references cited therein). It is also clearly true that the concept of a POVM plays an important role in the problem of arriving at a suitably extended concept of a physical observable. We only wish to point out that the question is not straightforward, and that a simple identification of the concept of a POVM with the concept of a generalized observable may be productive of confusion. Some of the pitfalls appear from the discussion in Uffink’s paper. Others will appear from the discussion in Section 8.
## 5. Approximate Joint Measurements of the Projections $`(𝐞_r\widehat{𝐋})^2`$
We now specialise the theory presented in Sections 2 and 3 to the case which will be discussed in the next two sections, of an approximate joint measurement of the projections $`(𝐞_1\widehat{𝐋})^2`$, $`(𝐞_2\widehat{𝐋})^2`$, $`(𝐞_3\widehat{𝐋})^2`$ where $`\widehat{𝐋}`$ is the angular momentum operator for a spin 1 system, and where the unit vectors $`𝐞_1`$, $`𝐞_2`$, $`𝐞_3`$ are approximately, but perhaps not exactly orthogonal.
Let us start by considering an exact measurement of the single projection $`\widehat{P}=(𝐧\widehat{𝐋})^2`$, for an arbitrary unit vector $`𝐧`$. The system state space $`_{\mathrm{sy}}`$ is thus 3-dimensional. To measure $`\widehat{P}`$ we couple the system to a single pointer observable $`\widehat{\alpha }`$ which has the two (non-degenerate) eigenvalues $`0`$ and $`1`$. The apparatus state space $`_{\mathrm{ap}}`$ is thus 2-dimensional. Let $`|0`$, $`|1`$ be the eigenvectors of $`\widehat{\alpha }`$ with eigenvalues $`0`$ and $`1`$ respectively. Let $`\widehat{\sigma }:_{\mathrm{ap}}_{\mathrm{ap}}`$ be the operator defined by
$$\widehat{\sigma }|0=i|1\widehat{\sigma }|1=i|0$$
(14)
Let $`\widehat{U}:_{\mathrm{sy}}_{\mathrm{ap}}_{\mathrm{sy}}_{\mathrm{ap}}`$ be the unitary operator defined by
$$\widehat{U}=\mathrm{exp}\left[i\frac{\pi }{2}\widehat{P}\widehat{\sigma }\right]=(1\widehat{P})+i\widehat{P}\widehat{\sigma }$$
Let the initial apparatus state be $`|\varphi _0=|0`$. Then
$$\widehat{U}\left(|\psi |0\right)=\{\begin{array}{cc}|\psi |0\text{if }\widehat{P}|\psi =0\hfill & \\ |\psi |1\text{if }\widehat{P}|\psi =|\psi \hfill & \end{array}$$
from which we see that $`\widehat{U}`$ describes a completely ideal measurement of $`\widehat{P}`$.
In order to obtain a joint measurement of the three projections $`\widehat{P}_1=(𝐞_1\widehat{𝐋})^2`$, $`\widehat{P}_2=(𝐞_2\widehat{𝐋})^2`$, $`\widehat{P}_3=(𝐞_3\widehat{𝐋})^2`$ we can chain together three ideal measurements of the kind just described, so that $`\widehat{P}_1`$ is measured first, $`\widehat{P}_2`$ second and $`\widehat{P}_3`$ third, as illustrated in Figure 1. We then have three commuting pointer observables $`\widehat{\alpha }_1`$, $`\widehat{\alpha }_2`$, $`\widehat{\alpha }_3`$ acting on the 6-dimensional space $`_{\mathrm{ap}}`$. Let $`|\alpha _1,\alpha _2,\alpha _3`$ be the joint eigenvector of $`\widehat{\alpha }_1`$, $`\widehat{\alpha }_2`$, $`\widehat{\alpha }_3`$ with eigenvalues $`\alpha _1,\alpha _2,\alpha _3`$. The unitary operator describing the measurement interaction is
$$\widehat{U}=\widehat{U}_3\widehat{U}_2\widehat{U}_1=\left((1\widehat{P}_3)+i\widehat{P}_3\widehat{\sigma }_3\right)\left((1\widehat{P}_2)+i\widehat{P}_2\widehat{\sigma }_2\right)\left((1\widehat{P}_1)+i\widehat{P}_1\widehat{\sigma }_1\right)$$
(15)
where the operators $`\widehat{\sigma }_r`$ are defined by
$`\widehat{\sigma }_1|\alpha _1,\alpha _2,\alpha _3`$ $`=(1)^{\overline{\alpha }_1}i|\overline{\alpha }_1,\alpha _2,\alpha _3`$
$`\widehat{\sigma }_2|\alpha _1,\alpha _2,\alpha _3`$ $`=(1)^{\overline{\alpha }_2}i|\alpha _1,\overline{\alpha }_2,\alpha _3`$
$`\widehat{\sigma }_3|\alpha _1,\alpha _2,\alpha _3`$ $`=(1)^{\overline{\alpha }_3}i|\alpha _1,\alpha _2,\overline{\alpha }_3`$
and where we have employed the notation $`\overline{0}=1,\overline{1}=0`$. Referring to Eq. (10) we see that
$`\widehat{T}_{\alpha _1\alpha _2\alpha _3}={\displaystyle \underset{m,m^{}}{}}\left(m|\alpha _1,\alpha _2,\alpha _3|\right)\widehat{U}\left(|m^{}|0,0,0\right)|mm^{}|`$
where $`|m`$ is any orthonormal basis for $`_{\mathrm{sy}}`$. Defining $`\widehat{P}_r^{(0)}=1\widehat{P}_r`$, $`\widehat{P}_r^{(1)}=\widehat{P}_r`$ this becomes
$$\widehat{T}_{\alpha _1\alpha _2\alpha _3}=\widehat{P}_3^{(\alpha _3)}\widehat{P}_2^{(\alpha _2)}\widehat{P}_1^{(\alpha _1)}$$
Using Eq. (11), and the fact that the $`\widehat{P}_3^{(\alpha _3)}`$ are projections, the POVM describing the measurement outcome is
$$\widehat{E}_{\alpha _1\alpha _2\alpha _3}=\widehat{P}_1^{(\alpha _1)}\widehat{P}_2^{(\alpha _2)}\widehat{P}_3^{(\alpha _3)}\widehat{P}_2^{(\alpha _2)}\widehat{P}_1^{(\alpha _1)}$$
We may assume that the basis in $`^3`$ has been chosen in such a way that
$$𝐞_1=\left(\begin{array}{c}1\\ 0\\ 0\end{array}\right)𝐞_2=\left(\begin{array}{c}\mathrm{sin}\psi \\ \mathrm{cos}\psi \\ 0\end{array}\right)𝐞_3=\left(\begin{array}{c}\mathrm{sin}\theta \mathrm{cos}\varphi \\ \mathrm{sin}\theta \mathrm{sin}\varphi \\ \mathrm{cos}\theta \end{array}\right)$$
where the angles $`\psi ,\theta `$ (but not necessarily $`\varphi `$) are small. There is no loss of generality in assuming that the basis is right handed (since we are free to adjust the signs of the $`𝐞_r`$). Making this assumption, and working to second order in $`\theta ,\psi `$, we find (after some rather lengthy algebra)
$`\widehat{E}_{111}`$ $`\psi ^2(1\widehat{L}_2^2)+\theta ^2(1\widehat{L}_3^2)\theta \psi \mathrm{cos}\varphi \{\widehat{L}_2,\widehat{L}_3\}`$
$`\widehat{E}_{110}`$ $`(1\theta ^2)(1\widehat{L}_3^2)+\theta \psi \mathrm{cos}\varphi \{\widehat{L}_2,\widehat{L}_3\}`$
$`\widehat{E}_{101}`$ $`(1\theta ^2\mathrm{sin}^2\varphi \psi ^2)(1\widehat{L}_2^2)`$
$`\widehat{E}_{011}`$ $`(1\theta ^2\mathrm{cos}^2\varphi \psi ^2)(1\widehat{L}_1^2)`$
$`\widehat{E}_{100}`$ $`\theta ^2\mathrm{sin}^2\varphi (1\widehat{L}_2^2)`$
$`\widehat{E}_{010}`$ $`\theta ^2\mathrm{cos}^2\varphi (1\widehat{L}_1^2)`$
$`\widehat{E}_{001}`$ $`\psi ^2(1\widehat{L}_1^2)`$
$`\widehat{E}_{000}`$ $`0`$
where $`\{\widehat{L}_2,\widehat{L}_3\}`$ denotes the anti-commutator.
$`<\widehat{E}_{\alpha _1\alpha _2\alpha _3}>`$ is the probability that the measurements of $`\widehat{P}_1`$, $`\widehat{P}_2`$, $`\widehat{P}_3`$ will give the values $`\alpha _1`$, $`\alpha _2`$, $`\alpha _3`$ respectively. If $`\psi =\theta =0`$ then $`\widehat{E}_{\alpha _1\alpha _2\alpha _3}`$ reduces to the PVM (projection valued measure)
$$\widehat{E}_{110}=1\widehat{P}_3\widehat{E}_{101}=1\widehat{P}_2\widehat{E}_{011}=1\widehat{P}_1$$
$$\widehat{E}_{111}=\widehat{E}_{100}=\widehat{E}_{010}=\widehat{E}_{001}=\widehat{E}_{001}=0$$
and the probability of the outcome of the measurement being one of the “illegal” combinations $`111`$, $`100`$, $`010`$, $`001`$, $`000`$ is zero. If, however, $`\theta `$, $`\psi `$ are not both $`=0`$, then the probability of obtaining one of these combinations, though small, is not exactly zero—as was to be expected.
Using Eq. (12) we obtain (after some algebra) the following expressions for the retrodictive errors, to lowest order in $`\psi `$, $`\theta `$:
$`\mathrm{\Delta }_{\mathrm{ei}}P_1`$ $`0`$ (16)
$`\mathrm{\Delta }_{\mathrm{ei}}P_2`$ $`|\psi |`$ (17)
$`\mathrm{\Delta }_{\mathrm{ei}}P_3`$ $`|\theta |`$ (18)
For the sake of completeness we also give the formulae for the predictive errors, to lowest order in $`\psi `$, $`\theta `$:
$`\mathrm{\Delta }_{\mathrm{ef}}P_1`$ $`\left(2(\psi ^2+\theta ^2\mathrm{cos}^2\varphi )\right)^{\frac{1}{2}}`$ (19)
$`\mathrm{\Delta }_{\mathrm{ef}}P_2`$ $`|\sqrt{2}\theta \mathrm{sin}\varphi |`$ (20)
$`\mathrm{\Delta }_{\mathrm{ef}}P_3`$ $`0`$ (21)
It is not difficult to see that the equalities $`\mathrm{\Delta }_{\mathrm{ei}}P_1=\mathrm{\Delta }_{\mathrm{ef}}P_3=0`$ are actually exact. This is because the joint measurement is constructed by stringing together a sequence of measurements which are individually ideal. The errors are entirely attributable to the disturbance of the system caused by the successive measurements. The measurement of $`\widehat{P}_1`$ comes first, there has been no preceding measurement to alter the state of the system, and so it is retrodictively ideal: which is why $`\mathrm{\Delta }_{\mathrm{ei}}P_1=0`$. The measurement of $`\widehat{P}_3`$ comes last, there is no subsequent measurement to alter the state of the system, and so it is predictively ideal: which is why $`\mathrm{\Delta }_{\mathrm{ef}}P_3=0`$.
We have assumed that the measurements of the three operators $`\widehat{P}_r`$ are performed sequentially, one after the other, because that is the easiest case to analyse. However, it is perfectly possible to apply these methods to cases where the three measurements are performed all at once (so to speak). For instance, one might consider the evolution described by the Hamiltonian
$$\widehat{H}=\widehat{H}_{\mathrm{sy}}+\widehat{H}_{\mathrm{ap}}+\widehat{H}_{\mathrm{meas}}$$
$`\widehat{H}_{\mathrm{sy}}`$ and $`\widehat{H}_{\mathrm{ap}}`$ are the Hamiltonians describing the free evolution of the system and apparatus respectively. $`\widehat{H}_{\mathrm{meas}}`$ is the time-dependent Hamiltonian describing the measurement interaction, given by
$$\widehat{H}_{\mathrm{meas}}=\frac{h}{4}f(t)\left(\widehat{P}_1\widehat{\sigma }_1+\widehat{P}_2\widehat{\sigma }_2+\widehat{P}_3\widehat{\sigma }_3\right)$$
where $`f`$ is a “bump” function, which is zero outside the short time interval $`[0,\tau ]`$, and which satisfies the normalisation condition $`_0^\tau 𝑑tf(t)=1`$. If $`\tau `$ is sufficiently small then the unitary operator describing the evolution between $`t=0`$ and $`t=\tau `$ is approximately given by
$$\widehat{U}^{}\mathrm{exp}\left[i\frac{\pi }{2}\left(\widehat{P}_1\widehat{\sigma }_1+\widehat{P}_2\widehat{\sigma }_2+\widehat{P}_3\widehat{\sigma }_3\right)\right]$$
If the vectors $`𝐞_r`$ are exactly orthonormal, then $`\widehat{U}^{}`$ coincides with the operator $`\widehat{U}`$ given by Eq. (15). Otherwise it does not. However, it is easily seen that it still describes an approximate joint measurement of the operators $`\widehat{P}_r`$.
Every measurement occupies a finite time interval. There is no difference in principle between joint measurements which are performed sequentially, so that the measurement of each individual observable is allotted its own individual time-slice; and joint measurements which are performed contemporaneously, so that the measurement of each individual observable takes up the whole of the time which is allotted to all.
## 6. A Modified Kochen-Specker Argument
We now apply the concepts and methods developed in the last four sections to the questions posed in the Introduction. We consider a spin 1 system, with angular momentum $`\widehat{𝐋}`$; and we consider the problem of making a joint measurement of the observables $`(𝐞_r\widehat{𝐋})^2`$, for some triad $`𝐞_r`$.
The argument in this section was originally inspired by some of the points made by Mermin . However, it appears to us that our formulation is considerably sharper than Mermin’s. Moreover, Mermin does not remark on the need to assume that the alignment errors are statistically independent (see below).
Before proceeding further, it will be helpful to introduce some terminology. Suppose that an analyzer is designed to measure the observable $`(𝐧\widehat{𝐋})^2`$ but, due to the imprecision in the alignment of the analyzer, does in fact measure the slightly different observable $`(𝐧^{}\widehat{𝐋})^2`$, where $`𝐧`$ and $`𝐧^{}`$ are both unit vectors. Then we will refer to $`(𝐧\widehat{𝐋})^2`$ as the target observable, and to the unit vector $`𝐧`$ as the target alignment; while $`(𝐧^{}\widehat{𝐋})^2`$ will be referred to as the actual observable, and $`𝐧^{}`$ as the actual alignment.
Kochen and Specker make no allowance for the imprecision in any real measurement procedure. They consequently assume that the measurements they consider are all ideal (in the sense explained in Section 2), and they assume that the actual observables exactly coincide with the target observables. They further assume that, in a measurement of the three observables $`(𝐞_r\widehat{𝐋})^2`$, the vectors $`𝐞_r`$ can be any orthonormal triad contained in the real unit 2-sphere, $`S_2`$.
MKC, by contrast, recognise that the imprecision of any real experimental procedure means that the actual alignments $`𝐞_r^{}`$ may be slightly different from the target alignments $`𝐞_r`$. This permits them to make the crucial postulate, that the actual alignments are constrained to lie in a proper, dense subset $`S_2^{}S_2`$. However, MKC continue to assume that the triad $`𝐞_r^{}`$ is precisely orthonormal, and that the measurements of the observables $`(𝐞_r^{}\widehat{𝐋})^2`$ are all ideal. As we discussed in the Introduction, these assumptions are unduly restrictive. The argument of MKC has little force unless it can be extended to the case when the triad $`𝐞_r^{}`$ is not precisely orthonormal and when, in consequence, the measurements of the (non-commuting) observables $`(𝐞_r^{}\widehat{𝐋})^2`$ are not ideal. This is the problem we now address.
The argument which follows does not depend on any assumption regarding the specific manner in which the measurements of the observables $`(𝐞_r^{}\widehat{𝐋})^2`$ are performed. The measurements could be performed sequentially, by three separate analyzers, as illustrated in Fig. 1. However, the argument applies equally well to the case when the measurements are performed contemporaneously, by a single piece of apparatus, as discussed at the end of Section 5.
In order to proceed it is necessary to make a definite hypothesis as to the distribution of actual alignments corresponding to a given target alignment. The most straightforward hypothesis, and the assumption on which the argument of this section will be based, is that the actual alignments are distributed randomly. We will further assume that, in the case of an apparatus which is designed to measure several different target observables, the actual observables are distributed independently. In other words, we assume that, for each possible target alignment $`𝐧`$, there is a probability measure $`\mu _𝐧`$ defined on the set $`S_2^{}`$ (*i.e.*, the set to which MKC postulate that the actual alignment must belong) such that, in a measurement of the target observables $`(𝐞_1\widehat{𝐋})^2`$, $`(𝐞_2\widehat{𝐋})^2`$, $`(𝐞_3\widehat{𝐋})^2`$, the probability that the actual alignments $`𝐞_1^{}`$, $`𝐞_2^{}`$, $`𝐞_3^{}`$ lie in the set $`A_1\times A_2\times A_3S_2^{}\times S_2^{}\times S_2^{}`$ is $`\mu _{𝐞_1}(A_1)\mu _{𝐞_2}(A_2)\mu _{𝐞_3}(A_3)`$.
The probability measure $`\mu _𝐧`$ depends, not only on the vector $`𝐧`$, but also on the construction of the apparatus. An apparatus which was constructed differently, so as to permit the alignments to be fixed more precisely, would have a different associated probability distribution.
The argument of this section is crucially dependent on the assumption that the alignment errors are statistically independent. We will consider the hypothesis that the distributions are not independent in Section 7.
We will denote the set of possible target alignments $`S_2^{\prime \prime }`$. The question arises: what, exactly, is this set? What conditions must a vector satisfy in order to be a possible target alignment? One could argue that *every* vector $`S_2`$ is a possible target alignment. However, it might be thought that this view would be too extreme. The target observable is the observable which the analyzer is designed to measure; and it may reasonably be argued that it is possible to design an instrument to measure some observable if and only if it is possible unambiguously to describe that observable. We will accordingly take the view that $`S_2^{\prime \prime }`$, the set of possible target alignments, consists of those unit vectors $`S_2`$ which are finitely specifiable—that is, which can be specified in standard mathematical notation, by means of a string consisting of finitely many characters.
The set $`S_2^{\prime \prime }`$ so defined is countable, like the set $`S_2^{}`$. However, unlike the set $`S_2^{}`$, the set $`S_2^{\prime \prime }`$ includes Kochen-Specker (KS) uncolourable sets. For instance, it includes the uncolourable set given by Kochen and Specker themselves, and the one given by Peres (the vectors belonging to these sets manifestly are finitely specifiable, since the authors explicitly do so specify them). It follows that $`S_2^{\prime \prime }`$ is itself KS uncolourable.
The fact that $`S_2^{\prime \prime }`$ is KS uncolourable is crucial. It opens the way to a modified version of the KS theorem.
Now consider a valuation $`f:S_2^{}\{0,1\}`$. The fact that $`S_2^{}`$ is countable means that $`f`$ is automatically $`\mu _𝐧`$-measurable, for every $`𝐧S_2^{\prime \prime }`$. Consequently, we may define for each $`𝐧S_2^{\prime \prime }`$,
$$p(𝐧)=\mu _𝐧\left(\{𝐧^{}S_2^{}:f(𝐧^{})=1\}\right)$$
$`p(𝐧)`$ is the probability that, if the analyzer is designed to measure the target observable $`\widehat{P}_𝐧`$, then the vector $`𝐧^{}S_2^{}`$ characterising the actual alignment of the analyzer will have $`f`$-value $`=1`$.
Using the function $`p(𝐧)`$ we next define an induced valuation $`\stackrel{~}{f}:S_2^{\prime \prime }\{0,1\}`$ by
$$\stackrel{~}{f}(𝐧)=\{\begin{array}{cc}0\text{if }p(𝐧)<0.5\hfill & \\ 1\text{if }p(𝐧)0.5\hfill & \end{array}$$
The valuation $`f`$ is defined on the set $`S_2^{}`$, which is KS colourable. However, the induced valuation $`\stackrel{~}{f}`$ is defined on $`S_2^{\prime \prime }`$ which, as we have seen, is KS uncolourable. It follows that there must exist an orthonormal triad $`𝐞_1,𝐞_2,𝐞_3S_2^{\prime \prime }`$ which $`\stackrel{~}{f}`$-evaluates to one of the “illegal” combinations $`111`$, $`100`$, $`010`$, $`001`$, $`000`$.
Let $`𝐞_1,𝐞_2,𝐞_3S_2^{\prime \prime }`$ be such a triad, and suppose that the sequence of three analyzers illustrated in Fig. 1 is used to make a joint measurement of the corresponding target observables, $`(𝐞_1\widehat{𝐋})^2`$, $`(𝐞_2\widehat{𝐋})^2`$, $`(𝐞_3\widehat{𝐋})^2`$. Let $`𝐞_1^{},𝐞_2^{},𝐞_3^{}S_2^{}`$ represent the actual alignments of the analyzers. We will assume that the precision with which the analyzers can be aligned is very high, so that the triad $`𝐞_1^{},𝐞_2^{},𝐞_3^{}`$ is very nearly orthonormal.
It follows from the definitions of $`f`$, $`\stackrel{~}{f}`$ that, for each $`r`$, there is probability $`0.5`$ that $`f(𝐞_r^{})=\stackrel{~}{f}(𝐞_r)`$. Consequently, there is a non-negligible probability that $`𝐞_1^{},𝐞_2^{},𝐞_3^{}`$ $`f`$-evaluates to one of the “illegal” combinations $`111`$, $`100`$, $`010`$, $`001`$, $`000`$ (in fact it is straightforward, though somewhat tedious to show that the probability of obtaining one of these combinations is $`0.5`$). On the other hand, the fact that the triad $`𝐞_1^{},𝐞_2^{},𝐞_3^{}`$ is almost orthonormal, and the results proved in Section 5, together imply that the probability that the result of the measurement will be one of these combinations is $`0`$. It follows that there is non-negligible probability (in fact, probability $`0.5`$) that, for at least one of the observables $`(𝐞_1^{}\widehat{𝐋})^2`$, $`(𝐞_2^{}\widehat{𝐋})^2`$, $`(𝐞_3^{}\widehat{𝐋})^2`$, the measured value is not close to the $`f`$-value.
This establishes that, if the alignment errors are random, and statistically independent, then the model must exhibit a form of contextuality: for it means that the probable outcome of an approximate measurement must, in general, be strongly dependent, not only the observable which is being measured, but also on the particular way in which the measurement is carried out. It follows that, if the stated assumptions are true, then the model fails to satisfy clause 3 of the AOV principle (as stated in the Introduction).
## 7. The Assumption of Independence
It is easily seen that the assumption that the alignment errors are statistically independent is crucial to the argument in the last section. For instance, one can envisage a model in which the errors are correlated in such a way that the actual alignments are always orthogonal. In that case the situation would reduce to the one considered by MKC. It should, however, be noted that this assumption, besides being somewhat implausible, is empirically falsifiable. We have, until now, been following MKC in assuming that the discrepancies between target and actual alignments are a consequence of the limited precision of the measuring device. However, one can equally well consider a case where the alignment “errors” are artificially controlled, using a random number generator, to be much larger than the minimum attainable errors (and consequently measurable). The argument of the last section applies to this case just as well as to the case when the errors are due to the finite precision of the instrument. Consequently, if one wished to avoid the conclusion to that argument in the manner suggested, then one would have to postulate that it is physically impossible to set up the apparatus in such a way that the angles between the alignments are measurably different from $`90^{}`$—which (quite apart from the implausibility of the suggestion) is a definite empirical prediction.
Of course, the fact that, in the MKC models, for each of the values 0 and 1, the set of vectors which are assigned that value constitute a dense subset of $`S_2`$, means that these models are very flexible. Consequently, it may be that there exist other, rather more subtle, postulates regarding the distribution of actual alignments which are not empirically falsifiable, and which do have the property that clause 3 of the AOV principle is then satisfied. However, it would not be very easy to prove that this was the case. In a *complete* hidden variables theory, the probability measure describing the distribution of actual alignments is not a feature which one is free simply to postulate. It has to be derived, from the dynamics of the interacting system+apparatus+environment composite. The models discussed by MKC are incomplete, since they do not include a specification of the dynamics. It is a highly non-trivial question, as to whether there exists a dynamics which, in every situation, gives rise to a probability distribution having the desired properties—not only in situations where the alignment errors arise “naturally”, but also in situations where the errors are adjusted “by hand” (in the manner described in the last paragraph). In the absence of a solution to this problem, the question as to whether there exist hidden variables theories satisfying clause 3 of the AOV principle must be regarded as open.
Suppose, however, that a suitable dynamics could be constructed. Regarded from the perspective of classical intuition the carefully adjusted correlations between the different alignment errors in a complex apparatus which such a model must exhibit would seem very peculiar (they would have something of the flavour of a “conspiracy”). However, what is perhaps rather more to the point is the fact that this phenomenon would itself represent a kind of contextuality: for it would mean that the statistical fluctuations in an analyzer were *ineluctably* dependent on the overall experimental context in which the analyzer was used. A theory of this kind (supposing that it could be constructed) would not so much nullify the contextuality asserted in the conclusion to the Kochen-Specker theorem, as change the locus of the contextuality, from the system, onto the alignment errors.
The conclusion consequently seems to be that, although one may modify its precise form, some kind of contextuality must appear somewhere, in any hidden variables theory.
The problem of trying to nullify the non-classical elements in a hidden variables theory might be compared with the problem of trying to nullify the rucks in a badly fitted carpet. The carpet corresponds to the hidden-variables theory. The nails holding the carpet down correspond to the empirical data. One has a certain amount of freedom to move the rucks around; and with sufficient ingenuity one may succeed in making them less noticeable. However, so long as the nails remain in place, the rucks cannot actually be eliminated.
## 8. Clifton and Kent’s POVM Theorem
In addition to the part of their argument which we have been considering up to now, Clifton and Kent also prove a theorem which, according to them, “rule\[s\] out falsifications of non-contextual models based on generalized observables, represented by POV measures” (also see Kent ). As we have seen, an approximate joint measurement of non-commuting observables is most conveniently described in terms of a POVM. Clifton and Kent’s theorem 2 also concerns measurements whose outcome is described in terms of a POVM. Consequently, it may at first sight seem that their theorem 2 contradicts the result proved in Section 6 of this paper. One purpose of this section is to show that this is not in fact the case. The other purpose is to point out that, although Clifton and Kent’s theorem 2 is valid if regarded as a piece of pure mathematics, its physical significance is more questionable.
Clifton and Kent show, that given any finite dimensional Hilbert space $``$, there exists a set $`𝒜_\mathrm{d}`$ of positive operators acting on $``$, and a truth function $`t_𝒜:𝒜_\mathrm{d}\{0,1\}`$ such that
1. If $`\{A_i\}𝒜_\mathrm{d}`$ is a finite positive operator resolution of the identity (so that $`_iA_i=1`$), then
$$\underset{i}{}t_𝒜(A_i)=1$$
2. the set of finite positive operator resolutions of the identity contained in $`𝒜_\mathrm{d}`$ is a countable, dense subset of the set of all finite positive operator resolutions of the identity (“dense” relative to a topology defined in their paper).
They further argue that the set of all such functions $`t_𝒜`$ is sufficiently large for the theory to be able to recover the statistical predictions of any density operator (in so far as these are testable using finite precision instruments).
In considering this claim we note, first of all, that every projection is also a positive operator. Consequently, the set $`𝒫_\mathrm{d}`$ of admissible projections which features in Clifton and Kent’s theorem 1 should be contained in the set $`𝒜_\mathrm{d}`$ of admissible positive operators which features in their theorem 2. Also, for any given assignment of hidden variables, the truth function $`t_𝒫:𝒫_\mathrm{d}\{0,1\}`$ should be the restriction to $`𝒫_\mathrm{d}`$ of the truth function $`t_𝒜:𝒜_\mathrm{d}\{0,1\}`$. It is not entirely clear from Clifton and Kent’s paper that these requirements can be satisfied. For the sake of argument, we will assume that they are satisfied.
We next address the question, as to what, if any, connection exists between Clifton and Kent’s theorem 2 and the result which we proved in Section 6 of this paper. It is true that both of these results concern measurements whose outcome is described by a POVM (positive operator valued measure) which is not also a PVM (projection valued measure). However, there is a significant difference between the way in which the POVM is interpreted. As we stressed in Section 4, the measurements which we consider are approximate measurements of *ordinary* observables. The role of the POVM is simply to provide a convenient mathematical description of the measurement outcome. By contrast, Clifton and Kent take it that the POVM’s which feature in their theorem (and which are not also PVM’s) represent an entirely new species of *generalized* observable.
Clifton and Kent do not provide any explicit details, as to how these generalized observables are to be measured. However, they appear to assume that the function $`t_𝒜:𝒜_\mathrm{d}\{0,1\}`$ can be specified independently of the function $`t_𝒫:𝒫_\mathrm{d}\{0,1\}`$ (as we noted above, they do not even explicitly impose the requirement that $`𝒫_\mathrm{d}`$ should be contained in $`𝒜_\mathrm{d}`$, and that $`t_𝒫`$ should be the restriction of $`t_𝒜`$). This suggests that they are assuming that, corresponding to the new class of generalized observables, there exists a new class of generalized measurements.
In the approximate measurement procedures which we described in Sections 2, 3 and 5 the final result is obtained by recording the pointer positions. The primary mathematical construct describing the result of the measurement is thus the PVM which gives the distribution of pointer positions. The POVM is a secondary construct which is defined in terms of this PVM \[see Eqs. (10) and (11)\]. Let $`\widehat{E}_{a_1,\mathrm{},a_n}`$ be the element of the POVM which describes the probability that the pointer positions will be $`a_1,\mathrm{},a_n`$. Then the admissibility of $`\widehat{E}_{a_1,\mathrm{},a_n}`$, as an operator describing the outcome of a physically possible approximate measurement process, entirely depends on the admissibility of the corresponding projection operator, acting on the apparatus state space. That is, $`\widehat{E}_{a_1,\mathrm{},a_n}`$ is admissible if and only if the corresponding projection belongs to the set $`𝒫_\mathrm{d}^{\mathrm{ap}}`$, of admissible apparatus projections. Moreover, if $`\widehat{E}_{a_1,\mathrm{},a_n}`$ is admissible, then it should be assigned the same truth value as the corresponding apparatus projection. Clifton and Kent, on the other hand, because they interpret the POVM as the mathematical representation of a completely different kind of observable, assume that they are free to fix the set $`𝒜_\mathrm{d}`$ and truth function $`t_𝒜:𝒜_\mathrm{d}\{0,1\}`$ without having any regard for the apparatus set $`𝒫_\mathrm{d}^{\mathrm{ap}}`$ and truth function $`t_𝒫^{\mathrm{ap}}:𝒫_\mathrm{d}\{0,1\}`$. Consequently, their theorem 2 has no bearing on the result proved in Section 6 of this paper.
Until now we have been concerned with the relationship (or lack of relationship) between Clifton and Kent’s theorem 2 and the result which we proved in Section 6. However, similar considerations also show that there are certain obscurities regarding the significance of Clifton and Kent’s result, even when it is interpreted in their suggested terms. It is, of course, the case that not every POVM arises in the manner discussed in this paper, as a way of describing the outcome of an approximate measurement of an ordinary observable. In other contexts it may be appropriate to think of a POVM as representing a generalized observable. But, whatever the entity is called, one always needs to specify how it is to be measured. The usual answer to this question makes use of the Neumark extension theorem . In the variant of this approach which was proposed by Peres , the system of interest is combined with an ancilla, and an ideal measurement is performed on the composite. The outcome of this measurement is described by a PVM. It follows that, in this case too, once one has fixed the set $`𝒫_\mathrm{d}`$ of admissible projections for the system+ancilla composite, one is no longer free to omit the corresponding system-space positive operators from the set of admissible positive operators.
It can be seen from this that one has less freedom to choose the set $`𝒜_\mathrm{d}`$ than Clifton and Kent assume. The situation regarding the truth values which should be assigned to the members of this set is even more problematic. It is not simply that these values are already partially determined by the function $`t_𝒫`$ describing the pointer observables (in the case of an approximate measurement procedure of the kind discussed in Section 3), or the system+ancilla composite (in the case of Peres’ version of the Neumark construction). It is not even clear that the truth values are determined unambiguously; for it may happen that a given POVM can be physically realized in more than one way. Suppose, for example, that the POVM described in Section 5 is alternatively realized by Peres’ version of the Neumark construction. In both cases, the outcome is fixed once the relevant functions $`t_𝒫`$ are fixed; but it is far from clear that the outcomes will be the same. In short, it is questionable whether it is appropriate to think in terms of there being a well-defined truth function on the set $`𝒜_\mathrm{d}`$.
## 9. Conclusion
In this paper we have made a number of criticisms of the arguments of MKC. This should not be allowed to obscure the fact that their work is, in our view, both interesting and valuable. MKC have taken a question which seemed clear-cut, and shown that it is in fact much more subtle and intricate than had previously been appreciated. They have thereby significantly deepened our understanding of the conceptual implications of quantum mechanics.
Similar qualifications apply to our critical remarks (based on Uffink’s criticisms) concerning the use of the concept of an unsharp observable to describe approximate measurements of ordinary observables. As we stressed at the end of Section 4, these criticisms should not be taken to imply that we question the need for an extended concept of physical observable.
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# Effect of changes in meson properties in a nuclear medium: 𝐽/Ψ dissociation in nuclear matter, and meson-nucleus bound states11footnote 1ADP-00-23/T406, UGI-00-8 Work supported by the Australian Research Council, and the Forschungzentrum Jülich
## 1 Mean-field potentials for mesons and baryons in QMC
This report is based on the quark-meson coupling (QMC) model , which has been successfully applied to many problems in nuclear physics . A detailed description of the Lagrangian density and the mean-field equations of motion are given in Ref. . The Dirac equations for the quarks and antiquarks in hadron bags ($`q=u,\overline{u},d`$ or $`\overline{d}`$, hereafter, and including up to $`s,\overline{s},c,`$ and $`\overline{c}`$) neglecting the Coulomb force, are given by ($`|𝒙|`$ bag radius) :
$`\left[i\gamma _x(m_qV_\sigma ^q)\gamma ^0\left(V_\omega ^q+{\displaystyle \frac{1}{2}}V_\rho ^q\right)\right]\left(\begin{array}{c}\psi _u(x)\\ \psi _{\overline{u}}(x)\end{array}\right)`$ $`=`$ $`0,\left((V_\omega ^q{\displaystyle \frac{1}{2}}V_\rho ^q)\mathrm{for}\left(\begin{array}{c}\psi _d\\ \psi _{\overline{d}}\end{array}\right)\right),`$ (5)
$`\left[i\gamma _xm_{s,c}\right]\psi _{s,c}(x)(\mathrm{or}\psi _{\overline{s},\overline{c}}(x))`$ $`=`$ $`0.`$ (6)
The mean-field potentials for a bag in nuclear matter are defined by $`V_\sigma ^qg_\sigma ^q\sigma `$, $`V_\omega ^q`$ $`g_\omega ^q\omega `$ and $`V_\rho ^qg_\rho ^qb`$, with $`g_\sigma ^q`$, $`g_\omega ^q`$ and $`g_\rho ^q`$ the corresponding quark-meson coupling constants.
The normalized, static solution for the ground state quarks or antiquarks with flavor $`f`$ in the hadron, $`h`$, may be written, $`\psi _f(x)=N_fe^{iϵ_ft/R_h^{}}\psi _f(𝒙)`$, where $`N_f`$ and $`\psi _f(𝒙)`$ are the normalization factor and corresponding spin and spatial part of the wave function. The bag radius in medium, $`R_h^{}`$, will be determined through the stability condition for the mass of the hadron against the variation of the bag radius (see Eq. (8)). The eigenenergies, $`ϵ_f`$, in the wave function in units of $`1/R_h^{}`$, are given by
$$\left(\begin{array}{c}ϵ_u\\ ϵ_{\overline{u}}\end{array}\right)=\mathrm{\Omega }_q^{}\pm R_h^{}\left(V_\omega ^q+\frac{1}{2}V_\rho ^q\right),\left(\begin{array}{c}ϵ_d\\ ϵ_{\overline{d}}\end{array}\right)=\mathrm{\Omega }_q^{}\pm R_h^{}\left(V_\omega ^q\frac{1}{2}V_\rho ^q\right),ϵ_{s,c}=ϵ_{\overline{s},\overline{c}}=\mathrm{\Omega }_{s,c},$$
(7)
where $`\mathrm{\Omega }_q^{}=\sqrt{x_q^2+(R_h^{}m_q^{})^2}`$, with $`m_q^{}=m_qg_\sigma ^q\sigma `$ and $`\mathrm{\Omega }_{s,c}=\sqrt{x_{s,c}^2+(R_h^{}m_{s,c})^2}`$. The hadron masses in a nuclear medium are calculated by
$`m_h^{}`$ $`=`$ $`{\displaystyle \frac{(n_q+n_{\overline{q}})\mathrm{\Omega }_q^{}+(n_{s,c}+n_{\overline{s},\overline{c}})\mathrm{\Omega }_{s,c}z_h}{R_h^{}}}+{\displaystyle \frac{4}{3}}\pi R_h^3B,{\displaystyle \frac{m_h^{}}{R_h}}|_{R_h=R_h^{}}=0,`$ (8)
where $`n_q`$ ($`n_{\overline{q}}`$) and $`n_{s,c}`$ ($`n_{\overline{s},\overline{c}}`$) are the lowest mode light quark (antiquark) and strange, charm (antistrange, anticharm) quark numbers in the hadron, $`h`$, respectively, and the $`z_h`$ parametrize the sum of the center-of-mass and gluon fluctuation effects, and are assumed to be independent of density. The parameters are determined in free space to reproduce the corresponding masses. We chose the values, $`m_q`$=5 MeV, $`m_s`$=250 MeV and $`m_c`$=1300 MeV for the current quark masses, and $`R_N=0.8`$ fm for the bag radius of the nucleon in free space. Other input parameters and some of the quantities calculated are given in Refs. . The quark-meson coupling constants, $`g_\sigma ^q`$, $`g_\omega ^q`$ and $`g_\rho ^q`$, are adjusted to fit the nuclear saturation energy and density of symmetric nuclear matter, and the bulk symmetry energy . Exactly the same coupling constants, $`g_\sigma ^q`$, $`g_\omega ^q`$ and $`g_\rho ^q`$, are used for the light quarks in the mesons and hyperons as in the nucleon. However, in studies of the kaon system, we found that it was phenomenologically necessary to increase the strength of the vector coupling to the non-strange quarks in the $`K^+`$ (by a factor of $`1.4^2`$) in order to reproduce the empirically extracted $`K^+`$-nucleus interaction . We assume this also for the $`D`$ and $`\overline{D}`$ mesons . The scalar ($`U_s^h`$) and vector ($`U_v^h`$) potentials felt by the hadrons, $`h`$, in nuclear matter are given by:
$$U_s=m_h^{}m_h,U_v=(n_qn_{\overline{q}})V_\omega ^qI_3V_\rho ^q,(V_\omega ^q1.4^2V_\omega ^q\mathrm{for}K,\overline{K},D,\overline{D}),$$
(9)
where $`I_3`$ is the third component of isospin projection of the hadron, $`h`$. We show in Figs. 2 and 2 some of the calculated mean field potentials.
## 2 $`J/\mathrm{\Psi }`$ dissociation in nuclear matter
There is a great deal of interest in possible signals of Quark-Gluon Plasma (QGP) formation (or precursors to its formation) and $`J/\mathrm{\Psi }`$ suppression is a promising candidate. On the other hand, there may be other mechanisms which produce an increase in $`J/\mathrm{\Psi }`$ absorption in a hot, dense medium. We are particularly interested in the rather exciting suggestion that the charmed mesons, $`D`$, $`\overline{D}`$, $`D^{}`$ and $`\overline{D}^{}`$, should suffer substantial changes in their properties in a nuclear medium (see Fig. 2). This is expected to have a considerable impact on charm production in heavy ion collisions , although the mass of the $`J/\mathrm{\Psi }`$ is expected to change by a tiny amount in nuclear matter within QMC .
The suppression of $`J/\mathrm{\Psi }`$ production observed in relativistic heavy ion collisions, from $`p+A`$ up to central $`S+U`$ collisions, has been relatively well understood. But recent data from $`Pb+Pb`$ collisions shows a considerably stronger $`J/\mathrm{\Psi }`$ suppression . In an attempt to explain this “anomalous” suppression, many authors have studied one of two possible mechanisms, namely hadronic processes and QGP formation .
In the hadronic dissociation scenario the reactions involving the $`J/\mathrm{\Psi }`$, $`\pi +J/\mathrm{\Psi }D^{}+\overline{D},\overline{D}^{}+D`$ and $`\rho +J/\mathrm{\Psi }D+\overline{D}`$, are well known. The absorption of the $`J/\mathrm{\Psi }`$ through these reactions have been found to be important in general and absolutely necessary in order to fit the data on $`J/\mathrm{\Psi }`$ production. (See Refs. and references therein.) $`J/\mathrm{\Psi }`$ dissociation on comovers, combined with the absorption on nucleons, is the main mechanism proposed as an alternative to that of Matsui and Satz , i.e., the dissociation in a QGP.
Within the hadronic scenario the crucial point is the required dissociation strength. In particular, one needs a total cross section for the $`\pi ,\rho +J/\mathrm{\Psi }`$ interaction of around $`1.53`$ mb in order to explain the data in heavy ion simulations . Recent calculations of the reactions, $`\pi +J/\mathrm{\Psi }D+\overline{D}^{},\overline{D}+D^{}`$ and $`\rho +J/\mathrm{\Psi }D+\overline{D}`$, based on $`D`$ exchange, indicate a much lower cross section than this. The main uncertainty in the discussion of the $`J/\mathrm{\Psi }`$ dissociation on a meson gas is associated with the estimates of the $`\pi ,\rho +J/\mathrm{\Psi }`$ cross sections . According to the predictions for the $`\pi +J/\mathrm{\Psi }`$ cross section taking into account the pion kinetic energy and a thermal pion gas with average temperature of 150 MeV, one might conclude that the rate of this process is small independent of the $`\pi +J/\mathrm{\Psi }`$ dissociation model used .
However, this situation changes when the in-medium potentials of the charmed mesons are taken into account, because they lower the $`\pi +J/\mathrm{\Psi }\overline{D}+D^{}`$ and $`\rho +J/\mathrm{\Psi }\overline{D}+D^{}`$ reaction thresholds due to the effect of the vector and scalar potentials felt by the charmed $`D`$, $`D^{}`$ and $`\rho `$ mesons (see Fig. 2). The cross sections calculated for the $`\pi ,\rho +J/\mathrm{\Psi }`$ collisions with the in-medium potentials are shown in Fig. 4. Clearly the $`J/\mathrm{\Psi }`$ absorption cross sections are substantially enhanced for both the $`\pi +J/\mathrm{\Psi }`$ and $`\rho +J/\mathrm{\Psi }`$ reactions, not only because of the downward shift of the reaction threshold, but also because of the in-medium effect on the reaction amplitude. Moreover, now the $`J/\mathrm{\Psi }`$ absorption on comovers becomes both energy and density dependent – a crucial finding given the situation in actual heavy ion collisions. These effects have never been considered before. The absorption cross section has hitherto been taken as a constant.
We found that the thermally averaged, in-medium $`\pi +J/\mathrm{\Psi }`$ and $`\rho +J/\mathrm{\Psi }`$ absorption cross sections, $`\sigma v`$, depend very strongly on the nuclear density . Even for $`p_{J/\mathrm{\Psi }}=0`$, with a pion gas temperature of 120 MeV, which is close to the saturation pion density, the thermally averaged $`J/\mathrm{\Psi }`$ absorption cross section on the pion at $`\rho _B=3\rho _0`$ is about a factor of 7 larger than that at $`\rho _B=0`$ . As for the $`\rho +J/\mathrm{\Psi }`$, the thermally averaged dissociation cross section at $`\rho _B=3\rho _0`$ becomes larger than 1 mb . Thus, the $`J/\mathrm{\Psi }`$ absorption on $`\rho `$ mesons should be also appreciable, although it is expected the $`\rho `$ meson density is only half of the pion density in $`Pb+Pb`$ collisions .
In order to compare our results with the NA38/NA50 data on $`J/\mathrm{\Psi }`$ suppression in $`Pb+Pb`$ collisions, we have adopted the heavy ion model proposed in Ref. with the $`E_T`$ model from Ref. . We introduce the absorption cross section on comovers as function of the density of comovers, while the nuclear absorption cross section is taken as 4.5 mb . Our calculations are shown in Fig. 4 by the solid line, using the density dependent, thermally averaged cross section, $`\sigma v`$, for $`J/\mathrm{\Psi }`$ absorption on comovers . The dashed line in Fig. 4 shows the calculations with the phenomenological constant cross section for $`J/\mathrm{\Psi }`$ absorption on comovers, $`\sigma v`$1 mb given in Ref. . Both curves clearly reproduce the data quite well, including most recent results from NA50 on the ratio of $`J/\mathrm{\Psi }`$ over Drell-Yan cross sections, as a function of the transverse energy up to $`E_T100`$ GeV. It is important to note that if one neglected the in-medium modification of the $`J/\mathrm{\Psi }`$ absorption cross section the large cross section, $`\sigma v`$$``$1 mb, could not be justified by microscopic theoretical calculations and thus the NA50 data could not be described. Furthermore, our calculations with in-medium modified absorption provide a significant improvement in the understanding of the data compared to the models quoted by NA50 .
The basic difference between our results and those quoted by NA50 is that in previous heavy ion calculations the cross section for $`J/\mathrm{\Psi }`$ absorption on comovers was taken as a free parameter to be adjusted to the data and was never motivated theoretically.
## 3 $`𝝎`$-, $`𝜼`$\- and $`𝜼^{\mathbf{}}`$-nuclear bound states
Next, we discuss the $`\omega `$, $`\eta `$-, and $`\eta ^{}`$-nuclear bound states . We have solved the Klein-Gordon equation using the calculated potentials (see Fig. 2):
$`\left[^2+E_j^2\stackrel{~}{m}_j^2(r)\right]\varphi _j(\stackrel{}{r})=0,E_j^{}E_j+m_ji\mathrm{\Gamma }_j/2,(j=\omega ,\eta ,\eta ^{}),`$ (10)
$`\stackrel{~}{m}_j^{}(r)=m_j^{}(r){\displaystyle \frac{i}{2}}\left[(m_jm_j^{}(r))\gamma _j+\mathrm{\Gamma }_j\right]m_j^{}(r){\displaystyle \frac{i}{2}}\mathrm{\Gamma }_j^{}(r),`$ (11)
where $`E_j^{}`$ is the complex valued, total energy of the meson, and we included the widths of the mesons in a nucleus assuming a specific form using $`\gamma _j`$, which are treated as phenomenological parameters. According to the estimates in Refs. , the widths of the mesons in nuclei and at normal nuclear matter density are $`\mathrm{\Gamma }_\eta ^{}3070`$ MeV and $`\mathrm{\Gamma }_\omega ^{}3040`$ MeV , respectively. Thus, we calculate the single-particle energies for the values $`\gamma _\omega =0.2`$, and $`\gamma _\eta =0.5`$, which are expected to correspond best with experiment, while for the $`\eta ^{}`$, $`\mathrm{\Gamma }_\eta ^{}^{}=0`$ is assumed. For a comparison we give also the results for the $`\omega `$ calculated using the potential obtained in Quantum Hadrodynamics (QHD) .
Our results suggest that one should expect to find bound $`\omega `$-, $`\eta `$\- and $`\eta ^{}`$-nuclear states for all nuclei investigated and relatively wide range of the in-medium meson widths . (For the predictions made by the other approaches, see Refs. .)
## 4 Summary
We have presented two possible signals for a change of meson properties in a nuclear medium. First, we discussed the impact on the observed $`J/\mathrm{\Psi }`$ suppression in relativistic heavy ion collisions, which has received a lot of interest recently because of the Quark-Gluon Plasma. Second, we discussed the prediction of $`\omega `$-, $`\eta `$\- and $`\eta ^{}`$-nuclear bound states. Both of these signals, for which many experiments are currently being perfomed or planned, certainly will provide important information about the in-medium properties of hadrons.
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# Observational implications of precessing protostellar discs and jets
## 1 Introduction
Recent observations of young stellar objects (YSOs) provide evidence for the existence of binary systems with circumstellar discs that are misaligned with the binary’s orbital plane. The Hubble Space Telescope (HST) and adaptive optics images of the HK Tau pre-main-sequence system provide the most striking and direct evidence, assuming that it is indeed a binary system (Stapelfeldt et al. 1998; Koresko 1998). In addition, there are several examples of pre-main-sequence binaries or unresolved YSOs from which two protostellar jets are seen to emanate in different directions (e.g. Davis, Mundt & Eislöffel 1994). It is assumed that these twin jets originate from binary systems with two circumstellar discs that are misaligned with each other. Thus, one or both of the discs must be misaligned with the binary’s orbital plane. Furthermore, observations of solar-type main-sequence binary systems indicate that the stellar rotational equatorial planes in wide binaries ($`\text{ }>40`$ a.u.) are frequently misaligned with the orbital plane, while for closer systems there is a tendency for alignment (Hale 1994). Assuming the equatorial plane of a star is determined by the plane of its original circumstellar disc, this indicates that circumstellar discs are frequently misaligned with the orbital plane in wide binaries.
In binary systems with a circumstellar disc whose plane is misaligned with the orbital plane, the circumstellar disc is expected to precess due to the tidal interaction of the companion (e.g. Papaloizou & Terquem 1995). Disc precession has been used to explain long-period variations in the light curves of a number of X-ray emitting binary systems (Gerend & Boynton 1976; Katz et al. 1982; Wijers & Pringle 1999). For YSOs, the major application of disc precession has been to predict the precession of protostellar jets, since many Class 0 and I objects are observed to emit jets (Bontemps et al. 1996). Several hydrodynamic investigations into the appearance of precessing jets have been performed (Raga, Cantó & Biro 1993; Biro, Raga & Cantó 1995; Cliffe et al. 1996; Völker et al. 1999) and precession has been used to explain the changes in the flow directions of several jets (Eislöffel et al. 1996; Davis et al. 1997; Mundt & Eislöffel 1998), although Eislöffel & Mundt (1997) stress that alternative explanations exist. Recently, Terquem et al. (1999) discussed the orbital periods that are required to give visibly precessing jets and the expected precession periods for systems that are observed to have misaligned jets.
In this paper, we present a simple review of the behaviour of a circumstellar disc which is misaligned with the orbital plane in a binary system, and discuss the implications of this behaviour for observations of discs and jets in binary protostellar systems. In Section 2, along with simple precession, we also consider the effect of the oscillating torque produced by the companion, which could cause the disc and jet to ‘wobble’ with a period of half the orbital period. In Section 3, we consider the effect of dissipation in the disc and calculate the timescale for a disc to realign itself with the binary’s orbital plane due to dissipation. The general results of this study, along with the expected observational consequences and the implications for models of binary star formation are discussed in Section 4. Finally, we give our conclusions in Section 5. The reader more interested in the results and conclusions of the paper than the details of the analysis may care to move directly to Section 4 from this point.
## 2 Precession induced by the binary companion
We consider here in simple terms the effect of a binary companion (the secondary) on a gaseous disc around the primary, whose plane is not aligned with the plane of the binary orbit.
### 2.1 Ring precession
We first consider the effect of the secondary on a circular ring of material orbiting the primary at radius $`a`$. We consider the binary star system with component masses $`M_p`$ (primary) and $`M_s`$ (secondary), with a circular orbit, and with stellar separation $`D`$. Now also consider a ring of material, of negligible mass, $`m_r`$, in orbit about the primary star. The ring has mean radius $`aD`$, and is tilted to the plane of the binary orbit at an angle $`\delta `$. We can regard the effect of the secondary star on the dynamics of the ring as a perturbation. We use a set of non-rotating coordinates centred on the primary star, with the OZ axis parallel to the rotation axis of the binary. In these coordinates, the acceleration caused by the secondary star at position vector r is given by
$$𝐅_s=(GM_s/D^3)𝐫+(3GM_s/D^5)(𝐫.𝐃)𝐃,$$
(1)
where $`𝐃`$ is the position vector of the secondary relative to the primary. We should remark that the fictitious force arising from the acceleration of the coordinate system has been included, cancelling the $`(GM_s/D^3)𝐃`$ term, and that terms of relative order $`(a/D)`$ have been neglected. Here the first term simply represents an augmentation of the effect of the gravity of the primary due to the presence of the secondary. It affects the centrifugal balance of the ring, but otherwise has no effect on its dynamics. The second term is the tidal term and represents a force in a direction parallel (or anti-parallel) to the instantaneous vector joining the two stars, and with magnitude proportional to the distance of the position r from the plane passing through the primary perpendicular to the line joining the two stars. From a physical point of view, this term can be thought of as comprising two components. The first is a force which acts similarly to a centrifugal force in that it is everywhere directed perpendicular to and away from the OZ axis, and in that its magnitude is proportional to the distance from that axis. In terms of azimuthal Fourier components $`\mathrm{exp}(im\varphi )`$, where $`\varphi `$ is the azimuth about the OZ axis, this term is derived from an axially symmetric potential, and thus has $`m=0`$. The second is similar except that the magnitude of the force is proportional to distance from the OZ axis multiplied by $`\mathrm{cos}(2\varphi )`$, and thus has $`m=2`$.
#### 2.1.1 The $`m=0`$ term
Breaking the tidal term into these two parts enables us to consider their effects on the dynamics of the ring in a more straightforward manner. Suppose that the coordinates of the secondary are given by
$$𝐃=(D\mathrm{sin}(\mathrm{\Omega }_bt),D\mathrm{cos}(\mathrm{\Omega }_bt),0),$$
(2)
where $`2\pi /\mathrm{\Omega }_b`$ is the binary period, and thus
$$\mathrm{\Omega }_b=(G(M_s+M_p)/D^3)^{1/2}.$$
(3)
Suppose that the ring is tilted in our system of coordinates about the OX axis. Thus the axis of the ring is given by
$$𝐤=(0,\mathrm{sin}\delta ,\mathrm{cos}\delta ).$$
(4)
Then the effect of the $`m=0`$ part of the force is to exert a torque $`𝐆_0=(G_{0x},0,0)`$ on the ring about the negative OX axis, of magnitude
$$G_{0x}=(3/4)(GM_sm_ra^2/D^3)\mathrm{sin}\delta \mathrm{cos}\delta ,$$
(5)
where $`m_r`$ is the mass of the ring. The direction of the $`m=0`$ torque, $`𝐆_\mathrm{𝟎}`$, lies along the line of nodes of the ring (the line along which the ring intersects the OXY-plane). If $`𝛀=\mathrm{\Omega }𝐤`$ is the angular velocity of the ring, then $`𝐆_0𝛀=0`$. Thus the effect of the torque on the ring is to cause the plane of the ring to precess retrogradely about the OZ axis with a precession rate given by $`\omega _p`$, where
$$\omega _p=G_{0x}/(\mathrm{sin}\delta m_ra^2\mathrm{\Omega }).$$
(6)
Note that $`\mathrm{\Omega }`$ is given by
$$\mathrm{\Omega }=(GM_p/a^3)^{1/2},$$
(7)
and $`m_ra^2\mathrm{\Omega }`$ represents the angular momentum of the ring.
Thus we find that the mean precession rate is
$$\omega _p=(3/4)\mathrm{cos}\delta (GM_s/D^3\mathrm{\Omega }),$$
(8)
or, equivalently, that
$$\omega _p/\mathrm{\Omega }=(3/4)\mathrm{cos}\delta (M_s/M_p)(a/D)^3.$$
(9)
We should also note that the above treatment, which essentially assumes that the ring maintains its shape as if it were a solid body, will be valid provided that $`\mathrm{\Omega }\mathrm{\Omega }_b\omega _p`$, which is satisfied in this case (see below).
#### 2.1.2 The $`m=2`$ term
The effect of the $`m=2`$ term is to give rise to an oscillating torque
$$𝐆_2=(G_{2x},G_{2y},G_{2z})$$
(10)
on the ring which has a frequency of $`2\mathrm{\Omega }_b`$. The $`m=2`$ torque is of the form
$$G_{2x}=(3/4)(GM_sm_ra^2/D^3)\mathrm{sin}\delta \mathrm{cos}\delta \mathrm{cos}(2\mathrm{\Omega }_bt),$$
(11)
$$G_{2y}=(3/4)(GM_sm_ra^2/D^3)\mathrm{sin}\delta \mathrm{cos}\delta \mathrm{sin}(2\mathrm{\Omega }_bt),$$
(12)
and
$$G_{2z}=(3/4)(GM_sm_ra^2/D^3)\mathrm{sin}^2\delta \mathrm{sin}(2\mathrm{\Omega }_bt).$$
(13)
We note also that $`𝐆_2𝛀=0`$. Indeed, the total torque on the ring may be expressed simply as
$$𝐆=(3/2)(GM_sm_ra^2/D^5)(𝐃𝐤)(𝐃\times 𝐤).$$
(14)
Because of the oscillating nature of the $`m=2`$ torque, the mean precession rate is unaffected by the $`m=2`$ term. However, the instantaneous precession rate oscillates with a period of (approximately) one half of the orbital period (Katz et al. 1982). Since the amplitude of the torque due to the $`m=2`$ term is equal to the torque generated by the $`m=0`$ term the modulation is not a negligible effect.
Perhaps the simplest way to envisage what is happening is to consider the motion of the symmetry axis $`𝐤`$ of the ring. The $`m=0`$ term causes the axis to move in a cone, semi-angle $`\delta `$, around the OZ axis at a rate $`\omega _p`$. Superimposed on this, the $`m=2`$ term causes a wobble of the symmetry axis in the direction of precession, and a nodding motion perpendicular to this direction. The amplitude of the wobble (in radians) is roughly $`\omega _p/(2\mathrm{\Omega }_b)`$, and the amplitude of the nodding motion is the same but multiplied by a factor of $`\mathrm{tan}\delta `$ (see Katz et al. 1982). We should remark that this does not lead to a divergence as $`\delta \pi /2`$ because $`\omega _p`$ is proportional to $`\mathrm{cos}\delta `$. We should note that, strictly, because the ring is precessing in a retrograde direction because of the $`m=0`$ term, the period of the wobble is in fact $`2\pi /(2\mathrm{\Omega }_b+\omega _p)`$.
### 2.2 Disc precession
We have seen above that the precession rate for a ring of radius $`a`$ is proportional to $`a^{3/2}`$. Thus if we regard a disc as being made up of a succession of concentric rings, the effect of the $`m=0`$ potential is to cause the disc to precess differentially. However, if the disc is able to hold itself together in some way, either by means of wave-like communication (Larwood et al. 1996), or by viscous communication (Wijers & Pringle 1999), then it may be able to respond coherently to the $`m=0`$ term in the potential by precessing at the same rate at all radii, that is, as if it were a solid body. The precession rate can then be calculated straightforwardly from the above, by rewriting the ring mass, $`m_r`$, as the mass of an annulus in the disc, $`2\pi \mathrm{\Sigma }ada`$, where $`\mathrm{\Sigma }`$ is the disc surface density at radius $`a`$.
The net precession rate is then given by
$$\omega _p=K\mathrm{cos}\delta (GM_s/D^3\mathrm{\Omega }_d),$$
(15)
where (Terquem 1998)
$$K=\frac{3}{4}R^{3/2}\left[_0^R\mathrm{\Sigma }a^3𝑑a\right]/\left[_0^R\mathrm{\Sigma }a^{3/2}𝑑a\right].$$
(16)
Here $`\mathrm{\Omega }_d`$ is the angular velocity of disc material at the outer edge of the disc which is at radius $`R`$, and the lower limit in the integrals can be simply a small radius rather than zero.
Thus, for a disc,
$$\omega _p/\mathrm{\Omega }_d=K\mathrm{cos}\delta (M_s/M_p)(R/D)^3.$$
(17)
For a constant surface density disc $`K=15/32`$ (Larwood et al. 1996); for a surface density proportional to $`R^{3/2}`$, $`K=3/10`$.
### 2.3 Application to protostellar discs
In protostellar discs the ratio of disc semi-thickness to radius, $`H/R`$, is of order 0.1 (Burrows et al. 1996; Stapelfeldt et al. 1998), which is larger than the typical dimensionless viscosity $`\alpha 0.01`$ (Hartmann et al. 1998). Under these circumstances bending waves can propagate through the disc (assumed Keplerian) (Papaloizou & Lin 1995, Pringle 1997), and thus the disc can communicate with itself in a wave-like manner. The propagation speed for such bending waves is of order $`\frac{1}{2}c_s`$, where $`c_s`$ is the disc sound speed (Papaloizou & Lin 1995, Pringle 1999). Thus the condition that the disc be able to precess as a solid body subject to forced differential precession might be expected to be that the wave crossing timescale be less than the precession timescale, that is
$$R/c_s\text{ }<\omega _p^1.$$
(18)
This is the criterion given by Larwood et al. (1996), and by Papaloizou & Terquem (1995). Using the standard result for accretion discs (e.g. Pringle 1981) that $`c_s/(R\mathrm{\Omega }_d)H/R`$, we find that the condition may be written
$$\omega _p/\mathrm{\Omega }_d\text{ }<H/R.$$
(19)
For a disc in a binary system which is truncated by tidal forces, the outside disc edge is typically a substantial fraction of the binary separation (e.g Papaloizou & Pringle 1977; Paczyński 1977; Artymowicz & Lubow 1994; Larwood et al. 1996). As a typical value, for reasonable mass ratios, we shall take $`R/D0.3`$. Thus, taking $`K0.4,q=M_s/M_p1`$, and $`\mathrm{cos}\delta 1`$, we find that
$$\omega _p/\mathrm{\Omega }_d0.011(K/0.4)(R/0.3D)^3q\mathrm{cos}\delta ,$$
(20)
We note that condition (19) is readily satisfied for protostellar discs. However, as we shall note below, this is not the whole story.
Using the above estimates, we note that
$$\begin{array}{cc}\mathrm{\Omega }_d/\mathrm{\Omega }_b\hfill & =[M_p/(M_p+M_s)]^{1/2}(D/R)^{3/2}\hfill \\ & \\ & 4.3[2/(1+q)]^{1/2}(R/0.3D)^{3/2}.\hfill \end{array}$$
(21)
Thus typically
$$\omega _p/\mathrm{\Omega }_b0.05(K/0.4)q[2/(1+q)]^{1/2}(R/0.3D)^{3/2}\mathrm{cos}\delta .$$
(22)
This implies that the amplitude of the wobble/nodding motion about steady precession (discussed in Section 2.1.2) is of order a degree or so at the outer edge of the disc although, as we note in Section 4.2, the amplitude inside the disc depends on the properties of the disc.
Some protostellar discs may be sufficiently massive that a noticeable deviation from Keplerian rotation may occur. This effect, and also the self-gravitation of the disc, could change the nature of the propagating bending waves, making them dispersive (Papaloizou & Lin 1995). However, assuming that self-gravity is not dominant (i.e. the Toomre parameter $`Q1`$), these effects are expected to be small because the fractional deviation from Keplerian rotation is at most of order $`H/R`$, and our estimates here are probably still valid. The character of wave propagation also differs between vertically isothermal and thermally stratified discs \[Lubow & Ogilvie 1998\], however this is relevant only when the radial wavelength is comparable to the disc thickness, and therefore does not affect the propagation of long-wavelength bending waves as considered here.
## 3 The effect of dissipation
As we have seen, the tidal effect of the secondary star can be regarded being due to the $`m=0`$ and $`m=2`$ parts of the tidal potential independently. Thus the effect of dissipation within the disc on the motions caused by these two parts of the potential can also be calculated separately.
### 3.1 The $`m=0`$ term
The effect of the $`m=0`$ term on the misaligned disc is to produce a torque $`𝐆_0`$ on the disc which is in the direction of aligning it with the binary orbit. This comes about because there is a potential energy, $`\mathrm{\Phi }`$, associated with the misalignment of the disc which is given by
$$\mathrm{\Phi }=(3GM_s/8D^3)\mathrm{sin}^2\delta _0^Ra^2\mathrm{\Sigma }2\pi a𝑑a.$$
(23)
This has minima at $`\delta =0`$ (alignment) and $`\delta =\pi `$ (anti-alignment), and a maximum at $`\delta =\pi /2`$. Thus we would normally expect loss of energy associated with the motions induced by the $`m=0`$ part of the tidal potential to lead to the disc becoming aligned with the orbital plane.
However, because the $`m=0`$ term is symmetric about the OZ-axis, $`\widehat{𝐳}`$, any torque produced by the $`m=0`$ term must lie in the OXY-plane. We have seen that in the absence of dissipation the torque is such that $`𝐆_0`$ is parallel to $`𝐤\times \widehat{𝐳}`$, where $`𝐤`$ is the unit vector in the direction of the angular momentum $`𝐉`$ of the disc. The effect of viscosity is to give rise to a phase delay in the torque about the OZ axis, which has the effect of producing an additional (viscous) torque $`𝐆_{0\nu }`$, which is perpendicular to $`\widehat{𝐳}`$, which lies in the plane defined by $`𝐉`$ and $`\widehat{𝐳}`$, and which is directed such that $`𝐆_{0\nu }𝐉<0`$, so that its effect is dissipative. Thus the net effect of this torque is to align the disc with the OZ-axis, but in such a way that $`𝐉\widehat{𝐳}`$ is conserved. This means that in the process of alignment, angular momentum is removed from the disc, giving rise to an enhanced accretion rate (see Papaloizou & Terquem 1995). We compute the timescale on which this torque leads to disc alignment below.
### 3.2 The $`m=2`$ term
In the absence of dissipation, the effect of the $`m=2`$ term is to induce oscillations in the disc which have a frequency of $`2\mathrm{\Omega }_b`$, and which therefore cause the angular momentum vector $`𝐉`$ to oscillate about some mean value. Thus we would naively expect that the effect of dissipation on the motions induced by the $`m=2`$ part of the tidal potential would be to reduce the amplitude of the oscillations, but to have no net effect on the mean disc plane, that is, to have no long-term time-averaged effect on the value of $`\delta `$.
Papaloizou & Terquem (1995) and Terquem (1998) have argued that the effect of dissipation on the $`m=2`$ induced motions is to give rise to a torque which may increase the inclination of the disc. This possibility has been investigated in detail by Lubow & Ogilvie \[Lubow & Ogilvie 2000\], who formulated the problem in terms of the linear stability of an initially coplanar disc in the presence of the tidal field. It was confirmed that the $`m=2`$ component of the tidal potential causes the inclination to grow, but the effect is usually negligible and is outweighed by the effect of the $`m=0`$ potential, so the net outcome is that the inclination decays in time. An exception occurs if there is a coincidence between the frequency of the oscillating torque, $`2\mathrm{\Omega }_\mathrm{b}`$, and the natural frequency of a global bending mode of the disc. However, this resonance can occur only if the disc is very thick ($`H/R0.4`$) or much smaller than the standard tidal truncation radius ($`R/D0.3`$). We therefore neglect the effect of the $`m=2`$ term on the long-term time-averaged value of $`\delta `$ and restrict our attention to the alignment effects of the $`m=0`$ term.
### 3.3 Application to protostellar discs
In the analysis so far we have assumed that the disc is able to communicate internally sufficiently fast that the disc can precess as a rigid body. This communication takes the form of bending waves which propagate through the disc at a speed of order $`\frac{1}{2}c_\mathrm{s}`$. However, these bending waves are associated with nearly resonant horizontal epicyclic motions that are strongly shearing, being proportional to the distance above the mid-plane (Papaloizou & Pringle 1983). Viscosity in the disc can then act on these shearing motions, leading to dissipation of energy in the precessional motion, and alignment of the disc with the orbital plane.
We have seen that the torques exerted at different radii in the disc by the $`m=0`$ component of the tidal potential would naturally result in differential precession. In order for the disc to resist this, hydrodynamic stresses must be established within the disc so that the net torque on each ring is such as to maintain a uniform, global precession rate. The required internal torque, although it varies with radius and vanishes at the edges of the disc, is therefore generally comparable to the total tidal torque on the disc, and is given approximately by
$$G_{\mathrm{int}}2\pi \mathrm{\Sigma }R^4\mathrm{\Omega }_\mathrm{d}\omega _\mathrm{p}\mathrm{sin}\delta .$$
(24)
These hydrodynamic stresses are associated with the horizontal epicyclic motions mentioned above, which take the form
$$v_r^{}2v_\varphi ^{}Az,$$
(25)
referred to cylindrical polar coordinates $`(r,\varphi ,z)`$ based on the mean disc plane, and where $`A`$ is independent of $`z`$. As a result of the hydrodynamic stress $`\rho v_r^{}r\mathrm{\Omega }`$ there is a net horizontal angular momentum flux
$$2\pi r\rho v_r^{}r\mathrm{\Omega }z𝑑z2\pi A\mathrm{\Sigma }R^2H^2\mathrm{\Omega }_\mathrm{d}.$$
(26)
Equating this with $`G_{\mathrm{int}}`$ gives
$$A(R/H)^2\omega _\mathrm{p}\mathrm{sin}\delta ,$$
(27)
and so
$$v_r^{}/c_\mathrm{s}2v_\varphi ^{}/c_\mathrm{s}(R/H)^2(\omega _\mathrm{p}/\mathrm{\Omega }_\mathrm{d})\mathrm{sin}\delta (z/H).$$
(28)
Note that for these velocities to be subsonic we require that
$$\omega _p/\mathrm{\Omega }_d\text{ }<(H/R)^2/\mathrm{sin}\delta .$$
(29)
For $`\mathrm{sin}\delta >H/R`$, this is a stronger condition than the one derived by Papaloizou & Terquem (1995), and if $`\delta 1`$, it is only marginally satisfied in protostellar discs. The fact that the internal disc velocities are likely to be close to sonic, means that the dissipation may well be strongly enhanced. We discuss this further below.
For an internal disc viscosity $`\nu `$, these velocity perturbations lead to a rate of energy dissipation in the disc of order
$$dE/dtm_dc_s^2(\nu /H^2)(R/H)^4(\omega _p/\mathrm{\Omega }_d)^2\mathrm{sin}^2\delta ,$$
(30)
where $`m_d`$ is the mass of the disc.
As we have seen above, the effect of this dissipation is to lead to an alignment of the disc with the orbital plane. The means by which this is accomplished is through the viscous torque $`𝐆_{0\nu }`$, which leads to alignment by removing the component of disc rotation which lies in the orbital plane (recall that $`𝐆_{0\nu }\widehat{𝐳}=0`$). Thus the amount of energy to be dissipated in order to bring about alignment is
$$E_{\mathrm{kin}}m_dR^2\mathrm{\Omega }_{d}^{}{}_{}{}^{2}\mathrm{sin}^2\delta .$$
(31)
From these estimates we deduce the alignment timescale for the disc $`\mathrm{t}_{\mathrm{align}}`$, which is given by
$$\mathrm{t}_{\mathrm{align}}E_{\mathrm{kin}}/(dE/dt)(R^2/\nu )(H/R)^4(\mathrm{\Omega }_d/\omega _p)^2.$$
(32)
We note further that the viscous evolution timescale for a disc $`t_\nu `$ is given by
$$t_\nu R^2/\nu .$$
(33)
Thus, from this analysis we would conclude that for precessing discs in which the induced velocities are subsonic, the alignment timescale for an misaligned protostellar disc is typically of order, or somewhat longer than, the viscous evolution timescale, and that the additional accretion rate caused by the process of alignment is typically at most comparable to the accretion rate already present in the disc.
The viscous evolution timescale can be written in terms of the dimensionless measure of viscosity, $`\alpha `$, as
$$t_\nu \mathrm{\Omega }_{d}^{}{}_{}{}^{1}(R/H)^2\alpha ^1.$$
(34)
Using this we find that
$$t_{\mathrm{align}}\omega _{p}^{}{}_{}{}^{1}\alpha ^1(H/R)^2(\mathrm{\Omega }_d/\omega _p).$$
(35)
This expression for the alignment timescale and the expression for the precession rate (equation 22) have been verified by the numerical calculations of Lubow & Ogilvie \[Lubow & Ogilvie 2000\].
#### 3.3.1 The effect of sonic induced velocities
However, since the velocities associated with the induced shearing in these discs are close to sonic, the parametric instabilities discussed by Gammie, Goodman & Ogilvie (2000) are likely occur. These authors concluded that the internal shear motions induced by the disc bending are, for sonic shearing motions, unstable on a local (disc) dynamical timescale, $`\mathrm{\Omega }^1`$. The local turbulence induced by these instabilities increases the damping rate for the shearing motions. For example, if $`v^{}c_s`$, then we may estimate the dissipation rate in the disc due to these instabilities as
$$dE/dtm_dc_{s}^{}{}_{}{}^{2}\mathrm{\Omega }_d,$$
(36)
which implies an alignment timescale of
$$t_{\mathrm{align}}\omega _{p}^{}{}_{}{}^{1}/\mathrm{sin}\delta .$$
(37)
This implies that alignment occurs almost as fast as precession.
Furthermore, even when the induced velocities, $`v^{}`$ are subsonic, the instabilities grow on a timescale $`H/v^{}`$ (Gammie et al. 2000), provided that the growth time is less than the viscous damping timescale $`H^2/\nu `$, i.e. provided that
$$\alpha \text{ }<v^{}/c_s,$$
(38)
or, alternatively, that
$$\mathrm{sin}\delta \text{ }>(\mathrm{\Omega }_d/\omega _p)(H/R)^2\alpha $$
(39)
i.e. that
$$\mathrm{sin}\delta \text{ }>0.01\frac{(\alpha /0.01)(H/0.1R)^2}{(K/0.4)(R/0.3D)^3q\mathrm{cos}\delta }.$$
(40)
In this case the parametric instabilities again give enhanced dissipation which is likely to lead to an enhanced effective viscosity, $`\alpha _{\mathrm{eff}}`$, such that the enhanced damping rate is of order the instability growth rate, viz.
$$\alpha _{\mathrm{eff}}v^{}/c_s.$$
(41)
This enhanced damping implies an alignment timescale of
$$\begin{array}{cc}t_{\mathrm{align}}\hfill & \omega _{p}^{}{}_{}{}^{1}(H/R)^4(\mathrm{\Omega }_d/\omega _p)^2/\mathrm{sin}\delta \hfill \\ & \\ & \omega _{p}^{}{}_{}{}^{1}\frac{(H/0.1R)^4(K/0.4)^2(R/0.3D)^6q^2}{\mathrm{cos}^2\delta \mathrm{sin}\delta }.\hfill \end{array}$$
(42)
We should stress that this estimate is probably a lower limit to the alignment timescale since the efficiency of the damping due to the parametric instability may not be as perfect as suggested by equation 41 (Gammie et al. 2000). We also note that although Gammie et al. (2000) considered an isothermal, Keplerian disc, the growth rate of the parametric instability is not sensitive to these assumptions. The instability would also act in a thermally stratified disc, or in a self-gravitating disc with slightly non-Keplerian rotation.
## 4 Discussion
### 4.1 General results
We have presented a simple discussion of the dynamics of a misaligned disc in a binary star system, and have applied the results to the parameters of discs and binary parameters which are relevant to young stellar objects. Our general conclusions are as follows:
For typical protostellar disc parameters, and for circular binary orbits, tidal forces cause a misaligned disc to precess like a solid body about the angular momentum vector of the binary orbit, with a precession period of order $`P_p20P_{orb}`$, where $`P_{orb}`$ is the orbital period (equation 22; Papaloizou & Terquem 1995; Larwood et al. 1996; Terquem 1998; Terquem et al. 1999). In addition, the outer disc plane is forced to wobble with a period of $`\frac{1}{2}P_{orb}`$ (Katz et al. 1982).
We have presented order of magnitude estimates for the timescale on which dissipation within the disc leads to alignment with the orbital plane. (For the reasons described above, we overlook the possibility raised by Papaloizou & Terquem (1995) that dissipation might lead to disc misalignment). If the disc evolution is determined by a simple isotropic viscosity, then we find, in line with the estimates given by Papaloizou & Terquem (1995) and by Terquem et al. (1999) that the alignment timescale is, for protostellar disc parameters, of order the normal viscous evolution timescale.
However, we have also pointed out that the velocities induced within the disc by the action of tidal forces on the tilt are, for typical protostellar disc parameters, transonic, and that the criterion used by Papaloizou & Terquem (1995) to justify the use of their linearization procedure is incorrect if the disc tilt exceeds the opening angle of the disc (equation 29). The induced velocities take the form of a horizontal epicyclic motion proportional to the distance above the mid-plane within the disc which oscillates in a frame rotating with the fluid with a period approximately equal to that of the orbital period of the disc material. It has long been suspected that such a shear flow is unstable (Kumar & Coleman 1993), and recent work by Gammie et al. (2000) has demonstrated that the flow is indeed unstable to a parametric hydrodynamic instability which has a growth rate of order the shearing timescale and leads to rapid dissipation. We have estimated the effect of such instabilities. For large disc tilts, and for typical protostellar disc parameters, we find that the disc alignment timescale is comparable to the precession timescale. However, for smaller tilt angles $`\delta `$, we find that the alignment timescale varies as $`(\mathrm{sin}\delta )^1`$. It is worth noting that the enhanced dissipation will also result in a greater disc luminosity and a larger mass accretion rate than those provided by standard viscous evolution. These effects should be considered when modelling protostellar discs in binary systems. However, given the large range of observed accretion rates from protostellar discs (e.g. Hartmann et al. 1998) and the uncertainties in current models of protostellar discs, it would be difficult to detect an unambiguous signature of the enhanced luminosity or accretion rate.
We should draw attention to the fact that, in common with previous authors, all the results presented above are for binary stars with circular orbits. If, as is the case for many binary stars, the orbit is non-circular, then the analysis is similar but rather more complicated. However, the main effects of orbital eccentricity on the results reported above are to decrease the ratio of disc radius to orbital semi-major axis and to modify the time-averaged potential due to the companion. Tidal truncation of the disc now takes place at periastron, so that $`R0.3a(1e)`$, while the modified time-averaged potential can be approximated by replacing the orbital separation $`D`$, by $`a(1e^2)^{1/2}`$ (Holman, Touma & Tremaine 1997), where $`a`$ is the semi-major axis and $`e`$ is the eccentricity. Thus, to a first approximation, the major effects of orbital eccentricity can be taken into account by replacing $`R/0.3D`$ in the above formulae by the quantity
$$\sqrt{\frac{1e}{1+e}}\left(\frac{R}{0.3a(1e)}\right).$$
(43)
### 4.2 Precession and wobbling of protostellar jets
Disc precession has been used as an ingredient in the explanation of long-period variations in the light curves of a number of X-ray emitting binary systems (Gerend & Boynton 1976; Katz et al. 1982; Wijers & Pringle 1999). However, for protostellar systems the major application for tidally induced disc precession has been to provide an explanation for changes in flow direction of protostellar jets (e.g. Eislöffel & Mundt 1997; Terquem et al. 1999). Although changes in flow direction are often discussed in terms of precession, there are, as discussed by Eislöffel & Mundt (1997), various alternative explanations for such phenomena, and there is as yet no convincing case of a jet which has been steadily precessing for many precession periods. In the light of this we discuss, in general terms, the kind of effects which tidally induced disc precession might be expected to lead to from a theoretical point of view.
Protostellar jets appear to be produced during the major accretion phase in the star-formation process, that is during the Class 0 and Class I phases of the life of a protostellar core/young star (Bontemps et al. 1996). This phase of the star formation process is thought to last about $`10^5`$ years (Lada 1999). In addition, given that the jet velocities are of order the escape velocities from the central stars, and that the mass outflow rates are a non-negligible fraction of the likely accretion rates, it seems a reasonable assumption that the jets are formed close to the centre of the disc (e.g. Pringle 1993), and therefore that the jet direction is governed by the disc axis in the central regions of the disc. Thus, since for a strongly misaligned disc, the alignment timescale is of order the precession timescale, which is of order $`20`$ orbital periods, we expect that, regardless of how the misalignment might have come about, strongly misaligned discs and jets are only likely to occur in binaries with orbital periods longer than $`5000`$ years, that is with separations larger than about $`100`$ a.u. Examples of systems which appear to have misaligned jets include: Cep E \[Eislöffel et al. 1996\]; T Tau \[Böhm & Solf 1994\]; HH 1,2,144 VLA 1/2 \[Reipurth et al. 1993\]; and HH 111/121 \[Gredel & Reipurth 1993\]. These systems are either known to be wide binaries ($`\text{ }>100`$ a.u.), in agreement with expectations, or have not yet been resolved.
Since the alignment and precession timescales are comparable for strongly misaligned discs, one would not expect to observe the multiple large scale wiggles which might be the result of a jet undergoing many precession cycles at a large angle to the orbital axis. One might, however, see the results of such a jet whose direction starts at large angle to the binary axis, and then changes direction on the sky as it simultaneously precesses and aligns (with the binary axis) on a timescale of $`20`$ orbital periods. For such a jet, with jet axis at a large angle to the binary axis, provided that the inner and outer parts of the disc are in good communication, the jet direction is forced to wobble by the $`m=2`$ tidal component with a period of half the orbital period. The amplitude of the wobble at the outside of the disc is only about a degree or so, but since the communication to the disc centre is wave-like, the amplitude of the wobble at the disc centre must depend on the wave amplitude induced there, which in turn depends on the properties of the disc, and in particular on the dependences with radius of the angular momentum density ($`\mathrm{\Sigma }r^{1/2}`$) and the group velocity ($`c_s`$) (e.g. Terquem 1998). Such a regular wobbling of the jet direction may already have been observed in one of the jets emanating from Cep E \[Eislöffel et al. 1996\]. The jet appears to undergo a regular wobble with an amplitude of $`4^{}`$ and a period of $`400`$ years (both dependent on the inclination). The mechanism driving this oscillation may be able to be tested observationally using the VLA or adaptive optics in the infrared. If the oscillation of the jet is due to a wobble of the disc, the period of the binary should be $`800`$ years (separation $`80`$ a.u. or $`0.12^{\prime \prime }`$), whereas if the oscillation is due to precession of the disc the binary’s period should be much shorter ($`\text{ }<20`$ years) and the binary would not be resolvable ($`\text{ }<0.01^{\prime \prime }`$). Many YSOs with single jets also show evidence for direction changes, and a high-resolution survey to determine the binarity of such sources would be invaluable for testing the theory of wobbling and precessing jets.
For jets that are weakly misaligned, we note that the alignment timescale is inversely proportional to $`\delta `$, (implying that the misalignment angle decreases with time as $`t^1`$) and scales, for typical protostellar disc parameters, approximately with the orbital period. Furthermore, the ratio of the precession timescale to the alignment timescale scales with $`\delta `$. In order to see multiple wiggles caused by jet precession in a jet whose length corresponds typically to a dynamical timescale of $`10^4`$ years, we require a precession period of less than or of order a few thousand years, and thus binary periods of less than or of order a few hundred years, and binary separations less than or of order a few tens of a.u. For example, an initially strongly misaligned binary with a separation of a few tens of a.u. will rapidly evolve into a weakly misaligned system, but even after $`10^5`$ years will still show a misalignment angle of $`0.03`$ radians, that is, a few degrees, and thus, with a precession period of a few thousand years, such a system would be marginally observable. Wider binaries than this would have precession periods too long to show wiggling of a jet whose dynamical age is $`10^4`$ years, whilst closer binaries with smaller precession periods would have alignment angles too small to be observed. We note that for small misalignment angles, the amplitude of the wobble of the outer disc caused by the $`m=2`$ tidal term (which is approximately $`\delta \omega _\mathrm{p}/(2\mathrm{\Omega }_\mathrm{b})\delta /40`$ as $`\delta 0`$) is probably too small to be detected.
In summary, one expects jets from strongly misaligned discs to be a rarity amongst binaries closer than $`100`$ a.u. In wider binaries, jets in strongly misaligned systems are wiggled at twice the binary orbital frequency with an amplitude which depends on the details of the disc structure. Weakly misaligned systems are longer-lived, and therefore potentially observable in closer binaries. As $`\delta 0`$, precession should give rise to observable jet wiggling (with amplitude of angle $`\delta `$ decreasing as $`\delta ^1`$) in binaries with separation of order a few tens of a.u.
### 4.3 Implications for binary star formation
We now consider briefly the implications that evidence for disc misalignment might have for the various theories of binary star formation, noting that in order for a misaligned system to be observed, the system as a whole must be assembled over a timescale which is shorter than the alignment timescale, $`t_{\mathrm{align}}`$.
#### 4.3.1 Binary fragmentation
Most published fragmentation calculations result in discs which are coplanar with the resulting fragments. Of these, there are typically two cases: fragmentation due to initial density perturbations (e.g. Boss & Bodenheimer 1979; Boss 1986), or centrifugally-supported fragmentation (e.g. fragmentation of a massive protostellar ring or disc; Norman & Wilson 1978; Bonnell 1994; Bonnell & Bate 1994; Burkert & Bodenheimer 1996; Burkert, Bate & Bodenheimer 1997) with some calculations exhibiting both types of fragmentation (e.g. Bonnell et al. 1991; Bate, Bonnell & Price 1995). The fragmentation of centrifugally-supported material leads to the orbit(s) and discs of all fragments occupying the same plane. In most calculations where the fragmentation occurs due to initial density perturbations, the initial conditions have a single axis of rotation and there is typically an $`m=2`$ density perturbation which is perpendicular to the rotation axis. With these simple initial conditions, a binary forms from the initial $`m=2`$ density perturbation in a plane perpendicular to the rotation axis and passing through the centre of the cloud, while the discs that form around these fragments have rotation axes that are parallel to the rotation axis of the cloud. Thus, the protostellar discs lie in the same plane as the orbit.
In order to produce discs that are misaligned with the orbital plane, it is natural to expect that the angular momentum distribution of the molecular gas must have strong spatial variations to enable the discs to have different rotation axes from the larger-scale orbit. However, this is not necessary. All that is required is to force the binary to form in a plane that is not perpendicular to the rotation axis. This can be achieved simply by rotating the $`m=2`$ density perturbation slightly so that it is no longer perpendicular to the rotation axis (or visa versa). Such a lack of correlation between the axis of rotation and the initial density perturbation(s) is expected if the collapse of a molecular gas is triggered by an external source (e.g. a shock wave or gravitational interaction with a passing object), or if the pre-collapse clump of gas is formed dynamically within a turbulent molecular cloud. In fact, given that the rotational energy of observed molecular clumps is generally insignificant compared with their gravitational energy (Goodman et al. 1993), it is hard to find a reason why there should be any correlation. In this case, the two fragments form above and below the plane that is perpendicular to the rotation axis and passes through the cloud’s centre and, hence, the orbital plane is no longer perpendicular to the rotation axis of the initial cloud. The rotation axes of the discs that form around the fragments, on the other hand, are still parallel to the rotation axis of the initial cloud. Thus, even though the pre-collapse gas is all rotating around a single axis, possibly in solid-body rotation, the discs are misaligned with the orbital plane of the binary. Such fragmentation calculations have been performed by Bonnell et al. (1992) and Bonnell & Bastien (1992). In their particular case, the $`m=2`$ density perturbation was in the form of the initial molecular cloud core being prolate. We note that in such calculations, although the discs are misaligned with the orbital plane, they are still aligned with each other. To produce the added complexity of discs which are initially misaligned with each other would require spatial variations in the angular momentum distribution. However, even if the discs are initially aligned, they are almost certain to undergo precession at different rates (equation 17) and, therefore, will very rapidly become misaligned with each other.
#### 4.3.2 Misaligned systems from dynamical interactions
The formation of wide binaries with misaligned discs via fragmentation that was described above is the simplest example of ‘prompt initial fragmentation’ (Pringle 1989) where each stellar component and its concomitant disc forms from a spatially-distinct region of the collapsing cloud, the accretion of large amounts of material on to the binary as a whole is avoided (see below), and, thus, the discs need not be initially aligned.
More complicated initial conditions than those described above can lead to the formation of a small, three-dimensional, cluster of stars (e.g. Larson 1978; Chapman et al. 1993; Klessen et al. 1998). The stars in such a group are expected to undergo interactions with one another, such as dissipative star-disc encounters (Larson 1990; Clarke & Pringle 1991a, 1991b; Clarke & Pringle 1993; McDonald & Clarke 1993, 1995).
Highly-dissipative encounters may lead to the formation of binary (or multiple) systems, via capture of the passing object, typically with misaligned discs. However, it remains to be determined to what extent these dissipative interactions also lead to disc alignment through strong tidal interactions (Heller 1993; Hall, Clarke & Pringle 1996; Hall 1997), especially given that in reality the major infall phase onto the separate protostellar nuclei occurs contemporaneously with (Bonnell et al. 1997; Klessen et al. 1998), rather than prior to, the dissipative binary formation process (see the next section).
Another possibility is that the gravitational interaction of a passing object, although not leading to capture, may tilt the disc(s) within a binary system or around a single star. If such an interaction occurred during the main accretion phase, the tilting of the disc would likely lead to a change in the direction of an emanating jet which could be mistaken for precession or realignment of a misaligned disc in a binary system. Thus, to unambiguously identify a precessing jet we emphasize that the jet should exhibit several oscillations.
#### 4.3.3 Subsequent accretion
Whether a binary is formed directly by fragmentation or more complicated interactions between several protostars, the above discussion ignores the effect of accretion of material after the binary has formed.
If a protobinary forms with discs that are initially misaligned with the orbital plane, but subsequently accretes the majority of its mass, the misalignment will generally be diminished. For binaries formed directly via fragmentation, the mass of the protobinary immediately after its formation is less for binaries with smaller initial separations (Boss 1986). Therefore, to obtain the same final total mass, closer binaries must accrete more, relative to their initial mass, than wider binaries and, hence, close binaries ($`\text{ }<100`$ a.u.) are less likely to exhibit misaligned discs than wider binaries. Binaries formed via dissipative interactions in small clusters are also likely to accrete material subsequently (Bonnell et al. 1997; Klessen et al. 1998). Finally, along with these direct effects of infalling material on misalignment, as mentioned earlier, binaries with separations $`\text{ }<100`$ a.u. are also unlikely to have strongly misaligned discs because the alignment timescale for close binaries ($`\text{ }<100`$ a.u.) is short in comparison to typical protostellar accretion timescales.
Conversely, if a binary is formed with discs that lie in the orbital plane, there remains the possibility that later infall of material with a different angular momentum vector to the binary’s orbit might lead to a degree of misalignment. In particular, it is simple for the spin of a disc to be affected by late infall of a small amount of material once the main accretion phase is over, and the disc masses have been substantially reduced. On the other hand, the magnitude of the misalignment which could occur during the main accretion phase of the stars in the process of formation, when the jet-like outflows are at their strongest, and when tidal forces too are at their greatest, has yet to be determined. Further numerical investigation of accretion and fragmentation scenarios are required to test this further.
Finally, as with the above mentioned effects of an encounter on the disc and/or jet of a single star, a change in the angular momentum vector of material being accreted by a single star could also lead to a change in the orientation of its disc and, thus, a wandering of an emanating jet. Once again, therefore, to unambiguously identify a precession, several oscillations must be observed.
## 5 Conclusions
We have investigated the dynamics of a protostellar disc in a binary system where the disc is misaligned with the orbital plane of the binary. The disc is found to wobble with a period approximately equal to half the binary’s orbital period and to precess with a period of order 20 binary periods. We also determine the characteristic timescale for realignment of the disc with the orbital plane due to dissipation. If the dissipation is determined by a simple isotropic viscosity then we find, in line with previous studies, that the alignment timescale is of order the viscous evolution timescale (of order 100 precession periods). However, for typical protostellar disc parameters, if the disc tilt exceeds the opening angle of the disc, then tidally induced shearing within the disc is transonic. In general, hydrodynamic instabilities associated with the internally driven shear result in extra dissipation which is expected to drastically reduce the alignment timescale. For large disc tilts the alignment timescale is then comparable to the precession timescale, while for smaller tilt angles $`\delta `$, the alignment timescale increases by a factor as $`(\mathrm{sin}\delta )^1`$.
These general results lead to several observational consequences. Since the alignment timescale in strongly misaligned discs is so short, such discs are only likely to occur in binaries with periods $`\text{ }>5000`$ years (i.e. separations $`\text{ }>100`$ a.u.). This expectation is in good agreement with the separations of binary systems which are observed to have misaligned jets. In addition, because the alignment timescale is of order the precession period, multiple large wiggles are not expected to be seen. At best one might observe a ‘bending’ of the jet as the disc simultaneously precesses and aligns to the orbital plane. However, such bending could also result from other processes (e.g. a single star whose disc, and thus jet, orientation is altered by interaction with a passing object or by the accretion of material with a different angular momentum vector than that of the original disc). For this and other reasons, we emphasize that to unambiguously identify a precessing jet several oscillations must be observed. Finally, although multiple large oscillations are not expected to be seen from systems with strongly misaligned discs, the jet direction may be forced to wobble by a few degrees on a timescale of half the orbital period, with an amplitude that depends on the details of the disc structure. Such a wobble may have been observed in one of the jets emanating from the Cep E system and we strongly encourage efforts to determine the binarity of sources which emit double jets or display evidence of wobbling jets so that the theoretical expectations can be tested.
For discs which are misaligned by small angles the alignment timescale is much longer. Therefore, precession of jets from such systems may be detectable for systems with orbital periods of order one hundred years (separations of a few tens of a.u.) in order to produce multiple wiggles in the $`10^4`$ year dynamical timescale of a jet. Wobbling of the jet on a timescale of half the orbital period is unlikely to be large enough to be detected in this case.
Finally, we discuss the implications of the existence of discs that are misaligned with the orbital plane of a binary for mechanisms for binary star formation. Although most published fragmentation calculations have not resulted in the formation of discs which are misaligned with the orbital plane of the binary, it is trivial to produce initial conditions where this is the case (e.g. Bonnell et al. 1992; Bonnell & Bastien 1992). Furthermore, these calculations do not require spatial variation of the direction of the angular momentum vector in the initial clouds; there is a single axis of rotation and even solid-body rotation is permitted. Binary systems with misaligned discs may also be formed directly via dissipative interactions in small clusters of protostars formed via ‘prompt initial fragmentation’. However, in either case, formation of the binary is likely to be contemporaneous with the accretion phase and, thus, strongly misaligned discs are unlikely for binaries with separations $`\text{ }<100`$ a.u. both because of the rapid realignment timescale and because the subsequent accretion of a large amount of material by the binary would tend to align the orbital and disc planes. Alternatively, binary systems with misaligned discs could be formed through the gravitational interaction of a passing object with a previously aligned disc, or by the infall of a small amount of material with a different angular momentum vector to that of the binary’s orbit near the end of accretion phase. Any of these mechanisms could explain the misaligned disc which may be present in HK Tau (Stapelfeldt et al. 1998; Koresko 1998) provided that the time that has elapsed since the misaligned disc was formed is less than the current alignment timescale. Thus, if HK Tau is a binary with a misaligned disc, this does not necessarily mean that it was formed this way initially or even that the discs were misaligned during the main accretion phase when the system would be expected to have produced jets. Larger surveys of the alignment of discs in binary systems are required for us to draw conclusions on formation mechanisms.
## Acknowledgments
We thank Jochen Eislöffel for his comments on the manuscript. This work was supported in part by NASA grant NAG5-4310 and the STScI visitor program.
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# Untitled Document
HUTP-00/A014 hep-th/0005059
From Noncommutative Bosonization to S-Duality
Carlos Nuñez, Kasper Olsen and Ricardo Schiappa
Department of Physics
Harvard University
Cambridge, MA 02138, U.S.A.
nunez, kolsen, ricardo@lorentz.harvard.edu
Abstract
We extend standard path–integral techniques of bosonization and duality to the setting of noncommutative geometry. We start by constructing the bosonization prescription for a free Dirac fermion living in the noncommutative plane $`𝐑_\theta ^2`$. We show that in this abelian situation the fermion theory is dual to a noncommutative Wess–Zumino–Witten model. The non–abelian situation is also constructed along very similar lines. We apply the techniques derived to the massive Thirring model on noncommutative $`𝐑_\theta ^2`$ and show that it is dualized to a noncommutative WZW model plus a noncommutative cosine potential (like in the noncommutative Sine–Gordon model). The coupling constants in the fermionic and bosonic models are related via strong–weak coupling duality. This is thus an explicit construction of $`S`$–duality in a noncommutative field theory.
May 2000
1. Introduction and Discussion
Quantum field theories on noncommutative spaces has been a subject of renewed interest since the recent discovery of its connections to string and $`M`$ theories, see e.g. and references therein. From a string theory point of view, it was realized in these works that one can translate the effects of a large background magnetic field into a deformation of the $`D`$–brane world–volume. Still, one can envisage studying such theories from a purely quantum field theoretic point of view. For example, perturbative aspects of such noncommutative field theories have been studied and have revealed a surprising mixing of the IR and the UV . These phenomena are directly related to the string theoretic origins of these theories, but one would also like to know to which extent properties of quantum field theories on commutative spaces also arise in quantum field theories on noncommutative spaces. This may be of some interest given that it is not always simple to extract quantum results from string theory, while we are used to do so in field theory.
One important feature of many conventional quantum field theories is that of duality. As an example, it follows from bosonization in 1+1 dimensions that the Sine–Gordon model of a single scalar field,
$$d^2x\{\frac{1}{2}_\mu \varphi ^\mu \varphi +\frac{\alpha _0}{\beta ^2}(\mathrm{cos}\beta \varphi 1)\},$$
is dual to the massive Thirring model of a fermion field,
$$d^2x\{\overline{\psi }(i\gamma ^\mu _\mu +m)\psi \frac{\lambda }{2}j_\mu j^\mu \}.$$
One can relate bosonic composite operators to fermionic ones and vice–versa using the standard bosonization machinery. Of particular interest to us in here is that the Sine–Gordon/Thirring model duality is a strong/weak coupling duality since the coupling constants of the two theories are related according to:
$$\frac{4\pi }{\beta ^2}=1+\frac{\lambda }{\pi }.$$
The purpose of this paper is to study the analog of this duality on a noncommutative spacetime. In order to do this, we begin by considering bosonization on the noncommutative plane and will see how the bosonization rules get generalized to this situation. This is done in section 2, where we study the abelian bosonization of a free fermion field in two noncommuting dimensions, employing path–integral techniques . We shall learn that the free fermion action is bosonized to a noncommutative $`U(1)`$ WZW–action. That the WZW term in the action is nonvanishing for a $`U(1)`$ valued field is simply due to the noncommutativity of spacetime. In fact, the procedure follows much as for the conventional non–abelian bosonization , and the rules are very similar both in the non–abelian and noncommutative cases. Because of this, the non–abelian noncommutative bosonization will be a simple standard extension of the abelian noncommutative bosonization. In particular, the non–abelian free fermion action bosonizes to a noncommutative $`U(N)`$ WZW model.
A question that immediately arises is the following. In the abelian case, the free fermion has a quadratic kinetic action and therefore noncommutative and commutative descriptions should match. On the other hand the commutative abelian fermion is dualized to a free scalar field theory, apparently very different from a noncommutative $`U(1)`$ WZW model. The same phenomena happens in the non–abelian situation, where the commutative abelian fermion is dualized to a commutative $`U(N)`$ WZW model, again not the same as a noncommutative non–abelian WZW model. As we shall see, the noncommutativity in the free fermion action makes its appearance when we gauge the global symmetries in order to implement the path–integral duality techniques of .
The fact that this happens raises some interesting possibilities for future research. Indeed this is apparently establishing some sort of relation between commutative and noncommutative WZW models, and one could interpret this as a different version of the Seiberg–Witten map between commutative and noncommutative descriptions of the Born–Infeld action . It would be very interesting to find a string theory realization of this field theoretic scenario, and try to understand this relation between WZW models from a kind of $`B`$–field point of view.
After understanding the free fermion we proceed in section 3 to interacting theories, with the goal of realizing $`S`$–duality for noncommutative field theories. We shall see that the massive Thirring model on noncommutative space is dual to a WZW model plus a noncommutative cosine potential. The usual relation (1.1) between the coupling constants of the dual theories continues to hold in the noncommutative case, thus realizing an example of $`S`$–duality. Observe that in here the knowledge of the bosonization rules for the noncommutative free fermion plays a central role, as they allow us to derive the noncommutative duality in very simple steps. Indeed they allow for a full and explicit quantum construction of $`S`$–duality in this noncommutative setting.
Understanding duality in noncommutative field theory, one could hope to gain some starting grounds in order to try to match these results to a string theory description. This is not a clear task, however. On one hand the results in are derived at the CFT disk level, so that one can not assume that an $`S`$–duality in noncommutative field theory will translate to a string theory $`S`$–duality. On the other hand, we are dealing with simple bosonic theories which do not have immediate brane realizations. After this paper was concluded, two pre–prints appeared that describe $`S`$–duality in noncommutative gauge theories , and therefore have a closer connection to the string theory description . What we would like to stress from our work is that not only it provides a construction which does not rely on any string theory connection, but it also allows for an explicit and exact treatment.
Some words on notation. To study noncommutative bosonization in $`𝐑_\theta ^2`$, the underlying $`𝐑_\theta ^2`$ will be labeled by noncommuting coordinates satisfying $`[x^\mu ,x^\nu ]=i\theta ^{\mu \nu }`$. Here $`\theta ^{\mu \nu }`$ is real and antisymmetric and so in two dimensions one has $`\theta ^{\mu \nu }=\theta ϵ^{\mu \nu }`$. The algebra of functions on noncommutative $`𝐑_\theta ^d`$ can be viewed as an algebra of ordinary functions on the usual $`𝐑^d`$ with the product deformed to the noncommutative, associative star product,
$$\left(\varphi _1\varphi _2\right)(x)=e^{\frac{i}{2}\theta ^{\mu \nu }_\mu ^y_\nu ^z}\varphi _1(y)\varphi _2(z)|_{y=z=x}.$$
Thus, we shall study theories whose fields are functions on ordinary $`𝐑^2`$, with actions of the usual form $`S=d^2x[\varphi ]`$, except that the fields in $``$ are multiplied using the star product. Moreover, for any noncommutative theory the quadratic part of the action is the same as in the commutative theory, since if $`f`$ and $`g`$ are functions that vanish rapidly enough at infinity,
$$d^dxfg=d^dxfg.$$
2. Bosonization on Noncommutative Space
In this section we derive the bosonization rules of the free fermion action on a two–dimensional noncommutative space , using the path–integral approach described in . Our derivation is carried out for the abelian case (and will largely follow ) but as we shall see it generalizes immediately to the non–abelian case . In it was shown that a conventional free fermionic theory in 1+1 dimensions is equivalent to a bosonic theory which is the WZW model. In our case the free fermionic theory is equivalent to a noncommutative version of the WZW model, even in abelian case.
2.1. The Noncommutative WZW Model
Before looking at our specific problem, we start by taking the usual WZW model and define the noncommutative extension in the obvious way:
$$S[g]=\frac{1}{8\pi }_\mathrm{\Sigma }d^2x(_\mu g_\mu g^1)\frac{i}{12\pi }_Bd^3xϵ^{ijk}(g^1_igg^1_jgg^1_kg).$$
If one would like to extend this action to the non–abelian situation, one simply needs to include a trace over the algebra. This theory has been discussed in a number of recent papers . The manifold $`\mathrm{\Sigma }`$ is parametrized by $`(x^0,x^1)`$ and is the boundary of the three–dimensional manifold $`B`$: $`B=\mathrm{\Sigma }`$. The $``$–product on $`B`$ is the trivial extension of the product on $`\mathrm{\Sigma }`$, i.e. the extra dimension $`x^2`$ is taken to be commutative. The commutative non–abelian WZW model obeys the important Polyakov–Wiegmann identity :
$$S[gh^1]=S[g]+S[h^1]\frac{1}{4\pi }_\mathrm{\Sigma }d^2x\mathrm{Tr}(g^1_+gh^1_{}h).$$
The same identity holds with a $``$–product in the noncommutative case since the identity follows from using the cyclic property of the integral $`AB=BA`$ and from $`gg^1=1`$. Here the group element $`g`$ of ”noncommutative” $`U(1)`$ is
$$g=e_{}^{i\alpha }=1+i\alpha \frac{1}{2}\alpha \alpha \frac{i}{6}\alpha \alpha \alpha +\mathrm{}.$$
In the abelian commutative case, where $`g=e^{i\alpha }`$ without any $``$–product, the action (2.1) of course reduces trivially to a free boson action $`d^2x(\alpha )^2`$. However, the WZW action is nontrivial even for the abelian noncommutative case, where $`g=e_{}^{i\alpha }`$.
2.2. Definition of the Current
Using the path–integral approach to bosonization one starts with the partition function of the free (abelian) fermion theory:
$$Z=𝒟\overline{\psi }𝒟\psi e^{{\scriptscriptstyle \overline{\psi }i/\psi }}.$$
We want to show that this theory is equivalent to a given bosonic theory. To prove that, one needs to show that the correlation functions obtained from the two theories are equal. We therefore consider the generating functional for these correlators as,
$$Z[s]=𝒟\overline{\psi }𝒟\psi e^{{\scriptscriptstyle }\overline{\psi }(i/\psi +s/)\psi },$$
where $`s_\mu `$ is an external source. Due to the noncommutativity of our problem, one might argue on where should one insert the source term, as for instance $`\overline{\psi }s/\psi s/\overline{\psi }\psi `$. As we shall see later, (2.1) is the definition we need in order to be able to carry out the dualization procedure using the results for the fermionic determinant in . So, we need to address the question of whether (2.1) is generating the correct conserved noncommutative current.
In the commutative fermi theory, the corresponding current is simply $`j^\mu =\overline{\psi }\gamma ^\mu \psi `$ and this is conserved because of the equations of motion, which are $`/\psi =0`$ and $`_\mu \overline{\psi }\gamma ^\mu =0`$. In the noncommutative case one finds similarly that the conserved current is:
$$j^\mu =\overline{\psi }\gamma ^\mu \psi =\overline{\psi }\gamma ^\mu \psi +\frac{i}{2}\theta ^{\mu \nu }_\mu \overline{\psi }\gamma ^\mu _\nu \psi +\mathrm{}.$$
Now consider the following source–term in the partition function:
$$d^2x\overline{\psi }s/\psi ,$$
where $`s_\mu `$ is a source. It is not immediately obvious that this will generate the correct correlation function, i.e. that
$$\frac{\delta }{\delta s_\mu (x)}Z[s]=j^\mu (x),$$
because of the infinite number of derivatives in the Moyal product. Let us see that this is actually true. Because of (1.1) one can remove for example the last $``$–product in (2.1) and then this integral equals,
$$d^2x\left(e^{\frac{i}{2}\theta ^{\mu \nu }_\mu ^y_\nu ^z}\overline{\psi }(y)s/(z)|_{y=z=x}\right)\psi (x).$$
The first order term in $`\theta `$ can be written as,
$$\frac{i}{2}\theta ^{\mu \nu }_\mu \overline{\psi }_\nu s/\psi =\frac{i}{2}\theta ^{\mu \nu }_\nu [_\mu \overline{\psi }s/\psi ]\frac{i}{2}\theta ^{\mu \nu }_\mu \overline{\psi }s/_\nu \psi .$$
The first term on the RHS is a total derivative and so vanishes under the integral sign. The second term on the RHS seems to have the wrong sign (and the same appears to be the case with all higher–order odd terms in $`\theta `$), since taking the functional derivative of the partition function with respect to $`s_\mu `$ will not lead to the current in (2.1). However, the second term in (2.1) also vanishes identically because of antisymmetry of $`\theta ^{\mu \nu }`$ and because of the following identity for Dirac fermions:
$$\overline{\chi }\gamma ^\mu \psi =(\overline{\psi }\gamma ^\mu \chi )^{},$$
which ensures that the current is real. Namely,
$$\frac{i}{2}\theta ^{\mu \nu }_\mu \overline{\psi }s/_\nu \psi =\frac{i}{2}\theta ^{\mu \nu }(_\nu \overline{\psi }s/_\mu \psi )^{}=\frac{i}{2}\theta ^{\mu \nu }(_\mu \overline{\psi }s/_\nu \psi )^{}=\frac{i}{2}\theta ^{\mu \nu }_\mu \overline{\psi }s/_\nu \psi ,$$
with the source $`s_\mu `$ being real. All higher–order terms with odd number of $`\theta `$’s vanish for the same reason, e.g. the third order terms is – up to total derivatives – of the form,
$$\theta ^{\mu \nu }\theta ^{\alpha \beta }\theta ^{\gamma \delta }_\mu _\alpha _\gamma \overline{\psi }s/_\delta _\beta _\nu \psi $$
and vanishes identically. This shows that the term in (2.1) does indeed generate the correct current.
2.3. Path-Integral Derivation
In the following we will use the fact that the generating functional in equation (2.1) is gauge invariant, i.e.
$$Z[s]=Z[s^g],$$
where under a gauge transformation the source transforms according to
$$s_\mu s_\mu ^g=g^1s_\mu g+g^1_\mu g.$$
This invariance follows from the invariance of the measure under local transformations of the fermion fields, i.e. transformations $`\psi e_{}^{i\alpha }\psi =g\psi `$, together with $`\overline{\psi }\overline{\psi }e_{}^{i\alpha }=\overline{\psi }g^1`$. From (2.1) we have:
$$Z[s]=𝒟\overline{\psi }𝒟\psi 𝒟ge^{{\scriptscriptstyle }\overline{\psi }(i/\psi +s/^g)\psi }=𝒟g\mathrm{det}_{}(i/+s/^g),$$
where the last equality is obtained after integrating out fermions. Note that the determinant is evaluated with respect to the $``$–product (we will shortly use the fact that this determinant was computed in ). Introduce the connection
$$b_\mu =s_\mu ^g,$$
such that the field strengths of $`b_\mu `$ and $`s_\mu `$ are related according to:
$$f_{\mu \nu }[b]=g^1f_{\mu \nu }[s]g.$$
We will choose a gauge where $`b_+=s_+`$, with $`\mathrm{\Delta }_{FP}`$ being the corresponding Faddeev–Popov determinant (we have $`\mathrm{\Delta }_{FP}=\mathrm{det}_{}D_+[s_+]`$, where $`D_+=_++i[s_+,]_{}`$). This allows us to write (2.1) in the form,
$$Z[s]=𝒟b_\mu \mathrm{det}_{}(i/+b/)\delta [ϵ_{\mu \nu }(f_{\mu \nu }[b]f_{\mu \nu }[s])]\delta [b_+s_+]\mathrm{\Delta }_{FP}.$$
By introducing a Lagrange–multiplier field $`\widehat{a}`$ that lives in the “noncommutative” $`U(1)`$ group with gauge transformation $`\widehat{a}\widehat{a}^g=g^1\widehat{a}g`$ one can write,
$$Z[s]=𝒟b_\mu 𝒟\widehat{a}\mathrm{det}_{}(i/+b/)e^{\xi {\scriptscriptstyle \widehat{a}(f_{\mu \nu }[b]f_{\mu \nu }[s])}}\delta [b_+s_+]\mathrm{\Delta }_{FP},$$
where $`\xi `$ is a constant which will be conveniently determined later. Now, make the following change of variables:
$$\begin{array}{cc}& s_+=i\stackrel{~}{s}^1_+\stackrel{~}{s},\hfill \\ & s_{}=is_{}s^1,\hfill \\ & b_+=i(\stackrel{~}{b}\stackrel{~}{s})^1_{}(\stackrel{~}{b}\stackrel{~}{s}),\hfill \\ & b_{}=(sb)_{}(b^1s^1).\hfill \end{array}$$
As we stated before, the fermion determinant for the noncommutative $`U(1)`$ theory has been calculated in with the result that it is:
$$\mathrm{det}_{}(i/+a/)=\mathrm{exp}S_{WZW}[hg],$$
where $`a_+=h^1_+h`$ and $`a_{}=g_{}g^1`$. The action for the noncommutative WZW model on the RHS is given in equation (2.1). With this result we can express the fermion determinant in terms of the variables in (2.1):
$$\mathrm{det}_{}(i/+b/)=\mathrm{exp}S_{WZW}[\stackrel{~}{b}\stackrel{~}{s}sb].$$
The Jacobian for the change of variables $`(b_+,b_{})(b,\stackrel{~}{b})`$ gives
$$𝒟b_+𝒟b_{}=\mathrm{det}_{}D_+[\stackrel{~}{b}\stackrel{~}{s}]\mathrm{det}_{}D_{}[sb]𝒟b𝒟\stackrel{~}{b}=\mathrm{exp}(\eta S_{WZW}[\stackrel{~}{b}\stackrel{~}{s}sb])𝒟b𝒟\stackrel{~}{b},$$
where we recall that the covariant derivatives $`D_\pm `$ are now in the adjoint representation. Therefore the result for their determinant is the same as for the fundamental representation but with an extra factor, $`\eta `$, that accounts for the change in representation . This factor can actually be computed to be related to the Casimir in the commutative case, but as we shall never need it we simply leave it as $`\eta `$. Furthermore, with this change of variables, one can write the $`\delta `$–function in (2.1) as:
$$\delta [b_+s_+]=\frac{1}{\mathrm{det}_{}D_+[s_+]}\delta [b1].$$
Combining these two results one obtains,
$$Z[s]=𝒟b𝒟\stackrel{~}{b}𝒟\widehat{a}\mathrm{exp}\left(S_{WZW}[\stackrel{~}{b}\stackrel{~}{s}sb]\right)\mathrm{exp}\left[\xi d^2x\widehat{a}(f_+[b]f_+[s])\right]\delta [\stackrel{~}{b}1].$$
In the gauge $`b_+=s_+`$ we have $`f_+[b]f_+[s]=D_+[s_+](b_{}s_{})`$ and so this gives,
$$\begin{array}{cc}& Z[s]=𝒟b𝒟\widehat{a}\mathrm{exp}\left((1+\eta )S_{WZW}[\stackrel{~}{s}sb]\right)\hfill \\ & \mathrm{exp}(\xi d^2xD_+[s_+]\widehat{a}(isb_{}b^1s^1)).\hfill \end{array}$$
Note that the expression for the generating functional (2.1) is gauge invariant, the transformation laws for $`s,\stackrel{~}{s}`$ being $`\stackrel{~}{s}\stackrel{~}{s}g`$ and $`sg^1s`$. One further change of variables, from $`\widehat{a}`$ to a group valued variable $`a`$ is defined as follows: $`D_+[\stackrel{~}{s}]\widehat{a}=i\stackrel{~}{s}^1(a^1_+a)\stackrel{~}{s}`$ (note that $`a`$ is the bose field equivalent to the original fermi field and will be invariant under gauge transformations). The Jacobian for the change of variables from $`\widehat{a}`$ to $`a`$,
$$\frac{\mathrm{det}_{}D_+[a\stackrel{~}{s}]}{\mathrm{det}_{}D_+[\stackrel{~}{s}]},$$
is, however, not gauge invariant. The trick is to use instead the following Jacobian obtained from the above by multiplying with (formally) one:
$$\frac{\mathrm{det}_{}D_+[a\stackrel{~}{s}]}{\mathrm{det}_{}D_+[\stackrel{~}{s}]}\frac{\mathrm{det}_{}D_{}[s]}{\mathrm{det}_{}D_{}[s]}=\mathrm{exp}(\eta S_{WZW}[a\stackrel{~}{s}s])\mathrm{exp}(\eta S_{WZW}[\stackrel{~}{s}s]).$$
From this one finally obtains,
$$\begin{array}{cc}& Z[s]=𝒟a𝒟b\mathrm{exp}((1+\eta )S_{WZW}[\stackrel{~}{s}sb]+\eta S_{WZW}[a\stackrel{~}{s}s]\eta S_{WZW}[\stackrel{~}{s}s]\hfill \\ & +\xi d^2x\stackrel{~}{s}^1a^1_+a\stackrel{~}{s}sb_{}b^1s^1).\hfill \end{array}$$
We can now apply the Polyakov-Wiegmann identity (2.1) for the noncommutative WZW model. Also we choose the up to now arbitrary value of $`\xi `$ to be $`\xi =\frac{1}{4\pi }(1+\eta )`$. This gives the following result for the partition function:
$$Z[s]=𝒟a𝒟b\mathrm{exp}\left((1+\eta )S_{WZW}[a\stackrel{~}{s}sb]S_{WZW}[a\stackrel{~}{s}s]+S_{WZW}[\stackrel{~}{s}s]\right).$$
Now, make a change of variables $`b\widehat{b}=a\stackrel{~}{s}sb`$ with trivial Jacobian. Then the $`𝒟\widehat{b}`$–integration factors and it is just a normalization contribution, so that one obtains the simpler expression,
$$Z[s]=𝒟a\mathrm{exp}\left(S_{WZW}[a\stackrel{~}{s}s]+S_{WZW}[\stackrel{~}{s}s]\right).$$
One final change of variables with trivial Jacobian is $`a\stackrel{~}{s}s\stackrel{~}{s}as`$; together with the Polyakov-Wiegmann identity this leads to our final result for the bosonization of the free fermion action in equation (2.1) (renaming $`a`$ as $`g`$):
$$\begin{array}{cc}& Z[s]=𝒟g\mathrm{exp}[S_{WZW}[g]_{}\hfill \\ & \frac{1}{4\pi }d^2x(s_+s_{}s_+gs_{}g^1ig^1_+gs_{}is_+g_{}g^1)].\hfill \end{array}$$
Our prescription for the noncommutative currents becomes:
$$\begin{array}{cc}& \overline{\psi }\gamma _+\psi \frac{i}{4\pi }g^1_+g,\hfill \\ & \overline{\psi }\gamma _{}\psi \frac{i}{4\pi }g_{}g^1.\hfill \end{array}$$
This derivation shows that – in the abelian case – the free fermion action on $`𝐑_\theta ^2`$ is bosonized to the noncommutative $`U(1)`$ WZW model. From this result it also follows what should be done for the non–abelian system. In the non–abelian case, $`g`$ in equation (2.1) belongs to the ”noncommutative” $`U(N)`$ group, i.e. $`g=e_{}^{i\alpha ^aT^a}`$ and one should therefore include the ordinary trace over the $`N\times N`$ matrices in the appropriate places, as in (2.1). When this is done one immediately realizes that the above derivation goes through without any changes, except of course that the group elements now belong to $`U(N)`$ (note that the evaluation of the noncommutative fermion determinant in also applies to the $`U(N)`$ case). This shows that in the non–abelian case the free fermionic theory is dual to a noncommutative $`U(N)`$ WZW model.
With these results in hand, one is lead to make the following observation. The free fermionic action (in both abelian and non–abelian cases) in the noncommutative plane is the same as in the commutative plane, due to its quadratic nature. On the other hand, the standard commutative free fermionic action is equivalent to the commutative WZW model , and in the abelian case in particular it is equivalent to a theory of a free scalar field. This shows an equivalence between commutative and noncommutative WZW models, and the map between these two models might be some version of the Seiberg-Witten map between ordinary Yang–Mills theory and noncommutative Yang–Mills theory. It would be very interesting to further explore and make more precise this relation.
3. Noncommutative S–Duality
In this section we discuss a noncommutative version of the well–known duality between the Sine–Gordon model and the massive Thirring model. As we have just seen, the free fermion theory discussed in section 2 is dual to a noncommutative WZW model. The next step is to study an interacting fermionic system, where the natural candidate is the Thirring model. We shall see that we can dualize this theory in a straightforward manner, given the results of the previous section. Moreover, we will unravel a strong/weak coupling duality in the procedure.
3.1. The Thirring Model
We consider the Thirring model with the usual quartic coupling:
$$S_\lambda =\frac{\lambda }{2}d^2x(\overline{\psi }\gamma ^\mu \psi )(\overline{\psi }\gamma _\mu \psi )=\frac{\lambda }{2}d^2xj^\mu j_\mu .$$
In order to bosonize this term one can either use the bosonization prescription directly (as we know how to bosonize the currents) or do it via a ”completing the square” type of prescription at the path–integral level. The two methods obviously yield the same result and we simply use the first. Using the bosonization recipe of section 2 one immediately obtains from equation (2.1) that the four–fermion interaction (3.1) corresponds to the following term in the bose theory:
$$\frac{1}{2}\lambda d^2x2(\overline{\psi }\gamma _+\psi )(\overline{\psi }\gamma _{}\psi )\frac{\lambda }{16\pi ^2}d^2xg^1_+gg_{}g^1.$$
This term is quartic and cannot be made quadratic. This is unlike the commutative theory, where the bosonized current coupling term becomes an extra contribution to the quadratic kinetic term in the scalar action. In here, by including the quartic coupling (3.1) the resulting noncommutative theory becomes significantly different from the corresponding commutative theory by the introduction of an infinite series of derivative terms in the $``$–product. The bosonized Lagrangian therefore becomes:
$$S_{WZW}+\frac{\lambda }{16\pi ^2}d^2xg^1_+gg_{}g^1.$$
To identify the bosonized fermion theory with a certain bose theory, one still needs a canonically–normalized scalar variable . Recall that we are dealing with noncommutative $`U(1)`$ group elements, so that one needs to look at scalar fields $`\mathrm{\Lambda }(x)`$ appearing as $`g(x)=e_{}^{i\mathrm{\Lambda }(x)}`$. From the WZW action (2.1) one has the term,
$$\frac{1}{4\pi }d^2x_+\mathrm{\Lambda }_{}\mathrm{\Lambda }+\mathrm{},$$
while from the bosonized Thirring coupling (3.1) we get a term,
$$\frac{\lambda }{16\pi ^2}d^2x_+\mathrm{\Lambda }_{}\mathrm{\Lambda }+\mathrm{}.$$
By adding these two contributions we get the following kinetic term for the scalar field:
$$\frac{1}{4\pi }(1+\frac{\lambda }{4\pi })d^2x_+\mathrm{\Lambda }_{}\mathrm{\Lambda },$$
and so the canonically normalized scalar variable is:
$$\varphi =\left[\frac{1}{4\pi }(1+\frac{\lambda }{4\pi })\right]^{1/2}\mathrm{\Lambda }.$$
In particular, stability of the bosonic theory requires $`\lambda >4\pi `$. This result (3.1) will shortly turn out to be important in determining the relation between the couplings of the noncommutative “Sine–Gordon”<sup>1</sup> Observe that we will not actually have a simple noncommutative extension of the Sine–Gordon model due to the extra term appearing in (3.1). and Thirring models.
3.2. The Massive Thirring Model
We shall now turn to the fermion mass coupling. The relevant term is
$$S_m=md^2x\overline{\psi }\psi .$$
In order to extend the previous discussion to the massive Thirring model we follow the procedure outlined in , as it applies to our situation. Indeed one can bosonize the mass term by considerations of chiral symmetry alone, and the discussion in directly applies in here as well (as one is considering global chiral symmetry for which the $``$–product collapses to the standard product).
Free fermions are invariant under the global $`U_A(1)`$ axial symmetry $`\psi e^{i\alpha \gamma _5}\psi `$ and $`\overline{\psi }\overline{\psi }e^{i\alpha \gamma _5}`$, and one expects that the bosonic theory will share such a symmetry . On the other hand, under the duality procedure of the previous section (see also ) one is gauging the global vector symmetry of the free fermions and the axial symmetry will not survive quantization due to the presence of the background gauge field. This axial anomaly in the noncommutative plane was computed in ,
$$_\mu j_5^\mu =\frac{1}{2\pi }ϵ^{\mu \nu }\widehat{F}_{\mu \nu },$$
where the axial current is $`j_5^\mu =\overline{\psi }\gamma ^\mu \gamma ^5\psi `$ and $`\widehat{F}_{\mu \nu }`$ is the noncommutative gauge field strength. Because both initial and final theories (under duality) share the same symmetry, it may seem odd that this symmetry is broken under the duality procedure. The solution to make it manifest throughout the dualization procedure is to include a transformation on the Lagrange multiplier as well . Indeed, recall that the Lagrange multiplier field $`\mathrm{\Lambda }`$ appears in the initial gauged action as,
$$\mathrm{exp}\left(d^2x\frac{1}{2\pi }\mathrm{\Lambda }ϵ^{\mu \nu }\widehat{F}_{\mu \nu }\right).$$
The observation is that if the field $`\mathrm{\Lambda }`$ transforms as $`\mathrm{\Lambda }\mathrm{\Lambda }\alpha `$ under axial transformations, then the Lagrange multiplier term (3.1) will cancel the axial anomaly term (3.1)as it appears in the path–integral, and chiral symmetry is made manifest. In summary, what we have done is to nail down what is the correct transformation of the new field $`\mathrm{\Lambda }`$ under chiral rotations, and this is uniquely defined by the previous considerations (for details see ).
The bottom line is that one can now proceed to deduce the bosonization of fermionic mass terms from these transformation properties. Indeed, in the presence of a mass term the fermi theory is no longer axial symmetric, but the axial transformation rules nevertheless remain the same. As we shall see in the following, these rules alone are enough information for a unique determination of the bosonized composite operator that corresponds to the fermionic quadratic term. For instance, we would like to bosonize a term as $`\psi _R^{}\psi _L`$, which amounts to finding an appropriate bosonic functional $`(\mathrm{\Lambda })`$ such that $`(\mathrm{\Lambda })\psi _R^{}\psi _L`$. Under global chiral transformations the chiral mass term transforms as,
$$\psi _R^{}\psi _Le^{2i\alpha }\psi _R^{}\psi _L,$$
and due to the $`\mathrm{\Lambda }`$ chiral transformation rule just deduced, it follows that the bosonic functional must satisfy,
$$(\mathrm{\Lambda }\alpha )=e^{2i\alpha }(\mathrm{\Lambda }),$$
where we recall that due to the global character of the rotation, the exponential involves no $``$–product, even though the functional $``$ will be defined in terms of the field $`\mathrm{\Lambda }`$ through a $``$–product. Indeed, it immediately follows that
$$(\mathrm{\Lambda })e_{}^{2i\mathrm{\Lambda }}$$
uniquely solves the functional equation (where the exponential is properly defined with the $``$–product). This whole procedure naturally follows from as one is dealing with global chiral rotations. What all this amounts to, is that the mass term (3.1) bosonizes to:
$$S_m=d^2xm\alpha _0\mathrm{cos}_{}2\mathrm{\Lambda },$$
where $`\alpha _0`$ is a constant associated to the zero point energy (as in the commutative case ), and the noncommutative cosine is defined naturally by $`\mathrm{cos}_{}\phi =\frac{1}{2}(e_{}^{i\phi }+e_{}^{i\phi })`$. This story is very similar to the commutative Sine–Gordon/Thirring duality. In particular, the coupling constant in the bosonic theory is defined as $`\beta `$ and appears in the action exactly through a cosine potential. In the noncommutative case, this would be,
$$d^2x\alpha \mathrm{cos}_{}\beta \varphi $$
where the field $`\varphi `$ is the canonically normalized field. Finally, since the canonically normalized field was $`\varphi =\sqrt{\frac{1}{4\pi }(1+\frac{\lambda }{4\pi })}\mathrm{\Lambda }`$ we obtain the following $`S`$–duality property:
$$\frac{16\pi }{\beta ^2}=1+\frac{\lambda }{4\pi }.$$
This is exactly the same type of relation as that of the commutative theory (1.1). At first glance, one could worry that when we plug this relation (3.1) back in (3.1) there could be terms going like inverse powers of $`\beta `$ which would spoil the strong/weak coupling duality. However, this does not happen as one should also recall that the scalar fields have to be canonically normalized according to (3.1). What therefore happens is that an expansion of (3.1) in powers of the canonically normalized field $`\varphi `$ will only produce interaction terms proportional to positive powers of the bosonic coupling constant $`\beta `$. From this we see that the strong/weak coupling duality between the Sine–Gordon model and the massive Thirring model survives on a noncommutative space, as we have just shown.
It is known that the $`T`$–duality of string theory can be interpreted as Morita equivalence in noncommutative geometry . One can wonder if there could be a similar geometrical interpretation of string theory $`S`$–duality in terms of noncommutative geometry. In here we have given some first steps, by showing that one can construct quantum field theoretical models in noncommutative space displaying $`S`$–duality.
Acknowledgements: We have benefitted from discussions/correspondence with C-S. Chu, L. Cornalba, K. Hori, A. Matusis, S. Minwalla, E. Moreno, B. Pioline, F. Schaposnik and H. J. Schnitzer. CN is supported by CONICET. KO is supported by the Danish Natural Science Research Council. RS is supported by the Fundação para a Ciência e Tecnologia, under the grant Praxis XXI BPD-17225/98 (Portugal).
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# Shear response of a frictional interface to a normal load modulation
## I Introduction
Friction between solids carrying a time-dependent normal load is a subject of interest in different fields, from mechanical engineering, where the “friction-lowering” effect of external vibrations is well known and commonly used in applications, to geophysical studies of the effect of rapid stress changes on static and dynamic friction of rocks, aiming at a better understanding of the coupling between normal and tangential stress states on slipping faults.
These studies involve multicontact interfaces (MCI’s), i.e. interfaces between macroscopic solids with rough surfaces. The real area of contact thus consists of a large number of small contacts with sizes on the micrometer scale.
In a situation of constant normal load on the MCI, the phenomenological state- and rate-dependent friction (SRF) model, formulated by Rice and Ruina, successfully describes the details of the low-velocity dynamics (typically in the 0.01 – 100 $`\mu `$m.$`s^1`$ range) of such systems, such as the bifurcation between steady sate and stick-slip oscillations. The model states that the dynamic friction force $`F_{fr}`$ depends on the instantaneous sliding velocity $`\dot{x}`$ and on a dynamic state variable $`\mathrm{\Phi }`$ as:
$$F_{fr}(\dot{x},\mathrm{\Phi })=W\left[\mu _0+A\mathrm{ln}\left(\frac{\dot{x}}{V_0}\right)+B\mathrm{ln}\left(\frac{V_0\mathrm{\Phi }}{D_0}\right)\right],$$
(1)
where $`\mu _0`$ is the dynamic friction coefficient in steady sliding at the reference velocity $`V_0`$, and A and B are measured to be positive and of typical order $`10^2`$ (with $`B>A`$).
The state variable $`\mathrm{\Phi }`$ can be interpreted as the “age” of the MCI, i.e. the average duration of the transient contact between load bearing asperities. For example, in stationary sliding at velocity $`V`$, the set of microcontacts is destroyed and replaced by a new one over a characteristic sliding length $`D_0`$, and the state variable thus expresses as $`\mathrm{\Phi }=D_0/V`$.
More generally, the model specifies the time evolution law of $`\mathrm{\Phi }`$ as:
$$\dot{\mathrm{\Phi }}=1\frac{\dot{x}\mathrm{\Phi }}{D_0}$$
(2)
In Eq. (1) the two corrections to $`\mu _0`$ have distinct physical meanings: the first term describes an instantaneous velocity-strengthening of the interface, while the second expresses strengthening of the interface with its “age”, which in stationary sliding, where $`\mathrm{\Phi }=D_0/V`$, leads to a velocity-weakening effect.
In the spirit of the Bowden and Tabor analysis, one can write, for a MCI, the friction force as:
$$F_{fr}=\sigma _s\left(\dot{x}\right)\mathrm{\Sigma }_r(\mathrm{\Phi },W)$$
(3)
where $`\sigma _s`$ defines an interfacial shear strength, $`\mathrm{\Sigma }_r`$ is the real area of contact between the solids, and $`W`$ is the normal load carried by the multicontact interface. The age-strengthening effect is associated with the creep growth of the microcontact area under normal load:
$$\mathrm{\Sigma }_r(\mathrm{\Phi },W)=\mathrm{\Sigma }_0\left(W\right)\left[1+m\mathrm{ln}\left(\frac{\mathrm{\Phi }V_0}{D_0}\right)\right],$$
(4)
$`\mathrm{\Sigma }_0`$ — the real area of contact at $`\mathrm{\Phi }=D_0/V_0`$ — exhibits a linear dependence on $`W`$, as explained by Greenwood and Williamson’s model of contact between rough surfaces. That is, the friction force obeys the Amontons law $`F_{fr}W`$.
The velocity-dependent interfacial strength of the interface is described as:
$$\sigma _s\left(\dot{x}\right)=\sigma _{s0}\left[1+\eta \mathrm{ln}\left(\frac{\dot{x}}{V_0}\right)\right],$$
(5)
This form for the interface “rheology”, discussed in detail in reference, results from the thermally activated depinning of multistable nanometric units localized in a layer of nanometric thickness forming a junction between micrometric asperities.
Eqs. (4) and (5) yield Eq. (1) with $`\mu _0=\sigma _{s0}\mathrm{\Sigma }_0/W`$, $`m=B/\mu _0`$, $`\eta =A/\mu _0`$, and since $`m,\eta 1`$, non-linear logarithmic terms can be neglected.
The SRF model, and its physical interpretation presented above, have been validated by friction experiments on different classes of contacting materials, namely granite, paper, polymer glasses and elastomers, under constant normal load applied to the solids.
In the case of a time-dependent normal load, one can first note that in the Amontons-Coulomb description ($`F_{fr}=\mu W`$, with constant $`\mu `$), a change in $`W`$ would lead to a proportional change in $`F_{fr}`$, in particular a harmonic normal load modulation $`W=W_0\left(1+ϵcos\left(\omega t\right)\right)`$ would produce a harmonic frictional modulation about a non-modified average value $`\mu W_0`$.
In the SRF framework, the variations of $`\dot{x}`$ and $`\mathrm{\Phi }`$ are non-linearly coupled, through Eqs. (1) and (2), to the load modulation, thus resulting in non-trivial effects on the friction force (such as, for instance, an anharmonic response to a harmonic normal load). However, the model as expressed by Eqs. (1) and (2) may not be sufficient to describe correctly the frictional response for the following reasons:
(i) the interface rheology expressed by Eq. (5) may not hold for “fast” changes of $`W`$,
(ii) the load variation may modify the interface age strengthening process, thus leading to changes in the evolution of $`\mathrm{\Sigma }_r`$ with $`\mathrm{\Phi }`$, or in the evolution law of the state variable $`\mathrm{\Phi }`$ itself.
Based on their results on the response to normal stress steps and pulses in granite friction experiments, Linker and Dieterich suggested to modify the evolution law of $`\mathrm{\Phi }`$, while retaining the functional form (4) of the $`\mathrm{\Phi }`$ dependence of $`\mathrm{\Sigma }_r`$. Argueing that a sudden change in normal stress would result in a sudden change in $`\mathrm{\Phi }`$, they postulate:
$$\dot{\mathrm{\Phi }}=1\frac{\dot{x}\mathrm{\Phi }}{D_0}\frac{\alpha \dot{\sigma }}{B\sigma }\mathrm{\Phi },$$
(6)
where they infer $`\alpha =0.2`$, for granite, from their analysis of the response to sudden normal load steps.
In a recent study, Richardson and Marone investigated the influence of normal stress modulations on the so-called “frictional healing” effect in a granular material layer confined between rough granite blocks: starting from steady sliding, shear loading is stopped and the subsequent shear stress relaxation is measured in presence of a 1 Hz normal load modulation (a modulation frequency close to the characteristic stick-slip oscillations frequency that can be inferred for their system).
Friction experiments with confined granular media have been successfully described by the SRF model in situations of constant normal load (although the physical meaning of the variable $`\mathrm{\Phi }`$ is not clear yet for these systems). However, the use of the constitutive law proposed by Linker and Dieterich to include time-dependence of the normal load did not account properly for the details of the results obtained by Richardson and Marone.
In this paper we present an extensive study of the effect of a harmonic modulation of the normal load, $`W=W_0\left(1+ϵ\mathrm{cos}\left(\omega t\right)\right)`$, on the dynamic frictional response of a multicontact interface. Experiments are conducted on an interface between two blocks of poly(methyl-metacrylate) (PMMA), at velocities $`V<100\mu m.s^1`$, load modulation frequency $`f=120HzV/(2\pi D_0)`$ and relative amplitude $`ϵ`$ in the range $`5.10^3\mathrm{\hspace{0.17em}0.5}`$ (so that no loss of contact between the surfaces occurs).
We study quantitatively the average $`\overline{F}`$ and the components at frequency $`f`$ and $`2f`$, $`F_1`$ and $`F_2`$, of the tangential pulling force:
$$F=K(Vtx)=m\ddot{x}+F_{fr}$$
(7)
for different values of $`V`$ and $`ϵ`$ ; these results are presented in section II. We find in particular that the modulation of $`W`$ induces a systematic decrease of the average dynamic friction coefficient $`\overline{\mu }=\overline{F}/W_0`$. This effect, which increases with higher $`ϵ`$, is quite substantial: a typical magnitude of this effect is a $`20\%`$ decrease of $`\overline{\mu }`$ for $`ϵ=0.5`$.
To analyze quantitatively our experimental data, we need to evaluate which fraction, $`ϵ_{eff}/ϵ`$, of the load modulation is effectively borne by the microcontacts. Indeed, the normal load modulation is too fast for air to be drained in and out of the micrometer-thick interfacial gap. We have studied this “leaking air cushion” effect by conducting similar experiments under primary vacuum. From these experiments we infer that $`ϵ_{eff}/ϵ0.4`$, and we use this in the subsequent analysis as a scaling factor for the modulation amplitude.
Section III is devoted to the analysis of these results in terms of the SRF model and its possible extensions to fast load modulations:
(i) We first test the unmodified SRF model by setting in Eq. (1) $`WW(t)`$ and using the evolution law (2) for the state variable $`\mathrm{\Phi }`$. Numerical integration of these equations leads to a quantitatively good prediction of the average friction force $`\overline{F}\left(ϵ\right)`$. However, the predicted oscillating tangential force components $`F_1`$ and $`F_2`$ strongly depart from the observed dependences on $`ϵ`$ and $`V`$.
(ii) We have then tested the proposition of Linker and Dieterich. Using their evolution law (6) and their proposed value of $`\alpha =0.2`$, we find that (i) the decrease of $`\overline{F}(ϵ)`$ is much smaller than the measured values and (ii) the agreement for $`F_1`$ and $`F_2`$ is not better than in the previous $`\alpha =0`$ test.
An attempt with a value of $`\alpha `$ small enough ($`\alpha =0.02`$) to describe correctly the $`ϵ`$ dependence of $`\overline{F}`$ leads to results close to those obtained with the basic SRF equations; this also holds for $`F_1`$ and $`F_2`$. That is, as confirmed by a perturbation calculation in $`ϵ`$, our experiments are not discriminating with respect to the Linker-Dieterich evolution law for such small values of $`\alpha `$.
So, this modified ageing law, even if valid, does not suffice to account properly for the details of the frictional response.
(iii) We propose to modify expression (5) for the following physical reason: we know from static measurements that a MCI exhibits, at shear forces much smaller than the static threshold, an elastic tangential response. One can deduce from this a shear stiffness $`\kappa _{asp}`$ with the particular feature $`\kappa _{asp}W`$. Now, in our interpretation of friction, the rate variable appearing in $`\sigma _s`$ must be the true rate of irreversible (plastic) strain of the interfacial junction of nanometer thickness $`h`$. When taking into account the asperity elasticity $`\kappa _{asp}`$, strictly speaking, this strain rate reads $`h^1d\left(xF/\kappa _{asp}\right)/dt`$. In quasi stationary motion, this reduces to the $`\dot{x}/h`$ strain rate, hence the usual $`\sigma _s(\dot{x})`$ expression. In the present experimental situation, $`\kappa _{asp}`$ is modulated as $`W`$ itself, and the difference between the total and plastic strain rates becomes relevant. Indeed, we show that this extended phenomenological elastoplastic generalization of interfacial dissipation leads to a very satisfactory description of the average and oscillating shear responses to fast normal load modulations.
## II Experiments and results
### A Experimental setup and methods
The tribometer is composed of a slider of mass $`M`$ driven along a track through a loading spring of stiffness $`K`$, one end of which is pulled at constant velocity $`V`$, as schematized in the inset of Fig. 1. The slider and track are made of PMMA with nominally flat surfaces lapped with SiC powder to a roughness of order $`1\mu `$m, thus forming a multicontact interface.
A detailed drawing of the setup is given in Fig. 1. We impose the velocity $`V`$ of the loading point, in the range $`0.1\mathrm{\hspace{0.17em}100}\mu `$m.$`s^1`$, by means of a translation stage driven by a stepping motor. The tangential load is applied on the slider through a leaf spring of stiffness $`K=0.2`$ N.$`\mu `$m<sup>-1</sup>, which is the more compliant part of the system. The dead weight of the slider is 16 N. The average normal load $`W_0`$ can be set in the range $`316`$ N with the help of a vertical spring attached to a remote point itself translated horizontally at the pulling velocity $`V`$ through a second translation stage, in order to prevent any tangential coupling.
The normal load modulation is achieved by means of a vibration exciter rigidly attached to the slider: a harmonic voltage input of given amplitude and frequency $`f`$ results in a harmonic vertical motion of the moving element of the exciter on which an accelerometer is fixed. An acceleration of amplitude $`\gamma `$ of this moving element of mass $`m`$ induces a normal load modulation on the slider of amplitude $`m\gamma `$ at frequency $`f`$. We thus obtain a normal load $`W=W_0\left(1+ϵ\mathrm{cos}(\omega t)\right)`$ with $`\omega =2\pi f`$ and $`ϵ=m\gamma /W_0`$ in the range $`5.10^3\mathrm{\hspace{0.17em}0.5}`$.
We use the loading leaf spring as a dynamometer by measuring its deflection $`\mathrm{\Delta }X`$ by means of an eddy current displacement gauge. The tangential force applied to the slider is thus $`F=K\mathrm{\Delta }X`$. We measure the average value of the output voltage of the gauge, and use a lock-in amplifier to measure the amplitude of the first and second harmonics of this output signal with respect to the harmonic excitation signal. We thus characterize the shear force through its average value $`\overline{F}`$ and its A.C. components at frequency $`f`$ and $`2f`$, of respective amplitudes $`F_1`$ and $`F_2`$.
The experiments are conducted according to the following protocole: for a fixed set of parameter values $`W_0`$ and $`V`$ leading to steady sliding when $`ϵ=0`$, we measure $`\overline{\mu }\left(0\right)=F\left(ϵ=0\right)/W_0`$. The normal load modulation is then set at amplitude $`ϵ`$, while sliding, and shear force measurements yield $`\overline{\mu }=\overline{F}/W_0`$, $`\mu _1=\left|F_1\right|/W_0`$ and $`\mu _2=\left|F_2\right|/W_0`$. The modulation is then switched off and $`F\left(ϵ=0\right)`$ is systematically remeasured before setting a new value of $`ϵ`$, in order to check that no drift occurred during the measurement. Moreover, we check that for $`ϵ0`$ the shear force signal does not exhibit low-frequency stick-slip oscillations.
The experimental results reported below have been obtained with an average load $`W_0=7`$ N and modulation frequencies $`f`$ of $`120`$ or $`200`$ Hz, chosen to be away from any mechanical resonance frequency of the setup.
### B Results
#### 1 Average dynamic friction
The effect of the normal load modulation on the average tangential force response is to systematically lower the dynamic friction coefficient. The ratio $`\overline{\mu }=\overline{F}/W_0`$ decreases as the modulation amplitude $`ϵ`$ is increased. The variation $`\mathrm{\Delta }\overline{\mu }\left(ϵ\right)=\overline{\mu }\left(ϵ\right)\overline{\mu }(0)`$, plotted on Fig. 2, becomes larger than the experimental noise for $`ϵ0.05`$, and is then quasi linear with $`ϵ`$, though it does not extrapolate to 0 at $`ϵ=0`$.
Fig. 3 displays measurements of $`\overline{\mu }\left(V\right)`$ for different values of $`ϵ`$. It appears that the only effect of an increase of the load modulation amplitude is to shift down the $`\overline{\mu }(V)`$ curve, without changing the slope $`\overline{\mu }/\mathrm{ln}(V)`$. Therefore, within experimental accuracy, $`\mathrm{\Delta }\overline{\mu }\left(ϵ\right)`$ is velocity-independent.
#### 2 A.C. components of the force response
The oscillating force response to a load modulation at frequency $`f`$ is found to be weakly anharmonic. We characterize it by the amplitudes of the first and second harmonics $`\mu _1`$ and $`\mu _2`$. The ratio $`\mu _2/\mu _1`$ lies typically in the range 0.1 — 0.2.
The reduced first harmonic $`\mu _1=|F_1|/W_0`$ increases monotonically with $`ϵ`$ and does not show any measurable dependence on the driving velocity, as presented on Fig. 4.a where we plot results at $`V=1`$ and $`10\mu `$m.s<sup>-1</sup>. $`\mu _1`$ is of order $`10^3`$ at $`ϵ=0.5`$, i.e. two orders of magnitude lower than the average shift $`\mathrm{\Delta }\overline{\mu }`$.
The amplitude of the second harmonic in the shear force response also exhibits a monotonic increase with the modulation amplitude. Moreover, $`\mu _2(ϵ)`$ depends significantly on velocity, the measured amplitude of this component being lower for smaller $`V`$, as presented on Fig. 4.b.
#### 3 Role of the interfacial air layer
All the above results correspond to experiments performed at atmospheric pressure. The PMMA surfaces in contact are nominally flat over typically $`\mathrm{\Sigma }_0=7\times 7`$ cm<sup>2</sup> but their roughness implies that air is trapped in an interfacial gap of micrometric thickness $`h_0`$. Any increase in normal load is borne in parallel by the microcontacts and by the interfacial air layer. This excess pressure leads the air to leak out of the edges of the sample, the rate of flow being limited by the air viscosity. For instance, when trying to lift the slider from the track, a strong succion is experienced. One may therefore expect that the air layer plays a non negligible role in the interfacial response to load modulation.
In order to quantify experimentally this “leaking air cushion” effect, we conducted a set of control experiments under vacuum. The setup described in section (II A) was placed in a vacuum chamber allowing to work at pressures down to $`1`$ mbar (a pressure at which the mean free path of air molecules becomes of order $`10\mu `$m, i.e. much larger than the interfacial gap, ensuring that the air effect has become negligible).
We first measure the average dynamic friction coefficient under constant normal load and find $`\mu _00.5`$, a value equal to the friction coefficient at atmospheric pressure. This confirms that when the interfacial air layer is simply sheared, the corresponding viscous force is negligible with respect to the solid friction one.
Then, following the protocole described in section II A, we measure $`\overline{\mu }\left(ϵ\right)`$, at $`V=10\mu `$m.s<sup>-1</sup> and $`f=200`$ Hz. We have not been able to use in this control experiment the frequency $`f=120`$ Hz at which all other data have been obtained, due to the presence of a spurious mechanical resonance of the vacuum chamber close to 120 Hz.
A comparison of the average friction coefficient variation $`\mathrm{\Delta }\overline{\mu }`$ measured at $`P=1`$ atm and at $`P=1`$ mbar is presented on Fig. 5. Note that for a given modulation amplitude, $`\left|\mathrm{\Delta }\overline{\mu }\right|`$ is larger in vacuum than in air. Moreover, when plotted as a function of $`ϵ_{eff}=ϵ/2.5`$, the results obtained at $`P=1`$ atm are found to collapse on those at $`P=1`$ mbar (see Fig. 5).
We present in the Appendix a model calculation of the elastohydrodynamic response of the air layer. We show that, in all the range of $`ϵ`$ used in our experiments, the normal response of the interface is linear, hence the ratio $`ϵ_{eff}/ϵ`$ does not depend on $`ϵ`$. Moreover, the estimated order of magnitude of this parameter at $`f=200`$ Hz is found to be compatible with the above measured value.
## III Discussion and model
In this section we analyse our data within the SRF framework. The three parameters $`A`$, $`B`$ and $`D_0`$ involved in the SRF laws are determined experimentally, at constant load $`W`$, using the velocity dependence of the friction coefficient $`\overline{\mu }\left(\mathrm{ln}(V)\right)`$ and the dynamic characteristics of the response close to the bifurcation threshold (this method has been described in detail in reference ). We measure for our system $`A=0.013\pm 0.005`$, $`B=0.026\pm 0.01`$ and $`D_0=0.4\pm 0.04\mu `$m. All the numerical integrations of SRF laws presented below are performed with this set of parameter values.
### A Rice and Ruina’s model
Before coming to the question of whether or not the Rice-Ruina (RR) equations themselves should be modified in the presence of load modulations, it is reasonable to study first which response is predicted by the RR model as such.
Replacing in Eq. (1) $`W`$ by its instantaneous value, the equation of motion of the center of mass of the slider reads:
$$M\ddot{x}=K\left(Vtx\right)W_0\left(1+ϵ\mathrm{cos}(\omega t)\right)\left[\mu _0+A\mathrm{ln}\left(\frac{\dot{x}}{V_0}\right)+B\mathrm{ln}\left(\frac{\mathrm{\Phi }V_0}{D_0}\right)\right]$$
(8)
where $`x(t)`$ is the instantaneous position of the center of mass of the slider with respect to the track. We assume the evolution law of $`\mathrm{\Phi }`$ (Eq. 2) to be unmodified:
$$\dot{\mathrm{\Phi }}=1\frac{\dot{x}\mathrm{\Phi }}{D_0}$$
(9)
#### 1 The perturbative regime
Let us first consider the case where $`ϵ1`$. We linearize Eqs. (8) and (9) about the steady sliding state at velocity $`V`$, $`ϵ=0`$:
$$\mathrm{\Phi }_{st}=D_0/Vx_{st}=VtW_0/K$$
(10)
Setting:
$$\delta \mathrm{\Phi }=\mathrm{\Phi }\mathrm{\Phi }_{st}=\mathrm{}e\left(ϵX_1\mathrm{exp}\left(i\omega t\right)\right)$$
(11)
$$\delta x=xx_{st}=\mathrm{}e\left(ϵ\mathrm{\Phi }_1\mathrm{exp}\left(i\omega t\right)\right)$$
(12)
we get to first order in $`ϵ`$
$$\left(M\omega ^2+K+iW_0\omega \frac{A}{V}\right)X_1+W_0\frac{BV}{D_0}\mathrm{\Phi }_1=W_0\overline{\mu }(V)$$
(13)
$$i\frac{\omega }{V}X_1+\left(i\omega +\frac{V}{D_0}\right)\mathrm{\Phi }_1=0$$
(14)
where $`\overline{\mu }\left(V\right)=\mu _0\left(BA\right)\mathrm{ln}\left(V/V_0\right)`$. We thus obtain
$$X_1=\frac{W_0\overline{\mu }}{\mathrm{\Delta }}\left(i\omega +\frac{V}{D_0}\right)$$
(15)
$$\mathrm{\Phi }_1=\frac{i\omega W_0\overline{\mu }}{V\mathrm{\Delta }}$$
(16)
where $`\mathrm{\Delta }`$ reads
$$\mathrm{\Delta }=\frac{KV}{D_0}\left(\frac{\omega }{\omega _c}\right)^2\left[\frac{K}{K_c}+\left(1\left(\frac{\omega }{\omega _0}\right)^2\right)\left(\frac{\omega _c}{\omega }\right)^2i\left(\frac{\omega _c}{\omega }\right)\sqrt{\frac{BA}{A}}\left(\left(\frac{\omega }{\omega _0}\right)^2\left(1\frac{K}{K_c}\right)\right)\right]$$
(17)
and (see reference)
$$K_c=\frac{(BA)W_0}{D_0}$$
(18)
$$\omega _c=\sqrt{\frac{BA}{A}}\frac{V}{D_0}$$
(19)
are, respectively, the critical stiffness and pulsation at the stick-slip bifurcation for the unmodulated system. $`\omega _0=\sqrt{K/M}360`$s<sup>-1</sup> is the inertial frequency.
In our experimental conditions, with $`V=10\mu `$m.s<sup>-1</sup> and $`(BA)/A1`$, $`\omega _c25`$s<sup>-1</sup>, so $`\omega _c\omega ,\omega _0`$. On the other hand $`K/K_c1`$. Then $`\mathrm{\Delta }\omega ^2AW_0/V`$, indicating in particular that inertia is negligible. Finally:
$$\mu _1=\frac{K\left|X_1\right|}{W_0}ϵ=\frac{ϵ\overline{\mu }}{A}\frac{V}{\omega }\frac{K}{W_0}$$
(20)
Similarly, a second order expansion in $`ϵ`$ yields the corrections at frequencies $`2\omega `$ and $`0`$, namely:
$$\mu _2\frac{V}{\omega }\frac{K\overline{\mu }^2}{2W_0A^2}ϵ^2$$
(21)
$$\mathrm{\Delta }\overline{\mu }\frac{\overline{\mu }\left(\overline{\mu }+2A\right)}{4A}ϵ^2$$
(22)
Note that Eq. (22) correctly predicts a decrease of the average friction coefficient.
It is interesting to compare the relative perturbative corrections on the age and velocity variables. One finds
$$\left|\frac{\delta \mathrm{\Phi }/\mathrm{\Phi }_{st}}{\delta \dot{x}/V}\right|\frac{\omega _c}{\omega }1$$
(23)
showing that in our regime, the modulation of the age variable contributes negligibly to the shear response.
Moreover, due to the smallness of $`A`$, the effective perturbation parameter is given by:
$$\left|\frac{\delta \dot{x}}{V}\right|=\frac{ϵ\overline{\mu }}{A}50ϵ$$
(24)
that is, the perturbative regime ($`ϵ10^2`$) is in practice out of experimental reach. We thus must resort to full integration of the above RR equations.
#### 2 Average friction coefficient decrease
Numerical results for $`\mathrm{\Delta }\overline{\mu }(ϵ)`$ are plotted on Fig. 6 at excitation frequency $`f=120`$ Hz and velocities $`V=1`$ and $`10\mu `$m.s<sup>-1</sup>. Note that a very weak dependence on $`V`$ is predicted, as observed experimentally.
In section II B we emphasised the role played by the interfacial air layer in our experiments, and pointed out that it should be taken into account through an effective modulation amplitude $`ϵ_{eff}`$, accounting for the fact that only a part of the excitation is borne by the contacting asperities. Therefore, the modulation parameter $`ϵ`$ introduced in Eq. (8) must be understood as $`ϵ_{eff}`$.
On the other hand, we have measured, at $`f=200`$ Hz, $`ϵ_{eff}=0.4ϵ`$. As explicited in the appendix, we expect $`ϵ_{eff}/ϵ`$ to exhibit some relatively weak frequency dependence. This effect depends crucially on the interfacial normal stiffness which is difficult to measure accurately. So, we have chosen to treat $`ϵ_{eff}/ϵ`$ as a free fitting parameter with an initial trial value 0.4.
Fig. 7 shows the best fit obtained for $`\mathrm{\Delta }\overline{\mu }(ϵ)`$ at $`V=1\mu `$m.s<sup>-1</sup>. It corresponds to $`ϵ_{eff}=0.48ϵ`$.
From now on, all experimental data will be plotted versus this effective modulation parameter.
#### 3 A.C. response
The computed first and second harmonics of the frictional response are plotted on Figs. 8.a and 8.b. One can first notice that the quasi-linear dependence of $`\mu _1`$ and $`\mu _2`$ on $`V`$ predicted by the perturbation calculation also holds here. Moreover, both computed harmonics saturate at large $`ϵ`$. None of these features agrees with the experimental behaviour.
We therefore conclude that, in spite of the excellent agreement between the predicted and observed $`\mathrm{\Delta }\overline{\mu }(ϵ,V)`$, the unmodified RR model is insufficient to describe the full response of the interface.
### B Linker and Dieterich’s ageing law
As mentioned in section I, Linker and Dieterich (LD) have proposed an extended version of the RR model in which the evolution law of the age variable $`\mathrm{\Phi }`$ is modified according to Eq. (6). We now study the predictions of this extended model.
A perturbation calculation to first order in $`ϵ`$, using the equation of motion (8) and the age law (6), leads to a first harmonic amplitude:
$$\mu _1=\frac{ϵ|\overline{\mu }\alpha |}{A}\frac{V}{\omega }\frac{K}{W_0}$$
(25)
This expression points to the fact that the dimensionless LD parameter $`\alpha `$ acts as a correction to the bare dynamic friction coefficient $`\overline{\mu }`$. LD propose for granite $`\alpha =`$0.2 – 0.3, i.e. a sizeable fraction of $`\overline{\mu }`$ ($`0.6`$ for that material). This leads one to expect that such a value should induce significant effects on the predicted shear response. However, we have estimated an order of magnitude of $`\alpha `$ for a sparse population of microcontacts (Greenwood interface) ageing under normal load. We have considered the two limits of (i) linear viscoelastic and (ii) fully developped plastic creep, using parameters compatible with the measured value of the RR parameter $`B`$. Both limits lead to the same estimate for $`\alpha `$, namely one order of magnitude smaller than the LD value.
In view of this discrepancy, we have chosen to perform numerical integrations of Eqs. (8) and (6) for various values of $`\alpha `$ in the range 0.02 – 0.2.
The magnitude of the load modulation effect on $`\mathrm{\Delta }\overline{\mu }(ϵ)`$ depends strongly on $`\alpha `$, as shown on Fig. 9.a. Whatever $`\alpha `$, $`\mathrm{\Delta }\overline{\mu }(ϵ)`$ remains quasi-independent of $`V`$, but for $`\alpha =0.2`$ it is significantly smaller than the experimental one.
Moreover, the dependences of $`\mu _1`$ and $`\mu _2`$ on $`ϵ`$ and $`V`$, shown on Figs. 9.b and 9.c, as for the RR model, clearly disagree with the experimental results.
$`\alpha =0.02`$ is found to provide a satisfactory fit for $`\mathrm{\Delta }\overline{\mu }\left(ϵ\right)`$. However, this $`\alpha `$ value is small enough for age effects to become negligible, as noticed in section (III A 1). We indeed check (Figs. 9.b and 9.c) that the corresponding $`\mu _1`$ and $`\mu _2`$ are very close to those obtained from the unmodified RR model.
We are thus led to conclude that:
(i) The LD evolution law with their proposed value of $`\alpha `$ does not agree with the experimental results.
(ii) The $`\mathrm{\Delta }\overline{\mu }`$ data permit to set an upper limit on $`\alpha `$ without, however, allowing to check the validity of the fuctional form of the LD model. Experiments at much lower frequencies (comparable to the stick-slip frequency $`\omega _c`$) would be needed to answer this question.
### C Extension of the RR model
The above analysis suggests that in our “high frequency” regime, where the response is controlled by the velocity modulations, it is the “rheological” factor $`\sigma _s(\dot{x})`$ which should be modified. As mentioned in section I, $`\sigma _s`$ describes the plastic dissipation occurring in a junction of nanometer thickness between contacting asperities, and the rate involved in $`\sigma _s`$ is a rate of irreversible strain of this junction.
It has been shown that when a multicontact interface is submitted to a shear much smaller than the static threshold, its response is elastic. Since the asperity “bodies” (which deform on a micrometric thickness, of the order of their diameter) are much more compliant than the nanometer-thick elastically pinned adhesive joint, it is their response which controls the interfacial shear stiffness $`\kappa _{asp}`$. This obeys an extented Amontons law: $`\kappa _{asp}=W/\lambda `$, with $`\lambda `$ a length of order $`1\mu `$m for our surfaces.
Sliding amounts to depinning of the adhesive joint which becomes dissipative, while the bodies of the asperities retain their elasticity. Therefore, we can schematically represent the sliding interface as an elastic element of stiffness $`\kappa _{asp}`$, accounting for the bulk elastic strain of the asperities, coupled in series to a (frictional) dissipative element (see Fig. 10). When this latter is sheared at velocity $`\dot{x}_{pl}`$, the corresponding force is $`F=f\left(\dot{x}_{pl}\right)`$.
$$F=\kappa _{asp}x_{el}=f(\dot{x}_{pl})$$
(26)
with $`x_{el}`$ and $`x_{pl}`$ respectively the elastic and irreversible displacements. The instantaneous velocity of the center of mass of the slider thus reads:
$$\dot{x}=\dot{x}_{el}+\dot{x}_{pl}=\frac{d\left(F/\kappa _{asp}\right)}{dt}+f^1\left(F\right)$$
(27)
and the tangential force on the slider finally reads:
$$F=f\left(\frac{d}{dt}\left(xF/\kappa _{asp}\right)\right)$$
(28)
We therefore express the external force using the same functional form as in Eq. (1), but the argument of the rate-dependent term becomes $`\dot{x}_{pl}`$. In stationary sliding under constant normal load, $`F/\kappa _{asp}`$ is constant, hence the usual dependence on $`\dot{x}`$. In the presence of a load modulation, both $`F`$ and $`\kappa _{asp}=W_0\left(1+ϵcos\left(\omega t\right)\right)/\lambda `$ are modulated, and the elastic strain term becomes significant.
We present hereafter the results obtained from numerical integration of the corresponding extended RR equations. Taking into account the above mentioned fact that in our experimental conditions, inertia can be neglected, Eq. (1) becomes:
$$F/W=\frac{K}{W}(Vtx)=\mu _0+A\mathrm{ln}\left[\frac{1}{V_0}\frac{d}{dt}\left(x\frac{K}{\kappa _{asp}}(Vtx)\right)\right]+B\mathrm{ln}\left(\frac{V_0\mathrm{\Phi }}{D_0}\right),$$
(29)
The parameters $`A`$, $`B`$ and $`D_0`$ are set to their experimentally determined values. The length $`\lambda `$ has been obtained from a quasi-static loading-unloading test at various normal loads. We find $`\lambda =0.62\pm 0.15\mu `$m. In view of the relatively large experimental uncertainty on this parameter, we have integrated Eqs. (29) and (2) with $`\lambda `$ as a free parameter. The best fit, performed on the most sensitive data, namely the $`\mu _1(ϵ_{eff})`$ ones, is found to correspond to $`\lambda =0.7\mu `$m, within the experimental uncertainty braket.
While $`\mathrm{\Delta }\overline{\mu }\left(ϵ\right)`$ is found to be only very weakly affected by the rheological correction, this extension of the model yields predictions for $`\mu _1`$ and $`\mu _2`$ markedly different from those of both the unmodified RR and LD models. Namely, their quasi-linear $`V`$-dependences are replaced by much weaker ones. On the other hand, neither $`\mu _1`$ nor $`\mu _2`$ exhibit any longer saturation within the relevant range $`ϵ_{eff}<0.3`$.
As appears from Figs. (11), the global agreement is now excellent, confirming the validity of the extended RR model.
### D Concluding remarks
This study leads us to the following conclusions:
On the one hand, from an engineering point of view, the most spectacular effect of modulating the normal load applied to a frictional system is to lower significantly the dynamic friction coefficient. This occurs as soon as the modulation is applied, even though its amplitude is low enough to ensure permanent interfacial contact between the sliding bodies.
On the other hand, an important aim of this study was to elucidate the question, relevant to seismology, of whether the RR model should be modified to describe the frictional response to fast variations of the normal stress. We have shown that, in order to study this, it is essential to measure and analyze not only the zero frequency component of the response to an oscillatory load, but also its harmonic content.
In the range of frequencies, much larger than the stick-slip one, that we have studied, the shear response is controlled by the velocity modulation, that is by the rate-dependent term of the RR constitutive law. However, the quantitative analysis of $`\mu _1`$ and $`\mu _2`$ data shows that the relevant displacement rate is not, for fast load modulations, the slider velocity, but the rate of plastic deformation of the adhesive frictional joint.
This confirms our picture of sliding friction as 2D plasticity prelocalized within a nanometer-thick adhesive joint coupled to the bulk of the slider through elastic asperities.
This enables us to extend correspondingly the expression of the rate-dependent part of the RR state- and rate-dependent model.
The question of the precise effect of a load modulation on interfacial age remains at this stage open. Indeed, we have concluded that, at least for our system, this effect is certainly much smaller than proposed by Linker and Dieterich. However, precisely for this reason, the “high frequency” response is not a good tool for investigating this question. This should be adressed through similar experiments at low frequencies, close to the stick-slip one.
## A Elasto-hydrodynamic response of the interfacial air layer
The aim of this appendix is to establish the equation for the vertical motion of the slider (i.e. along the $`z`$-direction normal to the interface) and to estimate the relative contributions of the forces that are involved. We will, as a result, justify the use of an effective amplitude of modulation of the normal load, to account for the fraction of the modulation which is borne by the air cushion trapped within the interfacial gap. The order of magnitude of this fraction, referred to as $`ϵ_{eff}/ϵ`$ in the text, and which is the only fitting parameter of our model, is checked independently in a control experiment, performed in vacuo, and described in the text.
The motion of the slider along the $`z`$-axis is assumed to be decoupled from the sliding motion along $`x`$. It is parametrized by the width $`h`$ of the “gap”, i.e. the separation between the average planes passing through the rough surfaces of, respectively, the track and the slider. When no modulation is superimposed to the bare normal load $`W_0`$, the width is $`h_0`$, a value fixed by the deformation of the load bearing asperities which are randomly distributed over the interface of nominal macroscopic area $`\mathrm{\Sigma }_0`$.
##### a Elastic response of the multicontact interface.
According to Greenwood and Williamson’s model for multicontact interfaces, the number of load bearing asperities and the real area of contact increase linearly with the load. This induces a non-linear dependence on load of the gap thickness. Experimentally, it has beeen found that $`h_0h\lambda _z\mathrm{ln}(W/W_0)`$, with $`\lambda _z`$ a length, the order of magnitude of which is given by the roughness of the surfaces in contact (the standard deviation of the surface heights, here $`1.3\mu `$m). In the small amplitude linear regime ($`\mathrm{\Delta }h\lambda _z`$), the stiffness $`\kappa _z=W/\lambda _z`$ is a constant and the emlastic restoring force in the $`z`$-direction reads:
$$F_{el}W_0\frac{(hh_0)}{\lambda _z}$$
(A1)
This expression has its exact counterpart for tangential motion, as discussed in the text. Shear elasticity involves a length $`\lambda `$ which is expected to be about $`1.7\lambda _z`$ for a Poisson ration $`\nu =0.44`$. Therefore, the measured value $`\lambda =0.7\mu `$m yields $`\lambda _z0.4\mu `$m.
##### b Elasto-hydrodynamic response of the interfacial air layer.
When the gap is e.g. narrowed, air is compressed until being drained out of the interfacial zone. The resulting pressure force on the slider will be denoted $`F_p`$. For the sake of evaluating $`F_p`$, we will simplify the problem and consider a thin layer of air, of viscosity $`\eta `$ and density $`\rho _0`$ at atmospheric pressure $`P_0`$, trapped between two perfectly smooth discs of radius $`R`$, parallel and distant of $`hR`$. The relative velocity is supposed to vanish at $`z=0`$ and $`z=h`$, an assumption which is legitimate if the roughness of the surfaces is much smaller than the gap width. Brown and Scholz have reported measurement of the gap width between macroscopic ground glass surfaces. At low average pressure corresponding to $`10^5`$ of the Young modulus of the glass, as encountered in our experiments, they have found that the gap width is typically 5 times larger than the rms roughness of the statistically identical surfaces. This figure is clearly too small for the “smoothed” model of the interfacial gap to be expected to provide a very accurate value of the hydrodynamic force, though it is certainly sufficient to estimate its order of magnitude.
An upper bound for the average pressure excess resulting from the motion of the disk is $`\mathrm{\Delta }P=ϵW_0/\mathrm{\Sigma }_0`$, with $`ϵW_0`$ the amplitude of the normal load modulation and $`\mathrm{\Sigma }_0=\pi R^2`$. The macroscopic loading pressure $`W_0/\mathrm{\Sigma }_0`$ remains of order 10 mbar in the reported experiments, while $`ϵ`$ is smaller than unity. As a result, $`\mathrm{\Delta }P`$ remains much smaller than the atmospheric pressure $`P_0`$. However, the compressibility of the air layer may be of paramount importance, as suggested by the following argument. For infinite plates, no leak occurs at the edge of the gap and the response of the layer, trapped under the mean pressure $`P_0`$, is elastic with an overall stiffness:
$$\kappa _{air}=P_0\mathrm{\Sigma }_0/h_0$$
(A2)
For $`P_0=10^5Pa`$, $`\mathrm{\Sigma }_0=49`$ cm<sup>2</sup> and $`h_0=6.5\mu `$m, one finds $`\kappa _{air}=\mathrm{7.5\hspace{0.17em}10}^7`$ N/m, namely one order of magnitude larger than the interfacial stiffness $`\kappa _z`$ originating from the load bearing asperities (see Eq.A1) at $`W_0=7`$ N. For finite radius $`R`$, edge flow will reduce the amount of air to be compressed in order to accomodate the change of gap volume. It is therefore necessary to compute the expression of $`F_p`$ by taking account the radial, viscosity controlled, Poiseuille flow which results from the density (hence pressure) gradient compatible with mass conservation.
The continuity equation for the radial flow reads:
$$\frac{}{t}(\rho h)+\frac{1}{r}\frac{}{r}(r\overline{v}\rho h)=0$$
(A3)
where $`\overline{v}(r)`$ is the mean velocity at radius $`r`$ (averaged across the gap along the z-direction). The pressure field is given by the equation of state of the air at pressures close to $`P_0`$ = 1 atm, which is assumed to be:
$$\frac{P}{\rho }=\frac{P_0}{\rho _0}$$
(A4)
The set of equations is closed by assuming that the flow is of the Poiseuille type, namely is parabolic along the $`z`$-direction and varies slowly along the radial direction according to:
$$\overline{v}=\frac{h^2}{12\eta }\frac{P}{r}$$
(A5)
As mentioned, the pressure modulation remains much smaller than $`P_0`$, and the gap modulation is smaller than $`h_0`$, hence linearization of Eqs. (A3A5) is legitimate. One therefore sets $`P=P_0+\delta P`$, $`\rho =\rho _0+\delta \rho `$, and $`h=h_0+\delta h`$, with $`\delta PP_0`$, $`\delta \rho \rho _0`$ and $`\delta hh_0`$. Eliminating $`\delta \rho `$ yields the following equation for the pressure field:
$$\frac{h^2}{12\eta }\frac{1}{r}\frac{}{r}\left(r\frac{(\delta P)}{r}\right)\frac{\dot{\delta P}}{P_0}=\frac{\dot{\delta h}}{h_0}$$
(A6)
where the dot indicates the partial derivative with respect to time.
Assuming that the normal elastic stiffness $`\kappa _z`$ of the asperities remains linear, the gap modulation resulting from the normal load one is harmonic and we therefore seek for a complex solution to Eq.A6 of the form: $`\delta P=\stackrel{~}{\delta P}\mathrm{exp}(i\omega t)`$ with $`\delta h=\stackrel{~}{\delta h}\mathrm{exp}(i\omega t)`$. Taking into account the boundary condition $`\delta P=0`$ at $`r=R`$ and the symetry requirement $`\overline{v}=0`$, hence $`P/r=0`$, at $`r=0`$, one obtains:
$$\stackrel{~}{\delta P}=P_0\frac{\stackrel{~}{\delta h}}{h_0}\left[1\frac{J_0(\gamma r)}{J_0(\gamma R)}\right]$$
(A7)
With $`J_0`$ the Bessel function of zeroth order and $`\gamma `$ a complex constant given by:
$$\gamma =\frac{1i}{\sqrt{2}}\sqrt{\frac{12\eta \omega }{P_0h_0^2}}$$
(A8)
Integration of the pressure field over the interface yields the complex amplitude of the pressure force:
$$\stackrel{~}{\delta F_p}=P_0\pi R^2\frac{\stackrel{~}{\delta h}}{h_0}\frac{J_2(\gamma R)}{J_0(\gamma R)}$$
(A9)
with $`J_2`$ the Bessel function of second order.
The asymptotic limits deserve comments. For $`\gamma R\mathrm{}`$, $`J_2/J_01`$, $`\stackrel{~}{\delta F_p}`$ is real and one recovers the purely elastic response (Eq.A2) predicted for large $`R`$. It also corresponds to the high frequency limit for which the air has no time to leak. For $`\gamma R0`$, $`J_2/J_0(\gamma R)^2/8=i\omega (3\eta R^2)/(2P_0h_0^2)`$, hence $`\stackrel{~}{\delta F_p}`$ is purely imaginary and reduces to a linear viscous damping force which could have been derivated by assuming a non compressive Poiseuille flow. For intermediate values of $`\gamma R`$, $`\stackrel{~}{\delta F_p}`$ has both a reactive component, which increases the interfacial stiffness, and a dissipative one.
##### c Prediction for the effective amplitude of load modulation.
The slider of mass $`M`$ oscillates in the normal $`z`$-direction under the combined action of the load modulation, the restoring elastic force resulting from deformation of the load bearing asperities and compression of the air cushion, and the damping force resulting from the air flow. The complex amplitude of modulation of the gap width $`\stackrel{~}{\delta h}`$ is therefore given by:
$$(M\omega ^2+\kappa _z)\stackrel{~}{\delta h}\stackrel{~}{\delta F_p}=ϵW_0$$
(A10)
with $`\stackrel{~}{\delta F_p}`$ given by Eq.A9.
The fraction of the load which is effectively borne by the microcontacts is $`ϵ_{eff}/ϵ=|\kappa _z\stackrel{~}{\delta h}/(ϵW_0)|`$. It reads:
$$\frac{ϵ_{eff}}{ϵ}=\left|1\frac{\omega ^2}{\omega _0^2}\frac{\pi R^2P_0}{W_0}\frac{\lambda _z}{h_0}\frac{J_2(\gamma R)}{J_0(\gamma R)}\right|^1$$
(A11)
with $`\omega _0=\sqrt{\kappa _z/M}`$.
The assumption that $`ϵ_{eff}/ϵ`$ does not depend on the amplitude of the modulation relies upon the fact that both the elasticity and the viscosity remain linear, namely, as previously discussed, that $`\mathrm{\Delta }h\lambda _z`$ and $`h_0`$. This reduces to $`ϵ_{eff}1`$, a criterion which is always fulfilled in our experiments.
For $`W_0=7`$ N, $`M=1.6`$ kg and $`\lambda _z0.4\mu `$m, $`\omega _0/(2\pi )=530`$ Hz. Hence, at 120 Hz, the inertia is $`\mathrm{5.2\hspace{0.17em}10}^2`$ of the elastic restoring force due to the asperities solely. It is clear from the above analysis that a key parameter for evaluating the viscoelastic response of the air is the gap width $`h_0`$. Taking, as discussed previously, the conservative value of 5 times the roughness, namely $`h_0=6.5\mu `$m, $`\eta =10^5`$ Pa.s, $`R=3.9`$ cm, one computes $`\gamma R=4(1i)`$ and $`|ϵ_{eff}/ϵ|0.24`$, a value of the same order of magnitude than the one, namely 0.48, which is found to provide the best agreement between the experimental data and the model prediction for $`\mathrm{\Delta }\overline{\mu }`$. The role of the interfacial air cushion is further confirmed by the control experiment performed in vacuo. At a remaining pressure of 1 mbar, the elastic stiffness of the air layer falls two orders of magnitude below the multicontact one. Moreover, since the mean free path (at 300 K) of the gas molecules is now of order several 10 $`\mu `$m, i.e. larger than the gap width, the viscosity of the layer should vanish. Consequently, the effective amplitude is essentially ruled by the slider inertia according to: $`ϵ_{eff}^{vacuum}/ϵ|1\omega ^2/\omega _0^2|^1`$ = 1.2 at 200 Hz. At atmospheric pressure, keeping the nominal value $`\mathrm{\Sigma }_0=49`$ cm<sup>2</sup>, $`ϵ_{eff}^{air}/ϵ=0.23`$. When bringing the data for $`\mathrm{\Delta }\overline{\mu }(ϵ)`$ performed in the air and in vacuo to collapse on a single curve, as explained in the text, one makes use of a scaling ratio which reads explicitely: $`ϵ_{eff}^{air}/ϵ_{eff}^{vacuum}0.20`$. The experimental value is 0.4.
The fact that, in both cases, the estimated value is smaller than the observed one by the same amount may be possibly attributed to some long wavelength modulations of the gap width $`h_0`$ which is likely to remain after the lapping process. Microcontacts may be distributed on patches of macroscopic area smaller than $`\mathrm{\Sigma }_0`$, separated by regions of much wider gap in which the air would play a negligible role. Typically, a patch radius of 2.5 cm, while keeping the other parameters unchanged, would account for the observed value $`ϵ_{eff}/ϵ0.48`$ at 120 Hz in the air. This would correspond to an effective area of $`0.4\mathrm{\Sigma }_0`$, a value still large enough for the microcontacts — the number of which does not depend on $`\mathrm{\Sigma }_0`$, according to Greenwood — to remain elastically independent. In addition, we have assumed a single degree of freedom for the slider, which is certainly a strong requirement since the slider is left free to find its own seat on the track. It is clear that a small amount of ”rocking” would promote the air flow and reduce the cushion effect.
Finaly, normal load modulation induces a tangential oscillating motion of the slider of amplitude $`\mathrm{\Delta }x`$, hence an air shear flow within the gap. The associated A.C. viscous force on the slider $`\eta \omega \mathrm{\Delta }x\mathrm{\Sigma }_0/h_0`$, which has been neglected in our models, must be compared to the leading term in the rate dependent friction force for oscillations about the sliding velocity $`V`$, namely $`AW_0\omega \mathrm{\Delta }x/V`$. The ratio of both terms is $`\eta \mathrm{\Sigma }_0V/(W_0h_0)10^7`$ for $`V=100\mu `$m/s, therefore shear viscosity of the layer is totally negligible.
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# 1 Introduction
## 1 Introduction
Hawking and Turok found that a family of gravitational instantons exist for a generic potential for the scalar field. In the Hartle–Hawking ‘no-boundary’ proposal the universe is supposed to be created from a compact instanton, so that Hawking–Turok instantons have a significant importance. The novelty is the allowance of singularities of special kind (Hawking–Turok (HT)-singularity) in the curvature and the scalar field. Hawking and Turok argued that the use of such singular instantons is justified because of the finiteness of the action. In addition, the singularity is mild and behaves as a reflecting boundary for scalar and tensor cosmological perturbations . Hence the Cauchy problem seems to be well posed and the model is well suited for the quantisation of small perturbations and for comparison with observations. As another justification, Wu argued that Hawking–Turok instantons may be regarded as ‘constrained gravitational instantons,’ whose action is required to be stationary under the constraint that the 3-geometry is given on the three-surface where quantum transition occurs.
On the other hand, there is a criticism by Vilenkin that asymptotically flat HT-singular instantons would destabilise the flat spacetime. Among such instantons especially the ones which do not have the potential for the scalar field are called as Vilenkin’s instantons. To this criticism it was argued that Vilenkin’s instanton can be taken as compactified five-dimensional Schwarzschild metric and the large extra dimension would metastabilises the flat spacetime. This mechanism was first discussed by Witten . Turok argued that with a careful definition of a constraint Vilenkin’s instantons possess no negative mode, so that they do not lead to the decay of flat spactime.
With such arguments it is desirable to examine the contribution of the HT-singularity to the action with much care. Most simply it is calculated as follows : at first we consider a manifold $`^{}`$ which consists of the whole spacetime $``$ except a four-sphere which contains the singular point and then make the excised four-sphere shrink. In the zero-volume limit of the excised four-sphere, the Gibbons–Hawking boundary term of the boundary of $`^{}`$ would give the contribution of the singularity. Other than this simple method, there are several attempts which do not agree with each other. Garriga calculated the contribution of the singularity by regularizing the singularity with a membrane which wraps the singularity and found that the contribution of the singularity is just $`1/3`$ of the Gibbons–Hawking term. But the method was criticised by Bousso and Chamblin because of allowing negative energy density. In ref. it is considered to regularise the singularity with a membrane of a positive energy density which couples to four-form field strength. Garriga calculated in ref. the action of Vilenkin’s instanton by showing that compactified five-dimensional Schwarzschild metric looks as Vilenkin’s instanton and obtained the value of the $`1/3`$ of the Gibbons–Hawking term. In ref. González-Díaz calculated the action of Vilenkin’s instanton which is replaced the singularity with an axionic wormhole and obtained the value zero. Turok considered to define Vilenkin’s instantons as a limit of constrained instantons which are non-singular and have no boundary, and obtained the value of the whole Gibbons–Hawking term as the value of the action.
Although the HT-singularity has been often regularised with new matter in these calculations, the singularity at the transition surface from Euclidean to Lorentzian metric would be allowed if we take HT-singular instantons as constrained gravitational instantons. On the other hand the simplest method mentioned above to calculate the contribution of the HT-singularity would not be justified because the contribution of the excised spacetime is overlooked. For example, in the case of the Euclidean Schwarzschild metric the limiting value of Gibbons–Hawking term does not contribute to the action although it has a non-zero value. For that reason let us consider to calculate the contribution of the excised HT-singularity. To do this we review the case of the Euclidean Schwarzschild metric in the next section. Then we apply it to the case of the HT-singular instantons using Kaluza–Klein technique.
## 2 The case of Euclidean Schwarzschild metric: A simple example
The Lorentzian Schwarzschild metric is given by
$$ds^2=\left(1\frac{2MG}{r}\right)dt^2+\left(1\frac{2MG}{r}\right)^1dr^2+r^2d\mathrm{\Omega }_2^2.$$
(1)
Euclideanising by $`t=i\tau `$ and putting $`x=4M\sqrt{12MG/r}`$, we obtain
$$ds^2=\left(\frac{x}{4MG}\right)^2d\tau ^2+\left(\frac{r^2}{4M^2G^2}\right)^2dx^2+r^2d\mathrm{\Omega }_2^2.$$
(2)
This metric will be regular at $`x=0`$, i.e. $`r=2MG`$, if $`\tau `$ is regarded as an angular variable with period $`8\pi MG`$. The Euclidean Schwarzschild metric is defined on the manifold given by $`x0`$, $`0\tau 8\pi GM`$. We denote this spacetime as $``$ in this section. The submanifold $`x=0`$ is called a bolt , which is the two-dimensional fixed point (FP) set of the periodic imaginary time isometry.
Let us consider to calculate the action of $``$ by using a spacetime $`^{}`$ which has an outer boundary three-surface $`_{\mathrm{𝑜𝑢𝑡}}^{}`$ at infinity and an inner boundary three-surface $`_{\mathrm{𝑖𝑛}}^{}`$ at $`x=x_\epsilon `$ in the vanishing limit of $`x_\epsilon `$. Since the metric is a vacuum solution, the Ricci scalar is zero everywhere and the action $`𝒮(^{})`$ comes entirely from the boundary terms. The Gibbons–Hawking boundary term is given by the integral over the boundary
$$𝒮_{GH}=\frac{1}{8\pi G}d^3\xi \sqrt{h}\left(KK_0\right),$$
(3)
where $`h`$ is the determinant of the induced metric on the boundary, $`K`$ and $`K_0`$ are the trace of the extrinsic curvature of the boundary and the boundary imbedded in flat spacetime, respectively. The second term is included when we consider a non-compact spacetime. For calculation of the Gibbons–Hawking term it is convenient to use the formula
$$d^3\xi \sqrt{h}K=_{normal}(\text{Volume of boundary}).$$
(4)
The boundary term of the outer boundary $`_{\mathrm{𝑜𝑢𝑡}}^{}`$ at infinity gives the contribution $`4\pi M^2G`$. In addition, the boundary term of the inner boundary $`_{\mathrm{𝑖𝑛}}^{}`$ is also non-vanishing:
$`𝒮(_{\mathrm{𝑖𝑛}}^{})`$ $`=`$ $`\underset{x_\epsilon 0}{lim}𝒮_{GH}`$ (5)
$`=`$ $`4\pi M^2G.`$ (6)
But this term is not expected to contribute to the action $`𝒮()`$ of the whole spacetime, because the metric is regular at the $`x=0`$ and so the spacetime $``$ has no boundary there. Therefore we expect that the contribution of the excised spacetime (in this case, the bolt) will be non-zero and will compensate the contribution of the new boundary which is added by considering $`^{}`$ instead of $``$.
To confirm this, we consider a four-sphere $`_B`$ which is centered around the bolt and whose boundary three-surface is at $`x=x_\epsilon `$. We define $`𝒮(\text{FP})`$ as the action of $`_B`$ in the limit of $`x_\epsilon 0`$. Because the Euclidean Schwarzschild metric is regular at the bolt, we can evaluate this term easily: In the vanishing limit of $`x_\epsilon `$ the bulk term vanishes and the action comes entirely from the Gibbons–Hawking term of $`_B`$, and we need not to calculate it because it is given by the Gibbons–Hawking term of $`_{\mathrm{𝑖𝑛}}^{}`$ only with the reversal of the direction of the normal of the boundary. So we obtain
$`𝒮(\text{FP})`$ $`=`$ $`\underset{x_\epsilon 0}{lim}𝒮_{GH}`$ (7)
$`=`$ $`𝒮(_{\mathrm{𝑖𝑛}}^{}).`$ (8)
The contribution of the bolt cancels with the inner boundary contribution of $`^{}`$, as expected. Therefore using $`^{}`$ we obtain the correct action of $``$ as
$`𝒮()`$ $`=`$ $`𝒮(^{})+𝒮(\text{FP})`$ (9)
$`=`$ $`4\pi M^2G.`$ (10)
## 3 The case of singular instantons
In this section we calculate the fixed point contribution of HT-singularities. The method to calculate it presented in the previous section cannot be applied directly because the excised spacetime is not regular in this case. But we show that we can go around this obstruction by using Kaluza–Klein technique.
Introduction of an extra dimension in HT-singular instantons was first discussed by Garriga . He showed that Vilenkin’s instanton may be reinterpreted à la Kaluza–Klein as the Euclidean Schwarzschild solution of five-dimensional gravity, and calculated the action of Vilenkin’s instanton as the action of the Euclidean Schwarzschild solution. He argued that if the fifth dimension is physical and we can interpret the HT-singularity as such, then the large extra dimension as in the M-theory metastabilises the flat spacetime. But at the same time he also argued that the five-dimensional interpretation would not work if a general potential term is present, because HT-singular instantons can be represented as compactified five-dimensional regular metrics only when the initial value of the scalar field takes a particular value. He exemplified this by using a model of five-dimensional gravity with cosmological constant and argued that this may indicate that the family of HT-singular instantons would not exist.
But if we take the HT-singular instantons as ‘constrained gravitational instantons ,’ the singularity does not need the justification by five-dimensional regularity. Hence in this paper we use the five-dimensional theory as a mere trick only to calculate the contribution of HT-singularity. Then the difficulty mentioned above disappears (see also the discussions in the final section). We show that near HT-singularity any HT-singular instantons can be approximated as five-dimensional regular metrics on compactification. It turns out that the HT-singularity corresponds to a bolt in five dimensions. Thus we can calculate the fixed point contribution of HT-singularities by evaluating the bolt contribution in five dimensions. The use of an extra dimension may be seen as a method to constrain HT-singularities.
First we review the behaviour of HT-singular instantons. The action of the HT-singular instantons is given by
$$𝒮=d^4x\sqrt{g}\left(\frac{}{16\pi G}+\frac{1}{2}(\varphi )^2+V(\varphi )\right)+,$$
(11)
where $``$ represents the boundary contribution $`𝒮(\text{boundary})`$ and the fixed point contribution of the HT-singularity $`𝒮(\text{FP})`$. We seek solutions of $`O(4)`$ symmetric metric
$$ds^2=d\sigma ^2+b^2(\sigma )d\mathrm{\Omega }_3^2=d\sigma ^2+b^2(\sigma )\left(d\psi ^2+\mathrm{sin}^2(\psi )d\mathrm{\Omega }_2^2\right).$$
(12)
The equations of motion read
$$\varphi ^{\prime \prime }+3\frac{b^{}}{b}\varphi ^{}=\frac{V}{\varphi },b^{\prime \prime }=\frac{8\pi G}{3}b\left(\varphi ^{\mathrm{\hspace{0.17em}2}}+V\right)$$
(13)
where primes denote derivatives with respect to $`\sigma `$. There are two choices of boundary conditions:
1) Hawking–Turok’s case . The metric and the scalar field are non-singular at $`\sigma =0`$, i.e.
$$b(\sigma )\sigma (\sigma 0),\varphi ^{}(0)=0.$$
(14)
We can set the value $`\varphi (0)=\varphi _0`$ arbitrarily except values which correspond to regular instanton solutions.
2) Asymptotically flat case (including Vilenkin’s case ).
$$b(\sigma )\sigma ,\varphi (\sigma )\varphi _{\mathrm{}}(\sigma \mathrm{}),$$
(15)
where $`\varphi _{\mathrm{}}`$ represents the value which the field $`\varphi `$ takes at the extremum of the potential, $`\frac{V}{\varphi }(\varphi _{\mathrm{}})=0`$ and $`V(\varphi _{\mathrm{}})=0`$.
Because $`b`$ goes to zero as we approach the singularity, the (anti-)damping term $`3b^{}\varphi ^{}/b`$ in the equation of motion overwhelms the potential term $`V/\varphi `$ near singularity, if the potential is not too steep there. Hence HT-singular instantons have the similar behaviours near the singularity $`\sigma =\sigma _{}`$, regardless of the potential and the choice of the boundary conditions, as
$`b^3(\sigma )`$ $``$ $`\kappa C\left|\sigma _{}\sigma \right|`$ (16)
$`\varphi (\sigma )`$ $``$ $`\pm {\displaystyle \frac{1}{\kappa }}\mathrm{log}\left|\sigma _{}\sigma \right|+\mathrm{𝑐𝑜𝑛𝑠𝑡}.,`$ (17)
where $`\kappa =\sqrt{12\pi G}`$ and $`C`$ is a constant. We see that the approximation we used is consistent if the potential is
$$\text{not so steep as }V(\varphi )e^{2\kappa \left|\varphi \right|}\text{near singularity.}$$
(18)
Using Eq.(16), (17), (18) and the relation $`=6(bb^{\prime \prime }b^{\mathrm{\hspace{0.17em}2}}+1)/b^2`$, we can check that the bulk terms of the action Eq.(11) do not contribute at the singularity.
Secondly let us calculate the boundary contribution $`𝒮(\text{boundary})`$ at the singularity. From the whole spacetime $``$ we excise a four-sphere which is centered around the singularity and of coordinate radius $`\sigma _\epsilon `$ and call the remaining spacetime as $`^{}`$. We consider to take the limit of $`\sigma _\epsilon 0`$. Hereafter we only consider the boundary at singularity. The contribution of the boundary at infinity is restored, if any, in the end of the calculation, Eq.(34). Then the action contribution of the boundary of $`^{}`$ becomes as
$`𝒮(^{})`$ $`=`$ $`\underset{\sigma _\epsilon 0}{lim}𝒮_{GH}`$ (19)
$`=`$ $`{\displaystyle \frac{1}{8\pi G}}\left(ϵ{\displaystyle \frac{}{\sigma }}2\pi ^2b^3(\sigma )\right)_{\sigma \sigma _{}}`$ (20)
$`=`$ $`\sqrt{{\displaystyle \frac{3\pi ^3}{4G}}}C,`$ (21)
where $`^{}`$ denotes the boundary of $`^{}`$ and $`ϵ`$ is defined by
$$ϵ=\{\begin{array}{cc}+1\hfill & \text{Hawking–Turok’s case}\hfill \\ 1\hfill & \text{asymptotically flat case.}\hfill \end{array}$$
(22)
We see that the constant in the behaviour of $`\varphi `$ (Eq.(17)) is unimportant for the boundary contribution. This is because the Gibbons–Hawking term involves only $`b`$, and only $`\varphi ^{}`$ and $`b`$ are important in the equations of motion near singularity.
Next, we calculate the contribution of the excised spacetime. In order to find it, we at first calculate the action contribution of bolts in five dimensions, and then we show the correspondence of the bolts in five-dimensions to HT-singularities in four dimensions. The five-dimensional action is given by
$$𝒮=d^5x\sqrt{\stackrel{~}{g}}\left(\frac{1}{16\pi \stackrel{~}{G}}\stackrel{~}{}+\text{(other terms)}\right)\frac{1}{8\pi \stackrel{~}{G}}d^4\xi \sqrt{\stackrel{~}{h}}\stackrel{~}{K},$$
(23)
where the tildes distinguish five-dimensional quantities from their four-dimensional counterparts. We consider $`O(4)\times U(1)`$ symmetric solutions. As it turns out below, the submanifold where the fifth dimension closes becomes the HT-singularity in four dimensions. So we approximate the five-dimensional regular metric there as
$$d\stackrel{~}{s}^2d\tau ^2+R_0^2d\mathrm{\Omega }_3^2+\tau ^2d\theta ^2(\tau 0),$$
(24)
where $`R_0`$ is a constant and the regularity at $`\tau =0`$ requires that $`\theta `$ is identified with period $`2\pi `$. The submanifold $`\tau =0`$ is the bolt in five dimensions (three-dimensional fixed point set). Its contribution $`𝒮(\text{FP})`$ becomes
$`𝒮(\text{FP})`$ $`=`$ $`\underset{\tau 0}{lim}𝒮_{GH}`$ (25)
$`=`$ $`{\displaystyle \frac{1}{8\pi \stackrel{~}{G}}}{\displaystyle \frac{}{\tau }}\left(2\pi ^2R_0^32\pi \tau \right)_{\tau 0}`$ (26)
$`=`$ $`{\displaystyle \frac{\pi ^2R_0^3}{2\stackrel{~}{G}}}.`$ (27)
Then we show the correspondence of the bolts in five dimensions to HT-singularities in four dimensions. We take fields to be independent of the fifth coordinate $`\theta `$ and use the ansatz
$$\stackrel{~}{g}_{AB}=e^{\frac{2\kappa }{3}(\pm \varphi +a)}\left(\begin{array}{cc}g_{\mu \nu }& 0\\ 0& e^{2\kappa (\pm \varphi +a)}\end{array}\right)=\left(\begin{array}{cc}e^{\frac{2\kappa }{3}(\pm \varphi +a)}g_{\mu \nu }& 0\\ 0& e^{\frac{4\kappa }{3}(\pm \varphi +a)}\end{array}\right),$$
(28)
where $`a`$ is a constant. The addition of $`a`$ to $`\pm \varphi `$ is a mere shift of the zero point of $`\varphi `$ and does not affect the calculation of the contribution of singularity to the action, but it is added for clarity. We obtain the four-dimensional action which has appropriate coefficients
$$𝒮=d^4x\sqrt{g}\left(\frac{}{16\pi G}+\frac{1}{2}(\varphi )^2+\text{(other terms)}\right)\frac{1}{8\pi G}d^3\xi \sqrt{h}K.$$
(29)
$`G`$ and $`\stackrel{~}{G}`$ are related by $`\stackrel{~}{G}=G𝑑\theta =2\pi G`$. From Eq. (28), we see that the four-dimensional fields diverge when the fifth dimension closes. Namely the submanifold where the fifth dimension closes is identified as the singularity in four dimensions. The five-dimensional metric of Eq.(24) becomes
$`\tau ^3`$ $``$ $`{\displaystyle \frac{9}{4}}\left(\sigma _{}\sigma \right)^2`$ (30)
$`b^3(\sigma )`$ $``$ $`{\displaystyle \frac{3R_0^3}{2}}\left|\sigma _{}\sigma \right|`$ (31)
$`\varphi (\sigma )`$ $``$ $`{\displaystyle \frac{1}{\kappa }}\mathrm{log}\left|\sigma _{}\sigma \right|{\displaystyle \frac{1}{\kappa }}\mathrm{log}{\displaystyle \frac{3e^{\kappa a}}{2}},`$ (32)
near the bolt, $`\tau =0`$ i.e. $`\sigma =\sigma _{}`$. We see that any HT-singular behaviours can be obtained by choosing $`R_0`$ and $`a`$ properly.
Comparing Eq.(16) and Eq.(31), we have that $`3R_0^3/2=\kappa C`$, so the bolt contribution $`𝒮(\text{FP})`$ can be rewritten in terms of the four-dimensional quantities as $`𝒮(\text{FP})=\sqrt{\pi ^3C^2/3G}=\frac{2}{3}𝒮(^{})`$. Hence the total action $`𝒮()`$ becomes
$`𝒮()`$ $`=`$ $`𝒮(^{})+𝒮(\text{FP})`$ (33)
$`=`$ $`𝒮(\text{bulk})+{\displaystyle \frac{1}{3}}𝒮(^{})\left(+𝒮(\text{boundary at infinity})\right).`$ (34)
Thus the contribution of the HT-singularity to the action can be calculated as the one third of the Gibbons–Hawking term of the three-surface which wraps the singularity.
## 4 The Total Action
We can obtain a very simple formula for the total action of the system considered above. Interestingly, we have the same formula for both the regular and HT-singular instantons only if we use the above results, Eq.(34).
From the Einstein equations
$`b^{\prime \prime }`$ $`=`$ $`{\displaystyle \frac{8\pi G}{3}}b\left(\varphi ^{\mathrm{\hspace{0.17em}2}}+V\right)`$ (35)
$`b^{\mathrm{\hspace{0.17em}2}}`$ $`=`$ $`{\displaystyle \frac{8\pi G}{3}}b^2\left({\displaystyle \frac{1}{2}}\varphi ^{\mathrm{\hspace{0.17em}2}}V\right)+1,`$ (36)
we have
$$\left(b^3\right)^{\prime \prime }=3\left(8\pi Gb^3V+2b\right).$$
(37)
By integration, we obtain the relation
$$𝑑\sigma b^3V=\frac{1}{4\pi G}𝑑\sigma b\frac{1}{3}\left[\frac{1}{8\pi G}(b^3)^{}\right].$$
(38)
On the other hand, the trace of the Einstein equation reads
$$=8\pi G\left((\varphi )^2+4V(\varphi )\right).$$
(39)
Hence the total action Eq.(11) can be rewritten as
$`𝒮`$ $`=`$ $`{\displaystyle d^4x\sqrt{g}V(\varphi )}+`$ (40)
$`=`$ $`2\pi ^2{\displaystyle 𝑑\sigma b^3(\sigma )V(\varphi )}+`$ (41)
$`=`$ $`{\displaystyle \frac{\pi }{2G}}{\displaystyle 𝑑\sigma b(\sigma )}+{\displaystyle \frac{1}{3}}\left[{\displaystyle \frac{1}{8\pi G}}(2\pi ^2b^3)^{}\right]+.`$ (42)
In the case of compact regular instantons, the second and third terms vanish, so we obtain a simple formula
$$𝒮=\frac{\pi }{2G}𝑑\sigma b.$$
(43)
We observe that the total action is given by the area of the circle plotted by $`b(\sigma )`$ in the $`(\sigma ,b)`$-plane. In the case of Hawking–Turok instantons, the second term is non-vanishing at the singularity, but it is exactly cancelled by the third term
$$=\frac{1}{3}\left(\frac{1}{8\pi G}(2\pi ^2b^3)^{}\right)_{\sigma \sigma _{}}.$$
(44)
Therefore we have the same expression for regular and HT-singular instantons. Here the coefficient $`1/3`$ in front of the Gibbons–Hawking term which originates from the addition of $`𝒮(\text{FP})`$ is crucial. It is to be noted that by changing $`\varphi _0`$ we can obtain both regular and HT-singular instantons when the potential $`V(\varphi )`$ has additional extrema.
In the case of HT-singular asymptotically flat instantons, the total action becomes
$`𝒮`$ $`=`$ $`\underset{s\mathrm{}}{lim}\left\{{\displaystyle \frac{\pi }{2G}}{\displaystyle _\sigma _{}^s}𝑑\sigma b(\sigma )+{\displaystyle \frac{\pi }{12G}}(b^3)^{}(s)\right\}.`$ (45)
Here we have dropped the boundary contributions at $`s`$ in the $``$ term because they cancels in the limit of $`s\mathrm{}`$ on account of the asymptotically flat condition.
## 5 Discussions
In conclusion, we showed in this paper that the HT-singularity has a fixed point contribution to the action and it is given by $`2/3`$ of the Gibbons–Hawking term of the three-surface which wraps the singularity and whose normal vectors point to the singularity, in the zero-volume limit of the three-surface. By adding the contribution to the action, we also obtained a simple formula for the action which is feasible for both regular and HT-singular instantons.
Two comments are in order.
At the singularity the equations of motion are not satisfied and so the action is not stationary. But the probability of quantum transition is calculated from the path integral which is integrated over instantons which is constrained by a given 3-metric on the considering three-surface. Hence the singularity at the constrained surface would be allowed. Motivated from this observation by Wu , it seems plausible that the HT-singularity in the cases of the creation of open universe and decay of flat spacetime would be allowed. In these cases the HT-singularity is on the quantum transition surface, so that the parameter $`C`$ (in other words, $`\varphi _0`$ or $`\sigma _{}`$) is constrained and the dependence of the action on it would be justified. (However, in ref. it was suggested to throw away the hemisphere which contains the HT-singularity and the whole manifold should be made by joining the remaining regular hemisphere and its oriented reversal in order to avoid the trouble caused by the HT-singularity.) Recently, Turok introduced a constraint and argued that Vilenkin’s instantons may be defined as a limit of constrained instantons which are non-singular (albeit with discontinuous first derivatives) and have no boundary. Using such instantons it was argued that the action of Vilenkin’s instantons is given by the same value calculated by Vilenkin , i.e. the whole Gibbons–Hawking term. This differs with our result, $`1/3`$ of the Gibbons–Hawking term, and so it seems doubtful that the obtained instantons are Vilenkin’s instantons. The suggested constraint may be an unappropriate one so that HT-singularity does not appear from the limit of such instantons.
As a second comment, we examine the five-dimensional theory consisting of gravity and a cosmological constant $`\stackrel{~}{\mathrm{\Lambda }}`$ to show how the difficulty mentioned in Sec. III disappears. In ref. it was stated that the regular instantons exist in this theory only when $`\varphi _0`$ is zero. However, careful examination reveals that this is not the case. The five-dimensional metric with $`O(4)\times U(1)`$ symmetry is
$$d\stackrel{~}{s}^2=d\tau ^2+\widehat{A}^2\mathrm{sin}^2\left(\frac{\tau }{\widehat{A}}\right)d\mathrm{\Omega }_3^2+\widehat{A}^2\mathrm{cos}^2\left(\frac{\tau }{\widehat{A}}\right)d\theta ^2,$$
(46)
where $`\widehat{A}=\sqrt{6/\stackrel{~}{\mathrm{\Lambda }}}`$. The coordinate $`\tau `$ runs from $`0`$ to $`\widehat{A}\pi /2`$. This metric is regular everywhere for any $`\widehat{A}`$. The HT-singularity appears at $`\tau =\widehat{A}\pi /2`$ on compactification. The behaviour of the scalar field near $`\sigma =0`$ (in this case it corresponds to $`\tau =0`$) becomes
$$\pm \varphi (\sigma )\frac{3}{2\kappa }\mathrm{log}\widehat{A}+\frac{3\sigma ^2}{4\kappa \widehat{A}^3}a,$$
(47)
and the potential becomes
$$V(\varphi )=\frac{\stackrel{~}{\mathrm{\Lambda }}}{8\pi G}e^{\frac{2\kappa }{3}(\pm \varphi +a)}=\frac{3}{4\pi G}e^{(2\mathrm{log}\widehat{A}+\frac{2\kappa }{3}a)}e^{\pm \frac{2\kappa }{3}\varphi }.$$
(48)
From these equations we see that we can have any $`\varphi _0`$ without changing the four-dimensional potential term by choosing $`a`$ and $`\stackrel{~}{\mathrm{\Lambda }}`$ appropriately.
We may interpret this calculation in terms of constraints on instantons as follows. Four-dimensional singular instantons are one parameter family in constrast to regular instantons and the action depends on the parameter. This parameter enters the instanton solution as an integration constant of the equations of motion. We saw in Sec. III that regular instantons in five dimensions whose fifth dimension closes show the HT-singular behaviour on compactification. In general there would be only a finite number of such regular instantons in a five-dimensional theory. They have one-to-one correspondence to four-dimensional HT-singular instantons of different parameter, so that we have a constraint on the parameter of the HT-singular instantons. If we choose one five-dimensional theory (i.e. the value of $`\stackrel{~}{\mathrm{\Lambda }}`$, in the above example), then one or some of HT-singular instantons are picked out from the family of HT-singular instantons. The parameter of the HT-singular instantons can be later integrated over in path integral by integrating over the appropriate parameter space of the constraining five-dimensional theory.
## Acknowledgements
I am grateful to Ken-Iti Izawa and Hiroaki Terashima for reading the manuscript and for useful comments.
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# Condensate fluctuations in finite Bose-Einstein condensates at finite temperature
## I Introduction
Bose-Einstein condensation in a weakly interacting Bose-gas in three dimensions in the thermodynamic limit of an infinitely extended system is a second order phase transition in which an order parameter, the macroscopic wave-function, appears spontaneously with a fixed but arbitrary phase, turning the global $`U(1)`$-gauge-symmetry connected with particle-number conservation into a spontaneously broken or hidden symmetry. The rigidity of the phase of the order parameter against local perturbations and the absence of any phase diffusion gives rise to the Goldstone modes, which take the form of collision-less (zero)-sound or hydrodynamic sound, respectively, depending on whether the sound frequency is in the collision-less mean-field regime or in the collision-dominated regime .
In finite systems, and thus also in all trapped Bose-gases, sharp phase-transitions are impossible and hidden symmetries in a rigorous sense cannot appear. However, a macroscopic wave-function describing a Bose-Einstein condensate still exists . Its phase cannot be stable and must undergo a diffusion process, which restores the $`U(1)`$ gauge-symmetry over sufficiently long time intervals. This diffusion process is different from the Goldstone modes mentioned before, which are oscillations around a fixed value of the phase and do not restore the symmetry. Rather the Goldstone modes show up either as collision-dominated hydrodynamic phonons or as collision-less phonons, which have also been observed in the finite Bose-Einstein condensates. In the present paper I would like to discuss the dynamics of the complex amplitude of a Bose-Einstein condensate containing a finite number of particles, and in particular analyse the diffusion of its phase. My discussion will extend and correct in several respects the work published in .
The stability of the phase-difference between the macroscopic wave functions of two Bose-Einstein condensates in a trap has been measured. In the experimental set-up the relative phase was measured using a time-domain separated-oscillatory-field condensate interferometer. Over the time-interval of 100 ms scanned in the experiment the relative phase was found to be robust. This experimental result demonstrates that the macroscopic wave functions of the condensates cannot be considered as quantum mechanical wave-functions of many particle systems entangled with each other, whose decoherence would indeed be extremely rapid. Rather, the macroscopic wave functions are appropriately viewed as robust classical objects, their quantum mechanical origin (just like magnets, crystals etc) notwithstanding. This does, of course not preclude that there may be quantum effects, for finite condensates, which lead to corrections of the dynamics described by the underlying classical wave-equation, the well-known Gross-Pitaevskii equation . In a number of papers the dispersion of the phase of a trapped Bose-Einstein condensate at zero-temperature was considered, which is due to fluctuations $`\delta \mu `$ of the chemical potential $`\mu `$ in a finite system with fixed particle number. An extension of this mechanism to finite temperature has also been proposed . This effect is not an irreversible phase diffusion but corresponds to an effect of inhomogeneous broadening, similar to the de-phasing of precessing spins occurring in spin-systems due to inhomogeneous broadening. Like the decay of the magnetization can be reversed in spin-echos, the decay of the order parameter expectation value in Bose-Einstein condensates due to a finite variance of $`\delta \mu `$ is in principle reversible in ‘revivals’. Experiments in Bose-Einstein condensation are done at temperature $`k_BT\mathrm{}\overline{\omega }`$ and often even at $`k_BT\mu `$, where $`\overline{\omega }`$ is the geometrical mean of the three main trap frequencies. A process of phase-diffusion process should occur in such a regime due to the interaction of the condensate with a thermal bath of collective modes and quasiparticles. An estimate of this phase-diffusion is of interest for the theory of atom lasers, because the fundamental limit of the line-width of an atom laser for a given temperature depends on it similar to the ‘Schawlow-Townes’-formula for the line-width of a laser.
In this paper a theory of dissipation and thermal fluctuations of a trapped Bose-Einstein condensate will be formulated. First a phenomenological framework for the theory in the form of a Langevin equation will be given in which dissipation appears via a phenomenological parameter and the fluctuation-dissipation relation is invoked to relate it to maximally three intensity-coefficients of the fluctuations. The solution of the Langevin-equation then determines the relaxation of the condensate-number and the diffusion of the phase, quite similar to the dynamics of a laser-amplitude above threshold. Then the Langevin equation is derived from the microscopic theory and formulas for the phenomenological parameters are derived. These are evaluated for a box-like trap and an isotropic harmonic trap-potential as a function of temperature, particle-number and scattering length. The final section contains a discussion of our results and a comparison with earlier related work. The theory presented here makes no claim to apply to the critical regime, nor can we examine here to what extent it covers the regime below but close to $`T_c`$, where it may be important to take the dynamics of the thermal cloud of non-condensed atoms into account in addition to the excitations from the condensate.
## II Microscopic equations of motion
The weakly interacting Bose-gas in a trap in standard notation is described by the Hamiltonian
$$\widehat{H}=d^3x\widehat{\psi }^+\left\{\frac{\mathrm{}^2}{2m}^2+V(𝒙)\mu +\frac{U_0}{2}\widehat{\psi }^+\widehat{\psi }\right\}\widehat{\psi }.$$
(1)
The total number of atoms $`N`$ is fixed, i.e. the Hilbert space is the restriction of the Fock-space of $`\widehat{\psi }`$ to the subspace on which $`\widehat{N}=N`$ is satisfied. $`\mu `$ is the average of the chemical potential, which is a fluctuating quantity in a system where N is fixed. Later-on we shall denote the fluctuating part of the chemical potential with $`\mathrm{\Delta }\mu `$. The presence of a Bose-Einstein condensate in equilibrium means that many ($`N_01`$) particles occupy a single mode of a macroscopic classical matter wave, determined as the mode of lowest energy of the classical Hamiltonian corresponding to eq.(1). The latter is obtained by replacing in $`H`$ the field operator $`\widehat{\psi }(𝒙)`$ by the classical field $`\psi (𝒙)=\sqrt{N_0}\mathrm{exp}(i\varphi )\stackrel{~}{\psi }_0(𝒙)`$. We shall restrict our attention to sufficiently low temperatures below the critical temperature $`T_c`$ so that the interaction of the condensate with the mean field of the thermal cloud of non-condensed particles is negligible. In this way one finds that the condensate mode $`\stackrel{~}{\psi }_0(𝒙)`$, which we take to be normalized to 1, satisfies the Gross-Pitaevskii equation
$$(\mathrm{}^2/2m)^2\stackrel{~}{\psi }_0+\left(V(𝒙)+U_0N_0|\stackrel{~}{\psi }_0(𝒙)|^2\right)\stackrel{~}{\psi }_0=\mu \stackrel{~}{\psi }_0.$$
(2)
For given $`N_0`$ the average value of the chemical potential $`\mu `$ follows by imposing the normalization condition
$$d^3x|\stackrel{~}{\psi }_0(𝒙)|^2=1$$
(3)
on the solution of the Gross-Pitaevskii equation, and thereby $`\mu `$ like $`\stackrel{~}{\psi }_0(𝒙)`$ becomes a function of the mean atom number in the condensate $`N_0`$. As an important consequence of this fact the chemical potential of the system can be expressed as a function of the average number of atoms in the condensate alone. $`N_0`$ differs from $`N`$, the fixed total number of atoms, by the average number $`N^{}`$ of non-condensed atoms, which needs to be calculated for given $`N_0`$. The condition $`N=N_0+N^{}`$ then fixes $`N_0`$ self-consistently. In the experimentally realized Bose-Einstein condensates it is possible to measure $`N_0`$ directly with reasonable accuracy as a function of temperature, and in practice it is therefore reasonable to regard $`N_0`$ as an experimentally given and known function of temperature. The space-dependent mean number density of the condensate is $`n_0(𝒙)=N_0|\stackrel{~}{\psi }_0(𝒙)|^2`$. We shall take the mode function $`\stackrel{~}{\psi }_0(x)`$ in the Gross-Pitaevskii equation as real and positive. (This also means we are not considering condensates containing vortices). The physical phase of the condensate is not carried by its mode-function $`\stackrel{~}{\psi }_0`$ but by its complex amplitude denoted as $`\alpha _0`$, where $`\alpha _0=\sqrt{N_0}\mathrm{exp}i\varphi `$.
If $`|\alpha _0|^2`$ makes a small fluctuation away from its equilibrium value $`N_0`$ the condensate mode function $`\psi _0`$ will no longer satisfy eq.(2) but will change its form slightly. We shall assume that such fluctuations of $`N_0=|\alpha _0|^2`$ occur on a sufficiently large time-scale, that the new form is again determined by the Gross-Pitaevskii equation, but for the changed condensate number $`|\alpha _0|^2`$ and a correspondingly changed chemical potential $`\mu _0`$ determined uniquely by $`|\alpha _0|^2`$, i.e. in eq.(2) the replacements $`(\stackrel{~}{\psi }_0,N_0,\mu )(\psi _0,|\alpha _0|^2,\mu _0)`$ have to be made in this case,
$$(\mathrm{}^2/2m)^2\psi _0+\left(V(𝒙)+U_0|\alpha _0|^2|\psi _0(𝒙)|^2\right)\psi _0=\mu _0\psi _0.$$
(4)
We cannot expect, in general, that in any given nonequilibrium state the difference defined by $`\mathrm{\Delta }_0\mu =\mu _0\mu `$ is the total deviation of the chemical potential from its equilibrium value, because there may obviously be states with $`|\alpha _0|^2=N_0`$ which differ in other respects from the equilibrium state and may therefore have $`\mu \mu `$. Therefore we use the notation $`\mu _0`$ for the part of the nonequilibrium chemical potential determined by $`|\alpha _0|^2`$.
The presence of the highly occupied condensate mode makes the decomposition of the Heisenberg field-operator
$$\widehat{\psi }(𝒙,t)=\left(|\alpha _0|\mathrm{exp}(i\varphi )\psi _0(𝒙)+\widehat{\chi }(𝒙,t)\right)\mathrm{exp}(i\mu t/\mathrm{})$$
(5)
useful, where we follow Bogoliubov and describe the condensate classically. $`\widehat{\chi }(𝒙,t)`$ is taken to be the field operator for the particles outside the condensate. We shall assume that the temporal changes in $`\varphi `$ can be considered as slow on the time-scales of the dynamics of $`\widehat{\chi }`$. The phase $`\varphi `$ and amplitude $`|\alpha _0|`$ are additional c-number variables in (5). Therefore the taking of expectation values has from now on to include an integration over a distribution of $`|\alpha _0|`$ and in addition an integration over all values of $`\varphi `$. Since the total number $`N`$ is fixed $`\widehat{\psi }=0`$ must hold for all times. However, it will also be useful to consider expectation values in the Fock-space of the operators $`\widehat{\chi },\widehat{\chi }^+`$ alone without averaging over $`\varphi `$. Such expectation values will be denoted as $`\mathrm{}_\varphi `$.
Gauge invariance, strictly speaking, is lost by splitting off a c-number term from the field-operator. However, this symmetry is saved by adopting the rule that the phase $`\varphi `$ of the c-number term in the decomposition also changes under a gauge transformation according to $`\varphi \varphi +ϵ`$. By this device we take into account the fact that the same change of phase would have occurred automatically, if we had not replaced the condensate term by a c-number. The generator of gauge transformations is thus taken as
$$\widehat{N}=i\frac{}{\varphi }|_{\widehat{\chi },\widehat{\chi }^+}+d^3x\widehat{\chi }^+\widehat{\chi },$$
(6)
from which it is clear (cf. eq.(27)) that $`i\frac{}{\varphi }|_{\widehat{\chi },\widehat{\chi }^+}`$ is a representation of $`\widehat{N}_0`$.<sup>*</sup><sup>*</sup>*This operator with fixed $`\widehat{\chi },\widehat{\chi }^+`$ has to be well distinguished from the unrestricted derivative operator $`i(/\varphi )`$, which is a representation of the total particle number $`N`$ and has as formal canonical-conjugate $`\widehat{\varphi }`$ with $`\mathrm{exp}(i\widehat{\varphi })=\mathrm{exp}(/N)`$. The canonical-conjugate is the phase $`\widehat{\varphi }`$ with
$$\mathrm{exp}(i\widehat{\varphi })=\mathrm{exp}(/N_0)_{\widehat{\chi },\widehat{\chi }^+}.$$
(7)
Via (5) the Hamiltonian furthermore splits up according to $`\widehat{H}=H_0+\widehat{H_1}+\widehat{H_2}+\widehat{H_3}+\widehat{H_4}`$ where $`\widehat{H_n}`$ comprises the terms of $`\widehat{H}`$ which are of n-th order in $`\widehat{\chi },\widehat{\chi }^+`$. Explicitly
$$H_0=|\alpha _0|^2d^3x\psi _0\{\frac{\mathrm{}^2}{2m}^2+V(𝒙)\mu _0+\frac{U_0}{2}|\alpha _0|^2|\psi _0|^2\}\psi _0+(\mu _0\mu )|\alpha _0|^2$$
(8)
$$\widehat{H}_1=|\alpha _0|d^3x\{e^{i\varphi }\widehat{\chi }(\frac{\mathrm{}^2}{2m}^2+V(𝒙)\mu +U_0|\alpha _0|^2\psi _0^2)\psi _0+(h.c.)\}$$
(9)
$`\widehat{H}_2={\displaystyle }d^3x\{`$ $`\widehat{\chi }^+({\displaystyle \frac{\mathrm{}^2}{2m}}^2+V(𝒙)\mu _0)\widehat{\chi }`$ (12)
$`+{\displaystyle \frac{U_0}{2}}|\alpha _0|^2\psi _0^2(e^{2i\varphi }\widehat{\chi }^2+e^{2i\varphi }\widehat{\chi }^{+2}+4\widehat{\chi }^+\widehat{\chi })`$
$`+(\mu _0\mu )\widehat{\chi }^+\widehat{\chi }\}`$
$$\widehat{H}_3=U_0|\alpha _0|d^3x\psi _0\widehat{\chi }^+(e^{i\varphi }\widehat{\chi }+e^{i\varphi }\widehat{\chi }^+)\widehat{\chi }$$
(13)
$$\widehat{H}_4=\frac{U_0}{2}d^3x\widehat{\chi }^+\widehat{\chi }^+\widehat{\chi }\widehat{\chi }.$$
(14)
Using the Gross-Pitaevskii equation (4) and its derivative with respect to $`|\alpha _0|^2`$ we derive in the appendix
$$H_0=_{N_0}^{|\alpha _0|^2}𝑑N_0(\mu _0(N_0)\mu )$$
(15)
The term $`\widehat{H_1}`$ can be simplified using the Gross-Pitaevskii equation (4) and becomes then
$$\widehat{H}_1=|\alpha _0|(\mu _0\mu )d^3x(e^{i\varphi }\widehat{\chi }+e^{i\varphi }\widehat{\chi }^+)\psi _0$$
(16)
This expression will be seen to vanish below due to an orthogonality condition.
The first part of $`\widehat{H_2}`$ in (12) is diagonalized by introducing quasi-particle operators $`\widehat{\alpha }_\nu ,\widehat{\alpha }_\nu ^+`$ defined by the standard Bogoliubov transformation, with time-dependent $`\varphi (t)`$
$$\widehat{\chi }(𝒙)=e^{i\varphi }\underset{\nu }{}\left(u_\nu (𝒙)\widehat{\alpha }_\nu +v_\nu ^{}(𝒙)\widehat{\alpha }_\nu ^+\right).$$
(17)
The $`u_\nu `$, $`v_\nu `$ satisfy the usual Bogoliubov-Fetter equations
$$\left(\begin{array}{cc}\frac{\mathrm{}^2}{2m}^2+U_{\mathrm{eff}}(𝒙)\mathrm{}\omega _\nu & K(𝒙)\\ K(𝒙)& \frac{\mathrm{}^2}{2m}^2+U_{\mathrm{eff}}(𝒙)+\mathrm{}\omega _\nu \end{array}\right)\left(\genfrac{}{}{0pt}{}{u_\nu (𝒙)}{v_\nu (𝒙)}\right)=0,$$
(18)
with the abbreviations
$`U_{\mathrm{eff}}(𝒙)=V(𝒙)\mu _0+2U_0|\alpha _0|^2\psi _0(𝒙)^2`$ (19)
$`K(𝒙)=|\alpha _0|^2U_0\psi _0(𝒙)^2.`$ (20)
The Hamiltonian $`\widehat{H}_2`$ now takes the form
$$\widehat{H}_2=\underset{\nu }{}\mathrm{}\omega _\nu (\widehat{\alpha }_\nu ^+\widehat{\alpha }_\nu +|v_\nu |^2)+(\mu _0\mu +\mathrm{}\dot{\varphi })d^3x\widehat{\chi }^+\widehat{\chi }$$
(21)
The coefficients $`u_\nu ,v_\nu `$ and the mode frequencies $`\omega _\nu `$ become also functions of $`|\alpha _0|`$ and fluctuate (slowly) with that number. Their equilibrium values will be denoted by $`\stackrel{~}{u}_\nu ,\stackrel{~}{v}_\nu `$ and $`\stackrel{~}{\omega }_\nu `$ and the corresponding operator $`\widehat{\chi }`$ according to (17) as $`\widehat{\stackrel{~}{\chi }}`$.
Eq. (18) is consistent with the ortho-normality conditions
$`{\displaystyle d^3x(u_\nu u_\mu ^{}v_\nu v_\mu ^{})}=\delta _{\nu \mu }`$ (22)
$`{\displaystyle d^3r(u_\nu ^{}v_\mu u_\mu ^{}v_\nu )}=0`$ (23)
which guarantee the Bose commutation relations of the $`\alpha _\nu `$, $`\alpha _\mu ^+`$. A formal solution of eq. (18) at zero energy $`\mathrm{}\omega _\nu =0`$ is given by the condensate
$$u_\nu (𝒙)=v_\nu ^{}(𝒙)=\psi _0(𝒙),\omega _\nu =0,$$
(24)
This solution is obviously not normalizable in the required sense (22) to furnish an acceptable solution for the $`u_\nu ,v_\nu `$ and must therefore be excluded from the sum over the terms containing the operators $`\widehat{\alpha }_\nu ,\widehat{\alpha }_\nu ^+`$. The existence of this formal solution implies however that the properly normalizable solutions $`u_\nu ,v_\nu `$ and the condensate mode $`\psi _0`$ satisfy the important orthogonality relation
$$d^3x\psi _0(u_\nu +v_\nu )=0.$$
(25)
It follows with (17) that
$$d^3x\psi _0(e^{i\varphi }\widehat{\chi }+e^{i\varphi }\widehat{\chi }^+)=0$$
(26)
which in turn implies that the reduced expression (16) for $`\widehat{H}_1`$ vanishes. Using the property (26) one can verify that the decomposition (5) of $`\widehat{\psi }`$ implies
$$N=|\alpha _0|^2+\widehat{N^{}}$$
(27)
with
$`\widehat{N^{}}={\displaystyle d^3x\widehat{\chi }^+}`$ $`(𝒙)\widehat{\chi }(𝒙)`$ (28)
$`={\displaystyle \underset{\nu ,\mu }{}}{\displaystyle d^3x}`$ $`(\widehat{\alpha }_\nu ^+\widehat{\alpha }_\mu (u_\nu ^{}u_\mu +v_\nu ^{}v_\mu )+{\displaystyle \frac{1}{2}}\widehat{\alpha }_\nu \widehat{\alpha }_\mu (u_\nu v_\mu +v_\nu u_\mu )`$ (29)
$`+`$ $`{\displaystyle \frac{1}{2}}\widehat{\alpha }_\nu ^+\widehat{\alpha }_\mu ^+(u_\nu ^{}v_\mu ^{}+v_\nu ^{}u_\mu ^{})+\delta _{\nu \mu }|v_\nu |^2)`$ (30)
which serves as a definition of $`N_0=|\alpha _0|^2`$. The mean thermal density $`n^{}`$ in equilibrium can now be determined via
$$n^{}(𝒙)=\widehat{\stackrel{~}{\chi }}^+(𝒙)\widehat{\stackrel{~}{\chi }}(𝒙)=\underset{\nu }{}\left\{(|\stackrel{~}{u}_\nu (𝒙)|^2+|\stackrel{~}{v}_\nu (𝒙)|^2)\overline{n}_\nu +|\stackrel{~}{v}_\nu (𝒙)|^2\right\}$$
(31)
with $`\overline{n}_\nu =(\mathrm{exp}(\beta \mathrm{}\stackrel{~}{\omega }_\nu )1)^1`$.
The fluctuations of $`N_0`$ are similarly fixed by
$`\mathrm{\Delta }N_0^2=N_0^2`$ $`N_0^2=\mathrm{\Delta }\widehat{N}^2`$ (32)
$`={\displaystyle \underset{\nu }{}}{\displaystyle \underset{\nu ^{}}{}}\{`$ $`\overline{n}_\nu (\overline{n}_\nu ^{}+1)\left|{\displaystyle d^3x(\stackrel{~}{u}_\nu ^{}(x)\stackrel{~}{u}_\nu ^{}(x)+\stackrel{~}{v}_\nu ^{}(x)\stackrel{~}{v}_\nu ^{}(x))}\right|^2`$ (34)
$`+(\overline{n}_\nu \overline{n}_\nu ^{}+{\displaystyle \frac{1}{2}}(\overline{n}_\nu +\overline{n}_\nu ^{}+1))\left|{\displaystyle }d^3x(\stackrel{~}{u}_\nu (x)\stackrel{~}{v}_\nu ^{}(x)+\stackrel{~}{u}_\nu ^{}(x)\stackrel{~}{v}_\nu (x))|^2\right\}`$
They have been evaluated in and are also needed below (see eq.(210)). For work in the mathematical physics literature on number-fluctuations in the condensate of the ideal Bose-gas and models of the interacting Bose-gas see and references given there. For an alternative proposal to define and calculate the number-fluctuations in a Bose-condensate see .
After the transformation (17) the Hamiltonian is now in the form
$$\widehat{H}=H_0+\widehat{H}_2+\widehat{H}_3+\widehat{H}_4$$
(35)
with $`H_0,\widehat{H}_2,\widehat{H}_3,\widehat{H}_4`$ given by eqs.(15,21,13,14).
## III Langevin-equation of the condensate amplitude
Neither the Gross-Pitaevskii equation nor the Bogoliubov-Fetter equations furnish an equation for the condensate amplitude $`\alpha _0=\sqrt{N_0}\mathrm{exp}i\varphi `$. To find such an equation phenomenologically we first turn to a macroscopic quantity like the entropy $`S(|\alpha _0|^2,N)`$ for a fixed particle number N, but restricted to a fixed arbitrary value of $`\alpha _0=\sqrt{N_0}\mathrm{exp}(i\varphi )`$, where $`N_0`$ is the instantaneous number of particles in the condensate and different from the equilibrium value $`N_0`$ corresponding to the maximum of $`S(|\alpha _0|^2,N)`$. Thus $`N_0`$ is a function of N. The fluctuations of $`N_0`$ in the closed system formed by the trapped condensate after the evaporative cooling has been switched off are governed by a canonical Boltzmann-Einstein distribution
$`P(N_0)=\mathrm{\Omega }^1\mathrm{exp}\left(S(|\alpha _0|^2,N)/k_B\right).`$
We shall restrict ourselves to temperatures in the condensed regime outside the critical regime, where $`N_0(N)`$ is much larger than its root mean square $`\sqrt{\mathrm{\Delta }N_0^2(N)}=\sqrt{\mathrm{\Delta }\widehat{N^{}}^2}=(\widehat{N^{}}^2\widehat{N^{}}^2)^{1/2}`$, which is also a function of N. Then $`S(|\alpha _0|^2,N)`$, expanded to lowest order around its maximum, takes the form
$`S(|\alpha _0|^2,N)=S^{(eq)}(N)+\mathrm{\Delta }S(|\alpha _0|^2,N)`$
with
$$\mathrm{\Delta }S(|\alpha _0|^2,N)=k_B\frac{(|\alpha _0|^2N_0)^2}{2\mathrm{\Delta }N_0^2}$$
(36)
The entropy $`S(|\alpha _0|^2,N)`$ not only determines the equilibrium-distribution of the condensate amplitude, but appears also in its equation of motion, both in the conservative part of the dynamics as a conserved quantity, and in the dissipative part as a potential for the irreversible part of the dynamics. Let us first consider both parts separately.
The conservative part of the dynamics of $`\alpha _0`$ is connected with the dynamics of its phase $`\varphi `$. According to eqs.(5,17) a change of $`\varphi `$ changes the total phase of the field-operator $`\widehat{\psi }`$. For this reason the dynamics of $`\varphi `$ is given by the equation of motion
$$\dot{\varphi }=\frac{1}{\mathrm{}}\frac{\widehat{H}}{N}=\frac{1}{\mathrm{}}\mathrm{\Delta }\mu $$
(37)
where $`\mathrm{\Delta }\mu `$ is the deviation of the chemical potential from its equilibrium value. Such deviations may occur as a result of any fluctuations present in the system and, as discussed already, may in particular occur as a result of fluctuations of the value of $`N_0`$ away from its average $`N_0`$. This part of the fluctuation of $`\mu `$ we shall denote as $`\mathrm{\Delta }_0\mu `$. Expanding again to lowest order around the equilibrium $`N_0=N_0`$ we can write
$$\mathrm{\Delta }_0\mu =\frac{\mu }{N_0}(|\alpha _0|^2N_0)$$
(38)
The systematic part of the conservative part of the equation of motion of $`\alpha _0`$ can now be written in the form
$$\left(i\mathrm{}\dot{\alpha }_0\right)_{cons}=\mathrm{\Delta }_0\mu \alpha _0.$$
(39)
It is convenient to introduce the fluctuation of the free energy by $`\mathrm{\Delta }F=T\mathrm{\Delta }S`$. The dynamics (39) conserves $`|\alpha _0|^2`$ and $`\mathrm{\Delta }F`$. In equilibrium the right-hand side of this equation vanishes, because there $`\mathrm{\Delta }_0\mu =0`$, and the total phase of the condensate $`\varphi \mu t/\mathrm{}`$ changes only with a rate given by the average chemical potential $`\mu `$ in equilibrium.
The dissipative part of the equation of motion of $`\alpha _0`$ near thermal equilibrium is written with the help of $`\mathrm{\Delta }F`$ in the general form
$$\mathrm{}\left(\dot{\alpha }_0\right)_{diss}=\mathrm{\Gamma }_0\frac{\mathrm{\Delta }F(|\alpha _0|^2,N)}{\alpha _0^{}}$$
(40)
which contains the positive phenomenological parameter $`\mathrm{\Gamma }_0`$ and describes the relaxation of $`N_0=|\alpha _0|^2`$ to its equilibrium-value $`N_0`$.
According to general principles of statistical thermodynamics the relaxation process (40) must be accompanied by some form of noise. Adding a noise-term the total Langevin-equation of $`\alpha _0`$ can be written in the form
$$i\mathrm{}\dot{\alpha }_0=\mathrm{\Delta }_0\mu \alpha _0i\mathrm{\Gamma }_0\frac{\mathrm{\Delta }F(|\alpha _0|^2,N)}{\alpha _0^{}}+\xi (t)\mathrm{exp}(i\varphi ).$$
(41)
Since the condensate amplitude $`\alpha _0`$ is a collective quantity the noise $`\xi (t)`$ can be assumed to be Gaussian due to the central limit theorem. In addition we shall assume $`\xi (t)`$ to be a white noise force. This means that the actual correlation time $`\tau _m`$ of the noise $`\xi `$ is assumed to be much smaller than the time-scale on which the dynamics of $`\alpha _0`$ is observed, an assumption which must be checked for its validity in any concrete microscopic description. (In the microscopic theory we describe later it is a consistent assumption because the relaxation rate $`\gamma _c`$ of $`|\alpha _0|^2`$ turns out to be small compared to the time-scale of motion in the trap). Thus we assume that $`\xi (t)=0`$ and
$`\xi ^{}(t)\xi (0)=`$ $`\mathrm{}k_BT(2\mathrm{\Gamma }_0+\mathrm{\Gamma }^{})\delta (t)`$ (42)
$`\xi (t)\xi (0)=`$ $`\mathrm{}k_BT(\mathrm{\Gamma }^{}+i\mathrm{\Gamma }^{\prime \prime })\delta (t)`$ (43)
where $`\mathrm{\Gamma }_0`$ reappears in (42) because of the fluctuation-dissipation theorem. The form of the Langevin-equation (41) generalizes the work in by taking into account a possible correlation of the phase of the condensate and of the Langevin force, which may exist in condensates with finite particle numbers due to gauge invariance, i.e. particle-number conservation. (However, it will turn out later that the coefficient $`\mathrm{\Gamma }^{\prime \prime }`$ vanishes in condensates with a real condensate mode i.e. without vortices, which can be understood generally as a consequence of time-reversal symmetry.) Gauge-invariance implies that the Langevin-equation for $`\alpha _0`$, including the fluctuating term, must be invariant under the transformation $`\varphi \varphi +ϵ`$. This makes it useful to write the fluctuating term as $`\mathrm{exp}(i\varphi (t))\xi (t)`$ where $`\xi (t)`$ is a complex noise source independent of $`\varphi `$, which, physically, describes the scattering of particles in the condensate with those outside.A simpler ansatz (see) ignores the $`\varphi `$-dependence of the Langevin-force in (41). Then gauge-invariance of the Fokker-Planck equation which is stochastically equivalent to the Langevin equation implies $`\mathrm{\Gamma }^{}=\mathrm{\Gamma }^{\prime \prime }=0`$, i.e. the complex noise $`F_0(t)`$ then has random phase-fluctuations which are completely uncorrelated with and equi-distributed with respect to the condensate-phase $`\varphi `$. Note however that this achieves gauge-invariance only in an ensemble sense, not for each individual stochastic physical realization which together form the ensemble. In contrast the form of the Langevin equation considered here does enforce gauge invariance for each stochastic realization. The coefficients $`\mathrm{\Gamma }^{},\mathrm{\Gamma }^{\prime \prime }`$ describe a possible correlation of the phases of $`F_0`$ and $`\alpha _0`$, i.e. the existence of a squeezing in the thermal bath of uncondensed particles, caused by the constraint of total particle-number conservation. We shall see that this effect actually does occur in finite condensates, i.e. the condensate mode imprints its (slowly) fluctuating phase on the non-condensed ’environment’ due to particle number conservation in such a way that the lowest-lying modes are nearly totally squeezed.
The multiplicative nature of the noise in (41) raises the question in which stochastic calculus this equation should be interpreted: in the sense of Ito, or Stratonovich, or in some intermediate sense? This will be specified in a moment. Within the Gauss- and Markoff-assumption the form of the noise-force with the same positive coefficient $`\mathrm{\Gamma }_00`$ appearing in the dissipative part (40) and two further real coefficients $`\mathrm{\Gamma }^{},\mathrm{\Gamma }^{\prime \prime }`$, is fixed by the requirement that the Langevin-equation must be consistent with the correct equilibrium distribution $`P(\alpha _0,\alpha _0^{})=Z^1\mathrm{exp}(\mathrm{\Delta }F(|\alpha _0|^2,N)/k_BT)`$ for the condensate. Splitting into real and imaginary parts eqs.(42,43) become
$`\mathrm{}(\xi (t))\mathrm{}(\xi (0))=`$ $`\mathrm{}k_BT(\mathrm{\Gamma }_0+\mathrm{\Gamma }^{})\delta (t)`$ (44)
$`\mathrm{}(\xi (t))\mathrm{}(\xi (0))=`$ $`\mathrm{}k_BT\mathrm{\Gamma }_0\delta (t)`$ (45)
$`\mathrm{}(\xi (t))\mathrm{}(\xi (0))=`$ $`{\displaystyle \frac{1}{2}}\mathrm{}k_BT\mathrm{\Gamma }^{\prime \prime }\delta (t).`$ (46)
Eq.(41) may now be rewritten as
$`{\displaystyle \frac{N_0}{t}}=`$ $`2{\displaystyle \frac{\mathrm{\Gamma }_0}{\mathrm{}}}(N_0{\displaystyle \frac{\mathrm{\Delta }F}{N_0}}k_BT)+{\displaystyle \frac{2}{\mathrm{}}}\sqrt{N_0}\mathrm{}(\xi (t))`$ (47)
$`{\displaystyle \frac{\varphi }{t}}=`$ $`{\displaystyle \frac{1}{\mathrm{}}}\mathrm{\Delta }_0\mu {\displaystyle \frac{1}{\mathrm{}\sqrt{N_0}}}\mathrm{}(\xi (t))`$ (48)
and must in this form be interpreted as a stochastic differential equation in the sense of ItoThen the Fokker-Planck equation corresponding to eqs.(47,48) is $`\mathrm{}P/t=2\mathrm{\Gamma }_0/N_0[N_0(\mathrm{\Delta }F/N_0+k_BT/N_0)P]`$ and has the desired equilibrium distribution $`\mathrm{exp}(\mathrm{\Delta }F/k_BT)`$. This implies that eq.(41) must be interpreted in some intermediate sense, which we need not specify here further. In order to obtain its version in the sense of Ito it is best to bring (47),(48) into the form (41) using the Ito calculus..
Eq.(48) can be compared with (37). This comparison reveals that $`\mathrm{}(\xi (t))`$ must describe the fluctuations of the chemical potential not caused by deviations of $`|\alpha _0|^2`$ from its equilibrium value but by other fluctuations in the system. We shall come back to this point in section VI below.
The three phenomenological coefficients $`\mathrm{\Gamma }_0,\mathrm{\Gamma }^{},\mathrm{\Gamma }^{\prime \prime }`$, are dimensionless, temperature dependent numbers, which must be determined from a microscopic theory. Only one of these coefficients, $`\mathrm{\Gamma }_0`$, is connected with the fluctuations of the number of condensed atoms. If fluctuations of the chemical potential due to other processes are neglected, i.e. $`\mathrm{}\xi (t)=0`$, the remaining two coefficients are fixed at
$$\mathrm{\Gamma }^{}=\mathrm{\Gamma }_0,\mathrm{\Gamma }^{\prime \prime }=0$$
(49)
This corresponds to the case of maximal squeezing of the noise in the phase-direction.
Linearizing with respect to the small fluctuations $`\delta N_0N_0`$ we find
$$\mathrm{}\delta \dot{N}_0=\frac{2k_BT}{\mathrm{\Delta }N_0^2}N_0\mathrm{\Gamma }_0\delta N_0+2\sqrt{N_0}\mathrm{}(\xi (t))$$
(50)
$$\mathrm{}\dot{\varphi }=\frac{\mu }{N_0}\delta N_0\frac{1}{\sqrt{N_0}}\mathrm{}(\xi (t))$$
(51)
Eq.(50) describes the relaxation of the condensate to the equilibrium at $`N_0=|\alpha _0|^2`$ on the time-scale
$$\tau _c=\frac{\mathrm{}\mathrm{\Delta }N_0^2}{2\mathrm{\Gamma }_0N_0k_BT}$$
(52)
and the thermal fluctuations around it with the correlation function
$$\delta N_0(t)\delta N_0(t^{})=\mathrm{\Delta }N_0^2e^{|tt^{}|/\tau _c}$$
(53)
The correlation time $`\tau _c`$ is an important time-scale of the problem. The noise sources $`\mathrm{}(\xi (t)),\mathrm{}(\xi (t))`$ must have correlation times short compared to $`\tau _c`$ in order to be well described by white noise.
On a time-scale much larger than the correlation time $`\tau _c`$ the fluctuations $`\delta N_0(t)`$ in the equation for the phase can also be considered as Gaussian white noise with correlation function $`2\tau _c\mathrm{\Delta }N_0^2\delta (t)`$. Using this long-time approximation in the equation for the phase $`\varphi `$ and taking the correlation of the effective white noise $`\delta N_0(t)`$ with $`\mathrm{}(\xi (t))`$ properly into account, $`\varphi (t)`$ is found to satisfy the Langevin equation of a Wiener process
$`d\varphi (t)=\sqrt{D_\varphi }dw`$
with $`(dw)^2=dt`$ and diffusion constant
$$D_\varphi =\frac{\mathrm{\Delta }N_0^2}{\mathrm{}N_0}\frac{\mu }{N_0}\left(\frac{\mathrm{\Delta }N_0^2}{k_BT}\frac{\mu }{N_0}\frac{1}{\mathrm{\Gamma }_0}+\frac{\mathrm{\Gamma }^{\prime \prime }}{\mathrm{\Gamma }_0}\right)+\frac{k_BT}{\mathrm{}N_0}(\mathrm{\Gamma }_0+\mathrm{\Gamma }^{}),$$
(54)
i.e
$$\left(\varphi (t)\varphi (0)\right)^2=D_\varphi |t|.$$
(55)
Eq.(54) agrees with the result of if we assume as in no squeezing in the noise, i.e. $`\mathrm{\Gamma }^{}=\mathrm{\Gamma }^{\prime \prime }=0`$ and $`\mu /N_0=k_BT/\mathrm{\Delta }N_0^2^1`$. Both assumptions will not be made in the present work, however (cf. also the corresponding discussion in the final section).
The expectation value $`\alpha _0(t)`$ then decays exponentially according to $`\alpha _0(t)=\sqrt{N_0}e^{\mathrm{\Delta }\nu |t|}`$ with the line-width $`\mathrm{\Delta }\nu `$ given by the Schawlow-Townes-type formula $`\mathrm{\Delta }\nu =\frac{1}{2}D_\varphi .`$
It is not difficult to solve eqs.(50,51) for the phase-fluctuations also on time-scales of the order of $`\tau _c`$. The result for the mean square of the phase-increment in time is
$$(\varphi (t)\varphi (0))^2=D_\varphi |t|+2\frac{\mu }{N_0}\left(k_BT\mathrm{\Gamma }^{\prime \prime }+\mathrm{\Delta }N_0^2\frac{\mu }{N_0}\right)\tau _c^2\left(e^{|t|/\tau _c}1\right).$$
(56)
It interpolates between the diffusive long-time behavior (55) for $`t>>\tau _c`$ and the short-time behavior for $`t\tau _c`$
$$(\varphi (t)\varphi (0))^2=\frac{k_BT}{\mathrm{}N_0}(\mathrm{\Gamma }_0+\mathrm{\Gamma }^{})|t|+\frac{1}{\mathrm{}^2}\frac{\mu }{N_0}\left(\frac{\mu }{N_0}\mathrm{\Delta }N_0^2+k_BT\mathrm{\Gamma }^{\prime \prime }\right)t^2.$$
(57)
The first term describes phase diffusion due to thermal fluctuations of the chemical potential on time scales much shorter than $`\tau _c`$. The second term describes a non-diffusive and in principle reversible phase-collapse with the collapse rate
$$\gamma _{collapse}=\frac{1}{\mathrm{}}\sqrt{\frac{\mu }{N_0}\left(\frac{\mu }{N_0}\mathrm{\Delta }N_0^2+k_BT\mathrm{\Gamma }^{\prime \prime }\right)}$$
(58)
including a contribution from the cross-correlation between both types of fluctuations.
## IV Microscopic derivation of the Langevin equation
The microscopic derivation of the equation of motion for the condensate amplitude $`\alpha _0`$ can be carried out by using the Hamiltonian (35). As we did for the phenomenological equations in the preceding section we wish to derive here the microscopic equation of motion only to first order in the deviation $`(|\alpha _0|\sqrt{N_0})`$ from equilibrium.
$`H_0`$ given by eq.(15) is the Hamiltonian, in mean field approximation, of the pure condensate. Its free equation of motion is
$$i\mathrm{}\dot{\alpha }_0=\frac{H_0}{\alpha _0^{}}=(\mu _0\mu )\alpha _0,$$
(59)
from which
$$\frac{d\varphi (t)}{dt}=(\mu _0(|\alpha _0(t)|^2)\mu )/\mathrm{}$$
(60)
follows. Let us first use (60) in (21) to simplify $`\widehat{H}_2`$ and then eliminate $`\widehat{H}_2`$ by proceeding to the Heisenberg picture with respect to it. This changes $`\widehat{\chi },\widehat{\chi }^+`$ in $`\widehat{H}_3,\widehat{H}_4`$ according to
$$\widehat{\chi }\widehat{\chi }(t)=e^{i\varphi (t)}\underset{\nu }{}(u_\nu \widehat{\alpha }_\nu e^{i\omega _\nu t}+v_\nu ^{}\widehat{\alpha }_\nu ^+e^{i\omega _\nu t})$$
(61)
and its adjoint<sup>§</sup><sup>§</sup>§For simplicity we disregard here the slow time-dependence of the frequencies $`\omega _\nu `$.. The transformed time-dependent Hamiltonians will be denoted as $`\widehat{H}_3(t),\widehat{H}_4(t)`$, but $`\widehat{H}_4(t)`$ will not be needed in the following.
The equation of motion of the condensate amplitude $`\alpha _0`$ now takes the form,For the canonically conjugate pair $`N_0,\varphi `$ at fixed $`\widehat{\chi },\widehat{\chi }^+`$ cf. eqs.(6,7) with the notation $`\widehat{H}(t)=H_0+\widehat{H}_3(t)+\widehat{H}_4(t)`$,
$$\mathrm{}\frac{d\varphi }{dt}=\left(\frac{\widehat{H}(t)}{|\alpha _0|^2}\right)_{\varphi ,\widehat{\chi },\widehat{\chi }^+}$$
(62)
$$\mathrm{}\frac{d|\alpha _0|^2}{dt}=\left(\frac{\widehat{H}(t)}{\varphi }\right)_{|\alpha _0|,\widehat{\chi },\widehat{\chi }^+}$$
(63)
We obtain
$`\mathrm{}{\displaystyle \frac{d\varphi }{dt}}`$ $`=\mathrm{\Delta }_0\mu {\displaystyle \frac{1}{\sqrt{N_0}}}\mathrm{}(\widehat{\xi }^{}(t)){\displaystyle \frac{1}{\sqrt{N_0}}}\delta \mathrm{}(\widehat{\xi }^{}(t))`$ (64)
$`\mathrm{}{\displaystyle \frac{d|\alpha _0|^2}{dt}}`$ $`=2\sqrt{N_0}\mathrm{}(\widehat{\xi }(t))+2\sqrt{N_0}\delta \mathrm{}(\widehat{\xi }(t))`$ (65)
with
$`\mathrm{}(\widehat{\xi }^{}(t))=`$ $`{\displaystyle \frac{1}{2}}U_0{\displaystyle d^3x(\stackrel{~}{\psi }_0+2N_0\frac{\stackrel{~}{\psi }_0}{N_0})\widehat{\stackrel{~}{\chi }}^+(t)(e^{i\varphi }\widehat{\stackrel{~}{\chi }}(t)+e^{i\varphi }\widehat{\stackrel{~}{\chi }}^+(t))\widehat{\stackrel{~}{\chi }}(t)}`$ (66)
$`\mathrm{}(\widehat{\xi }(t))=`$ $`{\displaystyle \frac{1}{2i}}U_0{\displaystyle d^3x\stackrel{~}{\psi }_0\widehat{\stackrel{~}{\chi }}^+(t)(e^{i\varphi }\widehat{\stackrel{~}{\chi }}(t)e^{i\varphi }\widehat{\stackrel{~}{\chi }}^+(t))\widehat{\stackrel{~}{\chi }}(t)}`$ (67)
The complex noise $`\xi (t)`$ in eqs.(41-48) should be identified with
$$\widehat{\xi }(t)=\mathrm{}(\widehat{\xi }^{}(t))+i\mathrm{}(\widehat{\xi }(t)).$$
(68)
It is indeed independent of $`\varphi `$ as required by gauge-invariance of the Langevin-equation, as can be seen from (64), (65) with (61).We shall see in the next section that $`\widehat{\xi }^{}(t)`$ can be replaced by a c-number.
To describe fluctuations around equilibrium we have replaced in the preceding expressions the quantities $`|\alpha _0|^2,\psi _0,\widehat{\chi },\widehat{\chi }^+`$ by their equilibrium expressions $`\sqrt{N_0},\stackrel{~}{\psi }_0,\widehat{\stackrel{~}{\chi }},\widehat{\stackrel{~}{\chi }}^+`$ and represent the difference in the nonequilibrium state by $`\delta \mathrm{}(\widehat{\xi }^{}(t)),\delta \mathrm{}(\widehat{\xi }(t))`$ in eqs.(64,65). Omitting these differences altogether amounts to neglecting the back-action of the condensate on the thermal reservoir, which describes not only a modification of the fluctuating forces, which can indeed be neglected for fluctuations around a stable thermodynamic equilibrium, but also dissipation. To take the latter into account we need to calculate the averages $`\delta \delta \mathrm{}(\widehat{\xi }^{}(t))_\varphi ,\delta \delta \mathrm{}(\widehat{\xi }(t))_\varphi `$ to lowest order in the interaction between the condensate and the thermal cloud of atoms. The form which these quantities must take is prescribed completely by the fluctuation-dissipation theorem and symmetry:
For the reversible phase-dynamics the back-action can only lead to a shift in the average chemical potential. Such shifts due to the interaction $`\widehat{H}_3`$ will be small and are neglected here. For the irreversible amplitude dynamics the fluctuation-dissipation theorem requires in addition the appearance of a dissipation term. If
$$S_{JJ}(tt^{})=\mathrm{}(\widehat{\xi }(t))\mathrm{}(\widehat{\xi }(t^{})_\varphi $$
(69)
is the correlation function of the fluctuating force in(65), the back action must modify eq.(65) to the form
$`\mathrm{}{\displaystyle \frac{d|\alpha _0|^2}{dt}}={\displaystyle \frac{4N_0}{\mathrm{}k_BT}}{\displaystyle _{\mathrm{}}^t}𝑑t^{}S_{JJ}(tt^{}){\displaystyle \frac{H_0(t^{})}{|\alpha _0|^2}}+2\sqrt{N_0}\mathrm{}(\widehat{\xi }(t))`$ (70)
The derivation of this equation is given in appendix B.
This stochastic differential equation still differs from the phenomenological equation (50) in two respects:
1. The noise still has a finite correlation time $`\tau _{mic}`$. We shall consider these correlation functions in more detail below. Taking the Markovian limit $`\tau _{mic}0`$ with
$`S_{JJ}(tt^{})=\mathrm{}k_BT\mathrm{\Gamma }_0\delta (tt^{})`$ (71)
eq.(70) becomes
$`{\displaystyle \frac{d|\alpha _0|^2}{dt}}=2{\displaystyle \frac{\mathrm{\Gamma }_0}{\mathrm{}}}N_0{\displaystyle \frac{H_0}{|\alpha _0|^2}}+{\displaystyle \frac{2}{\mathrm{}}}\sqrt{N_0}\mathrm{}(\widehat{\xi }(t))`$ (72)
2. The mean-field Hamiltonian $`H_0(|\alpha _0|^2)`$ appears in eqs.(70),(72) instead of the free energy $`\mathrm{\Delta }F(|\alpha _0|^2)`$. This is due to the fact that the influence of the thermal excitations on the energy are not yet taken into account. Doing this under isothermal or closed-system boundary conditions we should replace the energy $`H_0(|\alpha _0|^2)`$ by the free energy $`\mathrm{\Delta }F(|\alpha _0|^2)`$ or $`T\mathrm{\Delta }S(|\alpha _0|^2)`$, respectively.
This completes our derivation of the Langevin equation for the complex amplitude of the condensate.
## V Green-Kubo expressions for the transport coefficients
Let us now analyse the fluctuating forces in more detail. Inserting the Bogoliubov transformation (61) in (66, 67) the fluctuating forces take the form
$`\mathrm{}(\widehat{\xi }^{}(t))=`$ $`{\displaystyle \frac{1}{4}}{\displaystyle \underset{\kappa \nu \mu }{}}(((M_{\kappa ,\nu \mu }^{(1)})^{}+M_{\nu \mu ,\kappa }^{(2)})\widehat{\alpha }_\nu ^+\widehat{\alpha }_\mu ^+\widehat{\alpha }_\kappa e^{i(\stackrel{~}{\omega }_\kappa \stackrel{~}{\omega }_\nu \stackrel{~}{\omega }_\mu )t}+h.c.)`$ (74)
$`+(nonresonantterms).`$
$`\mathrm{}(\widehat{\xi }(t))=`$ $`{\displaystyle \frac{1}{4i}}{\displaystyle \underset{\kappa \nu \mu }{}}(((M_{\kappa ,\nu \mu }^{(1)})^{}M_{\nu \mu ,\kappa }^{(2)})\widehat{\alpha }_\nu ^+\widehat{\alpha }_\mu ^+\widehat{\alpha }_\kappa e^{i(\stackrel{~}{\omega }_\kappa \stackrel{~}{\omega }_\nu \stackrel{~}{\omega }_\mu )t}h.c.)`$ (77)
$`+(nonresonantterms).`$
Terms are called ’nonresonant’ if the frequencies of the quasi-particles cannot add up to zero. Such terms have not been written out explicitly because later-on we shall restrict ourselves to the resonance or rotating wave approximation in which they don’t contribute. The relevant matrix elements $`M^{(1)},M^{(2)}`$ are
$`M_{\kappa ,\nu \mu }^{(1)}=`$ $`2U_0{\displaystyle }d^3x\stackrel{~}{\psi }_0\stackrel{~}{v}_\nu (\stackrel{~}{u}_\kappa ^{}\stackrel{~}{u}_\mu +{\displaystyle \frac{1}{2}}\stackrel{~}{v}_\kappa ^{}\stackrel{~}{v}_\mu )+(\nu \mu )`$ (78)
$`M_{\nu \mu ,\kappa }^{(2)}=`$ $`2U_0{\displaystyle }d^3x\stackrel{~}{\psi }_0\stackrel{~}{u}_\nu ^{}(\stackrel{~}{v}_\mu ^{}\stackrel{~}{v}_\kappa +{\displaystyle \frac{1}{2}}\stackrel{~}{u}_\mu ^{}\stackrel{~}{u}_\kappa )+(\nu \mu )`$ (80)
and very similarly
$`M_{\kappa ,\nu \mu }^{(1)}=`$ $`2U_0{\displaystyle }d^3x(\stackrel{~}{\psi }_0+2N_0{\displaystyle \frac{\stackrel{~}{\psi }_0}{N_0)}})\stackrel{~}{v}_\nu (\stackrel{~}{u}_\kappa ^{}\stackrel{~}{u}_\mu +{\displaystyle \frac{1}{2}}\stackrel{~}{v}_\kappa ^{}\stackrel{~}{v}_\mu )+(\nu \mu )`$ (81)
$`M_{\nu \mu ,\kappa }^{(2)}=`$ $`2U_0{\displaystyle }d^3x(\stackrel{~}{\psi }_0+2N_0{\displaystyle \frac{\stackrel{~}{\psi }_0}{N_0)}})\stackrel{~}{u}_\nu ^{}(\stackrel{~}{v}_\mu ^{}\stackrel{~}{v}_\kappa +{\displaystyle \frac{1}{2}}\stackrel{~}{u}_\mu ^{}\stackrel{~}{u}_\kappa )+(\nu \mu )`$ (83)
The matrix-elements $`M^{(1)},M^{(2)}`$ coincide with $`M^{(1)},M^{(2)}`$ if the dependence of $`\stackrel{~}{\psi }_0`$ on $`N_0`$ is negligible or vanishes, as e.g. in homogeneous systems.
$`M_{\kappa ,\nu \mu }^{(1)}`$ and similarly $`M_{\kappa ,\nu \mu }^{(1)}`$ describes a scattering process in which one atom is scattered out of the condensate by the absorption of the two quasiparticles $`\nu ,\mu `$ from - and the emission of the new quasiparticle $`\kappa `$ into - the thermal bath. Likewise $`M_{\nu \mu ,\kappa }^{(2)}`$ and similarly $`M_{\nu \mu ,\kappa }^{(2)}`$ describes a scattering process where an incoming thermal quasiparticle $`\kappa `$ is absorbed, again an atom is kicked out from the condensate, and two quasiparticles $`\nu ,\mu `$ are emitted into the thermal bath. The scattering amplitudes for both processes are linearly superposed due to the phase-coherence of the condensate which exists on the time-scale of the relaxation process induced by the scattering process even if it is destroyed on a much longer time scale.
We can now calculate the correlation functions of the fluctuating forces. Their averages over the bath of quasi-particles vanish, $`\mathrm{}(\widehat{\xi }(t))=0=\mathrm{}(\widehat{\xi }(t))`$. Their second-order correlation functions are obtained as
$`\mathrm{}(\widehat{\xi }^{}(t))\mathrm{}(\widehat{\xi }^{}(t^{}))_\varphi `$ $`=`$ $`{\displaystyle \frac{1}{8}}{\displaystyle \underset{\kappa ,\nu ,\mu }{}}|\left((M_{\kappa ,\nu \mu }^{(1)})^{}+M_{\nu \mu ,\kappa }^{(2)}\right)|^2`$ (86)
$`\{\overline{n}_\kappa (\overline{n}_\nu +1)(\overline{n}_\mu +1)e^{i(\stackrel{~}{\omega }_\kappa \stackrel{~}{\omega }_\nu \stackrel{~}{\omega }_\mu )(tt^{})}`$
$`+\overline{n}_\nu \overline{n}_\mu (\overline{n}_\kappa +1)e^{i(\stackrel{~}{\omega }_\kappa \stackrel{~}{\omega }_\nu \stackrel{~}{\omega }_\mu )(tt^{})}\}`$
$`\mathrm{}(\widehat{\xi }(t))\mathrm{}(\widehat{\xi }(t^{}))_\varphi `$ $`=`$ $`{\displaystyle \frac{1}{8}}{\displaystyle \underset{\kappa ,\nu ,\mu }{}}|(M_{\kappa ,\nu \mu }^{(1)})^{}M_{\nu \mu ,\kappa }^{(2)}|^2`$ (89)
$`\{\overline{n}_\kappa (\overline{n}_\nu +1)(\overline{n}_\mu +1)e^{i(\stackrel{~}{\omega }_\kappa \stackrel{~}{\omega }_\nu \stackrel{~}{\omega }_\mu )(tt^{})}`$
$`+\overline{n}_\nu \overline{n}_\mu (\overline{n}_\kappa +1)e^{i(\stackrel{~}{\omega }_\kappa \stackrel{~}{\omega }_\nu \stackrel{~}{\omega }_\mu )(tt^{})}\}`$
$`\mathrm{}(\widehat{\xi }^{}(t))\mathrm{}(\widehat{\xi }(t^{}))_\varphi `$ $`=`$ $`{\displaystyle \frac{1}{8i}}{\displaystyle \underset{\kappa ,\nu ,\mu }{}}\{((M_{\kappa ,\nu \mu }^{(1)})^{}M_{\nu \mu ,\kappa }^{(2)})(M_{\kappa ,\nu \mu }^{(1)}+(M_{\nu \mu ,\kappa }^{(2)})^{})`$ (91)
$`(\overline{n}_\nu +1)(\overline{n}_\mu +1)\overline{n}_\kappa e^{i(\stackrel{~}{\omega }_\kappa \stackrel{~}{\omega }_\nu \stackrel{~}{\omega }_\mu )(tt^{})}`$
(94)
$`\left(M_{\kappa ,\nu \mu }^{(1)}(M_{\nu \mu ,\kappa }^{(2)})^{}\right)\left((M_{\kappa ,\nu \mu }^{(1)})^{}+M_{\nu \mu ,\kappa }^{(2)}\right)`$
$`(\overline{n}_\kappa +1)\overline{n}_\nu \overline{n}_\mu e^{i(\stackrel{~}{\omega }_\kappa \stackrel{~}{\omega }_\nu \stackrel{~}{\omega }_\mu )(tt^{})}\}`$
These correlation functions can be replaced by delta-functions provided that the frequency-sums contain a flat quasi-continuum of nearly resonant terms in a neighborhood of the resonance $`\stackrel{~}{\omega }_\kappa \stackrel{~}{\omega }_\nu \stackrel{~}{\omega }_\mu =0`$ which is broad compared to the damping rates we calculate here. This assumption will be satisfied in sufficiently large condensates. The strengths of the delta-functions can then be extracted from the expressions (86), (89), (94) by taking the time-averages $`_{\mathrm{}}^{\mathrm{}}d(tt^{})\widehat{\xi }(t)\widehat{\xi }(t^{})_\varphi `$ and $`_{\mathrm{}}^{\mathrm{}}d(tt^{})\widehat{\xi }^+(t)\widehat{\xi }(t^{})_\varphi `$.
$`\mathrm{}(\widehat{\xi }^{}(t)),\mathrm{}(\widehat{\xi }(t))`$ are here given as expressions involving operators. Provided the Markovian approximation is satisfied the average of their commutators over the quasi-particle bath are again given by delta functions in time. Explicitly we obtain for the coefficients of the delta-functions
$`{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑t[\mathrm{}(\widehat{\xi }^{}(t)),\mathrm{}(\widehat{\xi }^{}(0))]_\varphi =0=`$ $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑t[\mathrm{}(\widehat{\xi }(t)),\mathrm{}(\widehat{\xi }(0))]_\varphi `$ (95)
$`{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑t[\mathrm{}(\widehat{\xi ^{}}(t)),\mathrm{}(\widehat{\xi }(0))]_\varphi ={\displaystyle \frac{\pi }{8i}}{\displaystyle \underset{\kappa ,\nu ,\mu }{}}`$ $`\left(\left((M_{\kappa ,\nu \mu }^{(1)})^{}M_{\nu \mu ,\kappa }^{(2)}\right)\left(M_{\kappa ,\nu \mu }^{(1)}+(M_{\nu \mu ,\kappa }^{(2)})^{}\right)c.c\right)`$ (97)
$`\{(\overline{n}_\nu +1)(\overline{n}_\mu +1)\overline{n}_\kappa (\overline{n}_\kappa +1)\overline{n}_\nu \overline{n}_\mu \}\delta (\stackrel{~}{\omega }_\kappa \stackrel{~}{\omega }_\nu \stackrel{~}{\omega }_\mu )=0`$
It can easily be verified that the bracket $`\{\mathrm{}\}`$ in the last line of (97) vanishes if it is multiplied by the delta-function expressing energy-conservation. As a result the fluctuating force $`\widehat{\xi }`$ in the Markoffian limit can indeed be treated as a c-number and will henceforth again be denoted by $`\xi `$. This also serves as a nice consistency-check that it is indeed possible to treat the condensate classically, even after taking its interaction with the quasi-particles into account.
Let us now proceed to derive formulas for the three transport parameters $`\mathrm{\Gamma }_0,\mathrm{\Gamma }^{}`$ and $`\mathrm{\Gamma }^{\prime \prime }`$. From
$`2\mathrm{\Gamma }_0+\mathrm{\Gamma }^{}=`$ $`{\displaystyle \frac{1}{\mathrm{}k_BT}}{\displaystyle _{\mathrm{}}^+\mathrm{}}𝑑t\xi ^{}(t)\xi (0)_\varphi `$ (98)
$`\mathrm{\Gamma }^{}+i\mathrm{\Gamma }^{\prime \prime }=`$ $`{\displaystyle \frac{1}{\mathrm{}k_BT}}{\displaystyle _{\mathrm{}}^+\mathrm{}}𝑑t\xi (t)\xi (0)_\varphi `$ (99)
implied by eqs.(42,43) we obtain the representations
$`\mathrm{\Gamma }_0=`$ $`{\displaystyle \frac{1}{\mathrm{}k_BT}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}d(tt^{})\mathrm{}(\xi (t))\mathrm{}(\xi (t^{}))_\varphi `$ (100)
$`\mathrm{\Gamma }_0+\mathrm{\Gamma }^{}=`$ $`{\displaystyle \frac{1}{\mathrm{}k_BT}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}d(tt^{})\mathrm{}(\xi ^{}(t))\mathrm{}(\xi ^{}(t^{}))_\varphi `$ (101)
$`\mathrm{\Gamma }^{\prime \prime }=`$ $`{\displaystyle \frac{2}{\mathrm{}k_BT}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}d(tt^{})\mathrm{}(\xi ^{}(t))\mathrm{}(\xi (t^{}))_\varphi .`$ (102)
which have the form of Green-Kubo relations for the transport coefficients. Using the explicit form (74,77) of the fluctuating forces the thermal averages can be taken and the time-integrals in eqs.(101,100,102) can be carried out which leads to the formulas
$`\mathrm{\Gamma }_0={\displaystyle \frac{\pi }{2\mathrm{}k_BT}}{\displaystyle \underset{\kappa ,\nu ,\mu }{}}`$ $`|(M_{\kappa ,\nu \mu }^{(1)})^{}M_{\nu \mu ,\kappa }^{(2)}|^2\overline{n}_\nu \overline{n}_\mu (\overline{n}_\kappa +1)\delta (\stackrel{~}{\omega }_\kappa \stackrel{~}{\omega }_\mu \stackrel{~}{\omega }_\nu )`$ (103)
$`\mathrm{\Gamma }_0+\mathrm{\Gamma }^{}={\displaystyle \frac{\pi }{2\mathrm{}k_BT}}{\displaystyle \underset{\kappa ,\nu ,\mu }{}}`$ $`|(M_{\kappa ,\nu \mu }^{(1)})^{}+M_{\nu \mu ,\kappa }^{(2)}|^2\overline{n}_\nu \overline{n}_\mu (\overline{n}_\kappa +1)\delta (\stackrel{~}{\omega }_\kappa \stackrel{~}{\omega }_\mu \stackrel{~}{\omega }_\nu )`$ (104)
$`\mathrm{\Gamma }^{\prime \prime }={\displaystyle \frac{i\pi }{2\mathrm{}k_BT}}{\displaystyle \underset{\kappa ,\nu ,\mu }{}}`$ $`\{((M_{\kappa ,\nu \mu }^{(1)})^{}M_{\nu \mu ,\kappa }^{(2)})(M_{\kappa ,\nu \mu }^{(1)}+(M_{\nu \mu ,\kappa }^{(2)})^{})`$ (107)
$`(M_{\kappa ,\nu \mu }^{(1)}(M_{\nu \mu ,\kappa }^{(2)})^{})((M_{\kappa ,\nu \mu }^{(1)})^{}+M_{\nu \mu ,\kappa }^{(2)})\}`$
$`\overline{n}_\nu \overline{n}_\mu (\overline{n}_\kappa +1)\delta (\stackrel{~}{\omega }_\kappa \stackrel{~}{\omega }_\mu \stackrel{~}{\omega }_\nu )`$
These expressions constitute our general results for the three transport parameters. They have to be evaluated separately for each individual trap geometry.
## VI Relation to the fluctuation and dissipation of the excitations
As was pointed out after eq.(48) by phenomenological arguments, the noise term $`\mathrm{}(\xi (t))`$ is not connected with the fluctuations of the number of particles in the condensate, but must be due to other fluctuations, which are then necessarily thermal fluctuations of the amplitudes of the excited states. In our microscopic results this can be seen from the fact that the fluctuating force $`\mathrm{}(\widehat{\xi }^{}(t))`$ according to eq.(66) contains precisely the same operator which also appears in $`\widehat{H}_3(t)`$ and couples the atoms in the thermal cloud to the condensate.
In the special case where the difference between the coupling matrix-elements $`M^{(1,2)}`$ and $`M^{(1,2)}`$ is negligible (which is exactly satisfied in box-like traps, cf.section VII) the intensity $`\mathrm{\Gamma }_0+\mathrm{\Gamma }^{}`$ of the noise source can be expressed entirely as a property of the excitations, as we shall now demonstrate.In the general case the coupling of the condensate to the non-condensate modes differs from the coupling between the non-condensate modes and the relation between $`\mathrm{\Gamma }_0+\mathrm{\Gamma }^{}`$ and the $`\gamma _\nu `$ is less direct. For their amplitudes $`\widehat{\alpha }_\nu (t),\widehat{\alpha }_\nu ^+(t)`$ a quantum-Langevin equation can be derived microscopically along the same lines employed here for the condensate amplitude. We have done this elsewhere , see also with the result, in the Markoffian limit,
$`{\displaystyle \frac{d\widehat{\alpha }_\nu (t)}{dt}}=i\omega _\nu \widehat{\alpha }_\nu (t)\gamma _\nu \widehat{\alpha }_\nu (t)+\widehat{\xi }_\nu (t)`$ (108)
with Gaussian fluctuating force-operators with vanishing average and
$`\widehat{\xi }_\nu ^+(t)\widehat{\xi }_\mu (t^{})`$ $`=2\gamma _\nu \overline{n}_\nu \delta (tt^{})\delta _{\nu \mu }`$ (109)
$`[\widehat{\xi }_\nu (t),\widehat{\xi }_\mu ^+(t^{})]`$ $`=2\gamma _\nu \delta (tt^{})\delta _{\nu \mu }`$ (111)
where the damping rates $`\gamma _\nu `$ are given by
$`\gamma _\nu ={\displaystyle \frac{\pi N_0}{\mathrm{}^2}}{\displaystyle \underset{\kappa ,\mu }{}}\{`$ $`|(M_{\kappa ,\nu \mu }^{(1)})^{}+M_{\nu \mu ,\kappa }^{(2)}|^2(\overline{n}_\mu \overline{n}_\kappa )\delta (\omega _\kappa \omega _\mu \omega _\nu )`$ (113)
$`+|M_{\nu ,\kappa \mu }^{(1)}+(M_{\kappa \mu ,\nu }^{(2)})^{}|^2(\overline{n}_\kappa +{\displaystyle \frac{1}{2}})\delta (\omega _\kappa +\omega _\mu \omega _\nu )\}.`$
The first term describes Landau-damping of the mode $`\nu `$ by scattering a quasi-particle from mode $`\mu `$ to mode $`\kappa `$ and is equivalent to a result derived in by the golden rule. The second term in eq.(113) describes Beliaev damping, where the mode $`\nu `$ decays into two modes $`\kappa ,\mu `$. It survives even for $`T0`$ where $`\overline{n}_\kappa 0`$ for all modes.
Let us now establish the connection between $`\mathrm{\Gamma }_0+\mathrm{\Gamma }^{}`$ and the damping rates $`\gamma _\nu `$ as given by (113). We shall show that the simple sum-rule
$`\mathrm{\Gamma }_0+\mathrm{\Gamma }^{}={\displaystyle \frac{\mathrm{}}{3N_0k_BT}}{\displaystyle \underset{\nu }{}}\overline{n}_\nu (\overline{n}_\nu +1)\gamma _\nu `$ (114)
holds. To see this we need to consider
$`{\displaystyle \underset{\nu }{}}\overline{n}_\nu (\overline{n}_\nu +1)\gamma _\nu `$ $`=`$ $`{\displaystyle \frac{\pi N_0}{\mathrm{}^2}}{\displaystyle \underset{\kappa \mu \nu }{}}|(M_{\kappa ,\nu \mu }^{(1)})^{}+M_{\nu \mu ,\kappa }^{(2)}|^2\delta (\omega _\kappa \omega _\nu \omega _\mu )`$ (116)
$`\left\{(\overline{n}_\mu \overline{n}_\kappa )\overline{n}_\nu (\overline{n}_\nu +1)+{\displaystyle \frac{1}{2}}(\overline{n}_\mu +\overline{n}_\nu +1)\overline{n}_\kappa (\overline{n}_\kappa +1)\right\}`$
The second term in the curly bracket arises from the second term in (113) by first exchanging the notations for the summation indices $`\nu `$ and $`\kappa `$ and then symmetrizing in $`\nu `$ and $`\mu `$, because the matrix-elements are already symmetric in these indices. The remainder of the proof then consists simply in noting that for $`\omega _\kappa =\omega _\nu +\omega _\mu `$ the identities
$`(\overline{n}_\mu \overline{n}_\kappa )\overline{n}_\nu (\overline{n}_\nu +1)`$ $`=\overline{n}_\mu \overline{n}_\nu (\overline{n}_\kappa +1)`$ (117)
$`(\overline{n}_\mu +\overline{n}_\nu +1)\overline{n}_\kappa (\overline{n}_\kappa +1)`$ $`=\overline{n}_\mu \overline{n}_\nu (\overline{n}_\kappa +1)`$ (119)
hold. Using this in eq.(116) and then comparing with (104) establishes the sum-rule. We can also note that the processes due to Landau scattering contribute to the sum-rule with precisely twice the strength of those due to Beliaev scattering.
Thus we see that, in general, the noise-amplitudes proportional to the combination of matrix-elements $`(M_{\kappa ,\nu \mu }^{(1)})^{}M_{\nu \mu ,\kappa }^{(2)}`$ drive the number-fluctuations in the condensate, while those proportional to $`(M_{\kappa ,\nu \mu }^{(1)})^{}+M_{\nu \mu ,\kappa }^{(2)}`$ are due to fluctuations of the occupation numbers in the excited states, couple in the Hamiltonian to the particle number in the condensate and therefore drive the phase-fluctuations in the condensate.
## VII Evaluation of the Transport Parameters for a Box-like Trap
For simplicity we consider now a trap consisting of a cube of length $`L`$ with cyclic boundary conditions. In the following equilibrium values of all parameters are implied, but we shall omit in this section the tilde and write $`\mu `$ for $`\mu `$ to simplify our notation. The normalized $`u`$ and $`v`$ coefficients are in this case
$`u_\nu `$ $`=`$ $`{\displaystyle \frac{E_\nu +p_\nu ^2/2m}{\sqrt{2E_\nu p_\nu ^2/m}}}{\displaystyle \frac{1}{\sqrt{V}}}e^{i\stackrel{}{p}_\nu \stackrel{}{x}/\mathrm{}}`$ (120)
$`v_\nu `$ $`=`$ $`{\displaystyle \frac{E_\nu p_\nu ^2/2m}{\sqrt{2E_\nu p_\nu ^2/m}}}{\displaystyle \frac{1}{\sqrt{V}}}e^{i\stackrel{}{p}_\nu \stackrel{}{x}/\mathrm{}}`$ (121)
with
$`E_\nu =\sqrt{\left({\displaystyle \frac{p_\nu ^2}{2m}}+\mu \right)^2\mu ^2}`$ (122)
and $`\stackrel{}{p}_\nu =\mathrm{}\frac{2\pi }{L}\stackrel{}{n}_\nu `$ with integer vector $`\stackrel{}{n}_\nu `$.
### A Transport coefficients
The squares of the relevant matrix elements for $`E_\kappa =E_\nu +E_\mu `$ become
$`\left|\left(M_{\kappa ,\nu \mu }^{(1)}\right)^{}M_{\nu \mu ,\kappa }^{(2)}\right|^2`$ $`=\left({\displaystyle \frac{U_0}{V}}\right)^2{\displaystyle \frac{G(E_\nu ,E_\mu ,\mu )}{E_\nu E_\mu E_\kappa }}\delta _{\stackrel{}{n}_\kappa ,\stackrel{}{n}_\nu +\stackrel{}{n}_\mu }`$ (123)
$`\left|\left(M_{\kappa ,\nu \mu }^{(1)}\right)^{}+M_{\nu \mu ,\kappa }^{(2)}\right|^2`$ $`=\left({\displaystyle \frac{U_0}{V}}\right)^2{\displaystyle \frac{G(E_\nu ,E_\mu ,\mu )}{E_\nu E_\mu E_\kappa }}\delta _{\stackrel{}{n}_\kappa ,\stackrel{}{n}_\nu +\stackrel{}{n}_\mu }`$ (124)
with
$`G(x,y,\alpha )=`$ $`\sqrt{\alpha ^2+(x+y)^2}\left(3\sqrt{\alpha ^2+x^2}\sqrt{\alpha ^2+y^2}xy+2\alpha (\sqrt{\alpha ^2+x^2}+\sqrt{\alpha ^2+y^2}+\alpha )\right)`$ (127)
$`+(x+y)\left(x(\sqrt{\alpha ^2+y^2}+\alpha )+y(\sqrt{\alpha ^2+x^2}+\alpha )\right)`$
$`+2\alpha (\sqrt{\alpha ^2+x^2}+\alpha )(\sqrt{\alpha ^2+y^2}+\alpha )+\alpha (x^2+y^2+\alpha ^2)`$
The transport coefficients are then expressed as the simple result
$`\mathrm{\Gamma }^\mathrm{"}=0`$ (128)
and
$`\mathrm{\Gamma }_0`$ $`={\displaystyle \frac{2}{\pi }}\left({\displaystyle \frac{a}{L}}\right)^2{\displaystyle \underset{\stackrel{}{n}_\nu ,\stackrel{}{n}_\mu }{}}{\displaystyle \frac{\delta (\epsilon _\nu +\epsilon _\mu \epsilon _{\nu +\mu })}{\epsilon _\nu \epsilon _\mu (\epsilon _\nu +\epsilon _\mu )}}G(\epsilon _\nu ,\epsilon _\mu ,\alpha )F(\epsilon _\nu ,\epsilon _\mu ,{\displaystyle \frac{k_BT}{\mathrm{}\omega _0}})`$ (129)
$`\mathrm{\Gamma }_0+\mathrm{\Gamma }^{}`$ $`={\displaystyle \frac{2}{\pi }}\left({\displaystyle \frac{a}{L}}\right)^2{\displaystyle \underset{\stackrel{}{n}_\nu ,\stackrel{}{n}_\mu }{}}{\displaystyle \frac{\delta (\epsilon _\nu +\epsilon _\mu \epsilon _{\nu +\mu })}{\epsilon _\nu \epsilon _\mu (\epsilon _\nu +\epsilon _\mu )}}G(\epsilon _\nu ,\epsilon _\mu ,\alpha )F(\epsilon _\nu ,\epsilon _\mu ,{\displaystyle \frac{k_BT}{\mathrm{}\omega _0}})`$ (131)
with
$`F(\epsilon _\nu ,\epsilon _\mu ,{\displaystyle \frac{k_BT}{\mathrm{}\omega _0}})={\displaystyle \frac{\mathrm{}\omega _0}{k_BT}}{\displaystyle \frac{e^{\beta \mathrm{}\omega _0(\epsilon _\nu +\epsilon _\mu )}}{\left(e^{\beta \mathrm{}\omega _0(\epsilon _\nu +\epsilon _\mu )}1\right)\left(e^{\beta \mathrm{}\omega _0\epsilon _\nu }1\right)\left(e^{\beta \mathrm{}\omega _0\epsilon _\mu }1\right)}}`$ (132)
Here we scaled the scattering length $`a=mU_0/4\pi \mathrm{}^2`$ with $`L`$ and the energies $`E_\nu ,E_\mu `$ and $`\mu ,k_BT`$ with the energy $`\mathrm{}\omega _0=(2\pi \mathrm{})^2/2mL^2`$, defining
$`\epsilon _\nu =\sqrt{(n_\nu ^2+\alpha )^2\alpha ^2}`$ (133)
$`\epsilon _{\nu +\mu }=\sqrt{\left((\stackrel{}{n}_\nu +\stackrel{}{n}_\mu )^2+\alpha \right)^2\alpha ^2}`$ (134)
with $`\alpha =\mu /\mathrm{}\omega _0.`$
The double sums over $`\stackrel{}{n}_\nu ,\stackrel{}{n}_\mu `$ start with $`\stackrel{}{n}`$-values with $`|\stackrel{}{n}|=1`$. They are approximated by integrals according to
$`{\displaystyle \underset{\stackrel{}{n}_\nu ,\stackrel{}{n}_\mu }{}}\delta \left(\epsilon _\nu +\epsilon _\mu \epsilon _{\nu +\mu }\right)(\mathrm{})=\pi ^2{\displaystyle \underset{\sqrt{1+2\alpha }}{\overset{\mathrm{}}{}}}{\displaystyle \underset{\sqrt{1+2\alpha }}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\epsilon _\nu \epsilon _\mu (\epsilon _\nu +\epsilon _\mu )(\mathrm{})d\epsilon _\nu d\epsilon _\mu }{\sqrt{(\epsilon _\nu ^2+\alpha ^2)(\epsilon _\mu ^2+\alpha ^2)((\epsilon _\nu +\epsilon _\mu )^2+\alpha ^2)}}}`$ (135)
Here (…) is any smooth function of $`\epsilon _\nu ,\epsilon _\mu `$. In all experiments so far $`\alpha 1`$ is satisfied, i.e. we can replace $`\sqrt{1+2\alpha }\sqrt{2\alpha }`$. This leaves us with the integral expressions
$`\mathrm{\Gamma }_0`$ $`=2\pi \left({\displaystyle \frac{a}{L}}\right)^2{\displaystyle \underset{\sqrt{2\alpha }}{\overset{\mathrm{}}{}}}{\displaystyle \underset{\sqrt{2\alpha }}{\overset{\mathrm{}}{}}}{\displaystyle \frac{G(\epsilon _\nu ,\epsilon _\mu ,\alpha )F(\epsilon _\nu ,\epsilon _\mu ,\frac{k_BT}{\mathrm{}\omega _0})d\epsilon _\nu d\epsilon _\mu }{\sqrt{(\epsilon _\nu ^2+\alpha ^2)(\epsilon _\mu ^2+\alpha ^2)((\epsilon _\nu +\epsilon _\mu )^2+\alpha ^2)}}}`$ (136)
$`\mathrm{\Gamma }_0+\mathrm{\Gamma }^{}`$ $`=2\pi \left({\displaystyle \frac{a}{L}}\right)^2{\displaystyle \underset{\sqrt{2\alpha }}{\overset{\mathrm{}}{}}}{\displaystyle \underset{\sqrt{2\alpha }}{\overset{\mathrm{}}{}}}{\displaystyle \frac{G(\epsilon _\nu ,\epsilon _\mu ,\alpha )F(\epsilon _\nu ,\epsilon _\mu ,\frac{k_BT}{\mathrm{}\omega _0})d\epsilon _\nu d\epsilon _\mu }{\sqrt{(\epsilon _\nu ^2+\alpha ^2)(\epsilon _\mu ^2+\alpha ^2)((\epsilon _\nu +\epsilon _\mu )^2+\alpha ^2)}}}`$ (138)
The expression for $`\mathrm{\Gamma }_0`$ and the asymptotic behavior for $`\epsilon _\nu ,\epsilon _\mu 0`$: $`G(\epsilon _\nu ,\epsilon _\mu ,|\alpha |)F(\epsilon _\nu ,\epsilon _\mu ,k_BT/\mathrm{}\omega _0)18(k_BT/\mathrm{}\omega _0)^21/[\epsilon _\nu \epsilon _\mu (\epsilon _\nu +\epsilon _\mu )]`$ make it amply clear that the states with the smallest energies $`\epsilon _\nu \alpha `$ make a large contribution to $`\mathrm{\Gamma }_0`$ (but not to $`\mathrm{\Gamma }_0+\mathrm{\Gamma }^{}`$). To calculate this contribution it is permitted <sup>\**</sup><sup>\**</sup>\**We actually need the additional condition $`\sqrt{\mu \mathrm{}\omega _0}k_BT`$ to use $`\beta \mathrm{}\omega _0\epsilon _\nu 1`$, $`\beta \mathrm{}\omega _0\epsilon _\mu 1`$ under the integral and to approximate
$`F(\epsilon _\nu ,\epsilon _\mu ,{\displaystyle \frac{k_BT}{\mathrm{}\omega _0}})=\left({\displaystyle \frac{k_BT}{\mathrm{}\omega _0}}\right)^2{\displaystyle \frac{1}{(\epsilon _\nu +\epsilon {}_{\mu }{}^{})\epsilon _\nu \epsilon _\mu }}`$ (139)
in addition to approximating $`\sqrt{\epsilon _\nu ^2+\alpha ^2}\sqrt{\epsilon _\mu ^2+\alpha ^2}\sqrt{(\epsilon _\nu +\epsilon _\mu )^2+\alpha ^2}\alpha ^3`$, and neglecting terms of order $`\epsilon _\nu ^2/\alpha ^2`$.
This contribution to $`\mathrm{\Gamma }_0`$, which we shall denote as $`\mathrm{\Gamma }_{00}`$, then reduces to
$`\mathrm{\Gamma }_{00}=36\pi \left({\displaystyle \frac{a}{L}}\right)^2\left({\displaystyle \frac{k_BT}{\mathrm{}\omega _0}}\right)^2{\displaystyle \underset{\sqrt{2\alpha }}{\overset{\mathrm{}}{}}}𝑑\epsilon _\nu {\displaystyle \underset{\sqrt{2\alpha }}{\overset{\mathrm{}}{}}}𝑑\epsilon _\mu {\displaystyle \frac{1}{\epsilon _\nu \epsilon {}_{\mu }{}^{}(\epsilon _\nu +\epsilon {}_{\mu }{}^{})}}`$ (140)
The double-integral can be evaluated as $`\sqrt{2/\alpha }\mathrm{ln}2`$ which yields the final result
$`\mathrm{\Gamma }_{00}=36\pi \sqrt{2}\mathrm{ln}2\left({\displaystyle \frac{a}{L}}\right)^2\left({\displaystyle \frac{\mathrm{}\omega _0}{\mu }}\right)^{1/2}\left({\displaystyle \frac{k_BT}{\mathrm{}\omega _0}}\right)^2=1.59..\left({\displaystyle \frac{T}{T_c}}\right)^2\left({\displaystyle \frac{N}{N_0}}\right)^{1/2}N^{1/3}\left({\displaystyle \frac{k_BT_ca^2m}{\mathrm{}^2}}\right)^{3/4}`$ (141)
The second form of this expression is obtained by eliminating $`V=L^3`$ in favor of the critical temperature of the equivalent ideal Bose-gas at the same density via
$$V=(N/\zeta (3/2))(2\pi \mathrm{}^2/k_BT_cm)^{3/2}.$$
(142)
There is yet another particularly important contribution to $`\mathrm{\Gamma }_0`$ due to the infrared-singularity of the integrand, we shall denote it as $`\mathrm{\Gamma }_{01}`$, where only one of the two excitation frequencies, say $`E_\nu `$, is small compared to $`\mu `$, while the other is larger, of order $`\mu `$ or even $`k_BT`$. With respect to $`ϵ_\nu `$ the low-energy asymptotics may then still be used. We then obtain
$$\mathrm{\Gamma }_{01}=\frac{4\pi }{\alpha }\left(\frac{a}{L}\right)^2\underset{\sqrt{2\alpha }}{\overset{q\alpha }{}}𝑑\epsilon _\nu \underset{\sqrt{2\alpha }}{\overset{\mathrm{}}{}}𝑑\epsilon _\mu \frac{G(0,\epsilon _\mu ,\alpha )G(0,0,\alpha )}{\epsilon _\nu (\epsilon _\mu ^2+\alpha ^2)}\frac{e^{\beta \mathrm{}\omega _0\epsilon _\mu }}{(e^{\beta \mathrm{}\omega _0\epsilon _\mu }1)^2},$$
(143)
where we have set an upper cut-off for the small energy at a fraction $`q`$ of the chemical potential. Most of the $`\epsilon _\mu `$-integral comes from a range around $`\alpha `$ and we may therefore replace the thermal function by its asymptotics for $`\beta \mathrm{}\omega _0\epsilon _\mu 0`$, which is $`(k_BT/\mathrm{}\omega _0\epsilon _\mu )^2`$. The integrals can then be performed with the result
$$\mathrm{\Gamma }_{01}=16\pi \left(\frac{a}{L}\right)^2\mathrm{ln}\left(\frac{q^2\mu }{2\mathrm{}\omega _0}\right)\frac{(k_BT)^2}{\mathrm{}\omega _0\mu }=1.22..\left(\frac{T}{T_c}\right)^2\frac{N}{N_0}\left(\frac{k_BT_ca^2m}{\mathrm{}^2}\right)^{1/2}\mathrm{ln}\left(\frac{q^2\mu }{2\mathrm{}\omega _0}\right).$$
(144)
We conclude that this contribution to $`\mathrm{\Gamma }_0`$ is smaller than the leading term $`\mathrm{\Gamma }_{00}`$ by the order of magnitude $`(\mathrm{}\omega _0/\mu )^{1/2}\mathrm{ln}(q^2\mu /2\mathrm{}\omega _0)`$.
Let us now turn to the expressions for $`\mathrm{\Gamma }_0\mathrm{\Gamma }_{00}\mathrm{\Gamma }_{01}`$ and $`\mathrm{\Gamma }_0+\mathrm{\Gamma }^{}`$. They can be simplified for $`\alpha 1`$ by rescaling $`\epsilon _\nu ,\epsilon _\mu `$ by $`\alpha `$ and taking the limit $`\sqrt{2/\alpha }0`$ for the lower boundaries of the rescaled integrals. We find in this way
$`\mathrm{\Gamma }_0\mathrm{\Gamma }_{00}\mathrm{\Gamma }_{01}=2\pi \left({\displaystyle \frac{a}{L}}\right)^2\left({\displaystyle \frac{\mu }{\mathrm{}\omega _0}}\right){\displaystyle \underset{0}{\overset{\mathrm{}}{}}}dx{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}dy[{\displaystyle \frac{G(x,y,1)F(x,y,\frac{k_BT}{\mu })}{\sqrt{(x^2+1)(y^2+1)((x+y)^2+1)}}}`$ (145)
$`\left({\displaystyle \frac{k_BT}{\mu }}\right)^2({\displaystyle \frac{18}{xy(x+y)}}`$ (146)
$`+8{\displaystyle \frac{x(x^2+1)(y^2+\sqrt{y^2+1}1)+y(y^2+1)(x^2+\sqrt{x^2+1}1)}{x^2y^2(x^2+1)(y^2+1)}}`$ $`\left)\right]`$ (147)
$`\mathrm{\Gamma }_0+\mathrm{\Gamma }^{}=2\pi \left({\displaystyle \frac{a}{L}}\right)^2\left({\displaystyle \frac{\mu }{\mathrm{}\omega _0}}\right){\displaystyle \underset{0}{\overset{\mathrm{}}{}}}𝑑x{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}𝑑y{\displaystyle \frac{G(x,y,1)F(x,y,\frac{k_BT}{\mu })}{\sqrt{(x^2+1)(y^2+1)((x+y)^2+1)}}}`$ (149)
where we used the scaling property
$`F(\epsilon _\nu ,\epsilon _\mu ,{\displaystyle \frac{k_BT}{\mathrm{}\omega _0}})={\displaystyle \frac{1}{\alpha }}F({\displaystyle \frac{\epsilon _\nu }{\alpha }},{\displaystyle \frac{\epsilon _\mu }{\alpha }},{\displaystyle \frac{k_BT}{\mathrm{}\omega _0\alpha }})`$ (150)
Eqs(146,149) are the complete result for the temperature dependent transport parameters of the condensate for box-like traps. In general the integrals have to be done numerically. We shall here consider some asymptotic results only.
First we consider these expressions asymptotically for $`k_BT/\mu 1`$. Then the integrals receive important contributions from x, y of the order of 1, i.e. from quasi-particle-energies of the order of the chemical potential, and also from values of x,y large compared to 1, i.e. quasi-particle energies of order $`k_BT`$. The contributions $`\mathrm{\Gamma }_0^{(>)}`$ and $`\mathrm{\Gamma }^{(>)}`$ from large energies can be determined in leading power in $`(k_BT/\mu )`$ by approximating
$`F(x,y,{\displaystyle \frac{k_BT}{\mu }}){\displaystyle \frac{\mu }{k_BT}}e^{\frac{\mu }{k_BT}(x+y)}`$ (151)
and rescaling $`x`$ and $`y`$ by $`k_BT/\mu `$. In the integrals<sup>††</sup><sup>††</sup>††In the high-energy regime the subtractions of the infrared-divergent terms in the integrand of (146) are of no importance. for $`\mathrm{\Gamma }_0`$ and $`\mathrm{\Gamma }_0+\mathrm{\Gamma }^{}`$ we can then let $`\mu /k_BT0`$ without any problem, using the property $`G(x,y,0)=4xy(x+y)`$.whereupon they are easily evaluated with the asymptotic results
$`\mathrm{\Gamma }_0^{(>)}+\mathrm{\Gamma }^{(>)}\mathrm{\Gamma }_0^{(>)}8\pi \left({\displaystyle \frac{a}{L}}\right)^2{\displaystyle \frac{k_BT}{\mathrm{}\omega _0}}=1.27..{\displaystyle \frac{T}{T_c}}{\displaystyle \frac{k_BT_ca^2m}{\mathrm{}^2}}`$ (152)
We see that the result for $`\mathrm{\Gamma }^{(>)}`$ vanishes to this order.
For the contributions $`\mathrm{\Gamma }_0^{(\mu )},\mathrm{\Gamma }^{(\mu )}`$ from quasi-particles with energies around $`\mu `$ we can approximate $`F(x,y,k_BT/\mu )`$ according to eq.(139) and find
$`\mathrm{\Gamma }_0^{(\mu )}`$ $`B_0^{(\mu )}\left({\displaystyle \frac{a}{L}}\right)^2{\displaystyle \frac{(k_BT)^2}{\mathrm{}\omega _0\mu }}=0.0243..B_0^{(\mu )}\left({\displaystyle \frac{T}{T_c}}\right)^2{\displaystyle \frac{N}{N_0}}\left({\displaystyle \frac{k_BT_ca^2m}{\mathrm{}^2}}\right)^{1/2}`$ (153)
$`\mathrm{\Gamma }_0^{(\mu )}+\mathrm{\Gamma }^{(\mu )}`$ $`B^{(\mu )}\left({\displaystyle \frac{a}{L}}\right)^2{\displaystyle \frac{(k_BT)^2}{\mathrm{}\omega _0\mu }}=0.0243..B^{(\mu )}\left({\displaystyle \frac{T}{T_c}}\right)^2{\displaystyle \frac{N}{N_0}}\left({\displaystyle \frac{k_BT_ca^2m}{\mathrm{}^2}}\right)^{1/2}`$ (154)
with the numbers $`B_0^{(\mu )},B^{(\mu )}`$ defined by the integrals
$`B_0^{(\mu )}=2\pi {\displaystyle \underset{0}{\overset{\mathrm{}}{}}}𝑑x{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}𝑑y`$ $`[{\displaystyle \frac{G(x,y,1)}{\sqrt{(x^2+1)(y^2+1)((x+y)^2+1)}xy(x+y)}}`$ (157)
$`{\displaystyle \frac{18}{xy(x+y)}}`$
$`8{\displaystyle \frac{x(x^2+1)(y^2+\sqrt{y^2+1}1)+y(y^2+1)(x^2+\sqrt{x^2+1}1)}{x^2y^2(x^2+1)(y^2+1)}}]`$
$`B^{(\mu )}=2\pi {\displaystyle \underset{0}{\overset{\mathrm{}}{}}}𝑑x{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}𝑑y`$ $`{\displaystyle \frac{G(x,y,1)}{\sqrt{(x^2+1)(y^2+1)((x+y)^2+1)}xy(x+y)}}`$ (158)
We can conclude that the contribution from quasi-particles at energies of order $`\mu `$ is larger (for $`\mathrm{\Gamma }_0`$ by an order of magnitude $`k_BT/\mu `$) than the contribution from energies of order $`k_BT`$, but $`\mathrm{\Gamma }_0^{(\mu )}`$ is, in large condensates, still subdominant to $`\mathrm{\Gamma }_{00}`$ by the order of magnitude $`\sqrt{\mathrm{}\omega _0/\mu }`$.
Now let us consider also the low-temperature limit, namely $`k_BT/\mu 1`$ or, equivalently, $`T/T_c\left(k_BT_ca^2m/\mathrm{}^2\right)^{1/2}`$. In this region it is not neccessary to distinguish $`N`$ and $`N_0`$. The integrals now receive their contributions for $`x,y`$ both small compared to 1, but we can still use the approximation (151). For small $`x,y`$ we can expand
$`G(x,y,1)18+3(x^2+y^2+xy)G(x,y,1){\displaystyle \frac{9}{32}}x^2y^2(x+y)^2`$ (159)
To obtain the leading term it is enough to keep only the smallest powers of $`x,y`$ in the integrands. The integrals are easily evaluated with the asymptotic low-temperature results
$`\mathrm{\Gamma }_0\mathrm{\Gamma }_{00}=36\pi \left({\displaystyle \frac{a}{L}}\right)^2{\displaystyle \frac{k_BT}{\mathrm{}\omega _0}}=5.72..{\displaystyle \frac{T}{T_c}}{\displaystyle \frac{k_BT_ca^2m}{\mathrm{}^2}}`$ (160)
$`\mathrm{\Gamma }_0+\mathrm{\Gamma }^{}={\displaystyle \frac{189}{2}}\pi \left({\displaystyle \frac{a}{L}}\right)^2{\displaystyle \frac{\mu }{\mathrm{}\omega _0}}\left({\displaystyle \frac{k_BT}{\mu }}\right)^7=0.366..\left({\displaystyle \frac{T}{T_c}}\right)^7\left({\displaystyle \frac{k_BT_ca^2m}{\mathrm{}^2}}\right)^2`$ (161)
As long as the temperature is high enough to satisfy $`k_BT\sqrt{\mu \mathrm{}\omega _0}`$ the part $`\mathrm{\Gamma }_{00}`$ still dominates the value of $`\mathrm{\Gamma }_0`$.
### B Particle-number fluctuations
We follow the procedure of Giorgini et al. and deduce the particle-number fluctuations in the condensate from the number fluctuation in the thermal cloud. This leads to eq.(34), which we evaluate using the expressions for $`u_\nu ,v_\nu `$ and $`E_\nu `$.
We obtain
$`\mathrm{\Delta }N_0^2={\displaystyle \underset{\nu }{}}{\displaystyle \frac{2\overline{n}_\nu (\overline{n}_\nu +1)(E_\nu ^2+2\mu ^2)+\mu ^2}{2E_\nu ^2}}`$ (162)
Approximated by an integral this becomes
$`\mathrm{\Delta }N_0^2=\pi {\displaystyle _{\sqrt{1+2\alpha }}^{\mathrm{}}}{\displaystyle \frac{d\epsilon }{\epsilon }}\sqrt{{\displaystyle \frac{\sqrt{\epsilon ^2+\alpha ^2}\alpha }{\epsilon ^2+\alpha ^2}}}\left(\alpha ^2+{\displaystyle \frac{\epsilon ^2+2\alpha ^2}{2(\mathrm{sin}h\frac{\beta \mathrm{}\omega _0\epsilon }{2})^2}}\right)`$ (163)
The dominant contribution comes from the lower boundary of the integration which contributes, for $`\alpha 1`$
$`\mathrm{\Delta }N_0^22\pi \left({\displaystyle \frac{k_BT}{\mathrm{}\omega _0}}\right)^2=A^{}\left({\displaystyle \frac{mk_BT}{\mathrm{}^2}}\right)^2V^{4/3}`$ (164)
with
$`A^{}={\displaystyle \frac{1}{2\pi ^3}}=0.0161..`$ (165)
More precisely the dominant contribution to $`\mathrm{\Delta }N_0^2`$ is given by the discrete sum
$`\mathrm{\Delta }N_0^2=2\mu ^2(k_BT)^2{\displaystyle \underset{\nu }{}}{\displaystyle \frac{1}{E_\nu ^4}}`$ (166)
which gives the same expression as (164) but with the prefactor<sup>‡‡</sup><sup>‡‡</sup>‡‡The formula (167) differs from the one given in by a factor $`2^4`$ whereas the numerical result differs by yet another factor; the formula (169) differs from the one in by a factor 2.
$`A={\displaystyle \frac{2}{(2\pi )^4}}{\displaystyle \underset{\stackrel{}{n}_\nu 0}{}}{\displaystyle \frac{1}{n_\nu ^4}}=0.021..`$ (167)
If we eliminate the volume in favour of the critical temperature of the ideal Bose-gas of the same density via eq.(142) we get
$`\mathrm{\Delta }N_0^2=A{\displaystyle \frac{(2\pi )^2}{\zeta (3/2)^{4/3}}}\left({\displaystyle \frac{T}{T_c}}\right)^2N^{4/3}.`$ (168)
At temperature $`T=0`$ a similar evaluation of (34) gives
$`\mathrm{\Delta }N_0^2|_{T=0}=2\sqrt{\pi }(aN)^{3/2}V^{1/2}.`$ (169)
### C Particle-number relaxation rate
We are now in a position to evaluate the rate $`\tau _c^1`$ from eq.(52) using the results for $`\mathrm{\Delta }N_0^2`$ (numbers are calculated with the prefactor $`A^{}`$) and $`\mathrm{\Gamma }_0\mathrm{\Gamma }_{00}`$. We obtain
$`\gamma _c={\displaystyle \frac{1}{\tau _c}}=18.0..{\displaystyle \frac{T}{T_c}}\sqrt{{\displaystyle \frac{N_0}{N}}}\left({\displaystyle \frac{k_BT_ca^2m}{\mathrm{}^2}}\right)^{3/4}{\displaystyle \frac{k_BT_c}{\mathrm{}}}`$ (170)
This result applies also in the low-temperature region, because it makes use only of the results for $`\mathrm{\Delta }N_0^2`$ and $`\mathrm{\Gamma }_0`$ which also hold in that region.
To get an idea of order of magnitudes we compare this and the following results with the damping rate $`\gamma _0`$ of the lowest lying modes, which is given by
$`\gamma _0={\displaystyle \frac{3\pi ^2}{4}}\left({\displaystyle \frac{a}{L}}\right){\displaystyle \frac{k_BT}{\mathrm{}}}=4.06..{\displaystyle \frac{T}{T_c}}N^{1/3}\left({\displaystyle \frac{k_BT_ca^2m}{\mathrm{}^2}}\right)^{1/2}{\displaystyle \frac{k_BT_c}{\mathrm{}}}`$ (171)
We see that the relaxation rate $`\gamma _c`$ is of the order
$`\gamma _cN^{1/3}\left({\displaystyle \frac{k_BT_ca^2m}{\mathrm{}^2}}\right)^{1/4}\gamma _0`$ (172)
The proportionality factor is of the order of $`\sqrt{\mu /\mathrm{}\omega _0}`$ and is large in large and strongly interacting condensates. I.e. then the relaxation of the condensate to its equilibrium is faster than the relaxation of the low-lying collective modes, but slower than the frequency of the lowest lying modes, which is $`\sqrt{2\omega _0\mu /\mathrm{}}`$.
### D Phase collapse rate
The phase collapse rate is given by eq.(58) and requires only the result for $`\mathrm{\Delta }N_0^2`$ and $`\mathrm{\Gamma }^{\prime \prime }=0`$. At zero temperature it reduces to
$`\gamma _{collapse}|_{T=0}={\displaystyle \frac{1}{\mathrm{}}}{\displaystyle \frac{\mu }{N_0}}\sqrt{\mathrm{\Delta }N_0^2|_{T=0}}`$ (173)
from which we get
$`\gamma _{collapse}|_{T=0}={\displaystyle \frac{23.6..}{\sqrt{V}}}(an_0)^{3/4}{\displaystyle \frac{\mathrm{}a}{m}}=2.50..{\displaystyle \frac{k_BT_c}{\mathrm{}}}\left({\displaystyle \frac{k_BT_ca^2m}{\mathrm{}^2}}\right)^{7/8}{\displaystyle \frac{1}{\sqrt{N}}}`$ (174)
For finite temperature we obtain
$`\gamma _{collapse}=0.876..{\displaystyle \frac{k_BT}{\mathrm{}}}N^{1/3}\left({\displaystyle \frac{k_BT_ca^2m}{\mathrm{}^2}}\right)^{1/2}`$ (175)
By comparison with (171) we see that $`\gamma _{collapse}`$ is of the order of $`\gamma _0`$ and is therefore in large condensates smaller than $`\gamma _c`$.
In summary, the phase-collapse is not effective in large condensates because it occurs with a rate $`\gamma _{collapse}<\gamma _c`$ and is at the same time restricted to a time interval $`\mathrm{\Delta }t<\frac{1}{\gamma _c}`$ since for larger times phase-diffusion takes over.
### E Phase-diffusion rate
The phase-diffusion coefficient is a somewhat complicated quantity because it receives contributions from several processes, which are physically distinct. We consider the different contributions separately and also distinguish the two temperature regimes of high temperature $`k_BT>\mu `$, for which we give the result first, and $`k_BT<\mu `$ (low temperature).
#### 1 Low frequency condensate number fluctuations
From eq.(54) we infer with $`\mathrm{\Gamma }_0=\mathrm{\Gamma }_{00}`$
$`D_\varphi ^{(\alpha )}={\displaystyle \frac{1}{\mathrm{}k_BTN_0\mathrm{\Gamma }_{00}}}\left(\mathrm{\Delta }N_0^2{\displaystyle \frac{\mu }{N_0}}\right)^2`$ (176)
which is evaluated as
$`D_\varphi ^{(\alpha )}=0.0853..{\displaystyle \frac{T}{T_c}}\left({\displaystyle \frac{k_BT_cma^2}{\mathrm{}^2}}\right)^{1/4}{\displaystyle \frac{1}{N_0^{1/2}N^{1/6}}}{\displaystyle \frac{k_BT_c}{\mathrm{}}}.`$ (177)
The same result holds in the low temperature regime $`k_BT<\mu `$. In comparison with $`\gamma _0`$ (171) it is of the order
$`D_\varphi ^{(\alpha )}N^{1/3}\left({\displaystyle \frac{k_BT_cma^2}{\mathrm{}^2}}\right)^{1/4}\gamma _0\sqrt{{\displaystyle \frac{\mathrm{}\omega _0}{\mu }}}\gamma _0`$ (178)
and is much smaller in large and strongly interacting condensates. Still this contribution to the phase-diffusion rate is always the dominant one at low temperatures and may dominate even at higher temperatures (see below).
#### 2 Condensate number fluctuations due to quasi-particles around energies $`\mu `$
Splitting $`\mathrm{\Gamma }_0=\mathrm{\Gamma }_{00}+(\mathrm{\Gamma }_0\mathrm{\Gamma }_{00})`$ and expanding to first order
$`{\displaystyle \frac{1}{\mathrm{\Gamma }_0}}={\displaystyle \frac{1}{\mathrm{\Gamma }_{00}}}{\displaystyle \frac{\mathrm{\Gamma }_0\mathrm{\Gamma }_{00}}{\mathrm{\Gamma }_{00}^2}}`$ (179)
we estimate as contribution $`D_\varphi ^{(\beta )}`$ from the higher frequency condensate number fluctuations described by $`\mathrm{\Gamma }_0\mathrm{\Gamma }_{00}`$ as given by (152)
$`D_\varphi ^{(\beta )}{\displaystyle \frac{T}{T_c}}{\displaystyle \frac{1}{N_0}}{\displaystyle \frac{k_BT_c}{\mathrm{}}}`$ (180)
which is in absolute value smaller than the contribution $`D_\varphi ^{(\alpha )}`$ from low-energy excitations by the order of magnitude factor $`\sqrt{\mathrm{}\omega _0/\mu }`$. This contribution is therefore negligible in very large condensates. In not so large condensates the complete integral in the result for $`\mathrm{\Gamma }_0`$ needs to be evaluated.
In the low-temperature regime $`k_BT<\mu `$ we get instead
$`D_\varphi ^{(\beta )}=0.306..{\displaystyle \frac{1}{N}}\left({\displaystyle \frac{k_BT_ca^2m}{\mathrm{}}}\right)^{1/2}{\displaystyle \frac{k_BT_c}{\mathrm{}}}`$ (181)
This is much smaller than $`D_\varphi ^{(\alpha )}`$ by the order of magnitude factor $`(\mu /k_BT_c)^{1/2}/N^{1/3}`$.
#### 3 Fluctuations in the thermal cloud at energies of order $`\mu `$
By eq.(54) this contribution is given by
$`D_\varphi ^{(\gamma )}={\displaystyle \frac{k_BT}{\mathrm{}N_0}}(\mathrm{\Gamma }_0+\mathrm{\Gamma }^{})`$ (182)
which is evaluated to
$`D_\varphi ^{(\gamma )}=0.0243..B^{(\mu )}\left({\displaystyle \frac{T}{T_c}}\right)^3{\displaystyle \frac{N}{N_0^2}}\left({\displaystyle \frac{k_BT_ca^2m}{\mathrm{}^2}}\right)^{\frac{1}{2}}{\displaystyle \frac{k_BT_c}{\mathrm{}}}.`$ (183)
This contribution differs from $`D_\varphi ^{(\alpha )}`$ by the order of magnitude factor $`(T/T_c)^2\sqrt{\mu \mathrm{}\omega _0}/k_BT_c`$ and is therefore much smaller.
For temperatures $`k_BT<\mu `$ we find instead
$`D_\varphi ^{(\gamma )}=0.366..\left({\displaystyle \frac{T}{T_c}}\right)^8{\displaystyle \frac{1}{N_0}}\left({\displaystyle \frac{k_BT_ca^2m}{\mathrm{}^2}}\right)^2{\displaystyle \frac{k_BT_c}{\mathrm{}}}`$ (184)
which is again negligibly small compared to $`D_\varphi ^{(\alpha )}`$.
In summary, the phase-diffusion is caused dominantly by the low-frequency particle number fluctuations in the condensate and the phase-diffusion constant is given by eq.(177). It is proportional to temperature, and scales proportional to $`N^{2/3}`$ for fixed $`T_c`$ or proportional to $`N^{1/2}`$ for fixed volume of the trap.
## VIII Evaluation of the transport parameters for an isotropic harmonic trap
In this section we consider the more realistic case of condensates in a parabolic trapping potential $`m\omega _0^2x^2/2`$, which we assume to be isotropic for simplicity. In order to analyse the noise $`\mathrm{}(\xi (t)`$ driving the fluctuations of $`|\alpha _0|^2`$ we must consider in detail the relevant linear combination of matrix elements
$$(M_{\kappa ,\nu \mu }^{(1)})^{}M_{\nu \mu ,\kappa }^{(2)}=2U_0d^3x\psi _0\left\{(\stackrel{~}{u}_\kappa \stackrel{~}{v}_\kappa )(\stackrel{~}{u}_\mu ^{}\stackrel{~}{v}_\nu ^{}+\stackrel{~}{v}_\mu ^{}\stackrel{~}{u}_\nu ^{})\stackrel{~}{u}_\kappa \stackrel{~}{u}_\mu ^{}\stackrel{~}{u}_\nu ^{}+\stackrel{~}{v}_\kappa \stackrel{~}{v}_\mu ^{}\stackrel{~}{v}_\nu ^{}\right\}$$
(185)
In the following we shall make use of the local density and Thomas-Fermi approximation, restricting ourselves to large condensates. For the high-lying states we can then use the local energies in Thomas-Fermi approximation
$$E(p,𝒙)=\sqrt{(\frac{p^2}{2m}+|U_0n_0(𝒙)|)^2U_0^2n_0^2(𝒙)\mathrm{\Theta }(\mu V(𝒙))}$$
(186)
with the condensate density
$$n_0(𝒙)=N_0|\stackrel{~}{\psi }_0(𝒙)|^2=(\mu /U_0)(1(x/r_{TF})^2)$$
(187)
and the Thomas-Fermi radius
$$r_{TF}=\sqrt{2\mu /m\omega _0^2}=(\frac{15U_0N_0}{8\pi \mu })^{1/3}.$$
(188)
The high-lying quasi-particle modes can be represented similarly to the spatially homogeneous case as
$`u_\kappa (𝒙)=`$ $`{\displaystyle \frac{E_\kappa +p_\kappa ^2/2m}{\sqrt{2E_\kappa p_\kappa ^2/m}}}e^{i𝒑_\kappa 𝒙/\mathrm{}}`$ (189)
$`v_\kappa (𝒙)=`$ $`{\displaystyle \frac{E_\kappa p_\kappa ^2/2m}{\sqrt{2E_\kappa p_\kappa ^2/m}}}e^{i𝒑_\kappa 𝒙/\mathrm{}}`$ (191)
The low-lying collective modes can be represented as
$`u_\nu (𝒙)=`$ $`\left(\sqrt{{\displaystyle \frac{U_0n_0(𝒙)}{2\mathrm{}\stackrel{~}{\omega }_\nu }}}+{\displaystyle \frac{1}{2}}\sqrt{{\displaystyle \frac{\mathrm{}\stackrel{~}{\omega }_\nu }{2U_0n_0(𝒙)}}}\right)\chi _\nu (𝒙)`$ (192)
$`v_\nu (𝒙)=`$ $`\left(\sqrt{{\displaystyle \frac{U_0n_0(𝒙)}{2\mathrm{}\stackrel{~}{\omega }_\nu }}}+{\displaystyle \frac{1}{2}}\sqrt{{\displaystyle \frac{\mathrm{}\stackrel{~}{\omega }_\nu }{2U_0n_0(𝒙)}}}\right)\chi _\nu (𝒙)`$ (194)
with
$$d^3x|\chi _\nu (𝒙)|^2=1$$
(195)
The mode-functions $`\chi _\nu (𝒙)`$ are known in the hydrodynamic (long-wavelength) and Thomas-Fermi approximation by analytic solutions of the Bogoliubov-equations. In spatially isotropic parabolic traps they have the form
$$\chi _\nu (𝒙)=\frac{1}{r_{TF}^{3/2}}P_\mathrm{}_\nu ^{(2n_\nu )}(x/r_{TF})(x/r_{TF})^{\mathrm{}}Y_\mathrm{}m(\theta ,\phi )\mathrm{\Theta }(1x/r_{TF})$$
(196)
The polynomials $`P_{\mathrm{}}^{(2n)}(x)`$ of degree $`2n`$ are the normalized solutions of the radial part of the Bogoliubov-Fetter equations in the Thomas-Fermi and long-wavelength limit given by
$$P_{\mathrm{}}^{(2n)}(x)=\frac{\sqrt{4n+2\mathrm{}+3}}{n!}x^{2\mathrm{}1}\frac{d^n}{d(x^2)^n}\left[x^{2n+2\mathrm{}+1}(1x^2)^n\right]$$
(197)
with the normalization
$$_0^1𝑑xx^{2\mathrm{}+2}\left[P_{\mathrm{}}^{(2n)}(x)\right]^2=1$$
(198)
In the phonon part of the excitation spectrum we have $`u_\lambda v_\lambda \stackrel{~}{\omega }_\lambda ^{1/2}`$. Furthermore, in that low- energy region the statistical factor in eqs.(103), (104), (107) is well approximated by $`\overline{n}_\nu \overline{n}_\mu \overline{n}_\kappa (k_BT)^3/\mathrm{}^3\stackrel{~}{\omega }_\kappa \stackrel{~}{\omega }_\nu \stackrel{~}{\omega }_\mu `$. Just as in the case of box-like traps the frequency factors in the denominator, together with similar further factors in the denominator coming from the matrix elements, make the phonon contribution to the sums in (103) the dominant one, at least in large condensates, and we shall therefore concentrate on this contribution in the following. This frequency range has a natural upper cut-off at $`\mu /\mathrm{}`$, where the collective phonons go over smoothly into particle-like excitations.
For $`E_\kappa ,E_\nu ,E_\mu \mu `$ the matrix elements $`(M_{\kappa ,\nu \mu }^{(1)})^{},M_{\nu \mu ,\kappa }^{(2)}`$ are given by the integral
$`(M_{\kappa ,\nu \mu }^{(1)})^{}`$ $``$ $`M_{\nu \mu ,\kappa }^{(2)}`$ (199)
$``$ $`\sqrt{{\displaystyle \frac{15}{8\pi }}}{\displaystyle \frac{3U_0\mu ^{3/2}\delta _{m_\kappa ,m_\nu +m_\mu }}{r_{TF}^3\sqrt{2E_\nu E_\mu (E_\nu +E_\mu )}}}J(n_\kappa ,n_\nu ,n_\mu ;\mathrm{}_\kappa ,\mathrm{}_\nu ,\mathrm{}_\mu )C(\mathrm{}_\kappa |\mathrm{}_\mu m_\mu ,\mathrm{}_\nu m_\nu )`$ (200)
where $`J`$ denotes the integral
$`J(n_\kappa ,n_\nu ,n_\mu ;\mathrm{}_\kappa ,\mathrm{}_\nu ,\mathrm{}_\mu )={\displaystyle _0^1}𝑑xx^2(1x^2)^2x^{\mathrm{}_\kappa +\mathrm{}_\nu +\mathrm{}_\mu }P_\mathrm{}_\kappa ^{(2n_\kappa )}(x)P_\mathrm{}_\nu ^{(2n_\nu )}(x)P_\mathrm{}_\mu ^{(2n_\mu )}(x)`$
and the Clebsch-Gordon coefficients $`C(\mathrm{}_\kappa |\mathrm{}_\mu m_\mu ,\mathrm{}_\nu m_\nu )`$ are given by the angle-integral
$`C(\mathrm{}_\kappa |\mathrm{}_\mu m_\mu ,\mathrm{}_\nu m_\nu )={\displaystyle 𝑑\mathrm{\Omega }Y_{\mathrm{}_\kappa ,m_\nu +m_\mu }^{}(\theta ,\phi )Y_{\mathrm{}_\mu ,m_\mu }(\theta ,\phi )Y_{\mathrm{}_\nu ,m_\nu }(\theta ,\phi )}`$
if $`|\mathrm{}_\mu \mathrm{}_\nu |\mathrm{}_\kappa \mathrm{}_\mu +\mathrm{}_\nu `$, otherwise they vanish. Later-on we shall have to calculate e.g. $`_{m_\nu ,m_\mu }|(M_{\kappa ,\nu \mu }^{(1)})^{}M_{\nu \mu ,\kappa }^{(2)}|^2`$ where we can make use of the sum-rule, for $`|\mathrm{}_\mu \mathrm{}_\nu |\mathrm{}_\kappa \mathrm{}_\mu +\mathrm{}_\nu `$
$`{\displaystyle \underset{m_\nu ,m_\mu }{}}|C(\mathrm{}_\kappa |\mathrm{}_\mu m_\mu ,\mathrm{}_\nu m_\nu )|^2=1`$
so that the Clebsch-Gordon coefficients need actually not be used explicitly.
In order to have well-defined expressions for the rate-coefficients we again need to smoothen the delta-function expressing energy-conservation, which is done physically by experimental imperfections or limitations in resolution. Here this can be done by replacing the discrete sum over the ’quantum-number’ $`\mathrm{}_\kappa `$ by an integral
$`{\displaystyle \underset{\mathrm{}_\kappa }{}}\delta (E_\kappa E_\nu E_\mu )(\mathrm{})`$ $``$ $`{\displaystyle 𝑑E_\kappa \frac{1}{dE_\kappa /d\mathrm{}_\kappa }\delta (E_\kappa E_\nu E_\mu )(\mathrm{})}`$ (201)
$`=`$ $`{\displaystyle 𝑑E_\kappa \frac{E_\nu +E\mu }{(\mathrm{}\omega _0)^2(n_\kappa +1/2)}\delta (E_\kappa E_\nu E_\mu )(\mathrm{})}`$ (202)
where we used the expression for the excitation-energies
$`E_\kappa `$ $`=`$ $`\mathrm{}\omega _0e_\kappa `$ (203)
$`e_\kappa `$ $`=`$ $`\sqrt{2n_\kappa ^2+2n_\kappa \mathrm{}_\kappa +3n_\kappa +\mathrm{}_\kappa }.`$ (204)
We introduced the dimensionless eigenvalues $`e_{\kappa ,\nu ,\mu }`$ which will appear in the ensuing expressions from now on. The integration over $`E_\kappa `$ with the delta-function then picks out the energy-value $`E_\kappa =E_\nu +E_\mu `$ so that $`\mathrm{}_\kappa `$ becomes a function $`\mathrm{}_\kappa ^{(0)}`$ of the other ’quantum-numbers’
$`\mathrm{}_\kappa ^{(0)}={\displaystyle \frac{(e_\nu +e_\mu )^22n_\kappa ^23n_\kappa }{2n_\kappa +1}}`$
The inequalities
$`|\mathrm{}_\nu \mathrm{}_\mu |\mathrm{}_\kappa ^{(0)}\mathrm{}_\nu +\mathrm{}_\mu `$
then imply that $`n_\kappa `$ must lie in the interval
$`n_\kappa n_\kappa n_{\kappa +}`$
with
$`n_{\kappa \pm }={\displaystyle \frac{1}{2}}\left(\sqrt{|\mathrm{}_\nu \pm \mathrm{}_\mu |^2+|\mathrm{}_\nu \pm \mathrm{}_\mu |+9/4+2(e_\nu +e_\mu )^2}|\mathrm{}_\nu \pm \mathrm{}_\mu |3/2\right)`$
Using all this we obtain from eq.(103)
$`\mathrm{\Gamma }_0=B_{00}\left({\displaystyle \frac{a}{d_0}}{\displaystyle \frac{k_BT}{\mathrm{}\omega _0}}\right)^2=B_{00}{\displaystyle \frac{T^2}{T_c^2}}{\displaystyle \frac{k_BT_ca^2m}{\mathrm{}^2}}\left({\displaystyle \frac{N}{\zeta (3)}}\right)^{1/3}`$ (205)
where $`\mathrm{}\omega _0`$ is now eliminated in favour of $`k_BT_c`$ via the relation $`\mathrm{}\omega _0=k_BT_c(\zeta (3)/N)^{1/3}`$, and where the temperature- and particle-number-independent positive real number $`B_{00}`$ is defined by the multiple sums
$$B_{00}=\frac{135\pi ^2}{2}\underset{n_\nu }{}\underset{n_\mu }{}\underset{\mathrm{}_\nu }{}\underset{\mathrm{}_\mu }{}\underset{n_\kappa =n_\kappa }{\overset{n_{\kappa +}}{}}\frac{(2\mathrm{}_\kappa ^{(0)}+1)J^2(n_\kappa ,n_\nu ,n_\mu ;\mathrm{}_\kappa ^{(0)},\mathrm{}_\nu ,\mathrm{}_\mu )}{e_\nu ^2e_\mu ^2(e_\nu +e_\mu )(2n_\kappa +1)}$$
(206)
The result for $`\mathrm{\Gamma }_0`$ agrees, except for the numerical prefactor, with the result of which was evaluated there using the local-density approximation and imposing a lower cut-off for the excitation-frequencies at the geometrical mean trap-frequency $`\overline{\omega }`$. It can also be compared with the corresponding result (141) for the box-like trap, which shows the same dependence on temperature and particle-number (if we stipulate $`N_0N`$), but the comparison of the prefactor is problematic because the condensate in the parabolic trap has two length-scales $`d_0`$ and $`r_{TF}`$, whereas in the box-like trap only the length-scale L is relevant.
The property (199) of the matrix-elements implies that the low-lying excitations don’t contribute to $`\mathrm{\Gamma }_0+\mathrm{\Gamma }^{}`$. The reality of the matrix-elements furthermore implies that $`\mathrm{\Gamma }^{\prime \prime }`$ vanishes. These remarkably simple results mean that the noise-source $`\xi (t)`$ introduced in eq.(41) is purely imaginary, corresponding to total squeezing in the direction of the phase $`\varphi `$. In other words, the coupling of the condensate to the collective excitations introduces a direct Langevin noise-source only for the number-fluctuations $`\delta N_0`$, not the phase-variable $`\varphi `$.<sup>\**</sup><sup>\**</sup>\**The latter is of course affected by the noise-source indirectly, because the fluctuations of $`\delta N_0`$ driven by the latter cause fluctuations in the chemical potential.
The fact that $`\mathrm{\Gamma }_0+\mathrm{\Gamma }^{}=0`$ for the contribution from the low-lying states implies that the contributions from the higher lying states must also be considered in order to evaluate the small (compared to $`\mathrm{\Gamma }_0`$) but finite value of this quantity. For this purpose we need to consider the matrix-element $`(M_{\kappa ,\nu \mu }^{(1)})^{}+M_{\nu \mu ,\kappa }^{(2)}`$. It differs from the matrix-elements we have considered so far by the replacement $`\stackrel{~}{\psi }_0(x)\stackrel{~}{\psi }_0(x)+2N_0\stackrel{~}{\psi }_0(x)/N_0`$ in the matrix-element. In Thomas-Fermi approximation this is tantamount to the replacement
$$\stackrel{~}{\psi }_0(x)(2/5)\stackrel{~}{\psi }_0(x)/(1x^2/r_{TF}^2).$$
(207)
Physically this implies a reduced coupling of the thermal fluctuations with the center of the condensate and a strongly enhanced coupling at its boundary, as one would expect for fluctuations located in the thermal cloud. A mathematical consequence is the fact that the integrals defining these matrix-elements diverge at the boundary in the Thomas-Fermi approximation, meaning that we encounter here the limitations of that approximation. Instead of a full-fledged extension of the theory beyond the Thomas-Fermi approximation it will be sufficient for our purposes here to cure its deficiencies by substituting as a cut-off the finite thickness of the boundary-layer given by
$`d={\displaystyle \frac{1}{2}}r_{TF}\left({\displaystyle \frac{\mathrm{}\omega _0}{\mu }}\right)^{2/3}`$
The matrix-element itself is then evaluated in the local-density approximation , where we can make use to good purpose of the analysis already performed in the predeeding section. The finte volume $`V=L^3`$ (and the associated $`\mathrm{}\omega _0=(2\pi \mathrm{})^2/2mL^2`$ which is not to be confused with the trap frequency called $`\omega _0`$ in the present section) is then an arbitrary local subvolume of the condensate, introduced merely as a technical device like a quantization-volume. It must be sufficiently small so that the condensate within it can be treated as homogeneous, and sufficiently large that we can replace sums over local momenta by integrals. At the end we have to check for consistency whether the result is indeed independent of the choice of this volume. The result obtained in this way is the local average of the result (154) for the homogneous case, which now becomes space-dependent, because we have to substitute a space-dependent chemical potential $`\mu \mu (1x^2/r_{TF}^2)`$. This local result can be written as
$`\mathrm{\Gamma }_0(x)+\mathrm{\Gamma }^{}(x)={\displaystyle \frac{B^{(\mu )}}{2\pi ^2}}{\displaystyle \frac{(k_BT)^2a^2m}{\mathrm{}^2\mu (1x^2/r_{TF}^2)}}`$
and is indeed independent of the choice of $`V`$. The local average has to be performed with the weight $`(\stackrel{~}{\psi }_0(x)+2N_0\stackrel{~}{\psi }_0(x)/N_0)^2`$ determined from (207). Doing the average and regulating the divergency of the integral at the boundary of the condensate by the physical cut-off we obtain
$$\mathrm{\Gamma }_0+\mathrm{\Gamma }^{}=\frac{3}{10}\frac{2^{1/3}}{15^{2/15}\pi ^2}B^{(\mu )}\left(\frac{k_BT}{\mathrm{}\omega _0}\right)^2N_0^{2/15}\left(\frac{a}{d_0}\right)^{28/15}=0.024..B^{(\mu )}\left(\frac{T}{T_c}\right)^2\left(\frac{N}{N_0}\right)^{2/15}N^{2/9}\left(\frac{k_BT_ca^2m}{\mathrm{}^2}\right)^{14/15}.$$
(208)
In order to extract results for the relaxation-rate of the condensate number and the phase-diffusion rate it is necessary to know also the mean square of the number-fluctuations $`\mathrm{\Delta }N_0^2`$. This can be evaluated from eq.(34), using the fact that these fluctuations are also dominated by the low-lying modes . The result of this calculation to leading order in $`(\mathrm{}\omega _0/k_BT)`$ is
$`\mathrm{\Delta }N_0^2=`$ $`A\left({\displaystyle \frac{N_0a}{d_0}}\right)^{4/5}\left({\displaystyle \frac{k_BT}{\mathrm{}\omega _0}}\right)^2`$ (209)
$`=`$ $`{\displaystyle \frac{A}{(\zeta (3))^{8/15}}}\left({\displaystyle \frac{T}{T_c}}\right)^2\left({\displaystyle \frac{N_0}{N}}\right)^{4/5}N^{4/3}\left({\displaystyle \frac{k_BT_ca^2m}{\mathrm{}^2}}\right)^{2/5}`$ (210)
with the number $`A`$ given by the multiple sums
$$A=\frac{(15)^{4/5}}{2}\underset{n}{}\underset{n^{}}{}\underset{\mathrm{}}{}\frac{2\mathrm{}+1}{(e(n,\mathrm{})e(n^{},\mathrm{}))^2}\left|_0^1𝑑x(1x^2)x^{2(\mathrm{}+1)}P_{\mathrm{}}^{(2n)}(x)P_{\mathrm{}}^{(2n^{})}(x)\right|^2$$
(211)
In order to find the scaling of $`\mathrm{\Delta }N_0^2`$ in the thermodynamic limit $`N\mathrm{},\omega _00,k_BT_c=\mathrm{}\omega _0(N/\zeta (3))^{1/3}fixed`$ it is necessary to use the form of the preceding results in which $`\mathrm{}\omega _0`$ is eliminated in favor of $`k_BT_c`$ and to use $`N_0N`$.. Then the scaling $`\mathrm{\Delta }N_0^2N^{4/3}`$ derived in is recovered. The particle-number relaxation rate now follows from eqs.(52) and (205) as
$$\gamma _c=\frac{2B_{00}}{A}N_0^{1/5}\left(\frac{a}{d_0}\right)^{6/5}\frac{k_BT}{\mathrm{}}=\frac{2(\zeta (3))^{1/5}B_{00}}{A}\frac{T}{T_c}\left(\frac{N_0}{N}\right)^{1/5}\left(\frac{k_BT_ca^2m}{\mathrm{}^2}\right)^{3/5}\frac{k_BT_c}{\mathrm{}}$$
(212)
It is the largest of the various rates we calculate here but is still small compared to $`\omega _0`$, the inverse time-scale of motion in the trap, by the order of magnitude $`N^{2/3}(Na/d_0)^{6/5}`$.
The phase-collapse rate is obtained from (58). At $`T0`$ ( more precisely above a cross-over temperature of order $`\mathrm{}\omega _0`$) we find
$$\gamma _{collapse}=\frac{15^{2/5}A^{1/2}}{5}N_0^{1/5}\left(\frac{a}{d_0}\right)^{4/5}\frac{k_BT}{\mathrm{}}=\frac{15^{2/5}(\zeta (3))^{2/15}A^{1/2}}{5}\frac{T}{T_c}\left(\frac{N}{N_0}\right)^{1/5}N^{1/3}\left(\frac{k_BT_ca^2m}{\mathrm{}^2}\right)^{2/5}\frac{k_BT_c}{\mathrm{}}.$$
(213)
Apart from the numerical prefactor this is the same asymptotic expression as obtained for the damping-rate $`\gamma _0`$ of the low-lying collective modes (see e.g. ). It is smaller than $`\gamma _c`$ by the order of magnitude $`(N_0a/d_0)^{2/5}`$, i.e. the phase-collapse remains inefficient before phase-diffusion takes over.
The phase-diffusion constant $`D_\varphi ^{(\alpha )}`$ due to the exchange of particles between the condensate and low-lying excitations is gotten by inserting the results for $`\mathrm{\Delta }N_0^2`$ and $`\mathrm{\Gamma }_{00}`$ in eq.(54):
$$D_\varphi ^{(\alpha )}=\frac{(15)^{4/5}A^2}{25B_{00}}N_0^{3/5}\left(\frac{a}{d_0}\right)^{2/5}\frac{k_BT}{\mathrm{}}=\frac{15^{4/5}(\zeta (3))^{1/15}A^2}{25B_{00}}\frac{T}{T_c}\left(\frac{N}{N_0}\right)^{3/5}N^{2/3}\left(\frac{k_BT_ca^2m}{\mathrm{}^2}\right)^{1/5}\frac{k_BT_c}{\mathrm{}}.$$
(214)
It is smaller than $`\gamma _{collapse}`$, again by the order of magnitude of $`(N_0a/d_0)^{2/5}`$.
Finally, the contribution of the fluctuations in the thermal cloud to the phase-diffusion is also obtained from (54) by inserting the result (208) for $`\mathrm{\Gamma }_0+\mathrm{\Gamma }^{}`$:
$`D_\varphi ^{(\gamma )}`$ $`={\displaystyle \frac{3}{10}}{\displaystyle \frac{2^{1/3}B^{(\mu )}}{15^{2/15}\pi ^2}}N_0^{17/15}\left({\displaystyle \frac{a}{d_0}}\right)^{28/15}\left({\displaystyle \frac{k_BT}{\mathrm{}\omega _0}}\right)^2{\displaystyle \frac{k_BT}{\mathrm{}}}`$ (216)
$`={\displaystyle \frac{3}{10}}{\displaystyle \frac{2^{1/3}(\zeta (3))^{16/45}B^{(\mu )}}{15^{2/15}\pi ^2}}\left({\displaystyle \frac{T}{T_c}}\right)^3\left({\displaystyle \frac{N}{N_0}}\right)^{17/15}N^{7/9}\left({\displaystyle \frac{k_BT_ca^2m}{\mathrm{}^2}}\right)^{14/15}{\displaystyle \frac{k_BT_c}{\mathrm{}}}.`$
It differs from the previous rates, which were all proportional to temperature, by the stronger temperature-dependence $`T^3`$. However, this contribution to $`D_\varphi `$ remains smaller than $`D_\varphi ^{(\alpha )}`$ by the order of magnitude $`N^{1/9}(k_BT_ca^2m/\mathrm{}^2)^{11/15}(T/T_c)^2`$.
## IX Discussion and Conclusion
In this paper we have put forward a detailed theory of fluctuations and relaxation processes of the condensate in thermal equilibrium with the cloud of its excitations. For a given number of particles $`N_0`$ in the condensate, we have defined the condensate mode as the corresponding normalized solution of the Gross-Pitaevskii equation, defining at the same stroke the $`N_0`$-dependent part of the chemical potential. The equilibrium value of $`N_0`$ is distinguished as the value of $`N_0`$ for which the number of particles in the thermal cloud in equilibrium with the condensate plus $`N_0`$ equals $`N`$. We have calculated the fluctuations of $`N_0`$ around its equilibrium value and also the fluctuations of the phase of the complex amplitude $`\alpha _0`$ of the condensate with $`|\alpha _0|^2=N_0`$. In a general phenomenological framework presented in the first part of this paper we were able to separate the fluctuations of the complex condensate amplitude into several contributions, which have different physical origin:
– The fluctuation of the atom-number in the condensate, which are driven by the exchange of atoms between the condensate and the thermal cloud.
– The fluctuation of the chemical potential with two different contributions, namely the fluctuations of $`\mu `$ due to number-fluctuations in the condensate, and the faster fluctuations of $`\mu `$ at constant $`N_0`$ caused by number-fluctuations in the excitations.
The importance of number-fluctuations in the condensate, assumed at first in the phenomenological approach due to the importance of $`N_0`$ for the value of the chemical potential, but later born out by the microscopic calculations, leads to the appearance of the linear relaxation-rates $`\gamma _c`$ of the condensate-number as an important characteristic inverse time-scale of the problem. At times much shorter than $`\gamma _c^1`$ phase-diffusion of the condensate-phase due to the fast number-fluctuations in the excitations can occur. In the same regime may also occur the process of collapse due to the reversible spreading of the phase caused by the static uncertainty in $`N_0`$ and the associated chemical potential.
At times large compared to $`\gamma _c^1`$ the number-fluctuations in the condensate are dynamical and irreversible, and lead to the replacement of the reversible collapse by an irreversible phase-diffusion with a larger diffusion-rate than in the short-time regime.
The second and larger part of this paper was devoted to microscopic theory. First we have provided a microscopic derivation of the phenomenological Langevin equation, established microscopic formulas for all phenomenological parameters and also exhibited the relation between the short-time diffusion rate and fluctuation rates of the population numbers of excitations via a sum-rule. Then the microscopic theory was used to evaluate the transport-parameters and the various rates as a function of temperature, particle-number and the scattering length of the interaction potential. The evaluation was done for two simple cases – the cubic box-like trap, where the form of the condensate mode does not depend on $`N_0`$ and the thermal cloud penetrates the condensate homogeneously, and the isotropic harmonic trap, where the form of the condensate-mode changes with $`N_0`$ and the thermal cloud is located preferentially near the boundaries of the condensate. The physically important results for both kinds of traps are similar, even though they have to differ, obviously, in the details of the scalings with the atom-numbers and the scattering length.
The calculation of the transport-parameters reveals some interesting physical results:
– The fluctuations driving the absolute value $`|\alpha _0|`$ and the phase $`\varphi `$ of $`\alpha _0`$ are quite different in strength, those driving $`|\alpha _0|`$ being the much stronger ones. The reason for this is a pronounced squeezing of the bath of thermal excitations with respect to the instantaneous phase of the condensate. This squeezing reaches nearly 100% for the lowest-lying modes, which is the reason that fluctuations of $`\varphi `$ are practically not driven by such modes. On the other hand, the contribution of the high-lying modes to the fluctuating forces driving $`|\alpha _0|`$ and $`\varphi `$ is nearly the same (after the obvious normalization with $`|\alpha _0|`$), i.e. there is no squeezing in this (much weaker) contribution to the noise.
– The cross-correlation between the fluctuations driving $`|\alpha _0|`$ and $`\varphi `$ are found to vanish exactly in a real condensate, where both the Gross-Pitaevskii equation and the Bogoliubov-Fetter equations are real and all solutions can (but need not) be taken real. This can also be understood as a general consequence of time-reversal symmetry: $`\varphi `$ is a velocity potential and therefore odd under time-reversal while $`|\alpha _0|`$ is even under time-reversal. Their fluctuating forces therefore transform oppositely. In a time-reversal symmetric condensate (no vortices) the cross-correlation between an even and an odd quantity under time-reversal must vanish.
It turns out that the relaxation rate $`\gamma _c`$ of the atom-number in the condensate is the largest of the calculated rates. In particular it is larger than the collapse-rate and the phase-diffusion rate which, like $`\gamma _c`$, are proportional to temperature in the regime $`k_BT>\mu `$. It is also larger than the decay rates of the lowest-lying collective modes $`\gamma _0`$, which might look surprising because at the same time the theory tells us that $`\gamma _c`$ is dominated by the particle transfer-rates between the condensate and the low-lying modes. However, it is clear that $`\gamma _c`$ ought to be larger than $`\gamma _0`$ because the condensate couples to all low-lying modes in parallel which increases the number of decay channels by a factor proportional to the ratio of the chemical potential and the lowest-lying mode frequency.
The next largest rate we find is the thermal phase-collapse rate $`\gamma _{collapse}`$. It turns out to have the same functional dependence on $`T,a,N_0`$ and $`N`$ as the decay-rate of the lowest-lying collective modes. I cannot see any fundamental reason for this coincidence and have to count it just as that. Physically the smallness of $`\gamma _{collapse}/\gamma _c`$ means that the phase-collapse will not be observable at finite temperature because it can only lead to a decay-factor $`exp(\frac{1}{2}(\gamma _{collapse}/\gamma _c)^2)`$ very close to 1 before phase-diffusion takes over.
Finally, the phase-diffusion rate $`D_\varphi `$ is the smallest of the rates calculated here. We find the simple nice result that the ratios $`\gamma _c/\gamma _0`$ and $`\gamma _0/D_\varphi `$ are of about equal order of magnitude, given by the ratio of $`\mu `$ to the smallest excitation energy, which is $`\mathrm{}\omega _0`$ for the harmonic and $`\sqrt{\mu /m}(2\pi \mathrm{}/L)`$ for the box-like trap. Instead of $`\gamma _0`$ we may also take $`\gamma _c`$ in these ratios with the same conclusion. $`D_\varphi `$ like the rate $`\gamma _c`$ is found to be dominated by the atom-number exchange between the condensate and the low-lying modes.
This observation actually explains the coincidence of the two ratios we have just indicated and turns them into a precise relation: In (54) for $`D_\varphi `$ we put $`\mathrm{\Gamma }^\mathrm{"}=0`$ which is exact for real condensate-modes and neglect $`\mathrm{\Gamma }_0+\mathrm{\Gamma }^{}`$, which comes from high-lying excitations. Then multiplying the resulting expression for $`D_\varphi `$, with $`\gamma _c=\tau _c^1`$ from (52), we readily find
$$\frac{1}{2}D_\varphi \gamma _c=\gamma _{collapse}^2$$
(217)
with $`\gamma _{collapse}`$ from (58) again with $`\mathrm{\Gamma }^\mathrm{"}=0`$.
Let us now compare our results with related ones found in the literature. Most closely related to the present work in goal and scope is a paper by Jaksch et al on the intensity and amplitude fluctuations of a Bose-Einstein condensate at finite temperature, which builds on extensive earlier work by Gardiner and Zoller with collaborators (cf. the references given in ). Unlike the present paper it also takes into account trap losses. The theory presented in is based on a conceptual division of the Bose gas into two energy regions called the condensate band and the noncondensate band. In this construction the boundary between the two regions is chosen in such a way that the noncondensate band is not significantly affected by the mean field of the condensate, while the influence of excitations in the condensate band is neglected. Thus the main physical difference of to the present work is that it neglects fluctuations of particles from the condensate mode to quasiparticle modes as well as to very low-lying one-particle excitations.
By contrast in the present work we avoid the division of the energy region into two parts. We find, as we have discussed, that the exchange of particles between the condensate and the low-lying modes makes not only an important but in fact the dominant contribution to the relaxation rate of the condensate-number and the phase-diffusion rate, determining their dependence on temperature, atom number, population of the condensate and scattering length.
The importance of the particle-exchange between the low-lying excitations for the phase-diffusion of the condensate and the number-relaxation rate $`\gamma _c`$ was first pointed out in , while for the intensity of the number-fluctuations in the condensate this had already been shown in . The theory put forward in already proceeded along essentially the same lines we follow here, but it had some short-comings which we overcome and correct in the present work: The squeezing of the noise from the thermal cloud with respect to the phase of the condensate was briefly remarked upon in , but was not taken into account in the calculation of the transport coefficients presented there and in the formula for phase-diffusion. Moreover, in the conservative part of the Langevin-equation (41) $`\mathrm{\Delta }_0\mu `$ was replaced by $`\mathrm{\Delta }F(|\alpha )|^2)/|\alpha _0|^2`$ in , which, on scrutiny, appears questionable when used in conjunction with the fluctuation formula (32,34) for $`\mathrm{\Delta }N_0^2`$. After all, neither $`\mathrm{\Delta }_0\mu `$ nor $`\mathrm{\Delta }F`$ are equilibrium quantities. The use of the aforementioned relation between them is therefore avoided here.
Even though in the present paper I have opted for the use of the fluctuation-formulas (32,34), which in my opinion have a firm basis, it is only fair to mention that they are still under debate in the current literature, see e.g. . In another recent paper with some bearing on this topic Bergeman et al. use as equilibrium distribution for the condensate number $`P(N_0)\mathrm{exp}[(\mu N_0\frac{5}{14}(15N_0a/d_0)^{2/5}N_0)/k_BT]`$, (cf. the discussion after their eq.(21)), which implies $`\mathrm{\Delta }N_0^2TN_0^{3/5}`$, a result which is rather different, both in the temperature-dependence and in the scaling with the particle-number, from the result (209) on which our present calculations have been based. It is clear that not the method but the details of our results on the dynamics of the fluctuations of the condensate would change, if the results on the statics would be changed. Needless to say that a resolution of the theoretical debate concerning the correct approach to the statics seems urgent and would be highly wellcome. Vice versa experimental results on the dynamics (i.e. on $`\gamma _c`$ and $`D_\varphi `$) would also help to decide, by applying the theory presented here, which of the approaches to the statics of the number-fluctuations in the condensate proposed in the literature describes the physics correctly.
A quantum kinetic theory of trapped atomic gases has also been formulated by Stoof . In the general coupled Fokker-Planck equations of the condensate and the excited modes are presented and applied to the kinetics of the formation of a condensate. This problem has also been studied by Gardiner and coworkers as well as Kagan and Svistunov , where also earlier work by further authors is quoted.
By contrast the present work has focussed on the fluctuations around the equilibrium state of the condensate, after it has been formed. However, the application of our approach to the kinetics of the formation of a condensate would be an interesting goal for future work.
Experimentally the rates $`\gamma _c`$ and $`D_\varphi `$ we have calculated should be measurable. The rate $`\gamma _c`$ may be observable as the relaxation rate of the condensate back to its equilibrium state after creating a non-equilibrium state by a sudden small change of temperature via evaporative cooling. The sum of the phase-diffusion rates of two condensates could be measured by monitoring the phase-difference between them after it was initially fixed by measurement or preparation at a reference-time $`t=0`$. Methods for measuring phase-differences in Bose-Einstein condensates have recently been demonstrated . It is to be hoped therefore that the phase-diffusion in Bose-Einstein condensates - a fundamental process intimately linked to the spontaneously broken gauge symmetry in a finite system - will be measured in the near future.
## Acknowledgement
Usefull discussions with Walter Strunz are gratefully acknowledged. This work has been supported by the Deutsche Forschungsgemeinschaft through the Sonderforschungsbereich 237 “Unordnung und große Fluktuationen”.
## A
Here we wish to derive the expression (15) for$`H_0`$. Using the Gross-Pitaevskii equation (4) we put (8) in the form
$`H_0=(\mu _0\mu )|\alpha _0|^2{\displaystyle \frac{U_0}{2}}|\alpha _0|^4{\displaystyle d^3x\psi _0^4}`$ (A1)
Taking the derivative with respect to $`|\alpha _0|^2`$ we get
$`{\displaystyle \frac{H_0}{|\alpha _0|^2}}=\mu _0\mu +|\alpha _0|^2\left({\displaystyle \frac{\mu _0}{|\alpha _0|^2}}U_0{\displaystyle d^3x\psi _0^4}|\alpha _0|^2U_0{\displaystyle d^3x\psi _0^2\frac{\psi _0^2}{|\alpha _0|^2}}\right)`$ (A2)
To evaluate this further we use eq.(4) and its derivative with respect to $`|\alpha _0|^2`$
$`({\displaystyle \frac{\mathrm{}^2}{2m}}^2+V\mu _0+3U_0|\alpha _0|^2\psi _0^2){\displaystyle \frac{\psi _0}{|\alpha _0|^2}}=({\displaystyle \frac{\mu _0}{|\alpha _0|^2}}U_0\psi _0^2)\psi _0`$ (A3)
Multiplying eq.(A3) with $`\psi _0`$ and integrating over space, using the Gross-Pitaevskii equation (4) after partial integration, we derive the identity
$`U_0|\alpha _0|^2{\displaystyle d^3x\psi _0^2\frac{\psi _0^2}{|\alpha _0|^2}}={\displaystyle \frac{\mu _0}{|\alpha _0|^2}}U_0{\displaystyle \psi _0^4d^3x}`$ (A4)
which is used in in (A2) to yield $`H_0/|\alpha _0|^2=\mu _0\mu `$, and upon integration results in (15).
## B
Here we wish to derive eq.(72). This is achieved if we succeed to show that the coupling of the condensate and the thermal cloud via
$$\widehat{H}_3=U_0\sqrt{N_0}d^3x\stackrel{~}{\psi }_0\widehat{\stackrel{~}{\chi }}^+(e^{i\varphi }\widehat{\stackrel{~}{\chi }}+e^{i\varphi }\widehat{\stackrel{~}{\chi }}^+)\widehat{\stackrel{~}{\chi }}$$
(B1)
gives rise to the systematic change of $`\mathrm{}(\widehat{\xi }(t))`$, to first order in the interaction, of
$$\delta \mathrm{}(\widehat{\xi }(t)_\varphi =\frac{2\sqrt{N_0}}{\mathrm{}k_BT}_{\mathrm{}}^tdt^{}S_{JJ}(tt^{})\frac{H_0(t^{})}{|\alpha _0|^2},$$
(B2)
because this can then be used in (65) to yield (72). In (B1) we could put $`|\alpha _0|=\sqrt{N_0}`$ since we linearize around equilibrium and only wish to calculate the dissipation in $`|\alpha _0|^2`$ which is conjugate to $`\varphi `$, the variable we kept in (B1). Standard first order perturbation theory with adiabatic switch-on of the interaction gives, with $`ϵ+0`$,
$$\delta \mathrm{}(\widehat{\xi }(t)_\varphi =\frac{i}{\mathrm{}}_{\mathrm{}}^tdt^{}[\mathrm{}(\widehat{\xi }(t)),\widehat{H}_3(t^{})]_\varphi e^{ϵt^{}}.$$
(B3)
We can rewrite this as
$$\delta \mathrm{}(\widehat{\xi }(t))_\varphi =2i\sqrt{N_0}_{\mathrm{}}^t𝑑t^{}\left(\chi _{J\widehat{\xi }}^{\prime \prime }(t,t^{})e^{i\varphi (t^{})}+\chi _{J\widehat{\xi }^+}^{\prime \prime }(t,t^{})e^{i\varphi (t^{})}\right)e^{ϵt^{}}$$
(B4)
where we introduced the response functions
$`\chi _{J\widehat{\xi }}^{\prime \prime }(t,t^{})`$ $`={\displaystyle \frac{1}{2\mathrm{}}}[\mathrm{}(\widehat{\xi }(t)),\widehat{\xi }(t^{})]_\varphi \mathrm{\Theta }(tt^{})`$ (B5)
$`\chi _{J\widehat{\xi }^+}^{\prime \prime }(t,t^{})`$ $`={\displaystyle \frac{1}{2\mathrm{}}}[\mathrm{}(\widehat{\xi }(t)),\widehat{\xi }^+(t^{})]_\varphi \mathrm{\Theta }(tt^{})`$ (B7)
with
$$\widehat{\xi }(t)=U_0d^3x\stackrel{~}{\psi }_0\widehat{\stackrel{~}{\chi }}^+(t)\widehat{\stackrel{~}{\chi }}(t)\widehat{\stackrel{~}{\chi }}(t).$$
(B8)
Here $`\mathrm{\Theta }(tt^{})`$ is the Heaviside step-function. We shall define $`\mathrm{\Theta }(0)=0`$ without loss of generality. The fluctuation-dissipation theorem (in the classical frequency domain $`\mathrm{}\omega k_BT`$) ensures the relations
$`\chi _{J\widehat{\xi }}^{\prime \prime }(t,t^{})`$ $`={\displaystyle \frac{i\mathrm{\Theta }(tt^{})}{2k_BT}}{\displaystyle \frac{}{t^{}}}S_{J\widehat{\xi }}(t,t^{})`$ (B9)
$`\chi _{J\widehat{\xi }^+}^{\prime \prime }(t,t^{})`$ $`={\displaystyle \frac{i\mathrm{\Theta }(tt^{})}{2k_BT}}{\displaystyle \frac{}{t^{}}}S_{J\widehat{\xi }^+}(t,t^{})`$ (B11)
with the correlation functions<sup>\*†</sup><sup>\*†</sup>\*†$`\mathrm{}(\widehat{\xi }(t))`$ according to eq.(67) contains an explicit external time-dependence via $`\varphi (t)`$, in addition to the internal time-dependence of $`\widehat{\stackrel{~}{\chi }}(t),\widehat{\stackrel{~}{\chi }}^+(t)`$ via their Heisenberg equations of motion. This explicit time-dependence has to be taken into account when applying the fluctuation-dissipation theorem. We avoid this additional step by applying the time-derivative in the fluctuation-dissipation relation (LABEL:B7) directly to the second time argument $`t^{}`$, of course with the appropriate extra minus-sign.
$`S_{J\widehat{\xi }}(t,t^{})`$ $`=\mathrm{}(\widehat{\xi }(t))\widehat{\xi }(t^{})_\varphi `$ (B12)
$`S_{J\widehat{\xi }^+}(t,t^{})`$ $`=\mathrm{}(\widehat{\xi }(t))\widehat{\xi }^+(t^{})_\varphi .`$ (B14)
We can use (LABEL:B7) in (B4) and apply a partial integration in $`t^{}`$ to obtain
$`\delta \mathrm{}(\widehat{\xi }(t))_\varphi =`$ $`{\displaystyle \frac{\sqrt{N_0}}{ik_BT}}{\displaystyle _{\mathrm{}}^t}𝑑t^{}\left(S_{J\widehat{\xi }}(t,t^{})e^{i\varphi (t^{})}S_{J\widehat{\xi }^+}(t,t^{})e^{i\varphi (t^{})}\right){\displaystyle \frac{d\varphi (t^{})}{dt^{}}}e^{ϵt^{}}`$ (B16)
$`{\displaystyle \frac{\sqrt{N_0}}{k_BT}}\left(S_{J\widehat{\xi }}(t,t)e^{i\varphi (t)}+S_{J\widehat{\xi }^+}(t,t)e^{i\varphi (t)}\right)`$
which can be rewritten as
$`\delta \mathrm{}(\widehat{\xi }(t))_\varphi =`$ $`{\displaystyle \frac{2\sqrt{N_0}}{k_BT}}{\displaystyle _{\mathrm{}}^t}𝑑t^{}S_{JJ}(t,t^{}){\displaystyle \frac{d\varphi (t^{})}{dt^{}}}{\displaystyle \frac{2\sqrt{N_0}}{k_BT}}S_{JR}(t,t)`$ (B17)
with $`S_{JJ}(t,t^{})`$ defined by eqs.(67,69) and
$$S_{JR}(t,t^{})=\frac{1}{4i}(\widehat{\xi }(t)e^{i\varphi (t)}\widehat{\xi }^+(t)e^{i\varphi (t)})(\widehat{\xi }(t^{})e^{i\varphi (t^{})}+\widehat{\xi }^+(t^{})e^{i\varphi (t^{})})_\varphi .$$
(B18)
The constant term with $`S_{JR}(t,t)`$ amounts to a small shift of the equilibrium value of $`\mu `$ in the final result which we shall neglect like other terms contributing to such shifts. Then using eq.(60) we put $`\mathrm{}d\varphi (t^{})/dt^{}=H_0(t^{})/|\alpha _0|^2`$ in eq.(B17) which establishes (B1) and hence eq.(72).
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# Adiabatic compression of a trapped Fermi gas
## Abstract
We propose a method to reach conditions of high degeneracy in a trapped Fermi gas, based on the adiabatic transfer of atoms from a magnetic to a tighter optical trap. The transformation yields a large increase of the Fermi energy, without a significant change of the temperature. The large enhancement of the central density emphasizes the role of the interactions and makes the system much closer to the BCS transition. An estimate of the time needed to achieve the conditions of adiabaticity is also given.
The experimental realization of a highly degenerate atomic Fermi gas confined in traps is a task of primary importance, especially in view of the perspective of approaching the BCS transition to the superfluid phase. The regime of quantum degeneracy has been already reached in a sample of potassium atoms , where first signatures of Fermi statistics, like the deviation of the velocity distribution from a Boltzmann profile and the increase of the kinetic energy with respect to the classical value, have been observed. The main difficulties in further lowering the temperature are due to the fact that the efficiency of the evaporative cooling process is strongly quenched . In fact Fermi statistics inhibits collisional processes at low temperature both directly, by reducing the phase space available for collisions, and indirectly by lowering the density of the sample because of Pauli repulsion. Procedures to optimize the evaporative process have recently permitted to reach lower temperatures, of the order of $`0.20.3`$ $`T_F`$ where $`T_F`$ is the Fermi temperature .
The purpose of this work is to propose a method to reach conditions of high degeneracy, based on an adiabatic compression of the gas. A similar method has proven quite successful in producing Bose-Einstein condensation in a reversible way, starting from a trapped Bose gas above the critical temperature . The main point is that, by changing the shape of the confinement from a harmonic to a non harmonic trap, one can increase the degree of quantum degeneracy by keeping the entropy of the total system constant . At the same time, the process of thermalization is not drastically quenched and in typical experimental conditions takes place over times much shorter than the lifetime of the cloud. In the following we consider a gas occupying two different spin states, initially confined by a harmonic trap. We then switch on adiabatically a second tighter trap (see Fig. 1). Experimentally this can be realized using a magnetic trap for the first confinement and an optical trap for the second one. As a consequence of the adiabatic process a fraction of atoms will move from the magnetic to the optical trap. This can provide several important advantages:
i) The gas in the optical trap becomes much more degenerate than the original one. In particular, if the number of atoms transferred to the optical trap is a small fraction, the temperature will not change significantly with respect to the initial value, but the Fermi energy will increase, in a way proportional to the depth of the optical trap.
ii) The gas in the optical trap is much denser due to the tighter confinement. This produces an increase of interaction effects and hence of the value of the critical temperature. Both effects i) and ii) favour the reachability of the BCS transition.
iii) The Fermi energies of the two spin components, which initially were different because of the different magnetic trapping, become closer in the optical trap, thereby favouring the mechanism of Cooper pairing.
Another important advantage of the proposed method is that the velocity distributions of the atomic clouds occupying the magnetic and optical traps can be measured separately. By releasing first the magnetic trap, one can measure the temperature of the sample. The effects of quantum degeneracy can then be investigated by measuring the velocity distribution of the gas confined in the optical trap.
In the first part of the work we assume that thermodynamic equilibrium is ensured during each step of the adiabatic compression and we explore the properties of the new gas produced in the tight confinement by imposing entropy conservation. In the second part of the work we provide estimates for the times required to ensure adiabaticity and discuss possible scenarios where faster adiabatic trasformations take place out of thermal equilibrium.
Let us consider a gas initially confined in a harmonic trap (hereafter called magnetic trap). We assume that the trapping frequencies and the number of atoms are the same for the two spin species. The Fermi energy is given by $`ϵ_F^0=\mathrm{}\omega _{mag}(6N)^{1/3}`$, where $`N`$ is the number of atoms of each species and $`\omega _{mag}`$ is the geometrical average of the frequencies characterizing the magnetic trapping potential $`V_{mag}`$. We will consider systems lying initially in configurations of moderate degeneracy corresponding to $`k_BT_i=0.20.5ϵ_F^0`$, where $`T_i`$ is the initial temperature of the gas. For simplicity, we will assume that also the optical trap can be approximated by a harmonic potential $`V_{opt}`$ having a tighter frequency $`\omega _{opt}\omega _{mag}`$, and depth $`V_{opt}(𝐫=0)=U`$. In our model the trapping potential is then defined as $`V_{ext}(𝐫)=V_{opt}(𝐫)`$ inside the optical trap ($`V_{opt}<V_{mag}`$), and $`V_{ext}(𝐫)=V_{mag}(𝐫)`$ outside (see Fig. 1). The maximum number of atoms that can be transferred in the optical trap is given by the value
$$N_{opt}=\frac{1}{6}\left(\frac{U}{\mathrm{}\omega _{opt}}\right)^3.$$
(1)
As we will see, if we start from moderately low temperatures Eq.(1) provides an accurate estimate of the number of atoms which are actually transferred by the adiabatic process, provided $`N_{opt}N`$. The relative number of atoms transferred in the optical trap is then given by the useful expression
$$\frac{N_{opt}}{N}=\left(\frac{U\omega _{mag}}{ϵ_F^0\omega _{opt}}\right)^3.$$
(2)
Typical values that will be considered are $`\omega _{mag}/\omega _{opt}=0.1`$ and $`U/ϵ_F^0=5`$, corresponding to $`N_{opt}/N`$ 10%. Since the number of transferred atoms is small one expects that the final temperature $`T_f`$ of the gas will not change significantly with respect to the initial value $`T_i`$. The final degree of degeneracy of the gas will be however significantly higher, since the final Fermi energy is approximately given by $`ϵ_FU+ϵ_F^0`$. At the same time the central density of the gas, that at small temperature is equal to $`n(0)=(2mϵ_F^0/\mathrm{}^2)^{3/2}/(6\pi ^2)`$ in the initial stage, will increase by the factor $`(ϵ_F/ϵ_F^0)^{3/2}`$ in the optical trap. The increase of the Fermi energy and of the central density has an important effect on the value of the BCS temperature which is expected to behave, for negative scattering lengths, as
$$T_{BCS}\frac{ϵ_F}{k_B}\mathrm{exp}\left[\frac{\mathrm{}\pi }{2p_Fa}\right],$$
(3)
where $`p_F=\mathrm{}(6\pi ^2)^{1/3}n(0)^{1/3}`$ is the Fermi momentum calculated in the center of the trap. The value of the BCS temperature can increase significantly with respect to its value in the magnetic trap. For example, by taking $`a2000a_0`$, a ratio $`U/ϵ_F^0=5`$ and an initial central density $`n(0)10^{12}`$ cm<sup>-3</sup>, we obtain an increase of the density by a factor $``$15 and the BCS temperature (3) increases by the huge factor $``$50, becoming comparable to the initial value of the Fermi temperature $`ϵ_F^0/k_B`$. The above discussion suggests that the proposed adiabatic mechanism might provide conditions of high degeneracy, not far from the transition to the BCS phase. With such a denser gas also the effects of the mean field on the density profile may be significant. An estimate is given by the ratio $`E_{int}/E_{ho}0.3p_Fa/\mathrm{}`$ between the interaction energy and the oscillator energy of a spherically symmetric trap. By using the values employed above one finds corrections of the order of 10 %.
The high degeneracy realized in the gas is expected to show up in the velocity distribution and in the release energy. After completing the adiabatic transfer one can release the magnetic trap. Measuring the velocity distribution of these atoms then provides information on the temperature of the system which is expected to be close to the initial value. The atoms of the optical trap can be imaged in a second step. For them one predicts that the ratio $`E_{kin}/k_BT_f3U/8k_BT_f`$ between the kinetic energy per particle and the thermal energy, should be enhanced in a significant way if $`Uk_BT_f`$ revealing the effects of quantum degeneracy.
To confirm the scenario emerging from the above discussion we have carried out a numerical calculation of the thermodynamic functions before and after the adiabatic transformation. The calculation is obtained by imposing that the initial and final configurations have the same entropy. This has been calculated using the semiclassical expression
$`{\displaystyle \frac{S}{k_B}}`$ $`=`$ $`{\displaystyle \frac{1}{(2\pi \mathrm{})^3}}{\displaystyle }d𝐫d𝐩[{\displaystyle \frac{ϵ(𝐩,𝐫)/k_BT\mathrm{log}z}{z^1e^{ϵ(𝐩,𝐫)/k_BT}+1}}.`$ (4)
$`+`$ $`.\mathrm{log}(1+ze^{ϵ(𝐩,𝐫)/k_BT})],`$ (5)
where $`z=\mathrm{exp}(\mu /k_BT)`$ is the gas fugacity and
$$ϵ(𝐩,𝐫)=\frac{p^2}{2m}+V_{ext}(𝐫)$$
(6)
are the semiclassical particle energies. The results are presented in Figs. 2-5. In Fig. 2 we show the relative number of atoms in the optical trap as a function of the depth of the optical trap $`U/ϵ_F^0`$ for two initial temperatures. For the chosen configuration with $`\omega _{mag}/\omega _{opt}=0.1`$ and final depth $`U=5ϵ_F^0`$, the optical trap can host about 10 % of atoms confirming the analytic prediction (2).
In Fig. 3 the final temperature of the gas is plotted as a function of $`U/ϵ_F^0`$. As already anticipated the final temperature of the gas does not change significantly from the original value, except for values of $`U/ϵ_F^0`$ of the order of $`\omega _{opt}/\omega _{mag}`$. Notice however that, due to the tight confinement, the temperature of the gas can become comparable to the optical oscillator temperature $`\mathrm{}\omega _{opt}/k_B`$ if $`N`$ is not large enough.
In Fig. 4 we show the kinetic energy per atom of the gas in units of the final temperature. The initial temperature is $`k_BT_i=0.25ϵ_F^0`$. In the same figure, we also show the kinetic energy per particle of the gas occupying separately the magnetic and the optical trap. These values are obtained by averaging separately the kinetic energy over the particles occupying the optical trap ($`ϵ(𝐩,𝐫)<0`$), and the particles occupying the magnetic trap ($`ϵ(𝐩,𝐫)>0`$). The kinetic energy provides an important indicator of the quantum degeneracy of the gas. The effect is very spectacular for the atoms of the optical trap, where at $`U=5ϵ_F^0`$ one finds $`E_{kin}/k_BT_f7`$. By choosing a higher initial temperature the effect is less pronounced, but still large ($`E_{kin}/k_BT_f3`$ for $`k_BT_i=0.5ϵ_F^0`$ and $`U=5ϵ_F^0`$). Because of the high degeneracy the velocity distribution of the atoms in the optical trap deviates significantly from a Boltzmann distribution (Fig. 5). We have also calculated the central density as a function of the depth $`U`$. For $`U=5ϵ_F^0`$ and $`k_BT_i=0.25ϵ_F^0`$, we find an increase by a factor 15 with respect to the initial value.
In the second part of the work we discuss the conditions needed to achieve the proposed adiabatic transformation. Let us recall that there are several time scales in the problem. A first scale is fixed by the periods of the harmonic wells. These times (of order of $`10^3`$ \- $`10^1`$ sec) are expected to be shorter than the relaxation times due to collisions. The condition of reversibility requires that the relaxation time be much shorter than the time over which the optical trap is switched on. We have simulated the process of thermalization during the gradual increase of the depth of the optical trap by solving the quantum Boltzmann equation. We assume equal distribution functions for the two spin components and that the phase-space distribution of particles is a function only of the single-particle energies $`ϵ(𝐩,𝐫)`$ (ergodic assumption) . This yields the following equation to solve
$`\rho (ϵ_1){\displaystyle \frac{f(ϵ_1)}{t}}`$ $`=`$ $`{\displaystyle \frac{m\sigma }{\pi ^2\mathrm{}^3}}{\displaystyle 𝑑ϵ_3𝑑ϵ_4\rho (ϵ_{min})}`$ (7)
$`\times `$ $`[f(ϵ_3)f(ϵ_4)(1f(ϵ_1))(1f(ϵ_2))`$ (8)
$``$ $`f(ϵ_1)f(ϵ_2)(1f(ϵ_3))(1f(ϵ_4))],`$ (9)
where $`\rho (ϵ)`$ is the density of states in the potential $`V_{ext}`$, which we model by $`\rho (ϵ)=(ϵ+U)^2/2(\mathrm{}\omega _{opt})^3`$ if $`ϵ<0`$, and $`\rho (ϵ)=ϵ^2/2(\mathrm{}\omega _{mag})^3`$ if $`ϵ>0`$. Furthermore, $`ϵ_{min}=`$min$`\{ϵ_1,ϵ_2,ϵ_3,ϵ_4\}`$ is the minimum value of the four single-particle energies involved in the collisions and $`ϵ_2=ϵ_3+ϵ_4ϵ_1`$ according to energy conservation. The cross-section for collisions between the two distinguishable spin states is $`\sigma =4\pi a^2`$ and is fixed by the $`s`$-wave scattering length $`a`$. In the numerical simulation the depth $`U`$ of the optical trap is ramped up by keeping the ratio $`\omega _{mag}/\omega _{opt}`$ constant, and at each step we let the gas thermalize in the new configuration. From Eq. (9), it turns out that the time scale of thermalization is fixed by $`\tau ^1=\omega _{mag}(|a|/a_{mag})^2N^{2/3}`$, where $`a_{mag}=(\mathrm{}/m\omega _{mag})^{1/2}`$ is the magnetic oscillator length. If we choose $`\omega _{mag}/\omega _{opt}=0.1`$, $`U=5ϵ_F^0`$ and $`k_BT_i=0.5ϵ_F^0`$, we find that the time required for the reversible transformation is $`t_{rev}100\tau `$. For a configuration with $`N=10^6`$, $`\omega _{mag}=2\pi \times 100`$ Hz and $`|a|/a_{mag}=5\times 10^3`$, this corresponds to $`t_{rev}1`$ sec. If one starts from the lower initial temperature $`k_BT_i=0.25ϵ_F^0`$ and keeps the other parameters unchanged, the time $`t_{rev}`$ turns out to be about a factor two shorter. The origin of this behavior is due to the mechanism of atom exchange between the magnetic and the optical trap. Notice, however, that $`t_{rev}`$ depends rather strongly on the value of $`N`$ and becomes longer if $`N`$ decreases.
The true shape of the optical trap differs from the harmonic potential employed above. Differences appear because a dipole trap can be safely approximated by a harmonic potential only at its center and, in general, can host more atoms than the corresponding harmonic trap with the same depth $`U`$ and frequency ratio $`\omega _{mag}/\omega _{opt}`$. This results in more heating and, as a consequence, in a higher final temperature $`T_f`$. For $`k_BT_i=0.25ϵ_F^0`$, we find that the fraction of atoms in the optical trap at $`U=5ϵ_F^0`$ is about 50% larger than the value obtained using the harmonic approximation, while the final temperature is only about 10% higher. The final degree of degeneracy of the gas in the optical trap is consequently slightly reduced, but the qualitative features discussed above remain unchanged.
If the time of the transformation is longer than the inverse oscillator frequency, but faster than the relaxation time (a relatively easy condition to realize experimentally), then the transformation will be adiabatic, in the sense of entropy conservation, but the system will not be in thermal equilibrium. Such a transformation corresponds to an adiabatic transfer of the lowest $`N_{opt}`$ single particle states from the magnetic to the optical trap. The transformation keeps the corresponding occupation numbers unchanged. In this case, the atoms in the optical trap form a cold, out of equilibrium gas. With the initial condition $`k_BT_i=0.5ϵ_F^0`$ this out-of-equilibrium transformation transfers approximately the same fraction of atoms into the optical trap as with the fully reversible transformation discussed above and the average kinetic energy of these atoms takes a similar value, revealing that the system is highly degenerate. By using the quantum Boltzmann equation (9) we have found that the time needed for this configuration to relax to equilibrium is significantly shorter, the final temperature of the gas being only sligthly higher than the one obtained in the reversible transformation.
In conclusion we have investigated the consequences of an adiabatic transfer of a gas of Fermions from a magnetic to a tighter optical trap. We have seen that it is possible to reach configurations involving a significant fraction of atoms which, as a result of the transformation, will occupy a Fermi sea in conditions of high degeneracy. The large enhancement of the release energy should be easily observable through time of flight measurements. This method could make it possible to produce highly degenerate Fermi gases, which, in the case of gases interacting with negative scattering length, are close to the transition to the BCS phase.
Useful discussions with Brian DeMarco are acknowledged. It is also a pleasure to thank Wolfgang Ketterle for useful comments. L. V. wishes to thank the Danish Ministry of Education for support. This research was supported by Ministero dell’Università e della Ricerca Scientifica e Tecnologica (MURST).
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# Mapping IR Enhancements in Closely Interacting Spiral-Spiral Pairs. I. ISO CAM and ISO SWS Observations1footnote 11footnote 1Based on observations made with ISO, an ESA project with instruments funded by ESA Member States and with the participation of ISAS and NASA.
## 1 Introduction
It is well established from IRAS studies that many interacting galaxies show enhanced far-infrared (FIR) emission compared to non-interacting galaxies (Lonsdale et al. 1984; Kennicutt et al. 1987; Telesco et al. 1988). Most of the emission is due to young massive stars formed in recent starbursts, supporting the idea that galaxy-galaxy interactions can stimulate active star formation (Larson & Tinsley 1978; Rieke et al. 1980; Condon et al. 1982; Balzano 1983; Gehrz et al. 1983; Joseph et al. 1984; Cutri & McAlary 1985); albeit, dust-obscured AGNs might be responsible for the emission in some cases (Sanders et al. 1988a; Surace & Sanders 1999; Genzel et al. 1998; Lutz et al. 1998). The most extreme FIR luminosities ($`\genfrac{}{}{0pt}{}{>}{}\mathrm{\hspace{0.17em}10}^{12}L_{}`$; H<sub>0</sub>=75 km/s Mpc<sup>-1</sup>) and highest dust temperatures are found in galaxy mergers where the identity of the component galaxies is often indeterminate (Sanders et al. 1988a; Melnick & Mirabel 1990; Mazzarella et al. 1991; Sanders & Mirabel 1996). After these ultraluminous IR galaxies (ULIRGs), violent starbursts are most likely to be found in closely interacting pairs of spiral galaxies with overlapping disks such as Arp 244 (the Antennae). These pairs dominate the FIR luminosity function between $`10^{11}`$$`10^{12}L_{}`$ (Xu & Sulentic 1991).
There is clear evidence in the recent deep surveys carried out at: 1) optical/UV/NIR wavelengths by HST (Williams et al. 1996; Williams et al. 1998; Thompson et al. 1999), 2) MIR/FIR wavelengths by ISO (Elbaz et al. 1999; Puget et al. 1999), and 3) submm wavelengths by SCUBA (Blain et al. 1999), that much of the evolution in the history of star formation in the Universe is related to starbursts in closely interacting/merging systems. Therefore it is of great interest to understand how the enhanced star formation (starburst) is stimulated in these system.
The question of how starbursts are stimulated in interacting galaxies, leads directly to the more fundamental question of how the star formation is modulated in galaxies (Kennicutt 1989, 1998; Hunter et al. 1998; Wyse & Silk 1987; Dopita 1985). Since stars are formed from gas, an obvious necessary condition for a galaxy to be star-formation active is that it should contain significant gas. On the other hand, it appears that interacting galaxies have in general higher star formation rates (SFRs) not because they have more gas than isolated galaxies, but because the star formation rate per unit gas mass (the so called ’star formation efficiency’) is much higher (Young et al. 1986, 1989; Solomon & Sage 1988; Sanders et al. 1991). As pointed out by Kennicutt (1998) in his discussion of the ’Schmidt law’ (SFR$`\sigma _g^N`$, N$`1.4`$, where $`\sigma _g`$ is the gas surface density), the higher efficiency in starbursts in interacting galaxies may simply be a consequence of their much higher gas densities (Scoville et al. 1994; Solomon et al. 1997; Gao & Solomon 1999; Bryant & Scoville 1999). Unusually high gas density, which occurs almost exclusively in the nuclei of interacting galaxies and mergers, is certainly an interaction related phenomenon. Interaction induced gas inflow is predicted by simulations and is due either to the higher gas viscosity (caused by more frequent cloud-cloud collisions in interacting galaxies: Olson & Kwan 1991a; Struck 1997), or to the torque imposed on the gas by an interaction induced stellar-bar (Barnes & Hernquist 1996).
If enhanced star formation in interacting galaxies is due to higher gas density, which has a simple physical interpretation related to the growth rate of gravitational perturbation (Kennicutt 1998), two predictions can be made:
1. Interaction induced starbursts should concentrate in the nuclear region.
2. Starbursts in interacting systems of later stages (before the majority of the gas is converted to stars) should be stronger than those in early stages, because the later the stage, the more gas will sink into the nuclei.
As far as the ULIRGs are concerned, these two predictions seem to be in good agreement with the observations (see Sanders & Mirabel 1996 for a review). All ULIRGs appear to be in or close to the final stage of the merger process (e.g. Sanders et al. 1991; Murphy et a. 1996; Clements et al. 1996; Clements and Baker 1996; Mihos & Bothun 1998). The starbursts they harbor are primarily in the nuclear regions. Mihos & Bothun (1998) noted a trend between dynamical age and $`H\alpha `$ concentration in the ULIRGs they observed which is consistent with the physical scenario behind the two predictions.
Statistical studies of a large pair sample (Xu & Sulentic (1991) demonstrated that a subsample of closely interacting spiral-spiral pairs (hereafter CLO SS) with: a) separations less than a component diameter and b) showing optical signs of interaction, exhibit higher mean FIR–to–optical luminosity and FIR color ratios. This FIR excess exists not only with respect to isolated galaxies but relative to other pairs as well (see also Telesco et al. 1988; Jones & Stein 1989; Mazzarella et al. 1991). This result is again in line with prediction (2), because the small separations (normalized to the primary component diameter) of CLO SS pairs suggests that they may often be in a more advanced evolutionary state of a merger sequence compared to other SS pairs (Hwang et al. 1999). It should be noted, though, that: 1) since final coallescence usually requires several encounters between a pair of galaxies (Barnes & Hernquest 1996, 1998) and 2) since the orbital geometry of interacting systems is complicated, there is no one-to-one mapping between component separation and interaction stage. However, there is a clear tendency for the apoapsis of the orbit to decrease rapidly after the first encounter (Barnes & Hernquest 1998). This suggests that the later the stage, the more chance for the two galaxies to remain close together. On the other hand, the lack of one-to-one mapping between component separation and interaction stage for individual galaxy pairs (coupled with the projection effect) may be the reason for the absence of a monotonic dependence between starburst strength and separation, especially for pairs with wider separations (e.g. Fig.12 and Fig14 of Xu & Sulentic 1991; see also Keel 1993).
Contradicting prediction (1), some famous examples of the CLO SS pairs such as Arp 244 (Hummel & van del Hulst 1986; Vigroux et al. 1996; Mirabel et al. 1998) and Arp 299 (Gerhz et al. 1980; Hibbard & Yun 1998) show bright extranuclear starbursts in the overlap region, possibly due to interpenetrating collisions between components (Jog & Solomon 1992). Indeed, the simulation of Mihos et al. (1993), which modeled star formation in Arp 244 with a prescription similar to the ’Schmidt law’, failed to reproduce the overlap starburst. This motivates the following two questions:
1. Are overlap and other extranuclear starbursts common in CLO SS pairs?
2. If yes, then which starburst mode is the more important in CLO SS pairs, nuclear or overlap starbursts?
One must map the star formation distribution in these systems with spatial resolutions much better than IRAS in order to answer these questions. On the other hand, since dust extinction in interacting galaxies can be very large (Gehrz et al. 1983; Kunze et al. 1996), optical star formation mapping may not provide reliable results. These considerations motivated us to map a sample of CLO SS pairs with overlapping disks using two instruments (ISO CAM and ISO PHOT) on board of the Infrared Space Observatory (Kessler et al. 1996). ISO-SWS spectroscopic observations were also made for some pairs (pointing both to the nuclei and to the overlap regions). The imaging observations provided higher angular resolutions than IRAS and orders of magnitude higher sensitivity than KAO (e.g. Bushouse et al. 1998). This paper is the first in a series that presents and analyzes the ISO observations. We report the first results of ISO CAM and ISO SWS observations, comparing them with results from new ground based H$`\alpha `$ observations. ISO PHOT results and a more quantitative multiwavelength analysis will be presented in a following paper.
## 2 Observations
### 2.1 Sample Selection
The pair sample was selected from two catalogues:
1. Catalogue of Isolated Pairs of Galaxies in the Northern Sky (Dec.$`>3`$ deg) by Karachentsev (1972), hereafter KPG;
2. Atlas of Peculiar Galaxies (Dec. $`>27`$ deg) by Arp (1966).
The sample selection criteria were:
(1) spiral–spiral pairs with overlapping disks;
(2) pair component redshift difference $`\mathrm{\Delta }`$V$`<`$ 500 km s<sup>-1</sup>;
(3) pairs showing one of the three interaction morphology classes defined in the KPG (LIN=bridges and/or tails, ATM=common halo, or DIS=distortion in one or both components);
(4) major axis diameter of the primary component D$`>1^{}`$;
(5) a pair luminosity ratio $`L_{fir}/L_B>1`$, where $`L_{fir}`$ is the integrated IRAS FIR (82.5 $`\mu m`$) luminosity (Helou et al. 1988) and $`L_B`$ is the combined monochromatic luminosity ($`\nu L_\nu `$) at 4400Å estimated from the photographic magnitudes.
Criterion (1) restricts the sample to galaxy pairs with overlapping disks so that we can assess the role of overlap starbursts. This differs from the criterion used to define CLO SS pairs in Xu & Sulentic (1991) which included systems with projected separations up to one primary component diameter. Criterion (4) restricts the sample to relatively nearby pairs and insuring that the chosen sample could be resolved by ISO PHOT (beamwidth $`\theta `$$`45^{\prime \prime }`$ at 60 and 100$`\mu m`$). The final criterion selects SS pairs with FIR/B $`>`$ 2$`\times `$ the isolated spiral mean ($`L_{fir}/L_B=0.5`$, Xu & Sulentic 1991). This is also at least 2$`\sigma `$ above the mean $`L_{fir}/L_B`$ for normal spiral galaxies studied by Corbelli et al. (1991). Our sample is biased towards FIR enhanced SS pairs in order to minimize the chance of including false CLO SS pairs (wide pairs or accordant chance alignments). It is worthwhile to mention that although FIR/B is a good measure of starburst strength (Xu & De Zotti 1989) it is affected by several factors relatively unrelated to starburst activity including B-band extinction and diffuse FIR cirrus emission. Our original sample contained ten CLO SS pairs but two (Arp 263 and KPG 536) were dropped because of poor ISO visibility parameters. Basic properties for our sample of eight CLO SS pairs are presented in Table 1.
### 2.2 ISO CAM Observations
The 9.7$`\mu m`$ CAM images were obtained with the LW7 filter. They give information about the (9.7$`\mu m`$) silicon feature and facilitate comparisons with higher resolution ground-based (10$`\mu m`$) maps from the literature (e.g. Gehrz et al. 1983; Bushouse et al. 1998). The 15 $`\mu m`$ CAM images were obtained with the LW3 filter except for the very bright source Arp 299 (IRAS 12$`\mu m`$ flux of 3.7 Jy), where the narrower LW9 filter (centered at 15$`\mu m`$ was used in order to avoid saturation. These maps were expected to be significantly more sensitive than the 9.7 $`\mu m`$ observations because the filter had a broader bandpass and because the sources were expected to be brighter at longer wavelength (less extinction and more emission). With a resolution of $`10^{\prime \prime }`$ these CAM images reveal structures as small as a few kpc in our sample where typical distances are $`40Mpc`$ (H<sub>0</sub>=75 km/s Mpc<sup>-1</sup>). KPG 347 was a GTO target with the LW3 filter (Boselli et al. 1998). Table 2 summarizes details of the ISO CAM observations. Basic data reduction employed the CAM Interactive Analysis (CIA) software<sup>2</sup><sup>2</sup>2CIA is a joint development by ESA Astrophysics Division and the ISO CAM Consortium led by the ISO CAM PI, C. Cesarsky, Direction des Sciences de la Matiere, C.E.A., France.. Special attention was given to the correction of transient effects using code based on the “Fouks and Schubert algorithm” developed by A. Abergel, A. Coulais, and H. Wozniak (Saclay) and provided to us by M. Sauvage (1999, private communication). This code did a significantly better job than the standard CIA algorithms.
### 2.3 ISO SWS Observations
We made spectroscopic observations of four emission lines (Br<sub>β</sub> \- 2.63$`\mu m`$; Br<sub>α</sub> \- 4.05$`\mu m`$; \[Ne II\] - 12.81$`\mu m`$ and H<sub>2</sub> S(1) - 17.03$`\mu m`$) using ISO SWS in the grating scan mode (AOT SWS02). Observations were obtained for three pairs in our sample (marked by stars in Table 1) with three positions observed (galaxy nuclei plus overlap region) in Arp 81 and 157 and two positions (both nuclei) in Arp 278 (Table 3.1). Parameters of these observations are given in Table 3.2. The purpose of these observations was to measure star-formation rates and extinctions in different locations in the pairs. ISO SWS reductions were performed using the SWS Interactive Analysis package (IA3) developed by the international ISO SWS Consortium.
### 2.4 H$`\alpha `$ Observations
H$`\alpha `$ and R-band observations were carried out at: 1) Palomar in February 1996 using a Tek 1024$`\times 1024`$ CCD mounted on the 1.5m telescope giving a scale of 0$`\stackrel{}{\mathrm{.}}`$62 pixel<sup>-1</sup> and 2) Calar Alto in May 1996 and August 1997 using a 2048$`\times 2048`$ CCD on the 2.2 m telescope giving a scale of 0$`\stackrel{}{\mathrm{.}}`$33 pixel<sup>-1</sup> and/or 0$`\stackrel{}{\mathrm{.}}`$53 pixel<sup>-1</sup>. Palomar narrow-band filters centered on redshifted $`H\alpha `$ (+ \[N II\]6548,6583) with FWHM $`100`$ Å and adjacent continuum broad-band filter (R-band) were used. Calar Alto narrow-band (FWHM 50 – 80 Å centered on H$`\alpha `$) and R-band filters were also employed. The standard stars HD 84937, Landolt 104-334 (Palomar), BD +28 4211 and BD +33 2642 Calar Alto were used for photometric calibration. Standard IRAF data reduction procedures were used to reduce this data. Continuum subtracted H$`\alpha `$ images were produced by scaling and subtracting the $`R`$ frames using field stars to match the frames. Subtraction was carried out interactively until stellar residuals were minimized. Matching R band images were not obtained for all Calar Alto targets. In those cases continuum-free $`H\alpha `$ and \[N II\] emission line images were produced by subtracting the off-centered adjacent narrow-band image from the one centered on redshifted $`H\alpha `$ narrow-band maps, using about half dozen stars in the field for normalization. Aperture photometry was done using both elliptical and circular apertures.
## 3 Results
### 3.1 MIR Continuum v.s. H$`\alpha `$ Emission
#### 3.1.1 Individual Pair Properties
Figures 1–4 present MIR and H$`\alpha `$ images for our sample. Three images are presented for each source. TOP PANEL: ISO CAM 15$`\mu m`$ contours (9.7$`\mu m`$ for KPG 347) are plotted over an optical image from the Digitized Sky Survey<sup>3</sup><sup>3</sup>3The Digitized Sky Survey was produced at the Space Telescope Science Institute under U.S. Government grant NAG W-2166. The images of these surveys are based on photographic data obtained using the Oschin Schmidt Telescope on Palomar Mountain and the UK Schmidt Telescope. The plates were processed into the present compressed digital form with the permission of these institutions. The complete acknowledgement can be found at http://archive.stsci.edu/dss/dss\_acknowledgements.html.. Contour levels are 2<sup>n</sup> (n=1,2,3,…) times the rms noise ($`\sigma _{15\mu m}`$ or $`\sigma _{9.7\mu m}`$ for KPG 347) as given in Table 2. MIDDLE PANEL: A log grayscale image of the 15$`\mu m`$ to $`9.7\mu m`$ MIR color ratio with ISO CAM 15$`\mu m`$ contours superposed (9.7$`\mu m`$ contours for KPG347). Contour levels are 3+3<sup>n</sup> (n=0,1,2,3,…) times the rms noise of the corresponding ISO CAM image, with the range of the grayscale image varied from source to source in order to maximize the dynamic range. BOTTOM PANEL: H$`\alpha `$ contours plotted over the R-band image. Contour levels are different for each source because of the different rms noise levels. 15 and 9.7$`\mu m`$ flux densities are given for the sources in Table 4. The $`f_{15\mu m}/f_{9.7\mu m}`$ ratio depends on both small grain emission and silicate absorption in the 9.7$`\mu m`$ band. Detailed modeling of this ratio will be presented in a future paper.
Detailed notes
Arp 299 = NGC 3690/IC 694<sup>4</sup><sup>4</sup>4Confusion exists over the name of the eastern component in Arp 299. We adopt IC 694 in this paper. See Appendix in Hibbard & Yun 1998 for details. (Figures 1abc): A pair of gas-rich galaxies IC 694 (east) and NGC 3690 (west) regarded as a local example of a merger in progress (Gehrz et al. 1983; Telesco et al. 1985; Wynn-Williams et al. 1991; Casoli et al. 1999; Hibbard & Yun 1999; Gallais et al. 1999). Arp 299 is the most IR luminous system in our sample with $`L_{fir}=2.7\times 10^{11}L_{}`$. Remarkably with $`f_{60\mu m}=108.9`$ Jy, it is one of the brightest (60$`\mu m`$) extragalactic point sources in the IRAS point source catalog (Soifer et al. 1987). The component disks are in contact suggesting a relatively advanced stage of coallescence, although the nuclei are well separated and are resolved by ISO CAM. The extranuclear starburst (Source C from Gehrz et al. 1983) is about $`8^{\prime \prime }`$ north of the NGC 3690 nucleus (Source B in Gehrz et al. 1983, the western galaxy) is only marginally resolved (seen as a plateau on the contour plot) from the latter by ISO CAM. Gallais et al. (1999) report ISO CAM CVF observations of Arp 299 which cover $`48^{\prime \prime }\times 48^{\prime \prime }`$ with pixel size of $`1.5^{\prime \prime }\times 1.5^{\prime \prime }`$. This can be compared to our observations (Table 2) which have a $`3.8^{}\times 3.8^{}`$ field of view and a $`3^{\prime \prime }\times 3^{\prime \prime }`$ pixel size. Their LW3 (15$`\mu m`$) map clearly resolves the overlap starburst (Source C) from the nucleus of NGC 3690.
The $`f_{15\mu m}/f_{9.7\mu m}`$ ratio map is generally smooth with values close to the ratio of the integrated fluxes ($`f_{15\mu m}/f_{9.7\mu m}=2`$: Table 4) across most of the disks. The nucleus of IC 694 (Source A in Gehrz et al. 1983) and a few locations in the outer disk show ratios as high as $`f_{15\mu m}/f_{9.7\mu m}>10`$. While high values for the ratio may have large uncertainty in the outer regions, the high ratio in the IC 694 nucleus implies a high silicate absorption which may severely depress the 9.7$`\mu m`$ flux (see also Gehrz et al. 1983). This is consistent with results of CO imaging (Sargent & Scoville 1991; Aalto et al. 1997; Casoli et al. 1999) showing that the nucleus of IC 694 has the highest surface density of molecular gas in the Arp 299 system and comparable to values found in ULIRGs (Downes & Solomon 1998). The inner part of the MIR emission region shows good correspondence with the H$`\alpha `$ emission allowing for differences in spatial resolution (Fig. 1, Panel c). This indicates that: 1) most of the MIR emission is due to dust associated with star formation regions, as suggested by Sauvage et al. (1996) in an ISO CAM study of M51, and 2) the H$`\alpha `$ extinction is rather smooth, consistent with the smooth $`f_{15\mu m}/f_{9.7\mu m}`$ ratio image. The outer envelope of the MIR emission is elongated in along PA$``$45 which is also the ISO CAM scan direction. The elongation is contrary to that of the H$`\alpha `$ emission suggesting that the ISO CAM elongation is likely to be an artifact of the transient behavior in the ISO CAM detectors (Cesarsky et al. 1996). Careful transient corrections were included in our reduction of this data, but the very high surface brightness (up to $`1`$ Jy/beam) in the central region of Arp 299 produce strong effects that are difficult or impossible to remove completely.
Arp 244 = NGC 4038/9 (Figures 1def): The “Antennae” are regarded as another local example of a galaxy merger (Rubin et al. 1970; Toomre & Toomre 1972; Schweizer 1978; Hummel & van der Hulst 1986; Vigroux et al. 1996; Mirabel et al. 1998; Evans et al. 1998; Whitmore et al. 1999). ISO CAM images at 6.7 and 15$`\mu m`$ have been published (Vigroux et al. 1996; Mirabel et al. 1998). Our 9.7 and 15$`\mu m`$ maps cover a larger area ($`6^{}\times 6^{}`$). The brightest region in this system at MIR wavelengths involves an extranuclear starburst (Source A) about 15<sup>′′</sup> northeast of the NGC 4039 nucleus (c.f. Vigroux. et al. 1996; Mirabel et al. 1998). There is an H$`\alpha `$ peak at Source A, but much fainter than the emission from the NGC 4039 nucleus. Apparently much of the optical emission associated with Source A is extinguished by dust. More intriguingly, the FIR (60$`\mu m`$, 100$`\mu m`$ and 160$`\mu m`$) KAO maps (Evans et al. 1998; Bushouse et al. 1998) show that the FIR peak is displaced north of Source A, at the position of a dark patch in the optical images (e.g. the HST WFPC2 images by Whitmore et al. 1999). New SCUBA maps (Haas et al. 2000) at 450 and 850$`\mu m`$ (15<sup>′′</sup> resolution) reveal large amounts ($`10^{67}`$ M) of cold dust ($`<20`$K) in the overlap region. There are corresponding radio continuum peaks at both the MIR and FIR peaks (Hummel & van der Hulst 1986). The CO observations (Stanford et al. 1990; Gao et al. 1998; Gruendl et al. 1998) demonstrate that most of the molecular gas in this system is extended throughout the overlap region. The \[C II\] 158$`\mu m`$ emission also peaks at the dark patch north of Source A (Nikola et al. 1998). All of these observations indicate that much of the star formation activity in Arp 244 is hidden by dust. Our MIR maps also show some weak emission at the beginning of the southern tail starts and a similar feature is observed in the radio continuum (Hummel & van der Hulst 1986).
The $`f_{15\mu m}/f_{9.7\mu m}`$ ratio map also peaks around Source A with a value of $`(f_{15\mu m}/f_{9.7\mu m})3`$. Given that most of the molecular gas is found near this region (Gao et al. 1998), and very strong extinction (A$`{}_{V}{}^{}70`$) is found from an SWS study of MIR lines in the same region (Kunze et al. 1996), the high $`(f_{15\mu m}/f_{9.7\mu m})`$ level is likely due to significant silicate absorption.
In a separate ISO CAM observation (Arp 244-02, Table 2), we obtained ISO CAM LW3 images further down along the southern tail, including the end of the tail where a dwarf galaxy was found (Schweizer 1978; Mirabel et al. 1992). No MIR emission is detected in these tidal features.
Arp 157 = NGC 520a/b (Figures 2abc): This is a very complex system apparently involving two (colliding) disk galaxies, one oriented southeast-northwest (NGC 520a) and another oriented east-west (NGC 520b) (Stanford 1991; Stanford & Balcells 1990, 1991; Bernlöhr 1993ab). The collision center appears to lie near the nucleus of NGC 520b where where the MIR emission peaks. The CO emission also concentrates at this position (Sanders et al 1988b). The H I gas shows the the kinematic signature of a rotating disk with the same orientation as NGC 520a, while the rotation center is clearly at the nucleus of NGC 520b (Hibbard & van Gorkom 1996). An interesting possibility is that the H I originally belonged to the former but has been captured by the latter (which might have an order of magnitude more mass than the former; Bernlöhr 1993b). Note also that both disks show rotation axes nearly parallel to the axis of the orbital motion (Bernlöhe 1993b), making the hypothesis for migration of the H I gas more reasonable. There is a second, much weaker peak in the MIR emission associated with the nucleus of NGC 520a, where new millimeter synthesis observations failed to detect any CO emission (Hibbard et al. 2000, private communication). The H$`\alpha `$ map (see also Hibbard & van Gorkom 1996; and Young et al. 1988) shows very different morphology from the MIR maps which is most likely due to dust extinction of the former. Near the nucleus of NGC 520b, heavy dust lanes are visible in optical images. The $`(f_{15\mu m}/f_{9.7\mu m})`$ ratio, as a rough indicator of silicate absorption, also peaks there. A NIR K-band image of Arp 157 (Bushouse & Werner 1990) shows similar morphology to the MIR images, indicating that the difference between the MIR and optical/H$`\alpha `$ morphologies is not due to different angular resolutions. This is consistent with results of Bernlöhr (1993b) whose model predicts that both NGC 520a and NGC 520b have been undergoing starbursts, with the starburst associated with NGC 520a about 2-3 $`10^8`$ years older and a factor of $`8`$ fainter than the starburst associated with NGC 520b.
Arp 81 = NGC 6621/2 (Figures 2def): This system was included by Toomre (1978), along with Arp 157 and 244, in proposed sequence of mergers. A recent HST WFPC B-band image (Keel, private communication) shows that NGC 6622 (the southern galaxy) is likely to be a high inclination S0 galaxy. H$`\alpha `$ and MIR emission from the center of NGC 6622 indicates that some star formation is occurring there but at much lower level than in the nuclear region of NGC6621. The HST image shows a star formation region between the two galaxies along with gas/dust features extend from NGC 6621 and near the center of NGC 6622. There starburst in the overlap region is seen in both H$`\alpha `$ and MIR maps. The tidal tail is a marginal (2$`\sigma `$) detection on ISO CAM images. The nucleus of NGC 6622 and the overlap starburst show rather low $`f_{15\mu m}/f_{9.7\mu m}`$ ratios ($`1)`$, while the nucleus of NGC 6621 shows a higher level $`f_{15\mu m}/f_{9.7\mu m}`$ ($`3`$). A KAO 100$`\mu m`$ observation (Bushouse et al. 1998) detects only the NGC 6621 nucleus. Given the large beam ($`50^{\prime \prime }`$) it is possible that the overlap starburst may have contributed significantly to the flux. A 1.426 GHz VLA map (Condon et al. 1996) peaks on the NGC 6621 nucleus with an extension that may be due to the overlap starburst or even emission associated with NGC 6622.
Arp 278 = NGC 7253a/b (Figures 3abc): This system looks like a more edge-on version of the famous “Taffy Galaxies” (UGC 12914/5) (Condon et al. 1993; Jarrett et al. 1999). Both galaxies are highly inclined, and much diffuse emission is observed between the two disks (as in case of “Taffy”). Arp 278 shows evidence of starbursts in both galaxies (Bernlöhr 1993a) with a possible time delay of 10-70 Myr between the components (the larger galaxy shows the younger burst). Both galaxies are detected in the radio continuum (Condon et al. 1996). MIR emission from the northwestern galaxy (NGC 7253a, the primary) peaks at the nucleus, while the peak is offset (toward the primary) in the other component. From deep optical image, the two galaxies appear to be in contact near the nucleus of the southeastern galaxy (NGC 7253b), where several H$`\alpha `$ knots exist (Fig.3c). Given the clear link between the star formation regions in the two galaxies as shown in the H$`\alpha `$ image, at least some of the emission in those H$`\alpha `$ knots is likely due to overlap starburst. In connection with the Taffy Galaxy analogy (Condon et al. 1993; Jarrett et al. 1999), the components of Arp 278 may have recently undergone a face-on collision. The plumes (i.e. the diffuse emission between the two disks) may be viewed as debris from that collision. Keel (1993) suggests that NGC 7253a is undergoing a direct (prograde) encounter. Given that NGC 7253a also dominates the MIR emission of the pair (Table 4), it is likely to be the more disturbed component that is the source of most of the debris. The $`(f_{15\mu m}/f_{9.7\mu m})`$ ratio map peaks at the nucleus of NGC 7253a although, in general, that ratio is not high ($`<2`$) indicating that the overall dust extinction is less severe than in Arp 244 and 299.
Arp 276 = NGC 935/IC 1801 (Figures 3def): Neither galaxy in this pair shows evidence for strong perturbation. Both the MIR and H$`\alpha `$ emissions extend over the galaxy disks. Most of the H$`\alpha `$ emission in the primary comes from H II regions in the spiral structure while the nucleus shows little or no emission. This $`H\alpha `$ “hole” might be due to over-subtraction of the continuum or, alternatively, may indicate a high extinction in the nuclear region where the MIR emission peaks. A foreground star is projected between the two component nuclei (but not on the overlap region) making H$`\alpha `$ measures unreliable there. H$`\alpha `$ emission from the smaller galaxy appears more concentrated. with the nucleus contributing significantly to the total flux. Two H II regions on either side of the nucleus and perpendicular to the major axis, contribute much of the remaining H$`\alpha `$ flux. No localized H$`\alpha `$ or MIR feature is detected in the overlap region. Given the differences in angular resolution, the $`f_{15\mu m}`$ and H$`\alpha `$ maps shows similar morphology (i.e. the emission in the primary is extended and that in the secondary is more centrally concentrated. The $`(f_{15\mu m}/f_{9.7\mu m})`$ ratio map is smooth with little dependence on MIR surface brightness. It is unlikely that dust extinction is high in this pair and the temperature of small grains, which are the major contributors to $`f_{9.7\mu m}`$ and $`f_{15\mu m}`$, does not depend sensitively on the radiation intensity (Désert et al. 1990).
KPG 347 = NGC 4567/8 (Figures 4abc): This pair in the Virgo cluster is known as the “Butterfly Galaxies”. The 9.7$`\mu m`$ contours are plotted in Fig. 4. The $`(f_{15\mu m}/f_{9.7\mu m})`$ ratio image used a 15$`\mu m`$ map from Boselli et al (1998). The two galaxies in this system also show little morphological disturbance. A VLA H I map of the system (Cayatte et al. 1994) shows that the H I distribution of the component disks to be reasonably intact. Star formation is widespread in both galaxies and unlike the CLO SS pairs discussed above, the $`(f_{15\mu m}/f_{9.7\mu m})`$ ratio shows dips rather than peaks at the positions of the nuclei with high ratio plateaus in the outer disk. The reason for the variations may be different in this pair, related perhaps to heating of small grains that are responsible for MIR continuum emission. In more active pairs (e.g. Arp 299) silicate absorption is the more likely to be the major cause of variations in the flux ratio.
KPG 426 = UGC 9376a/b (Figures 4def): Deep optical image show that the component of this pair are embedded in a common envelope. Both the MIR and H$`\alpha `$ morphologies look undistorted (symmetric). Most of the star formation in the northern galaxy is concentrated in the nucleus while that in the southern galaxy shows a nuclear and ring structure. Detailed comparison between the morphology and simulations (Toomre and Toomre 1972) suggest that the northern galaxy has undergone a retrograde encounter, while the southern galaxy a prograde but high inclination ($`>30^{}`$) encounter. Similar to KPG 347, the $`(f_{15\mu m}/f_{9.7\mu m})`$ ratio shows dips rather than peaks in the nuclei.
#### 3.1.2 Collective Properties
The eight systems studied in this paper provide a small but representative sample of CLO SS pairs with enhanced FIR emission. The MIR and H$`\alpha `$ observations provide direct information about the distribution of star formation, which may impose new constraints on models for interaction induced star bursts. On the other hand, the uncertainties associated with small sample statistics make the following inferences suggestive rather than conclusive.
We divide our sample into three subgroups based upon the star formation morphology:
I. Advanced mergers (Arp 157, 244 and 299): The nuclei in these systems are still well separated (or they would not have been cataloged as pairs by Arp and/or Karachentsev) however both the stellar and gaseous disks are “entangled” with each other. These pairs have the least separations ($`<8kpc`$) in our sample. They also show line of sight velocity differences $`\mathrm{\Delta }`$V$``$60 km s<sup>-1</sup>, consistent with a significant dissipation of orbital angular momentum. Simulations (e.g. Mihos et al. 1993; Hibbard & Yun 1999) suggest that these pairs will merge in a few times of 10<sup>8</sup> years.
II. Severely disturbed systems (Arp 81 and 278): The two disks remain separated but severe tidal distortions are present. Line of sight $`\mathrm{\Delta }`$V$``$200 km s<sup>-1</sup> for these pairs.
III. Less disturbed systems (Arp 276, KPG 347 and 426): Galaxies in these pairs show reasonably normal morphologies in both MIR and optical/H$`\alpha `$ images. Although all three pairs are classified ATM or DIS by Karachentsev (1972), our R-band images show nearly normal (Arp 276 and KPG 347) or only weakly distorted (KPG 426) morphologies in contrast to subgroups 1) and 2). $`\mathrm{\Delta }`$V ranges from 50-150 km s<sup>-1</sup> for these binary systems.
The following trends can be drawn from our observations:
All five CLO pairs (Arp 81, 157, 244, 278 and 299) in the first two subgroups show localized MIR and H$`\alpha `$ emission enhancements in the overlap region. Only in Arp 244 is the overlap region starburst more luminous than the nuclei burst while in Arp 81 and 299 the overlap region starbursts are significantly fainter than the most active nucleus in the system (in Arp 157 and 278 the overlap region is close to one of the nuclei). The inference is that star formation induced by the hydrodynamic collisions of gaseous disks is a common phenomenon in closely interacting galaxy pairs. However, it is apparently seldom the dominant source of starburst emission in these systems. The nuclear starburst, driven presumably by the gravitational tidal force which induces infall of large quantities of gas into galaxy nuclei, is perhaps the major mechanism for the enhanced star formation in these galaxy pairs.
CLO pairs in the third subgroup show star formation that is more evenly distributed in the component galaxies. This includes both disk and nuclear activity as is observed in isolated samples at a lower average intensity level. H II emission regions in these components closely follow the spiral arms, similar to the distribution of star formation in M51 (Sauvage et al. 1996). No overlap starbursts are observed in these systems. With a caution about small sample statistics, they also show relatively low FIR/B ratios ($`<L_{fir}/L_B>=1.30\pm 0.25`$) and cool FIR colors ($`<f_{60\mu m}/f_{100\mu m}>=0.36\pm 0.02`$) compared to pairs in the first two subgroups which show $`<L_{fir}/L_B>=4.93\pm 1.37`$ and $`<f_{60\mu m}/f_{100\mu m}>=0.65\pm 0.09`$.
### 3.2 MIR Line Emission
The in-orbit performance of SWS over the observed wavelength regions was much lower than expected at the time our observations were planned. No Br$`\gamma `$ or Br$`\alpha `$ features were detected while \[N II\] 12.81$`\mu m`$ and H<sub>2</sub> S(1) 17.03$`\mu m`$ were only detected in Arp 157. The ISO SWS results are presented in Table 6. The \[N II\] 12.81$`\mu m`$ line and H<sub>2</sub> S(1) 17.03$`\mu m`$ lines were only detected in the nucleus of NGC 520b (Fig. 2ab). A pointing at the overlap region in this pair also detected the \[N II\] 12.81$`\mu m`$ line (Fig. 2c) but this may be contaminated by emission from the NGC 520b nucleus because of the large beam. A Br$`\gamma `$ line (EW= 17.0$`\pm 0.5`$ Å) was previously reported by Vanzi et al. (1998).
SWS observations of Arp 244 and 299 exist in the literature. The MIR line ratios, assuming that gas and dust are well mixed, imply rather high dust extinctions: A$`{}_{V}{}^{}20`$ for Arp 299 and A$`{}_{V}{}^{}80`$ for Arp 244 (Kuntz et al. 1996; Genzel et al. 1998). It is interesting to compare our new data for Arp 157 with SWS observations of related objects. The \[N II\] 12.81$`\mu m`$ flux for NGC 520b is comparable to the values detected in the SWS survey of ULIRGs as well as template starburst and AGN sources (Genzel et al. 1998). A more physically meaningful comparison between sources can be made using the ratio of the \[N II\] 12.81$`\mu m`$ line to total FIR flux estimated from the IRAS observations using $`FIR=1.26\times 10^{14}(2.58f_{60\mu m}+f_{100\mu m})[Wm^2]`$ (e.g., Sanders & Mirabel 1996). The ratio of $`f([NII])/f(FIR)=0.001`$ for NGC 520b similar to the values observed in NGC 6240 and Arp 244, and about half value observed for the nuclei of Arp 299 in the more sensitive observations reported by Genzel et al. (1998). This is another direct indication that the physical conditions of active star formation in NGC 520b are similar to those observed in Arp 244 and 299, the other two sources in our first subgroup.
## 4 Discussion
### 4.1 Overlap Starbursts
What is special about the overlap region? It is the actual interface region of a disk-disk collision. An overlap starburst could then be triggered by direct collisions between giant molecular clouds (GMCs) (Noguchi 1991). At the same time this mechanism may not be very efficient because the galactic disk filling factor of GMCs is very low ($`>0.01`$). Jog & Solomon (1992) proposed a more efficient model where collisions between H I clouds lead to the formation of a hot ionized, high-pressure remnant gas. The over-pressure due to this hot gas causes a radiative shock compression of the outer layers of existing GMCs in the overlap region. These layers become gravitationally unstable and trigger a burst of massive star formation in the initially barely stable GMCs. This model can be tested with X-ray observations of the hot remnant gas. High resolution ROSAT observations of Arp 244 (Fabbiano et al. 1997) revealed X-ray emission in the overlap region. However, given the morphology of the X-ray emission (localized, point-source like), it is more likely due to supernova remnants associated with the on-going starburst rather than a hot remnant gas which would be more diffuse. The X-ray emission and star formation associated with the ongoing collision in Stephan’s Quintet (Pietsch et al. 1997; Xu et al. 1999; Sulentic et al. 2000) may be closer to this situation. Future higher resolution AXAF observations could provide more definite results.
If cloud-cloud collisions can indeed trigger star formation (Scoville et al. 1986; Olson & Kwan 1991ab) then this may be a mechanism (in addition to the gas-density dependence of star formation rate) for interaction induced star formation enhancements in general. Simulations (Olson & Kwan 1991ab; Noguchi & Ishibashi 1986; Noguchi 1988, 1991) have shown that cloud-cloud collisions are significantly enhanced throughout the disk due to the orbit-crossing of gas clouds triggered by gravitational perturbations. This may provide an interpretation for more widely distributed star formation, that shows more moderate enhancement, in “less disturbed” CLO pairs.
It is interesting to compare the overlap starbursts in Arp 299 (Source C-C’ in Gehrz et al. 1983) and Arp 244 (Source A in Vigroux et al. 1996). Table 6 gives the MIR fluxes of overlap starbursts and galactic nuclei for these two systems. The overlap starburst in Arp 244 is significantly brighter than both nuclei combined! On the other hand, the overlap starburst in Arp 299 is much fainter than either of the component nuclei. One possible reason for this difference is that the two pairs are in different stages of the merging process. Arp 299 may be in a later merging stage where more gas has fallen into the galactic nuclei powering intensive starbursts (especially in IC 694; Casoli et al. 1999). Interaction kinematics may also play a role in Arp 299 where IC 694 is undergoing a retrograde encounter (Augarde & Lequeux 1985; Hibbard & Yun 1999). This apparently allows it to retain most of its gas which is then available to the nucleus as fuel for the violent starburst. On the other hand, both disks in Arp 244 are undergoing prograde encounters (Toomre & Toomre 1972; Mihos et al. 1993), and the gas disks have suffered severe disruption (Barnes & Hernquist 1996, 1998) with much of the gas moved away from the central regions (van der Hulst 1979; Gao et al. 1998; Grundl et al. 1998). The highest gas concentration is, in fact, found in the overlap region (Stanford et al. 1990; Gao et al. 1998; Grundl et al. 1998). It is possible that such a gas distribution is unstable and transient (retardation of infall) which might be the reason why such bright overlap starbursts are rare.
### 4.2 ’Less Disturbed Systems’
Our third CLO SS subgroup shows a modest FIR enhancement but less structural distortion and more extended star formation. These properties are consistent with the hypothesis that they are in an earlier interaction stage than the first two subgroups. The absence of overlap starbursts in these may indicate that a disk-disk collision has not yet occurred. Although no firm conclusion can be reached about their physical separation, the calculation of Xu & Sulentic (1991, Appendix A) shows that components in most close pairs are physically proximate (physical separation $``$ diameter of the primary). Examination of component morphologies in KPG 426, the pair with components most clearly separated in our sample, and consideration of simulations (e.g. Toomre & Toomre 1972), suggest that the components are well separated. There is, at the same time, evidence that ’less disturbed systems’ may have unfavorable interaction geometries. Two out of three pairs in subgroup three may be undergoing either retrograde or highly inclination encounters (the situation for Arp 276 is less clear).
### 4.3 Variations in the Strength of Star Formation Activity
We suggest that the following sequence of sources of increasing star formation efficiency: isolated galaxies – wide pairs – closely interacting pairs – ULIRGs, may be viewed as an evolutionary sequence for mergers. Results in this paper show further that, within the population of closely interacting pairs (CLO SS), the ’less disturbed systems’ may be in earlier stages of interaction and have less enhanced star formation activity compared to ’severely disturbed systems’ and ’advanced mergers’.
At the same time, CLO SS pairs show a large scatter in star formation indicators such as $`L_{fir}/L_B`$ and $`f_{60\mu m}/f_{100\mu m}`$ ratios (Xu & Sulentic 1991). Even among sources from the same subgroups in this study there is a wide range in star formation activity indicators (such as FIR/B). Arp 244, one of the ’advance mergers’, shows FIR/B lower than the ’less disturbed systems’ and one of the lowest $`L_{fir}/M_{gas}`$ ratios in interacting systems (Gao et al. 1998). One explanation for the large scatter in star formation indicators involves variations in galaxy gas content. Galaxies with less gas have less fuel for starbursts which is consistent with previous results (Sulentic 1989; Xu & Sulentic 1991) which found that SS pairs show higher FIR emission than normal galaxies but little FIR emission from pairs of early type galaxies (E/S0) which are gas poor. Differences in gas content between CLO SS pairs with strong and weak star formation enhancement will be an interesting area of study for future H I and CO surveys of galaxy pairs. A surprising lack of H I gas depletion was found in a statistical study of an FIR enhanced sample of E+S pairs (Zasov & Sulentic 1994).
Even within a sample of gas rich pairs significant differences in their star formation rates are found. This is true even for galaxies with the same gas surface density where the “star formation efficiency” can differ by as much as an order of magnitude (see, e.g., Fig. 2 of Solomon & Sage 1988). Much of this scatter may reflect the episodic nature of starbursts. Several physical mechanisms may contribute to this episodic behavior. Feed-back from massive star formation (such as supernovae explosions and stellar winds) may quench a burst after a few 10<sup>7</sup> yr (Krügel & Tutukov 1993). This effectively breaks the interaction induced star formation into pulses which is confirmed by observations: All starburst durations derived from observation (Rieke et al. 1980; Gehrz et al. 1983; Bernlöhr 1993a) are on the order of a few times 10<sup>7</sup> years even though interactions typically last for several 10<sup>8</sup> years (e.g. Hibbard & Yun 1999). When an interacting galaxy is observed in the ’off’ stage, only a ’post-starburst’ is seen (examples can be found in Bernlöhr 1993a). A ’post-starburst’, with most of the OB stars already gone, will have a much reduced effect on star formation indicators (e.g. FIR and H$`\alpha `$ emissions).
As demonstrated by the simulations of Noguchi (1991), the periodic orbital motion which swings the two galaxies back and forth relative to each other several times before they eventually merge, will induce sharp peaks (corresponding to the passages of periapses) in the cloud-cloud collision rate which in turn may also cause pulsational star formation (see also Olson & Kwan 1991ab).
The episodic nature of starbursts could be a key to understanding why only a single component in many pairs shows enhanced star formation (Joseph et al. 1984). Since the star formation in the members of a pair may well be unsynchronized (Bernlöhr 1993a), there is a good chance for one of them to be in the ’on’ phase while the other is ’off’. At the same time there is a significant probability that both components will be ’on’ (Lutz 1992; Surace et al. 1993). This suggests that the duration of the ’on’ and ’off’ phases must be comparable (a few 10<sup>7</sup> yrs).
Another reason for the scatter in star formation activity could be the interaction geometry. As shown in the pioneering work of Toomre & Toomre (1972), retrograde and high inclination encounters cause much less distortion than direct (prograde) and low inclination encounters. If, as suggested, the ’less disturbed systems’ have unfavorable interaction geometries the lower star formation activity might result from this rather than from an earlier merger evolutionary stage. However the evidence that the galaxies in Arp 244 (lowest FIR/B in our sample) show prograde rotation (van der Hulst 1979) while IC 694 in Arp 299 (highest FIR/B in our sample) shows retrograde rotation (Hibbard & Yun 1999), argues against this possibility (see also Keel 1993; Lutz 1992). This is not meant to imply that orbital geometry plays no role in the star formation enhancement process. It suggests only that this role is likely to be complex. For example an unfavorable orbital geometry may preserve most of the original gas in a galaxy until a very late stage in the merger process, similar to the role played by a massive bulge in the simulation of Mihos & Hernquist (1996; see also Evans et al. 2000), making a ULIRG-like starburst more possible (see, e.g., the simulations of Noguchi 1991).
### 4.4 Galaxy Pairs and Mergers
Throughout this discussion we have adopted a scenario where galaxy pairs merge after a few close encounters, as implied by models (e.g. Barnes & Hernquist 1996, 1998). Recent studies on the cosmic evolution of merger rate using HST data (Wu & Keel 1998; Le Fèvre et al. 2000) also hint that the time scale of mergers is much shorter than the Hubble time. However, pair merger time scales as short as $`1`$ Gyr require an explanation for the large number of isolated binary galaxies ($`10\%`$ of all galaxies, Xu & Sulentic 1991), the majority of which must have been gravitationally bound systems for $``$10 Gyr (Chatterjee 1987). They must also account for the rarity of candidate (early-type) merger postcursors found in the the same environments as these large pair (and compact group) populations (Sulentic & Rabaca 1994).
A possible solution for this dilemma is that today’s galaxy pairs are evolved from galaxy groups. Through coalescence, within the context of hierarchical galaxy formation (White 1997), two giant galaxies observed now may be products of a long history of mergers/accretions of smaller galaxies which were formed as a bound system. One possibility is that some mixed pairs (E+S)) represent the last stages in the coallescence of a compact group (Rampazzo & Sulentic 1992). Such a picture may apply to the two giant galaxies in the Local Group (Milky Way and M31) which may represent the early stages in the formation of a CLO SS pair (e.g., Irion 2000). There is evidence that the Milky Way is still growing by accretion of smaller satellites, with the Sagittarius dwarf galaxy being the current “meal” (Buser 2000). Given the recent discoveries closely linking galaxy evolution and interactions (Williams et al. 1996; Williams et al. 1998; Thompson et al. 1999; Elbaz et al. 1999; Puget et al. 1999; Blain et al. 1999), any theory for the formation/evolution of galaxy pairs must play an integral part in the formation and evolution of galaxies in general.
## 5 Conclusion
We present MIR imaging and spectroscopic observations for a well defined sample of eight closely interacting pairs of spiral galaxies with overlapping disks and enhanced FIR emission. Our goal was to study the star formation distribution in these pairs with special emphasis on the importance of overlap starbursts. We identified three possible subgroups in our sample according to star formation morphology:
(1) advanced mergers;
(2) severely disturbed systems;
(3) less disturbed systems.
Overlap region starbursts are detected in all of the five pairs of subgroups (1) and (2), suggesting that they are a common property of colliding systems. On the other hand, except for Arp 244, the ’overlap starburst’ is less intense than the nuclear starbursts in such pairs. Star formation pairs of subgroup (3) is often more widely distributed in the disks of both components with no evidence for overlap starbursts. These systems also show a smaller FIR enhancement, implying weaker star formation and suggesting that they may be in earlier interaction stages where direct disk-disk collisions have not yet occurred. Only one pair (Arp 157) is detected in our ISO SWS observations. The F(\[Ne II 12.81$`\mu m`$\])/F(FIR) ratio in the NGC 520b component of Arp 157 is $`0.1`$, comparable to values for other luminous infrared galaxies.
C.X. thanks Marc Sauvage for valuable assistance in reducing the ISO CAM data, and Eckard Sturm and Alberto Noriega-Crespo for helping reducing the SWS data. This work was supported by NASA grant for ISO data analysis. This research has made use of the NASA/IPAC Extragalactic Database (NED) which is operated by the Jet Propulsion Laboratory, California Institute of Technology, under contract with the National Aeronautics and Space Administration. C.X., Y.G., J.M., and N.L. were supported by the Jet Propulsion Laboratory, California Institute of Technology, under contract with NASA.
Table 1. CLO SS Pairs Sample
| (1) | (2) | (3) | (4) | (5) | (6) | (7) | (8) | (9) | (10) | (11) |
| --- | --- | --- | --- | --- | --- | --- | --- | --- | --- | --- |
| Name | R.A. | DEC. | B | d | T | z | SEP | $`\frac{L_{fir}}{L_B}`$ | $`\frac{f_{60\mu m}}{f_{100\mu m}}`$ | L<sub>fir</sub> |
| | (J2000) | (J2000) | mag | (kpc) | | km/s | (kpc) | | | 10<sup>10</sup>L |
| Arp 157: | | | | | | | 0.65 (5.6) | 4.37 | 0.68 | 4.56 |
| NGC 520a | 01h24m32.8s | +03d47m56s | 12.59 | 1.66 (14.4) | Sc | 2105 | | | | |
| NGC 520b | 01h24m34.9s | +03d47m30s | 13.55 | 1.14 (9.9) | Sm | 2360 | | | | |
| Arp 276: | | | | | | | 1.06 (16.7) | 1.74 | 0.33 | 2.19 |
| NGC 935 | 02h28m11.1s | +19d35m58s | 13.74 | 1.76 (27.7) | Sb | 4142 | | | | |
| IC 1801 | 02h28m12.7s | +19d35m00s | 14.65 | 1.19 (18.7) | Sb | 3970 | | | | |
| Arp 299: | | | | | | | 0.38 (4.6) | 9.12 | 0.99 | 26.6 |
| NGC 3690 | 11h28m31.0s | +58d33m41s | 12.72 | 1.30 (15.8) | Sc | 3132 | | | | |
| IC 694 | 11h28m33.5s | +58d33m47s | 12.49 | 1.47 (17.8) | Sc | 3121 | | | | |
| Arp 244: | | | | | | | 1.21 (7.7) | 1.15 | 0.59 | 3.86 |
| NGC 4038 | 12h01m52.8s | -18d51m54s | 10.91 | 5.25 (33.4) | Sm | 1642 | | | | |
| NGC 4039 | 12h01m55.2s | -18d53m06s | 11.05 | 3.09 (19.7) | Sm | 1641 | | | | |
| KPG 347: | | | | | | | 1.30 (11.4) | 1.20 | 0.36 | 3.18 |
| NGC 4567 | 12h36m32.7s | +11d15m28s | 12.02 | 2.47 (21.7) | Sb | 2274 | | | | |
| NGC 4568 | 12h36m34.7s | +11d14m15s | 11.99 | 3.50 (30.7) | Sc | 2255 | | | | |
| KPG 426: | | | | | | | 0.88 (26.3) | 1.55 | 0.38 | 4.01 |
| UGC 9376a | 14h33m46.8s | +40d04m52s | 14.80 | 1.35 (40.3) | Sb | 7616 | | | | |
| UGC 9376b | 14h33m48.4s | +40d05m39s | 14.56 | 1.52 (45.3) | Sa | 7764 | | | | |
| Arp 81: | | | | | | | 0.73 (17.9) | 3.39 | 0.55 | 8.36 |
| NGC 6621 | 18h12m54.7s | +68d21m49s | 14.39 | 0.96 (23.6) | Sb | 6210 | | | | |
| NGC 6622 | 18h13m00.2s | +68d21m12s | 14.23 | 1.00 (24.6) | Sa | 6466 | | | | |
| Arp 278: | | | | | | | 0.84 (15.0) | 6.61 | 0.46 | 4.57 |
| NGC 7253a | 22h19m26.2s | +29d23m55s | 15.03 | 1.46 (26.1) | Sm | 4718 | | | | |
| NGC 7253b | 22h19m31.3s | +29d23m25s | 14.95 | 1.42 (25.4) | Sm | 4493 | | | | |
Col.(1): Names of galaxy pairs and of pair components. Pairs with ISO SWS observations (see Table 3) are marked by $``$.
Col.(2): 2000 epoch right ascencion.
Col.(3): 2000 epoch declination.
Col.(4): Blue magnitude.
Col.(5): Major axis.
Col.(6): Hubble type.
Col.(7): Redshift.
Col.(8): Component separation.
Col.(9): FIR-to-blue luminosity ratio.
Col.(10): $`\frac{f_{60\mu m}}{f_{100\mu m}}`$ color ratio.
Col.(11): Integrated FIR (40–120 $`\mu m`$) luminosity calculated from IRAS 60$`\mu m`$ and 100$`\mu m`$ fluxes (Helou et al. 1988) and the mean redshift of the two components (H<sub>0</sub>=75 km/s Mpc<sup>-1</sup>).
Table 2. ISO CAM Observations of CLO SS Pairs
| Name | Obs. | Filter | $`\lambda `$ | $`\lambda /\delta \lambda `$ | pixel | samp. | M$`\times `$N | map | rms noise |
| --- | --- | --- | --- | --- | --- | --- | --- | --- | --- |
| | date | | ($`\mu m`$) | | size | step | | size | (1$`\sigma `$) |
| Arp 157 | 12.29.96 | LW7 | 9.62 | 4 | $`3^{\prime \prime }\times 3^{\prime \prime }`$ | $`45^{\prime \prime }\times 45^{\prime \prime }`$ | $`4\times 4`$ | $`3.8^{}\times 3.8^{}`$ | 0.09 mJy/pix |
| | 12.29.96 | LW3 | 15.0 | 3 | $`3^{\prime \prime }\times 3^{\prime \prime }`$ | $`45^{\prime \prime }\times 45^{\prime \prime }`$ | $`4\times 4`$ | $`3.8^{}\times 3.8^{}`$ | 0.09 mJy/pix |
| Arp 276 | 8.26.96 | LW7 | 9.62 | 4 | $`6^{\prime \prime }\times 6^{\prime \prime }`$ | $`60^{\prime \prime }\times 60^{\prime \prime }`$ | $`3\times 3`$ | $`5.0^{}\times 5.0^{}`$ | 0.11 mJy/pix |
| | 8.26.96 | LW3 | 15.0 | 3 | $`6^{\prime \prime }\times 6^{\prime \prime }`$ | $`60^{\prime \prime }\times 60^{\prime \prime }`$ | $`3\times 3`$ | $`5.0^{}\times 5.0^{}`$ | 0.15 mJy/pix |
| Arp 299 | 4.27.96 | LW7 | 9.62 | 4 | $`3^{\prime \prime }\times 3^{\prime \prime }`$ | $`45^{\prime \prime }\times 45^{\prime \prime }`$ | $`4\times 4`$ | $`3.8^{}\times 3.8^{}`$ | 0.09 mJy/pix |
| | 4.27.96 | LW9 | 15.0 | 8 | $`3^{\prime \prime }\times 3^{\prime \prime }`$ | $`45^{\prime \prime }\times 45^{\prime \prime }`$ | $`4\times 4`$ | $`3.8^{}\times 3.8^{}`$ | 0.20 mJy/pix |
| Arp 244 | 7.16.96 | LW7 | 9.62 | 4 | $`6^{\prime \prime }\times 6^{\prime \prime }`$ | $`60^{\prime \prime }\times 60^{\prime \prime }`$ | $`4\times 4`$ | $`6.0^{}\times 6.0^{}`$ | 0.21 mJy/pix |
| | 7.16.96 | LW3 | 15.0 | 3 | $`6^{\prime \prime }\times 6^{\prime \prime }`$ | $`60^{\prime \prime }\times 60^{\prime \prime }`$ | $`4\times 4`$ | $`6.0^{}\times 6.0^{}`$ | 0.17 mJy/pix |
| Arp 244-02 | 7.28.96 | LW3 | 15.0 | 3 | $`6^{\prime \prime }\times 6^{\prime \prime }`$ | $`6^{\prime \prime }\times 6^{\prime \prime }`$ | $`6\times 6`$ | $`3.5^{}\times 3.5^{}`$ | 0.06 mJy/pix |
| KPG 347 | 7.4.96 | LW7 | 9.62 | 4 | $`6^{\prime \prime }\times 6^{\prime \prime }`$ | $`90^{\prime \prime }\times 90^{\prime \prime }`$ | $`3\times 3`$ | $`6.0^{}\times 6.0^{}`$ | 0.47 mJy/pix |
| KPG 426 | 8.11.96 | LW7 | 9.62 | 4 | $`6^{\prime \prime }\times 6^{\prime \prime }`$ | $`60^{\prime \prime }\times 60^{\prime \prime }`$ | $`3\times 3`$ | $`5.0^{}\times 5.0^{}`$ | 0.11 mJy/pix |
| | 8.11.96 | LW3 | 15.0 | 3 | $`6^{\prime \prime }\times 6^{\prime \prime }`$ | $`60^{\prime \prime }\times 60^{\prime \prime }`$ | $`3\times 3`$ | $`5.0^{}\times 5.0^{}`$ | 0.11 mJy/pix |
| Arp 81 | 8.19.96 | LW7 | 9.62 | 4 | $`3^{\prime \prime }\times 3^{\prime \prime }`$ | $`45^{\prime \prime }\times 45^{\prime \prime }`$ | $`4\times 4`$ | $`3.8^{}\times 3.8^{}`$ | 0.10 mJy/pix |
| | 8.19.96 | LW3 | 15.0 | 3 | $`3^{\prime \prime }\times 3^{\prime \prime }`$ | $`45^{\prime \prime }\times 45^{\prime \prime }`$ | $`4\times 4`$ | $`3.8^{}\times 3.8^{}`$ | 0.09 mJy/pix |
| Arp 278 | 11.16.96 | LW7 | 9.62 | 4 | $`6^{\prime \prime }\times 6^{\prime \prime }`$ | $`60^{\prime \prime }\times 60^{\prime \prime }`$ | $`3\times 3`$ | $`5.0^{}\times 5.0^{}`$ | 0.10 mJy/pix |
| | 11.16.96 | LW3 | 15.0 | 3 | $`6^{\prime \prime }\times 6^{\prime \prime }`$ | $`60^{\prime \prime }\times 60^{\prime \prime }`$ | $`3\times 3`$ | $`5.0^{}\times 5.0^{}`$ | 0.12 mJy/pix |
Table 3.1. ISO SWS Observations of CLO SS Pairs
| Name | Pointing | R.A. (2000) | Dec. (2000) | Obs. date |
| --- | --- | --- | --- | --- |
| Arp 157 | N1 | 1h24m34.9s | 3d $`47^{}`$ $`29.8^{\prime \prime }`$ | 6.24.97 |
| | N2 | 1h24m32.8s | 3d $`47^{}`$ $`55.9^{\prime \prime }`$ | 6.24.97 |
| | OV | 1h24m33.9s | 3d $`47^{}`$ $`42.8^{\prime \prime }`$ | 6.24.97 |
| Arp 81 | N1 | 18h12m55.9s | 68d $`21^{}`$ $`50.1^{\prime \prime }`$ | 4.28.97 |
| | N2 | 18h13m00.1s | 68d $`21^{}`$ $`12.3^{\prime \prime }`$ | 4.28.97 |
| | OV | 18h12m58.1s | 68d $`21^{}`$ $`31.2^{\prime \prime }`$ | 5.29.97 |
| Arp 278 | N1 | 22h19m26.2s | 29d $`23^{}`$ $`53.2^{\prime \prime }`$ | 5.9.97 |
| | N2 | 22h19m28.9s | 29d $`23^{}`$ $`11.3^{\prime \prime }`$ | 5.9.97 |
Centers of ISO SWS pointings:
N1 — Nucleus of component 1 of the galaxy pair.
N2 — Nucleus of component 2 of the galaxy pair.
OV — Overlap region of the two components.
Table 3.2. Parameters of ISO SWS Observations
| | $`\lambda `$ | SWS- | Aperture | Resolution | Sensitivity |
| --- | --- | --- | --- | --- | --- |
| Line | ($`\mu m`$) | band | ($`{}_{}{}^{\prime \prime }\times _{}^{\prime \prime }`$) | ($`\lambda /\delta \lambda `$) | ($`3\sigma `$, Jy) |
| Br<sub>β</sub> | 2.63 | SW-1B | $`14\times 20`$ | 1470–1750 | 0.057 |
| Br<sub>α</sub> | 4.05 | SW-2A | $`14\times 20`$ | 1540–2130 | 0.29 |
| $`[`$Ne II$`]`$ | 12.81 | LW-3A | $`14\times 27`$ | 1250–1760 | 0.80 |
| H<sub>2</sub> S(1) | 17.03 | LW-3C | $`14\times 27`$ | 1760–2380 | 0.69 |
Table 4. ISO CAM Fluxes of CLO SS Pairs Sample
| (1) | (2) | (3) |
| --- | --- | --- |
| Name | f<sub>9.7μ</sub> | f<sub>15μ</sub> |
| | (mJy) | (mJy) |
| Arp 157: | 358 | 766 |
| NGC 520a | 32 | 49 |
| NGC 520b | 327 | 715 |
| Arp 276: | 241 | 300 |
| NGC 935 | 183 | 228 |
| IC 1801 | 63 | 72 |
| Arp 299: | 1699 | 5977 |
| NGC 3690 | 1071 | 3494 |
| IC 694 | 550 | 2158 |
| Arp 244: | 1227 | 2124 |
| NGC 4038 | 489 | 700 |
| NGC 4039 | 738 | 1427 |
| KPG 347: | 914 | 1269 |
| NGC 4567 | 124 | 252 |
| NGC 4568 | 790 | 1017 |
| KPG 426: | 84 | 110 |
| UGC 9376a | 32 | 41 |
| UGC 9376b | 52 | 66 |
| Arp 81: | 171 | 270 |
| NGC 6621 | 129 | 216 |
| NGC 6622 | 22 | 15 |
| Arp 278: | 178 | 253 |
| NGC 7253a | 126 | 200 |
| NGC 7253b | 44 | 49 |
The 15$`\mu m`$ data are taken from Bosselli et al. 1998.
Col.(1): Names of galaxy pairs and of pair components.
Col.(2): Flux density at 9.7$`\mu m`$. The uncertainty is $`15\%`$, dominated by the calibration error.
Col.(3): Flux density at 15$`\mu m`$. The uncertainty is $`15\%`$, dominated by the calibration error.
Table 5. SWS line emission of CLO SS Pairs
| | | Flux ($`10^{20}`$ W/cm<sup>2</sup>) | | | |
| --- | --- | --- | --- | --- | --- |
| Name | Pointing | Br<sub>β</sub> | Br<sub>α</sub> | \[Ne II\] | H<sub>2</sub> S(1) |
| | | $`2.63\mu m`$ | $`4.05\mu m`$ | $`12.81\mu m`$ | $`17.03\mu m`$ |
| Arp 157 | N1 | $`<2.3`$ | $`<2.2`$ | 19.3 ($`\pm 1.5`$) | 2.60 ($`\pm 0.34`$) |
| | N2 | $`<0.6`$ | $`<2.2`$ | $`<1.4`$ | $`<1.4`$ |
| | OV | $`<0.6`$ | $`<2.2`$ | 4.88 ($`\pm 0.96`$) | $`<1.4`$ |
| Arp 81 | N1 | $`<0.6`$ | $`<3.0`$ | $`<1.4`$ | $`<1.4`$ |
| | N2 | $`<2.3`$ | $`<2.2`$ | $`<0.7`$ | $`<0.5`$ |
| | OV | $`<0.6`$ | $`<3.0`$ | $`<1.9`$ | $`<1.2`$ |
| Arp 278 | N1 | $`<0.4`$ | $`<3.0`$ | $`<1.2`$ | $`<0.9`$ |
| | N2 | $`<0.5`$ | $`<2.2`$ | $`<2.3`$ | $`<1.2`$ |
Table 6. ISO CAM Fluxes of Nuclear and Overlap Starburst Emission in Arp 299 and Arp 244
| (1) | (2) | (3) | (4) | (5) |
| --- | --- | --- | --- | --- |
| Name | R.A. | DEC. | f<sub>9.7μ</sub> | f<sub>15μ</sub> |
| | (J2000) | (J2000) | (mJy) | (mJy) |
| Arp 299: | | | | |
| NGC 3690 | 11h28m31.0s | +58d33m41s | 613 ($`\pm 65`$) | 1410 ($`\pm 150`$) |
| IC 694 | 11h28m33.5s | +58d33m47s | 221 ($`\pm 25)`$) | 1044 ($`\pm 110`$) |
| overlap starburst | 11h28m31.0s | +58d33m49s | 123 ($`\pm 25`$) | 471 ($`\pm 70`$) |
| Arp 244: | | | | |
| NGC 4038 | 12h01m52.8s | -18d51m54s | 74 ($`\pm 10`$) | 135 ($`\pm 20`$) |
| NGC 4039 | 12h01m55.2s | -18d53m06s | 39 ($`\pm 10`$) | 76 ($`\pm 20`$) |
| overlap starburst | 12h01m54.8s | -18d53d03s | 142 ($`\pm 18`$) | 359 ($`\pm 50`$) |
Fluxes estimated from point source extractions.
Fluxes estimated by summing up the counts in a region of size$`15^{\prime \prime }`$ after the source associated to the nucleus of NGC 3690 is extracted.
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# Semi analytic approach to understanding the distribution of neutral hydrogen in the universe
## 1 Introduction
The nature and evolution of the initial power spectrum of density fluctuations could be obtained by studying the distribution of objects at different scales and different epochs. The formalism for studying the formation of dark matter (DM) structures is well established, as they are collisionless particles interacting only through gravity, and has been extensively studied using the large cosmological N-body simulations. However, in order to model the evolution of baryonic structures like galaxies, groups of galaxies etc. one needs to incorporate all the hydrodynamical processes, heating, cooling, star formation etc., in the N-body simulations. Because of such complications, our understanding of the formation of baryonic structures has been limited.
Among the various baryonic structures, the regions where one can neglect the star formation are comparatively easier to study. Two such areas are (a) low amplitude fluctuations in the intergalactic medium (IGM), where the star formation rate is very low, and (b) the intracluster medium, where one studies processes over large scales and thus the star formation details can be neglected. Hence considerable effort has been given in understanding these two types of structures.
The baryonic matter distribution at $`z5`$ is well probed through the absorption signatures they produce on the spectra of the distant QSOs. It is widely believed that while the metal lines systems (detected through Mg ii or C iv doublets) seen in the QSO spectra could be associated with the halos of the intervening luminous galaxies (Bergeron & Boisse 1991; Steidel 1993), most of the low neutral hydrogen column density absorption lines (commonly called as ‘Ly$`\alpha `$’ clouds) are believed to be due to low amplitude baryonic fluctuations in the IGM.
Semi analytical as well as hydrodynamical simulations are consistent with the view that the Ly$`\alpha `$ clouds are small scale density fluctuations (Bond, Szalay & Silk 1988; Cen et al. 1994; Zhang, Anninos & Norman 1995; Hernquist et al. 1996; Miralda-Escudé et al. 1996; Bi & Davidsen 1997; Riediger, Petitjean & Mücket 1998; Theuns, Leonard & Efstathiou 1998; Theuns et al. 1998; Davé et al. 1999) that are naturally expected in any standard structure formation models. This idea is supported by the detection and the evolution of the weak clustering among the Ly$`\alpha `$ clouds in the redshift space (Cristiani et al. 1995; Srianand 1997; Khare et al. 1997). Subsequently it is realised that the thermal history of the Ly$`\alpha `$ line forming regions depends on (i) epoch of reionisation (equation of state) (ii) rate of photoionisation and (iii) adiabatic cooling. One can in principle neglect shocks and other processes that are important only in the highly non-linear regime. However a simple linear evolution of the densities will fail to produce the saturated Ly$`\alpha `$ systems and one needs to incorporate non-linearities in the model.
As a first step, one can model the non-linear evolution of the baryonic fluctuations that produce Ly$`\alpha `$ clouds using one of the several approximations like: (i) Zeldovich approximation (Doroshkevich & Shandarin 1977; McGill 1990; Hui, Gnedin & Zhang 1997), (ii) lognormal approximation (Bi 1993; Gnedin & Hui 1996; Bi & Davidsen 1997). or (iii) power law approximation (Bi, Ge & Fang 1995) (strictly speaking, the baryonic fluctuations are calculated here using the linear theory). In all these cases the baryon density is estimated from the DM density by some rule and the neutral fraction is estimated by considering the equilibrium between the rate of photoionisation due to background radiation and the rate of recombination estimated from the temperature defined through the equation of state. All these models depend on various IGM parameters such as intensity of the background radiation, equation of state and density-averaged temperature as well as the cosmological parameters like $`\mathrm{\Omega }_m,\mathrm{\Omega }_\mathrm{\Lambda },`$ etc.
Observationally the statistical properties of the Ly$`\alpha `$ absorption lines are quantified through the column density distribution, correlation functions and their dependence on the mean redshift. The clustering properties of the Ly$`\alpha `$ absorption lines are studied through two point correlation function obtained either (a) in the redshift space using the lines detected along a single line-of-sight (LOS) which we call “LOS correlation function”, or (b) among the absorption lines detected along the lines of sights toward a few closely spaced QSOs which we call “transverse correlation function”. In either case the observed spectra is decomposed into clouds using “Voigt” profile fits. Though this process smoothens the density field over the width of the lines the average effects due to thermal broadening is taken care of by the Voigt profiles. One can also compute the two-point correlation function of the observed flux in different pixels. As this process does not decompose the actual density fields into cloudlets, in order to analyse the data the models should incorporate the thermal broadening and blending of contribution from different density fluctuations (Croft et al. 1999; McDonald et al. 1999). Most of the existing studies concentrate on obtaining constraints on the cosmological parameters using the observed statistical properties. Comparatively very less effort is directed to understand how the observed quantities depend on the physical conditions in the IGM.
In this work we make a preliminary attempt to investigate the dependence of the observable quantities on various parameters of the models using a simple analytic approach. We derive analytic relations for the two-point correlation function among the Ly$`\alpha `$ clouds and the column density distribution using a lognormal approximation. These equations are used with the observed Voigt profile fitted data to get constraints on different IGM and cosmological parameters. In section 2, we treat the non-linear evolution with a simple ansatz proposed by Bi & Davidsen (1997) for the baryonic density fluctuations, and derive analytic expressions for the correlation function along the LOS and in the transverse direction and the column density distribution. The model parameters used to obtain various results are discussed in section 3. In section 4, we study the correlation function at different redshifts for different structure formation models and for different values of the IGM parameters such as the density averaged temperature and the equation of state. We compare some of our results with the existing observational data. We also present the results for the column density distribution and study its dependence on various cosmological and IGM parameters. The results are summarised in section 5.
## 2 Analytic Model
The linear density contrast for dark matter in comoving $`k`$-space, for a particular redshift $`z`$, is given by
$$\delta _{\mathrm{DM}}(𝒌,z)=D(z)\delta _{\mathrm{DM}}(𝒌,0),$$
(1)
where $`D(z)`$ is the linear growth factor for the density contrast, normalised such that $`D(0)=1`$. If we assume the linear density contrast to be a Gaussian random field, then the corresponding linearly extrapolated power spectrum $`P_{\mathrm{DM}}(k)`$ is defined by
$$\delta _{\mathrm{DM}}(𝒌,0)\delta _{\mathrm{DM}}(𝒌^{\mathbf{}},0)=(2\pi )^3P_{\mathrm{DM}}(k)\delta _{\mathrm{Dirac}}(𝒌𝒌^{\mathbf{}}).$$
(2)
The power spectrum is only a function of the magnitude of $`𝒌`$, because of the isotropy of the background universe.
The linear density contrast for baryons in the IGM can be obtained from the DM density contrast by smoothing over scales below the Jeans length. We use the relation (Fang et al. 1993)
$$\delta _\mathrm{B}(𝒌,z)=\frac{\delta _{\mathrm{DM}}(𝒌,z)}{1+x_b^2(z)k^2},$$
(3)
where
$$x_b(z)=\frac{1}{H_0}\left[\frac{2\gamma k_BT_m(z)}{3\mu m_p\mathrm{\Omega }_m(1+z)}\right]^{1/2}$$
(4)
is the Jeans length; $`T_m`$ and $`\mu `$ are the density-averaged temperature and mean molecular weight of the IGM respectively; $`\mathrm{\Omega }_m`$ is the cosmological density parameter of total mass and $`\gamma `$ is the ratio of specific heats. Strictly speaking, equation (3) is valid only for the case where $`x_b`$ is independent of $`z`$, but it is shown by Bi, Borner & Chu (1992) that equation (3) is a good approximation for $`\delta _\mathrm{B}(𝒌,z)`$ even when $`x_b`$ has a redshift dependence. The linear density contrast in real comoving space, $`\delta (𝒙,z)`$, is the Fourier transform of equation (3).
In principle, to study the properties of the IGM one has to take into account the non-linearities in the density distribution and various physical processes such as shocks, radiation field, cooling etc. However, detailed hydrodynamical modelling of IGM has shown that most of the low column density Ly$`\alpha `$ absorption (i.e. $`N_{\mathrm{HI}}10^{14}`$ cm<sup>-2</sup>) are produced by regions that are either in the linear or in the weakly non-linear regime (Cen et al. 1994; Zhang, Anninos & Norman 1995; Hernquist et al. 1996; Miralda-Escudé et al. 1996; Theuns, Leonard & Efstathiou 1998; Theuns et al. 1998; Davé et al. 1999). The lower envelope of the column density, $`N_{\mathrm{HI}}`$ versus the thermal velocity dispersion, $`b`$ (given by $`b=(2k_BT/m_p)^{1/2}`$) scatter plot (Schaye et al. 1999a; Schaye et al. 1999b) suggests that there is a well defined relationship between the density and the temperature of the IGM (Hui & Gnedin 1997). Thus it is possible to model low column density systems using simple prescription for the non-linear density field and an equation of state.
In this work, we take into account the effect of non-linearity by assuming the number density distribution of the baryons, $`n_\mathrm{B}(𝒙,z)`$ to be a lognormal random field
$$n_\mathrm{B}(𝒙,z)=A\mathrm{e}^{\delta _\mathrm{B}(𝒙,z)}$$
(5)
where $`A`$ is a constant to be determined. The mean value of $`n_\mathrm{B}(𝒙,z)`$ is given by
$$n_\mathrm{B}(𝒙,z)n_0(z)=A\mathrm{e}^{\delta _\mathrm{B}(𝒙,z)}$$
(6)
Since $`\delta _\mathrm{B}(𝒙,z)`$ is a Gaussian random field, one can write
$$\mathrm{e}^{\delta _\mathrm{B}(𝒙,z)}=\mathrm{e}^{\mathrm{\Delta }^2(z)/2}$$
(7)
where
$$\mathrm{\Delta }^2(z)=\delta _\mathrm{B}^2(𝒙,z)=D^2(z)\frac{\mathrm{d}^3k}{(2\pi )^3}\frac{P_{\mathrm{DM}}(k)}{(1+x_b^2(z)k^2)^2}.$$
(8)
Hence,
$$A=n_0(z)\mathrm{e}^{\mathrm{\Delta }^2(z)/2}$$
(9)
and
$$n_\mathrm{B}(𝒙,z)=n_0(z)\mathrm{exp}[\delta _\mathrm{B}(𝒙,z)\frac{\mathrm{\Delta }^2(z)}{2}].$$
(10)
The lognormal distribution was introduced by Coles & Jones (1991) as a model for the non-linear matter distribution in the universe. This ansatz has several interesting features:
(a) It can be seen that the matter density given by equation (10) is always positive, even when $`\delta _\mathrm{B}\mathrm{}`$, unlike any polynomial function of $`\delta _\mathrm{B}`$. When the density contrast is small ($`\delta _\mathrm{B}1`$), equation (10) reduces to $`n_\mathrm{B}/n_01+\delta _\mathrm{B}`$, which is just what we expect from linear theory.
(b) On small scales, equation (10) becomes the isothermal hydrostatic solution, which describes highly clumped structures like intracluster gas, $`n_\mathrm{B}\mathrm{exp}(\mu m_p\psi _{\mathrm{DM}}/\gamma k_BT)`$, where $`\psi _{\mathrm{DM}}`$ is the dark matter potential (Sarazin & Bahcall 1977). The lognormal function can be thought of as the simplest function which links these two extreme regions smoothly.
(c) One can also think of the lognormal distribution as the kinematic model for the density field. If one assumes that the initial density and velocity fields are Gaussian, and extrapolates the continuity equation into non-linear regimes, treating the velocity field as linear, it turns out that the non-linear density field obtained in such a manner follows the lognormal distribution (Coles & Jones 1991).
(d) Bi & Davidsen (1997) have tested the distribution against hydrodynamical simulations, and found a reasonable match between them. The lognormal assumption has also been used to model the IGM in numerical simulations (Bi 1993; Bi & Davidsen 1997) and is found to be working well in reproducing the observations. In particular, the simulation results matches well with the observed column density distribution and number density of the Ly$`\alpha `$ absorption lines, the probability distribution of the $`b`$ parameter etc (see Bi & Davidsen 1997).
We shall also discuss briefly a more general argument as to why the lognormal distribution should be natural choice in a large class of phenomena. There is a wide class of quantities, denoted by $`f`$, the time evolution of which can be characterised by the following property – the change in the value of $`f`$ at some instant $`t_i`$ is proportional to its value at that instant, with the proportionality factor being a random variable. In mathematical notation, this can be written as $`f(t_{i+1})=f(t_i)+\epsilon _if(t_i)`$, where $`\epsilon _i`$ is the random variable. (Some examples of such phenomena in sociological context are (i) rich getting richer through fluctuations in stock market, and (ii) more facilities being provided to people who already have them.) Similar situation can occur in structure formation scenario also. The regions which have high density, because of stronger gravitational attraction, have a better chance of acquiring more mass. Let us denote the density field at some particular point at a given epoch $`t_i`$ by $`n(t_i)`$ and postulate the evolution,
$$n(t_{i+1})=n(t_i)+\epsilon _in(t_i)=(1+\epsilon _i)n(t_i)$$
(11)
We can now write $`n(t_{i+1})`$ in terms of some initial density field $`n(t_0)`$
$$n(t_{i+1})=(1+\epsilon _i)(1+\epsilon _{i1})\mathrm{}(1+\epsilon _0)n(t_0)$$
(12)
Taking logarithm of both sides
$$\mathrm{ln}[n(t_{i+1})]=\underset{j=0}{\overset{i}{}}\mathrm{ln}(1+\epsilon _j)+\mathrm{ln}[n(t_0)]$$
(13)
It is clear that when the time intervals $`(t_{i+1}t_i)`$ are small, the mass acquired within that interval will also be very small. Hence we expect that $`\epsilon _i1`$. Then the above expression becomes
$$\mathrm{ln}\left[\frac{n(t_{i+1})}{n(t_0)}\right]=\underset{j=0}{\overset{i}{}}\epsilon _j$$
(14)
This means that $`\mathrm{ln}[n(t)/n(t_0)]`$ is a sum of a large number of uncorrelated random variables. Using the central limit theorem, we can conclude that it follows a Gaussian distribution or, equivalently, $`n(t)`$ follows a lognormal distribution. This suggests that it may be reasonable to try an ansatz that the distribution of the non-linear baryonic density field is lognormal.
As an aside, we just mention that our analysis described here can easily be carried out for any other local ansatz for the non-linear baryonic density. \[The results for power law assumption in which $`n_\mathrm{B}(1+\delta )^p`$ will be discussed in a later paper.\]
Once we have obtained the total baryonic density, the fraction of hydrogen in the neutral form, $`f`$, in the IGM can be obtained by solving the ionisation equilibrium equation
$$\alpha (z,T(z))n_pn_e=J(z)n_{\mathrm{HI}},$$
(15)
where $`\alpha (z,T(z))`$ is the radiative recombination rate and $`J(z)`$ is rate of photoionisation for hydrogen at redshift $`z`$ (Black 1981); $`n_p,n_e`$ and $`n_{\mathrm{HI}}`$ are the number densities of proton, electron and neutral hydrogen, respectively. For simplicity, we assume that hydrogen is the only element present in the IGM and neglect the presence of helium and other heavier elements. In such a case, we have $`n_e=n_p`$. (This relation is not valid in the presence of helium or other heavier elements. If we have taken their presence into account, we would have got $`n_e=\kappa n_p`$, where $`\kappa `$ is a constant. Usually, $`1\kappa 1.2`$, because the amount of helium and heavier elements in the IGM is small compared to hydrogen. Since we do not know $`J(z)`$ beyond an accuracy of 10–20 per cent, we can always absorb $`\kappa `$ into $`J(z)`$.) Let us define the neutral fraction of hydrogen, $`f`$ by
$$f=\frac{n_{\mathrm{HI}}}{n_\mathrm{B}}=\frac{n_{\mathrm{HI}}}{n_{\mathrm{HI}}+n_p}$$
(16)
Hence we get from equation (15)
$$\frac{(1f)^2}{f}=\frac{J(z)}{\alpha (z,T(z))n_\mathrm{B}}.$$
(17)
In general, one can solve this equation and determine $`f`$ as a function of $`n_\mathrm{B}`$. This expression simplifies for two extreme cases. For $`f1`$, we get
$$f=\frac{\alpha (z,T(z))n_\mathrm{B}}{J(z)}$$
(18)
and for $`f1`$,
$$f=1\sqrt{\frac{J(z)}{\alpha (z,T(z))n_\mathrm{B}}}.$$
(19)
Hence, we have
$$n_{\mathrm{HI}}(𝒙,z)=\{\begin{array}{cc}\frac{\alpha (z,T(z))}{J(z)}n_\mathrm{B}^2(𝒙,z)\hfill & \text{(if }n_{\mathrm{HI}}n_\mathrm{B}\text{)}\hfill \\ n_\mathrm{B}(𝒙,z)\sqrt{\frac{J(z)n_\mathrm{B}(𝒙,z)}{\alpha (z,T(z))}}\hfill & \text{(if }n_{\mathrm{HI}}n_\mathrm{B}\text{)}\hfill \end{array}$$
(20)
The ionisation conditions in the Ly$`\alpha `$ absorbers are similar to the of H ii regions with $`f10^4`$. Thus, from now on we concentrate only on the case $`n_{\mathrm{HI}}n_\mathrm{B}`$.
We take the temperature dependence of the recombination coefficient $`\alpha `$ to be given by (Rauch et al. 1997)
$$\alpha (z,T(z))=\alpha _0\left(\frac{T(z)}{10^4K}\right)^{0.7},$$
(21)
where $`\alpha _0=4.2\times 10^{13}`$ cm<sup>3</sup> s<sup>-1</sup>. This relation is a good approximation for $`\alpha `$ in the temperature range relevant for Ly$`\alpha `$ forest. The temperature $`T`$ is related to the baryonic density $`n`$ through the equation of state. We assume a polytropic equation of state $`p\rho ^\gamma n^\gamma `$, or equivalently
$$T(z)=T_0(z)[n_\mathrm{B}(z)/n_0(z)]^{\gamma 1},$$
(22)
where
$$n_0(z)=\frac{\mathrm{\Omega }_{\mathrm{baryon}}\rho _c}{\mu _bm_p}(1+z)^3$$
(23)
is the mean baryonic number density at redshift $`z`$. $`\rho _c`$ is the critical matter density at the present epoch, given by
$$\rho _c=1.8791\times 10^{29}h^2\mathrm{cm}^3$$
(24)
and $`\mu _bm_p`$ is the mass per baryonic particle. Then, the H i density becomes
$$n_{\mathrm{HI}}(𝒙,z)=F(z)\left(\frac{n_\mathrm{B}(𝒙,z)}{n_0(z)}\right)^\beta $$
(25)
where
$$F(z)=\alpha _0n_0^2(z)\left(\frac{T_0(z)}{10^4K}\right)^{0.7}J^1(z)$$
(26)
and
$$\beta =2.70.7\gamma .$$
(27)
(We note, in passing, that $`\beta `$ becomes negative if $`\gamma >3.86`$.) We can write the H i density in terms of the linear baryonic density contrast
$$n_{\mathrm{HI}}(𝒙,z)=F_1(z)\mathrm{exp}[\beta \delta _\mathrm{B}(𝒙,z)],$$
(28)
where
$$F_1(z)=F(z)\mathrm{e}^{\beta \mathrm{\Delta }^2(z)/2}.$$
(29)
It is clear from equation (28) that the H i distribution at a particular redshift is also described by a lognormal distribution. All the statistical quantities regarding H i can be derived from this in a straightforward manner.
### 2.1 Correlation Function for Neutral Hydrogen
One of our main interest is the correlation function
$`n_{\mathrm{HI}}(𝒙,z)n_{\mathrm{HI}}(𝒙^{\mathbf{}},z^{})=`$ (30)
$`F_1(z)F_1(z^{})\mathrm{exp}\{\beta [\delta _\mathrm{B}(𝒙,z)+\delta _\mathrm{B}(𝒙^{\mathbf{}},z^{})]\},`$
from which several useful quantities can be obtained. Since $`\delta _\mathrm{B}`$ is a Gaussian random field, we can write, using equation (7)
$`\mathrm{exp}\{\beta [\delta _\mathrm{B}(𝒙,z)+\delta _\mathrm{B}(𝒙^{\mathbf{}},z^{})]\}=`$ (31)
$`\mathrm{exp}[{\displaystyle \frac{\beta ^2}{2}}\{\mathrm{\Delta }^2(z)+\mathrm{\Delta }^2(z^{})+2Q(𝒙,𝒙^{\mathbf{}};z,z^{})\}],`$
where
$$Q(𝒙,𝒙^{\mathbf{}};z,z^{})=\delta _\mathrm{B}(𝒙,z)\delta _\mathrm{B}(𝒙^{\mathbf{}},z^{}).$$
(32)
Simple algebra gives
$`Q(𝒙,𝒙^{\mathbf{}};z,z^{})Q(𝒙𝒙^{\mathbf{}};z,z^{})=`$ (33)
$`D(z)D(z^{}){\displaystyle \frac{\mathrm{d}^3k}{(2\pi )^3}\frac{P_{\mathrm{DM}}(k)\mathrm{e}^{\mathrm{i}𝒌(𝒙𝒙^{\mathbf{}})}}{(1+x_b^2(z)k^2)(1+x_b^2(z^{})k^2)}}.`$
One also notes from equations (8) and (33) that $`\mathrm{\Delta }^2(z)=Q(0;z,z)`$. We can now write equation (30) as
$$n_{\mathrm{HI}}(𝒙,z)n_{\mathrm{HI}}(𝒙^{\mathbf{}},z^{})=F_2(z)F_2(z^{})\mathrm{e}^{\beta ^2Q(𝒙𝒙^{\mathbf{}};z,z^{})},$$
(34)
where
$$F_2(z)=F_1(z)\mathrm{e}^{\beta ^2\mathrm{\Delta }^2(z)/2}.$$
(35)
One needs to normalise the quantity $`n_{\mathrm{HI}}(𝒙,z)n_{\mathrm{HI}}(𝒙^{\mathbf{}},z^{})`$, to obtain the correlation function $`\xi _{\mathrm{HI}}(𝒙𝒙^{\mathbf{}};z,z^{})`$ for H i. A natural way of normalising the correlation would be to use the definition
$$1+\xi _{\mathrm{HI}}(𝒙𝒙^{\mathbf{}};z,z^{})=\frac{n_{\mathrm{HI}}(𝒙,z)n_{\mathrm{HI}}(𝒙^{\mathbf{}},z^{})}{n_{\mathrm{HI}}(𝒙,z)n_{\mathrm{HI}}(𝒙^{\mathbf{}},z^{})}.$$
(36)
Since $`n_{\mathrm{HI}}(𝒙,z)=F_2(z)`$, we get
$$\xi _{\mathrm{HI}}(𝒙𝒙^{\mathbf{}};z,z^{})=\mathrm{e}^{\beta ^2Q(𝒙𝒙^{\mathbf{}};z,z^{})}1$$
(37)
with $`Q`$ given by equation (33).
All the analysis above is valid if one can probe any scale with arbitrary accuracy. But it turns out that one cannot obtain information about scales smaller than some particular value, due to various observational constraints. While observing along a LOS, it will be impossible to probe the velocity scales less than the spectroscopic limit due to thermal broadening and the blending of spectral lines. Similarly, while observing across the transverse direction, the peculiar velocities of individual points will constrain the velocity resolution which we have not taken into account in the above analysis. If $`\mathrm{\Delta }v`$ is the smallest scale one can probe, then the corresponding limit in the redshift-space is
$$\mathrm{\Delta }z=\frac{\mathrm{\Delta }v}{c}(1+z).$$
(38)
This means that we will not be able to probe below a comoving length scale given by
$$\mathrm{\Delta }x(z)=d_H(z)\mathrm{\Delta }z,$$
(39)
where
$`d_H(z)`$ $`=`$ $`c\left({\displaystyle \frac{\dot{a}}{a}}\right)^1`$ (40)
$`=`$ $`{\displaystyle \frac{c}{H_0}}[\mathrm{\Omega }_\mathrm{\Lambda }+\mathrm{\Omega }_m(1+z)^3+\mathrm{\Omega }_k(1+z)^2]^{1/2},`$
$$\mathrm{\Omega }_k=1\mathrm{\Omega }_m\mathrm{\Omega }_\mathrm{\Lambda }.$$
(41)
This effect can be included in our calculation by smoothing over all the length scales smaller than $`\mathrm{\Delta }x(\overline{z})`$, where
$$\overline{z}=\frac{1}{2}(z+z^{})$$
(42)
is the average redshift. We use a Gaussian window of width $`\sigma _x(\overline{z})=\mathrm{\Delta }x(\overline{z})`$, and get a smoothed version of $`Q`$ in equation (33). In Fourier space, this smoothing will introduce an extra Gaussian term in the integrand, and our smoothed $`Q`$ will be
$`Q_{\mathrm{smooth}}(𝒙𝒙^{\mathbf{}};z,z^{})=`$ (43)
$`D(z)D(z^{}){\displaystyle \frac{\mathrm{d}^3k}{(2\pi )^3}\frac{P_{\mathrm{DM}}(k)\mathrm{e}^{k^2\sigma _x^2(\overline{z})/2}\mathrm{e}^{\mathrm{i}𝒌(𝒙𝒙^{\mathbf{}})}}{(1+x_b^2(z)k^2)(1+x_b^2(z^{})k^2)}}.`$
The angular integrations can be carried out trivially, and we get
$`Q_{\mathrm{smooth}}(𝒙𝒙^{\mathbf{}};z,z^{})={\displaystyle \frac{D(z)D(z^{})}{2\pi ^2}}\times `$ (44)
$`{\displaystyle _0^{\mathrm{}}}dk{\displaystyle \frac{P_{\mathrm{DM}}(k)k^2\mathrm{e}^{k^2\sigma _x^2(\overline{z})/2}}{(1+x_b^2(z)k^2)(1+x_b^2(z^{})k^2)}}{\displaystyle \frac{\mathrm{sin}kX}{kX}},`$
where $`X=|𝒙𝒙^{\mathbf{}}|`$. The final integration can be done numerically, once the DM power spectrum is given.
At this stage, the relations derived above can be used for any $`𝒙,𝒙^{\mathbf{}},z,z^{}`$. As we mentioned earlier, if one observes the H i along a particular LOS, then one is probing different regions of the IGM at different redshifts. The position $`x`$ will be related to the redshift $`z`$ by the relation
$$x(z)=_0^zd_H(z^{})dz^{}$$
(45)
where $`d_H(z)`$ is given by equation (40). Then the LOS correlation function is given by
$$\xi _{\mathrm{HI}}^{\mathrm{LOS}}(z,z^{})=\mathrm{e}^{\beta ^2Q_{\mathrm{LOS}}(l(z,z^{});z,z^{})}1,$$
(46)
where
$`Q_{\mathrm{LOS}}(l(z,z^{});z,z^{})=`$
$`{\displaystyle \frac{D(z)D(z^{})}{2\pi ^2}}{\displaystyle _0^{\mathrm{}}}dk{\displaystyle \frac{P_{\mathrm{DM}}(k)k^2\mathrm{e}^{k^2\sigma _x^2(\overline{z})/2}}{(1+x_b^2(z)k^2)(1+x_b^2(z^{})k^2)}}{\displaystyle \frac{\mathrm{sin}kl}{kl}},`$
and
$$l(z,z^{})=x(z)x(z^{}).$$
(48)
It should be stressed that $`\xi _{\mathrm{HI}}^{\mathrm{LOS}}(z,z^{})\xi _{\mathrm{HI}}^{\mathrm{LOS}}(zz^{})`$. This means that one cannot rigorously define a power spectrum from the LOS correlation function because the correlation is a function of two variables $`z`$ and $`z^{}`$. In other words, the LOS power spectrum does not exist in strict sense. However, one can get an approximate LOS power spectrum for a small redshift range around any mean redshift. We have already defined the average redshift in equation (42). We define a redshift difference
$$\mathrm{\Delta }z=zz^{}$$
(49)
and evaluate the correlation function for a particular value of $`\overline{z}`$ as a function of $`\mathrm{\Delta }z`$, i.e.,
$$\xi _{\mathrm{HI}}^{\mathrm{LOS}}(\overline{z},\mathrm{\Delta }z)=\mathrm{e}^{\beta ^2Q_{\mathrm{LOS}}(\overline{z},\mathrm{\Delta }z)}1,$$
(50)
where
$`Q_{\mathrm{LOS}}(\overline{z},\mathrm{\Delta }z)=`$
$`{\displaystyle \frac{D(\overline{z}+\mathrm{\Delta }z/2)D(\overline{z}\mathrm{\Delta }z/2)}{2\pi ^2}}{\displaystyle _0^{\mathrm{}}}\mathrm{d}k\{{\displaystyle \frac{\mathrm{sin}(kd_H(\overline{z})\mathrm{\Delta }z)}{kd_H(\overline{z})\mathrm{\Delta }z}}`$
$`\times {\displaystyle \frac{P_{\mathrm{DM}}(k)k^2\mathrm{e}^{k^2\sigma _x^2(\overline{z})/2}}{(1+x_b^2(\overline{z}+\mathrm{\Delta }z/2)k^2)(1+x_b^2(\overline{z}\mathrm{\Delta }z/2)k^2)}}\}`$
For small $`\mathrm{\Delta }z`$, one can use equation (39) to write the correlation as a function of $`\mathrm{\Delta }x`$, Fourier transform the correlation, and get the power spectrum. Such a power spectrum will depend on the value of $`\overline{z}`$. We stress again this power spectrum is approximate in the sense that it exists only for $`\mathrm{\Delta }z\overline{z}`$.
The transverse correlation is observed at some particular redshift ($`z=z^{}`$), along the transverse direction. Then
$$\xi _{\mathrm{HI}}^{\mathrm{trans}}(l_{};z)=\mathrm{e}^{\beta ^2Q_{\mathrm{trans}}(l_{};z)}1,$$
(52)
$`Q_{\mathrm{trans}}(l_{};z)=`$
$`{\displaystyle \frac{D^2(z)}{2\pi ^2}}{\displaystyle _0^{\mathrm{}}}dk{\displaystyle \frac{P_{\mathrm{DM}}(k)k^2\mathrm{e}^{k^2\sigma _x^2(z)/2}}{(1+x_b^2(z)k^2)^2}}{\displaystyle \frac{\mathrm{sin}kl_{}}{kl_{}}},`$
where $`l_{}`$ is the comoving distance along the transverse direction. For a given redshift, the transverse correlation is only a function of $`l_{}`$. Hence, one can obtain the power spectrum from $`\xi _{\mathrm{HI}}^{\mathrm{trans}}`$ following usual methods.
### 2.2 Column Density Distribution
One of the other statistics the observers use to quantify the distribution Ly$`\alpha `$ absorption lines is column density distribution. Indeed one can get the analytic expression for this using the formalism developed so far in this work. Note that Voigt profile fitting to the absorption lines are used to get the observed column density distribution. Here we use a method called ‘density-peak ansatz’ (DPA), discussed in Gnedin & Hui (1996) and Hui, Gnedin & Zhang (1997) to derive an analytic expression for the column density distribution.
Suppose we are looking at the IGM along any one direction, at some redshift $`z`$. Then the linear density field $`\delta _\mathrm{B}^{(1D)}(x,z)`$ along that LOS will be described by a one dimensional Gaussian random field. DPA assumes that each density peak in the comoving space is associated with an absorption line, and one can assign a definite column density to each of them. In the articles referred above, each density peak is fitted with a Gaussian, and the column density is calculated using
$$N_{\mathrm{HI}}_{\mathrm{peak}}n_{\mathrm{HI}}(x)dx.$$
(54)
In such a case, there is a definite correlation between the value of the density field at the peak, and the effective width of the absorber (which is determined by the correlation between the density field and its second derivative at the peak, and is fixed once the fitting function for the density peak is given). We, however, take a simpler approach in assigning the column density to a density peak, which is described below.
The coherence scale of the distribution is defined as (Bardeen et al. 1986)
$$R^{}\frac{\sigma _1}{\sigma _2},$$
(55)
where $`\sigma _1`$ and $`\sigma _2`$ are defined in equation (74) (see Appendix A). This length is a measure of the distance between two successive zeroes for the one dimensional Gaussian random field. Since this is the relevant scale for the distribution of zero-crossing, we expect the effective length scale of a peak to be a fraction of $`R^{}`$. Then the column density corresponding to a particular peak will be
$$N_{\mathrm{HI}}n_{\mathrm{HI}}[\mathrm{peak}]\mathrm{R}^{}=\mathrm{n}_{\mathrm{HI}}[\mathrm{peak}]\mathrm{R}^{}ϵ$$
(56)
where $`n_{\mathrm{HI}}`$\[peak\] is the H i number density at the peak and $`ϵ`$ is the proportionality constant, which can be used as a free parameter in comparing with observations. We have assumed $`ϵ`$ to be independent of $`N_{\mathrm{HI}}`$, which means that the column density is directly proportional to the peak density. Using this prescription for obtaining the column density from the H i density, we can easily obtain the relation between $`N_{\mathrm{HI}}`$ and the total baryonic over-density $`n_\mathrm{B}/n_0`$ using equation (20). For the case $`n_{\mathrm{HI}}n_\mathrm{B}`$, the relation is given by
$$\frac{n_\mathrm{B}[\mathrm{peak}]}{n_0}=\left(\frac{N_{\mathrm{HI}}}{R^{}ϵF(z)}\right)^{1/\beta }.$$
(57)
The relation between $`N_{\mathrm{HI}}`$ and $`\delta _\mathrm{B}^{(1D)}`$, is then given by
$$\delta _\mathrm{B}^{(1D)}[\mathrm{peak}]=\frac{1}{\beta }\mathrm{ln}\left(\frac{\mathrm{N}_{\mathrm{HI}}}{\mathrm{R}^{}ϵ\mathrm{F}(\mathrm{z})}\right)+\frac{\mathrm{\Delta }^2}{2}$$
(58)
Given this relation, it is straightforward to obtain the quantity $`\mathrm{d}N_{\mathrm{pk}}/(\mathrm{d}z\mathrm{d}N_{\mathrm{HI}})`$, defined as the number of clouds (peaks) per column density interval per redshift interval. For completeness, we give the relevant calculations in the appendix.
## 3 Model Parameters
In this section we discuss about the various model parameters used in obtaining the results. The parameters defining the model can be divided into two categories : (i) cosmological parameters, and (ii) parameters related to the IGM.
The first set of parameters are those which determine the background cosmology. We assume that the background universe is described by the FRW metric. We have considered four different cosmological models with the parameters listed below:
1. $`\mathrm{\Omega }_m=1,\mathrm{\Omega }_\mathrm{\Lambda }=0,h=0.5`$
2. $`\mathrm{\Omega }_m=0.35,\mathrm{\Omega }_\mathrm{\Lambda }=0,h=0.5`$
3. $`\mathrm{\Omega }_m=0.35,\mathrm{\Omega }_\mathrm{\Lambda }=0.65,h=0.5`$
4. $`\mathrm{\Omega }_m=0.35,\mathrm{\Omega }_\mathrm{\Lambda }=0.65,h=0.65`$
The next cosmological input that is required is the form of the DM power spectrum. We take the following form for $`P_{\mathrm{DM}}(k)`$ (Efstathiou, Bond & White 1992)
$$P_{\mathrm{DM}}(k)=\frac{Ak}{(1+[ak+(bk)^{1.5}+(ck)^2]^\nu )^{2/\nu }}$$
(59)
where $`\nu =1.13`$, $`a=(6.4/\mathrm{\Gamma })`$$`h^1`$ Mpc, $`b=(3.0/\mathrm{\Gamma })`$$`h^1`$ Mpc, $`c=(1.7/\mathrm{\Gamma })`$$`h^1`$ Mpc and $`\mathrm{\Gamma }=\mathrm{\Omega }_mh`$. The normalisation parameter $`A`$ is fixed through the value of $`\sigma _8`$ (the rms density fluctuation in spheres of radius 8 $`h^1`$ Mpc).We take the values of $`\sigma _8`$ to be given by (Eke, Cole & Frenk 1996)
$$\sigma _8=\{\begin{array}{cc}(0.52\pm 0.04)\mathrm{\Omega }_m^{0.46+0.10\mathrm{\Omega }_m}\hfill & \text{(if }\mathrm{\Omega }_\mathrm{\Lambda }=0\text{)}\hfill \\ (0.52\pm 0.04)\mathrm{\Omega }_m^{0.52+0.13\mathrm{\Omega }_m}\hfill & \text{(if }\mathrm{\Omega }_\mathrm{\Lambda }=1\mathrm{\Omega }_m\text{)}\hfill \end{array}$$
(60)
The next set of parameters are related to the physical conditions in the IGM. The parameters we need to describe the IGM are $`\gamma `$, $`T_m`$, $`\mathrm{\Omega }_{\mathrm{baryon}}`$, $`J(z)`$ and $`T_0(z)`$. All these quantities were defined in section 2. Besides these, we have also introduced a parameter $`ϵ`$ while modelling the column density distribution of the IGM. This is taken as a free parameter, to be fixed through observations.
It is known that the value of $`\gamma `$, at any given epoch, depends on the reionisation history of the universe (Hui & Gnedin 1997). The value of $`\gamma `$ and its evolution is still quite uncertain. Using Voigt profile fits to the observed Ly$`\alpha `$ absorption lines one can in principle obtain the value of $`\gamma `$. Available observations are consistent with $`\gamma `$ in the range $`1.2<\gamma <1.7`$ (Schaye et al. 1999b) for $`2z4.5`$. As far as the evolution of $`\gamma `$ is concerned, we shall treat it as independent of $`z`$. The density averaged temperature is defined as,
$$T_m=\frac{\rho (T)TdT}{\rho (T)dT}$$
(61)
Using the equation of state $`T\rho ^{\gamma 1}`$, we get
$$T_m=\frac{2\gamma 1}{\gamma }(T_{\mathrm{max}}T_{\mathrm{min}}).$$
(62)
We take $`T_m`$ to be in the range 10,000 K$`<T_m<`$60,000 K. The minimum value of $`T_m`$ corresponds to the minimum temperature of the IGM, which is determined by the photoionisation equilibrium. The maximum value of $`T_m`$ corresponds to $`b`$ parameter $`31.7`$ km s<sup>-1</sup> and is consistent with the minimum value of $`b`$ observed at higher H i column densities (i.e. $`10^{14.5}`$ cm<sup>-2</sup>). The evolution of $`T_m`$ depends on how $`T_{\mathrm{max}}`$ and $`T_{\mathrm{min}}`$ evolve with redshift. One possibility is to take the adiabaticity relation $`T_m(1+z)^{3\gamma 3}`$. However, people have argued that since there is no conclusive evidence of the evolution of the temperature in IGM, one should treat the mean temperature of the IGM as constant (Bi et al. 1995). Hence, in this paper, we also consider the second possibility where $`T_m`$ is independent of $`z`$. For the cases where we consider a small redshift range $`\mathrm{\Delta }z`$ around a mean redshift $`\overline{z}`$, the effect of the evolution of $`T_m`$ is not very significant. But, whenever we study the evolution of a quantity over a large redshift range, we have to take into account the various redshift dependences of $`T_m`$.
We note that when we normalise the correlation function for neutral hydrogen, the parameters $`\mathrm{\Omega }_{\mathrm{baryon}}`$, $`J(z)`$ and $`T_0(z)`$ cancel out. Hence the knowledge of these parameters are not necessary for modelling the correlation function. However, in the case of column density distribution, they appear as a combination $`(\mathrm{\Omega }_{\mathrm{baryon}}h^2)^2J^1(z)T_0^{0.7}(z)`$ through $`F(z)`$ (equations (26) and (23)). The values of these quantities are not known accurately, nor do we know how $`J(z)`$ and $`T_0(z)`$ evolve. In this work, we take $`J(z)`$ to be independent of $`z`$. We fix the value of the combination $`(\mathrm{\Omega }_{\mathrm{baryon}}h^2)^2J^1`$ to be $`(0.026)^2\times 10^{12}`$ s, which is consistent with the values given in McDonald et al. (1999). One should note that any change in the values of the above parameters can be compensated (to some extent) by changing the value of $`ϵ`$, which is a free parameter in our model for the column density distribution. For $`T_0`$ evolution, we shall consider two separate cases like $`T_m`$, i.e., (i) $`T_0(1+z)^{3\gamma 3}`$ and (ii) $`T_0`$ independent of $`z`$.
As we have discussed earlier, we need to smooth the power spectrum below some velocity because the blending of spectral lines makes it impossible to resolve the lines below a particular velocity. Typically this velocity is of the order of a few tens km s<sup>-1</sup>. For definiteness, we take the smoothing velocity to be $`\mathrm{\Delta }v=30`$ km s<sup>-1</sup>.
## 4 Results
### 4.1 LOS Correlation
We shall now compute the results for the H i correlation function along the LOS as a function of the velocity separation $`v`$, where $`v`$ is related to $`\mathrm{\Delta }z`$ by $`v=c\mathrm{\Delta }z/(1+\overline{z})`$.
The results for the LOS correlation function for different cosmological models are shown in Figure 1.
We have chosen typical values for $`T_m`$ as 40,000 K and $`\gamma =1.5`$, at a redshift of $`\overline{z}=2.5`$. It can be seen that the correlation curves tend to flatten at low velocities, and goes to zero at high velocities.
Miralda-Escudé et al. (1996) give the correlation function of the transmitted flux along a LOS (see the solid curve in their Figure 12a). We note that the correlation curves for the transmitted flux and the neutral hydrogen density need not be exactly the same. However, we expect that the broad features and the general trends should be alike. It turns out that the correlation curve for transmitted flux obtained from the simulations does have the same trend as our results.
We can compare how the shape of $`\xi _{\mathrm{LOS}}(v)`$ depends on various parameters. We found that the $`\xi _{\mathrm{LOS}}(v)`$ curve falls approximately like a power law, $`\xi _{\mathrm{LOS}}(v)v^p`$, at velocities within 100–1000 km s<sup>-1</sup>. At lower velocities the curve is practically independent of $`v`$. We have given the value of $`p`$ and the rms error on $`p`$ for different cosmological models and for different IGM parameters in Table 1.
The dependence of $`p`$ on the IGM parameters can be understood easily. Higher value of $`T_m`$ implies a larger $`x_b`$ which, in turn, implies more smoothing of the power spectrum at low scales. However, the larger scales are more or less unaffected by $`x_b`$. Consequently, the correlation curve becomes flatter as we increase $`T_m`$. Also, for a fixed value of $`\gamma `$, the effect of the cosmological models on the shape of $`\xi _{\mathrm{LOS}}`$ is large for low $`T_m`$.
The effect of $`\gamma `$ on $`\xi _{\mathrm{LOS}}`$ is twofold – increasing $`\gamma `$ introduces more smoothing at low scales just like $`T_m`$, and there is also a reduction in the neutral hydrogen density fluctuations for given baryonic fluctuations see (equation (28)). Both these effects make the correlation curve flatter, which is what we see from Table 1. Furthermore, because of this twofold effect, $`\gamma `$ affects $`p`$ much more than $`T_m`$ does. This point can be seen clearly in Figure 2. In the left figure, we have kept $`x_b`$ constant, and shown the effect of changing only the neutral hydrogen density fluctuations. The middle figure shows the effect of changing only $`T_m`$ or, equivalently, the Jeans scale $`x_b`$, without changing anything else. In the right figure, the full effect of $`\gamma `$ can be seen, as it changes both the Jeans scale and the neutral hydrogen density fluctuations. Since $`T_m`$ does not have much effect on $`p`$, we can determine $`\gamma `$ from observations of LOS correlation even with ill-constrained values of $`T_m`$, provided the cosmology is known from some other studies.
Just like in the case of $`T_m`$, the effect of the cosmological models on $`p`$ is large for low $`\gamma `$. The first reason for this is same as in the case of $`T_m`$ – low $`\gamma `$ implies less smoothing and hence DM fluctuations are more effective. The second reason is that for low $`\gamma `$, the neutral hydrogen density fluctuations are larger for given linear baryonic density fluctuations. Thus, a slight change in the DM fluctuations causes a significant change in the neutral hydrogen fluctuations (for a fixed $`x_b`$) which, in turn, makes the correlation function sensitive to the DM power spectrum.
As the universe evolves after the reionisation, it is possible that the value of $`\gamma `$ increases (Hui & Gnedin 1997) and the LOS correlation becomes less and less sensitive to the DM power spectrum. Hence, if the reionisation has occurred very early, it is extremely unlikely that one can fix the cosmological parameters from observations of $`\xi _{\mathrm{LOS}}`$. We not only need to know the value of $`\gamma `$ and $`T_m`$ accurately, but also the value of $`p`$ to an accuracy better than 10 per cent.
Finally, we comment on the effect of $`h`$ on the shape of $`\xi _{\mathrm{LOS}}`$. As one can clearly see from Table 1, the effect is quite significant ($`12`$ per cent). The parameter $`h`$ affects the LOS correlation function (equation (LABEL:eq:q\_los\_delz)) in three ways – (i) it changes the shape of $`P_{\mathrm{DM}}(k)`$ (see equation (59)), (ii) it affects the value of $`x_b`$ (equation (4)) and (iii) it affects the relation between distance and redshift. As a result, the values of $`\sigma _x(\overline{z})`$ and $`d_H(\overline{z})`$ in the integrand of equation (LABEL:eq:q\_los\_delz) get modified. However, it turns out that the last two effects can be scaled out. To understand this more clearly, we rewrite equation (LABEL:eq:q\_los\_delz) in a slightly modified form
$`Q_{\mathrm{LOS}}(\overline{z},\mathrm{\Delta }z)={\displaystyle \frac{D(\overline{z}+\mathrm{\Delta }z/2)D(\overline{z}\mathrm{\Delta }z/2)}{2\pi ^2}}\times `$
$`{\displaystyle _0^{\mathrm{}}}\mathrm{d}K\{{\displaystyle \frac{\mathrm{sin}(KD_H(\overline{z})\mathrm{\Delta }z)}{KD_H(\overline{z})\mathrm{\Delta }z}}\times `$
$`{\displaystyle \frac{h^3P_{\mathrm{DM}}(Kh)K^2\mathrm{e}^{K^2\mathrm{\Sigma }_x^2(\overline{z})/2}}{(1+X_b^2(\overline{z}+\mathrm{\Delta }z/2)K^2)(1+X_b^2(\overline{z}\mathrm{\Delta }z/2)K^2)}}\}`$
where $`Kk/h`$ and
$$D_H(\overline{z})d_H(\overline{z})h,\mathrm{\Sigma }_x(\overline{z})\sigma _x(\overline{z})h,X_bx_bh$$
(64)
One can easily verify that all the three quantities defined above $`(D_H,\mathrm{\Sigma }_x,X_b)`$ are independent of $`h`$. Hence $`Q_{\mathrm{LOS}}(\overline{z},\mathrm{\Delta }z)`$ depends on $`h`$ only through the combination $`h^3P_{\mathrm{DM}}(Kh)`$. It is very difficult to study this function analytically. We have studied it using numerical methods and found that for the power spectra used in this paper, its effect is to make the LOS correlation steeper when $`h`$ is larger.
To compare our results with observational data we take the data from Cristiani et al. (1997). The data consists of several QSO spectra at various redshifts, ranging from 1.7 to 3.7. This range is pretty large, and evolutionary effects will be significant in the data. We compare the observed LOS correlation (points with error bars) with the theoretical curve for the four cosmological models, and for various ranges of values of $`T_m`$ and $`\gamma `$ in Figure 3 for $`\overline{z}=2.5`$.
It should be noted that the observational data points were obtained using the Ly$`\alpha `$ clouds with log($`N_{\mathrm{HI}}/`$cm<sup>-2</sup>)$`>14`$. However, we have not used any such constraint while obtaining the analytical curves. As a preliminary check, we can see that the analytical curves have the broad features which are expected from the observational data. We hope to carry out a more detailed comparison with observations in a future publication.
We next check the redshift evolution of the LOS correlation function. For definiteness, we consider $`\xi _{\mathrm{LOS}}`$ at a particular velocity, $`v=100`$ km s<sup>-1</sup>, and study it as a function of $`\overline{z}`$. We have assumed that $`\gamma `$ does not evolve with redshift. Since we are studying the evolution over a large redshift range, we have to consider both the possibilities for the evolution of $`T_m`$ discussed in section 3. The evolution curve closely follows a power-law dependence, i.e., $`\xi _{\mathrm{LOS}}(\overline{z})(1+\overline{z})^q`$ in the range $`1.5<\overline{z}<4.5`$. We give the values of $`q`$ and the rms error for different cases in Table 2.
It is clear that the cosmology has maximum influence when $`\gamma `$ and $`T_m`$ are small. The reason for this is same as that discussed in the case of Table 1.
For a given cosmology and a given value of $`\gamma `$, the effect of $`T_m`$ (or, equivalently, $`x_b`$) on $`q`$ is insignificant with the relative change in $`q`$ being about 6–10 per cent. The reason for this is as follows: the effect of $`x_b`$ is significant only for scales $`<x_b`$. For the parameters we are considering, the velocity scale corresponding to $`x_b`$ is $``$ 10–30 km s<sup>-1</sup>. Hence, for the case where $`v=100`$ km s<sup>-1</sup>, the evolution will not be affected significantly by the Jeans length. We studied the evolution at a higher velocity scale ($`v=250`$ km s<sup>-1</sup>), and found that the effect of $`x_b`$ was even less (the relative change in $`q`$ was about 3 per cent at $`v=250`$ km s<sup>-1</sup>).
The effect of $`\gamma `$ is threefold here – it affects the value of $`x_b`$, the evolution of $`x_b`$ (more precisely, this also depends on whether we evolve $`T_m`$ or not) and the neutral hydrogen density fluctuations (through the value of $`\beta `$). We have already seen that the effect of changing the value of $`x_b`$ is not very important. Increasing the value of $`\gamma `$ will make the evolution of $`x_b`$ more rapid. Since $`x_b`$ appears in the denominator of the integrand in equation (LABEL:eq:q\_los\_delz), the higher the value of $`\gamma `$, the more rapid is the decrease of $`\xi _{\mathrm{LOS}}`$ with increasing redshift. This feature alone will increase the value of $`q`$ with increasing $`\gamma `$. But, actually $`q`$ decreases when we increase $`\gamma `$ because of the third effect of $`\gamma `$ – increasing $`\gamma `$ reduces the value of $`\beta `$ (equation (27)), and $`\beta `$ appears in the exponential in equation (37). Hence, the $`\xi _{\mathrm{LOS}}(\overline{z})`$ curves will decrease less rapidly when $`\beta `$ is small, i.e., $`\gamma `$ is large.
Finally, we note that the effect of whether $`T_m`$ evolves or not becomes appreciable for large values of $`\gamma `$. This is obvious because larger the value of $`\gamma `$, more rapid is the evolution of $`T_m`$. Furthermore, we have already argued that if the evolution of $`T_m`$ is more rapid, then the value of $`q`$ should increase (provided, of course, the value of $`\beta `$ remains unchanged). Hence, the values of $`q`$ are smaller when $`T_m`$ is kept constant than when $`T_m`$ has a redshift dependence. This can be verified from Table 2.
These effects are shown in Figure 4. The curves are normalised in such a way that $`\xi _{\mathrm{LOS}}(v=100`$ km s<sup>-1</sup>) = 0.21 at $`\overline{z}=3.85`$, which is taken from Cristiani et al.(1997). In the top row of Figure 4, the value of $`T_m`$ is fixed at a particular redshift (in this case, at $`\overline{z}=2.5`$) and the value of $`T_m`$ at other redshifts are calculated using the relation $`T_m(1+z)^{3\gamma 3}`$. In the bottom row, $`T_m`$ is kept constant. We do not plot the effect of $`T_m`$ because we find it to be very weak.
The point to be noted here is that the effects due to change in cosmology and $`\gamma `$ are of the same order, which means that we cannot constrain both the parameters simultaneously. Thus the above analysis clearly suggest that it will be very difficult to recover the power spectrum of density fluctuations uniquely from the Ly$`\alpha `$ absorption lines without knowing the IGM parameters. However, the cosmological parameters determined through studies such as CMBR (and other data), can be used to constrain the equation of state (using the plot in the centre in Figure 4), provided we have some idea about the evolution of the Jeans length or equivalently, $`T_m`$.
Throughout this paper we have treated $`\gamma `$ to be independent of $`z`$. However, there are indications that $`\gamma `$ could change with $`z`$. Schaye et al. (1999b) notice that the temperature of the IGM has a peak at $`z3`$ and it decreases with decreasing redshift afterwards. They also notice that the slope of the equation of state become close to one at $`z3`$ then increases with decreasing redshift. Theoretical calculations suggest that $`\gamma `$ increases with time and the rate of evolution depends on the reionisation epoch (Hui & Gnedin 1997). From Table 2 and Figure 4 we can infer that when $`\gamma `$ becomes larger the rate of growth of $`\xi _{\mathrm{LOS}}`$ at a given velocity decreases. Thus our study clearly suggests that the evolution of $`\xi _{\mathrm{LOS}}`$ at a given velocity can be used as probe of a reionisation and thermal history of the IGM once the cosmological model and the evolution of $`x_b`$ is fixed. We hope to study this in detail in a future publication.
### 4.2 Transverse Correlation
In this section we present the results for the transverse correlation. As before, we consider the same CDM power spectrum, and essentially the same range of the IGM parameters. The smoothing velocity is taken to be 30 km s<sup>-1</sup>, which is the typical peculiar velocity of a blob in the IGM.
Given $`z`$, we calculate $`\xi _{\mathrm{trans}}`$ as a function of the transverse comoving distance $`l_{}`$. One can then convert this length scale to an angular scale $`\theta `$ through the following relations.
$$\theta =\frac{l_{}}{d_a^{com}(z)}$$
(65)
$$d_a^{com}(z)=\frac{c}{H_0\sqrt{|\mathrm{\Omega }_k|}}S_k\left(x(z)\frac{H_0}{c}\sqrt{|\mathrm{\Omega }_k|}\right)$$
(66)
where $`\mathrm{\Omega }_k`$ is given by equation (41), $`x(z)`$ is given by equation (45) and
$$S_k(r)=\{\begin{array}{cc}\mathrm{sin}r\hfill & \text{(if }\mathrm{\Omega }_k<0\text{)}\hfill \\ r\hfill & \text{(if }\mathrm{\Omega }_k=0\text{)}\hfill \\ \mathrm{sinh}r\hfill & \text{(if }\mathrm{\Omega }_k>0\text{)}\hfill \end{array}$$
(67)
Instead of plotting the correlation function directly, we plot the quantity $`𝒫_\theta (\theta )`$, which is defined as follows. The excess probability, over random background, of finding two neutral hydrogen overdense regions separated by a comoving transverse distance $`l_{}`$ is
$$𝒫_l_{}(l_{})\mathrm{d}l_{}=\frac{\xi _{\mathrm{trans}}(l_{};z)2\pi l_{}\mathrm{d}l_{}}{4\pi [d_a^{com}(z)]^2}.$$
(68)
Using equation (65) we get the excess probability over random background of finding two neutral hydrogen overdense regions separated by an angle $`\theta `$ as
$`𝒫_\theta (\theta )\mathrm{d}\theta `$ $`=`$ $`𝒫_l_{}(l_{}){\displaystyle \frac{\mathrm{d}l_{}}{\mathrm{d}\theta }}\mathrm{d}\theta `$ (69)
$`=`$ $`{\displaystyle \frac{1}{2}}\xi _{\mathrm{trans}}(\theta )\theta \mathrm{d}\theta .`$
From Figure 5 it is clear that even the maximum excess probability of finding two H i overdense regions over an angular scale greater than few arc seconds is less than 1 per cent. Observationally the distribution of H i along the transverse direction is probed by studying the common absorbers along the LOS towards closely spaced QSOs. The angular scales probed varies between few arc seconds and few arc minutes (Shaver, Boksenberg & Robertson 1982; Shaver & Robertson 1983; Smette et al. 1992; Dinshaw et al. 1994; Bechtold et al. 1994; Crotts et al. 1994; Bechtold & Yee 1994; Smette et al. 1995; D’Odorico et al. 1997; Petitjean et al. 1998). Based on our analysis it is most likely that the common absorbers seen in the spectra of closely spaced QSOs are most likely probe the transverse extent of the same overdense region rather than the clustering length scale of separate regions.
### 4.3 Difference between $`\xi _{\mathrm{trans}}`$ and $`\xi _{\mathrm{LOS}}`$
It should be noted that for a given mean redshift, the values of the LOS and the transverse correlation functions need not be the same. This is because, when we observe along one LOS, we actually sample different points at different redshifts. In contrast to this, the transverse correlation is calculated at the same redshift. The effect of evolution in the LOS correlation makes it different from the transverse correlation.
To illustrate this point more clearly, let us first assume that $`x_b`$ does not evolve with $`z`$. Then from equations (LABEL:eq:q\_los) and (LABEL:eq:q\_sp), we see that for a given length scale $`l`$, the integrands in the two equations are identical. Hence, we get
$$\frac{Q_{\mathrm{LOS}}}{Q_{\mathrm{trans}}}=\frac{D(\overline{z}+\mathrm{\Delta }z/2)D(\overline{z}\mathrm{\Delta }z/2)}{D^2(\overline{z})},$$
(70)
where $`\overline{z}`$ and $`\mathrm{\Delta }z`$ are defined in equations (42) and (49) respectively. The difference in the two correlation functions is now entirely due to the evolution of the power spectrum. Thus the two correlation functions will be nearly equal for small $`\mathrm{\Delta }z`$ but will start differing from each other for large $`\mathrm{\Delta }z`$. In the general case when $`x_b`$ evolves with $`z`$, the difference will be much more prominent.
This is indeed true, as one can see from Figure 6. For scales below 200 $`h^1`$ Mpc, the two correlation functions are nearly the same. But above such scales the two functions start differing appreciably. For the observations made in the scales of 10–100 $`h^1`$ Mpc, our analytical calculation shows that one should not see any appreciable difference between LOS and transverse correlations. This can be used as a important tool determining the power spectrum (provided, of course, we know the IGM parameters and the correlation function completely). As we have argued earlier, one cannot get the power spectrum from the LOS correlation. But the power spectrum can be obtained from the transverse correlation in usual manner. Since the two correlations are identical for scales upto 100 $`h^1`$ Mpc, one can start from the LOS correlation, replace it with the transverse correlation, and obtain the power spectrum.
### 4.4 Column Density Distribution
In this section we study the results for the column density distribution. As we have discussed in section 3, we shall consider two separate cases for the evolution of $`T_0`$ and $`T_m`$, i.e., (i) $`(T_0`$ and $`T_m)(1+z)^{3\gamma 3}`$ and (ii) $`T_0`$ and $`T_m`$ are independent of $`z`$. The comparison between the two cases is shown in Figure 7, where we plot the quantity $`f(N_{\mathrm{HI}})`$ for the two different cases. $`f(N_{\mathrm{HI}})`$ is related to $`(\mathrm{d}N_{\mathrm{pk}}/\mathrm{d}z\mathrm{d}N_{\mathrm{HI}})`$ through the relation (Bi & Davidsen, 1997)
$$f(N_{\mathrm{HI}})=(\mathrm{d}N_{\mathrm{pk}}/\mathrm{d}z\mathrm{d}N_{\mathrm{HI}})/(1+z).$$
(71)
As we can see from Figure 7, the difference between the two cases becomes more significant at higher column densities. We have checked and found that $`T_m`$ has very little effect on the column density distribution. The difference between the two cases at $`z=3.5`$ is because of the fact that the value of $`T_0`$ is different for the two cases at that redshift. However, we have already mentioned that any uncertainty in the knowledge of $`T_0`$ can be compensated to some extent by changing $`ϵ`$. Hence, for studying the column density distribution we shall consider only the case where both $`T_m`$ and $`T_0`$ evolve as $`(1+z)^{3\gamma 3}`$.
We shall now discuss the dependence of the column density distribution on the following parameters : (i) cosmological models, (ii) $`ϵ`$ and (iii) $`\gamma `$.
In what follows we try to get constraints on our model parameters using the observed column density distribution obtained from Hu et al. (1995) and Kim et al. (1997) at three mean redshifts 2.31, 2.85 and 3.35. In Figure 8 we have plotted the quantity $`f(N_{\mathrm{HI}})`$ for those redshifts. The observational data points are the points with errorbars in the figure.
For all the plots in Figure 8, we have fixed $`T_m(z=2.5)=40,000`$ K and $`T_0(z=2.5)=20,000`$ K. In the left most panel of each row in the figure we plot the predicted column density distribution for various cosmological models for a given set of IGM parameters ($`\gamma =1.2`$) and $`ϵ=0.3`$. We did not plot the LCDM2 model because it completely overlaps with the LCDM1 model. In all the redshift bins it is clear that the SCDM curves fall steeply at the higher column density end compared to other models. This is consistent with the results noted by Gnedin & Hui (1996). However, in our method of obtaining $`f(N_{\mathrm{HI}})`$, the SCDM model can be made to fit the data by slightly increasing the value of $`ϵ`$. The OCDM and LCDM1 curves fit the observed distribution upto $`\mathrm{log}(N_{\mathrm{HI}}/\mathrm{cm}^2)=14.2`$ for $`ϵ=0.3`$.
In the middle panel of each row in Figure 8 we plot the predicted distribution for the three assumed values of $`ϵ`$ for LCDM1 with $`\gamma =1.2`$. Increasing the value of $`ϵ`$ enhances the column density of a cloud (with fixed peak density), and consequently, the value of $`f(N_{\mathrm{HI}})`$ increases at the high column density region. It is clear that the observed distribution is consistent with $`ϵ=0.3`$. This means that, in the case of LCDM, the effective length of the overdense region is about (1/3) of the coherence scale $`R^{}`$. The value of $`R^{}`$ depends on the baryonic power spectrum, its typical value is of the order of few hundreds of Kpc; for LCDM model with $`\gamma =1.2,T_m(z=2.5)=40,000`$K, we get $`R^{}60`$$`h^1`$ Kpc. Hence the typical length of an individual overdense region is $`20`$$`h^1`$ Kpc. Zhang et al. (1998) infer that the typical LOS scale length of an absorbing cloud is $``$ 15–35 $`h^1`$ Kpc through their hydrodynamical simulations. Our values are consistent with theirs.
The effect of $`\gamma `$ on the column density distribution can be seen from the right most panel in Figure 8 where we plot the results of LCDM1 model for three values of $`\gamma `$. In the models with higher value of $`\gamma `$, the $`f(N_{\mathrm{HI}})`$ curve falls sharply at higher column density end at all redshifts, which is consistent with the conclusions of Hui et al. (1997). The reason for this sharp fall is because the neutral hydrogen density is less (for a given baryonic density) when the value of $`\gamma `$ is more and, hence there are less number of clouds with high column densities. At low redshifts, it is clear that lower values of $`\gamma `$ are preferred for $`ϵ=0.3`$. However, one can get the consistent fit for higher values of $`\gamma `$ by increasing the value of $`ϵ`$. For example, at $`\overline{z}=2.31`$, one can fit the data reasonably well with $`ϵ`$ in the range 0.3 to 0.5 and $`\gamma `$ in the range 1.2 to 1.6. But it is impossible to fit the data with $`\gamma >1.6`$, for any choice of $`ϵ`$ at $`\overline{z}=2.31`$. This implies that we cannot accommodate $`\gamma `$ larger than 1.6 at this redshift even by changing the values of $`T_0`$ or other IGM parameters. At $`\overline{z}=2.85`$, one can fit the observations for $`\gamma =1.7`$ (the highest value of $`\gamma `$ we are considering in this paper), by taking $`ϵ=0.4`$. For the case where $`\overline{z}=3.35`$, we can see that we can marginally fit the observations for whole range of $`\gamma `$ considered in this paper with $`ϵ=0.3`$ itself.
Finally, we comment on how our results compare with some of the hydrodynamical simulations. It has been observed through numerical simulations that the peak baryonic over-density and the neutral hydrogen column density are strongly correlated (Zhang et al. 1998; Davé et al. 1999). We tested equation (57) against the scatter plots showing the correlations. The parameters were taken to be: LCDM1 cosmological model with $`\gamma =1.2,ϵ=0.3`$ at $`\overline{z}=3`$. These are the parameters which fit the observations for the column density distribution reasonably well (see Figure 8). It turns out that the relation matches well with the median of the scatter plot in Figure 8 of Zhang et al. (1998). Davé et al. (1999) give an analytical formula which relates $`n_\mathrm{B}[\mathrm{peak}]/n_0`$ and $`N_{\mathrm{HI}}`$ by a power-law (see equation (7)). We find that one needs a much higher value of $`\gamma `$ ($`1.82`$) in our model to match the power-law index at redshift 3. Also, one needs a $`ϵ1`$ to match the overall scaling factor. If we take such values of $`ϵ`$ and $`\gamma `$, then our results for the column density distribution will not be able to match the observations. We believe that such a discrepancy arises because of the difference in methods used for assigning a column density to a cloud.
We next test the column density distributions against some of the recent simulation results. In the column density range $`10^{12.8}`$$`10^{14.3}`$ cm<sup>-2</sup>, one can fit a power-law of the form $`f(N_{\mathrm{HI}})N_{\mathrm{HI}}^{\beta _{\mathrm{HI}}}`$. The values of $`\beta _{\mathrm{HI}}`$ for our LCDM1 model with $`\gamma =1.2`$ and $`ϵ=0.3`$ are $`\beta _{\mathrm{HI}}=1.70\pm 0.02`$ for $`\overline{z}=2.31`$, $`\beta _{\mathrm{HI}}=1.65\pm 0.04`$ for $`\overline{z}=2.85`$, $`\beta _{\mathrm{HI}}=1.56\pm 0.04`$ for $`\overline{z}=3.35`$. When compared with fits to the observational points, the corresponding power-law indices are $`1.35\pm 0.03,1.39\pm 0.26,1.59\pm 0.13`$, respectively (Kim et al. 1997). All errors given here are $`2\sigma `$. Simulations of Zhang et al. (1997) using the SCDM model produce a $`\beta _{\mathrm{HI}}=1.39\pm 0.06`$ in the range $`2\times 10^{12}<N_{\mathrm{HI}}/\mathrm{cm}^2<10^{14}`$ for $`z=3`$. Our curves are also in quite good agreement with the P<sup>3</sup>M–SPH simulations (using SCDM model) by Theuns, Leonard & Efstathiou (1998) and Theuns et al. (1998) in the same column density and redshift ranges. More recently, Machacek et al. (2000) have performed hydrodynamical simulations for various cosmological models. The LCDM model ($`\mathrm{\Omega }_m=0.4,\mathrm{\Omega }_\mathrm{\Lambda }=0.6,h=0.65,\sigma _8=1.0`$) with $`\mathrm{\Omega }_\mathrm{B}h^2=0.015`$ gives $`\beta _{\mathrm{HI}}=1.61\pm 0.04`$ for $`z=2`$ and $`\beta _{\mathrm{HI}}=1.48\pm 0.04`$ for $`z=3`$ in the range $`10^{12.8}<N_{\mathrm{HI}}/\mathrm{cm}^2<10^{14.3}`$. They note that the power-law index is higher for the SCDM model, and lower for the OCDM model. This is also consistent with what we get from our model.
## 5 Conclusions
We have presented a simple analytic expression for the correlation function and the column density distribution for the low H i column density systems seen in the spectra of high redshift QSOs. We have used our results to get constraints on various cosmological and IGM parameters. We summarise our main results below.
1. One cannot rigourosly define a power spectrum from the LOS correlation function. However, since the LOS and transverse correlations are identical below scales of $`100`$ $`h^1`$ Mpc, it is possible to obtain the power spectrum by a Fourier transform of the LOS correlation function (provided the IGM parameters are known).
Previous studies have attempted to recover the power spectrum of density fluctuations from the observations of the IGM (Croft et al. 1998, Hui 1999, Croft et al. 1999). We show that it is difficult to recover a unique power spectrum from H i correlation function without constraining the IGM parameters, especially $`\gamma `$. We have shown that the shape of the LOS correlation function at a particular mean redshift becomes less and less sensitive to the DM power spectrum as the universe evolves after reionisation. We feel that the correct approach in studying this issue is to constrain the cosmological models using CMBR or supernovae data, and apply those constraints to study the IGM parameters using H i correlation functions.
2. The LOS correlation function and its evolution is much more sensitive to $`\gamma `$ than $`T_m`$. Using observations which give the shape of the LOS correlation at a particular redshift, one can constrain the value of $`\gamma `$, even with ill-constrained values of $`T_m`$, provided the background cosmology is known. However, the more uncertain is the value of $`T_m`$, the less constrained is $`\gamma `$. Carrying out such an exercise for different redshifts, it will be possible to constrain the evolution of $`\gamma `$. However, for this study, one needs accurate observational data at different redshift bins which are not affected by evolutionary effects. Thus one can use the study of correlation function as an independent method to constrain the reionisation history of the universe.
3. The analytic column density distribution for H i, like $`\xi _{\mathrm{LOS}}`$, is less sensitive to $`T_m`$ than $`\gamma `$. The distribution, when compared with observations, favours a lower value of $`\gamma `$, although at redshifts $`>2.5`$, one can marginally fit the observations with higher values of $`\gamma `$. Our model clearly rules out $`\gamma >1.6`$ at redshift 2.31.
## Acknowledgment
We gratefully acknowledge the support from the Indo-French Centre for Promotion of Advanced Research under contract No. 1710-1. TRC is supported by the University Grants Commission, India. We would like to thank the anonymous referee for suggestions which helped us to improve the clarity of the paper. We also thank Patrick Petitjean for useful comments.
## Appendix A Detailed calculations for the column density distribution
Suppose we are looking at the IGM along any one direction, at some redshift $`z`$. Then the linear density field $`\delta _\mathrm{B}^{(1D)}(x,z)`$ along that axis will be described by a one dimensional Gaussian random field, with a power spectrum
$$P_\mathrm{B}^{(1D)}(k,z)=\frac{1}{2\pi }D^2(z)_k^{\mathrm{}}dk^{}k^{}\frac{P_{\mathrm{DM}}(k)}{(1+x_b^2(z)k^2)^2}$$
(72)
From now on we shall derive all the expressions at a particular redshift $`z`$, and we shall not write the explicit $`z`$-dependence on the quantities.
To define the column density, we associate each local maximum or peak in the linear density field to a Ly$`\alpha `$ cloud. The column density corresponding to such a cloud is given by equation (56). We have expressed $`\delta _\mathrm{B}^{(1D)}`$\[peak\] in terms of $`N_{\mathrm{HI}}`$ in equation (58).
Using the properties of a Gaussian random field, we can derive the joint probability distribution for the three Gaussian random fields $`\delta _\mathrm{B}^{(1D)},\delta _\mathrm{B}^{(1D)\prime \prime },\delta _\mathrm{B}^{(1D)}`$. The probability that the field and its second derivative have values $`\delta _\mathrm{B}^{(1D)}`$ and $`\delta _\mathrm{B}^{(1D)\prime \prime }`$, respectively at the peak $`\delta _\mathrm{B}^{(1D)}=0`$ is
$`𝒫[\delta _\mathrm{B}^{(1D)},\delta _\mathrm{B}^{(1D)\prime \prime },\delta _\mathrm{B}^{(1D)}=0]\mathrm{d}\delta _\mathrm{B}^{(1D)}\mathrm{d}\delta _\mathrm{B}^{(1D)\prime \prime }|\delta _\mathrm{B}^{(1D)\prime \prime }|\mathrm{d}x=`$ (73)
$`{\displaystyle \frac{1}{(2\pi )^{3/2}\sigma _1\mathrm{\Sigma }}}\mathrm{exp}[{\displaystyle \frac{1}{2\mathrm{\Sigma }^2}}(\sigma _2^2\delta _\mathrm{B}^{(1D)^2}+2\sigma _1^2\delta _\mathrm{B}^{(1D)}\delta _\mathrm{B}^{(1D)\prime \prime }`$
$`+\sigma _0^2\delta _\mathrm{B}^{(1D)^2})]\mathrm{d}\delta _\mathrm{B}^{(1D)}\mathrm{d}\delta _\mathrm{B}^{(1D)\prime \prime }|\delta _\mathrm{B}^{(1D)\prime \prime }|\mathrm{d}x,`$
where
$$\sigma _m^2\sigma _m^2(z)=\frac{1}{2\pi }_{\mathrm{}}^{\mathrm{}}dkk^{2m}P_\mathrm{B}^{(1D)}(k,z),$$
(74)
and
$$\mathrm{\Sigma }^2=\sigma _0^2\sigma _2^2\sigma _1^4.$$
(75)
Note that $`\sigma _0=\mathrm{\Delta }`$, defined in equation (8).
For convenience, let us define some dimensionless quantities
$$\nu \frac{\delta _\mathrm{B}^{(1D)}}{\sigma _0},\lambda \frac{\delta _\mathrm{B}^{(1D)\prime \prime }}{\sigma _2},\kappa \frac{\sigma _1^2}{\sigma _0\sigma _2}$$
(76)
$`\nu `$ and $`\lambda `$ measure the field and its second derivative, respectively; $`\kappa `$ is a measure of the width of the power spectrum. One can use these quantities to obtain the number of peaks (clouds) per unit length
$`{\displaystyle \frac{\mathrm{d}N_{\text{pk}}}{\mathrm{d}x}}=`$ (77)
$`𝒫[\delta _\mathrm{B}^{(1D)},\delta _\mathrm{B}^{(1D)\prime \prime },\delta _\mathrm{B}^{(1D)}=0]\mathrm{d}\delta _\mathrm{B}^{(1D)}\mathrm{d}\delta _\mathrm{B}^{(1D)\prime \prime }|\delta _\mathrm{B}^{(1D)\prime \prime }|`$
Using equation (39), one can convert the above expression to the number of clouds per unit redshift interval. After simplification, the relation becomes
$`{\displaystyle \frac{\mathrm{d}N_{\mathrm{pk}}}{\mathrm{d}z}}={\displaystyle \frac{d_H(z)}{(2\pi )^{3/2}\sqrt{1\kappa ^2}R^{}}}\times `$ (78)
$`\mathrm{exp}\left[{\displaystyle \frac{1}{2}}\left\{{\displaystyle \frac{(\nu \kappa \lambda )^2}{1\kappa ^2}}+\lambda ^2\right\}\right]\lambda \mathrm{d}\lambda \mathrm{d}\nu `$
The $`\lambda `$ integration can be carried out to obtain
$`{\displaystyle \frac{\mathrm{d}N_{\text{pk}}}{\mathrm{d}z\mathrm{d}\nu }}={\displaystyle \frac{d_H(z)}{(2\pi )^{3/2}R^{}}}[\sqrt{1\kappa ^2}\mathrm{exp}({\displaystyle \frac{\nu ^2}{2(1\kappa ^2)}})`$ (79)
$`+\kappa \nu \sqrt{2\pi }\mathrm{e}^{\nu ^2/2}`$
$`\sqrt{{\displaystyle \frac{\pi }{2}}}\kappa \nu \text{erfc}\left({\displaystyle \frac{\kappa \nu }{\sqrt{2(1\kappa ^2)}}}\right)\mathrm{e}^{\nu ^2/2}]`$
We are interested in the quantity
$$\frac{\mathrm{d}N_{\mathrm{pk}}}{\mathrm{d}z\mathrm{d}N_{\mathrm{HI}}}=\frac{\mathrm{d}N_{\mathrm{pk}}}{\mathrm{d}z\mathrm{d}\nu }\frac{\mathrm{d}\nu }{\mathrm{d}N_{\mathrm{HI}}}$$
(80)
which is straightforward to obtain from equation (79), provided we know $`\nu `$ as a function of $`N_{\mathrm{HI}}`$. Equations (58) and (76) give $`\nu `$ in terms of $`N_{\mathrm{HI}}`$, and they can be used to calculate $`\mathrm{d}\nu /\mathrm{d}N_{\mathrm{HI}}`$ (for $`n_{\mathrm{HI}}n_\mathrm{B}`$)
$$\frac{\mathrm{d}\nu }{\mathrm{d}N_{\mathrm{HI}}}=\frac{1}{\beta \mathrm{\Delta }N_{\mathrm{HI}}}.$$
(81)
Thus we get an analytic expression for the number of clouds per unit redshift interval per unit column density range $`(\mathrm{d}N_{\mathrm{pk}}/\mathrm{d}z\mathrm{d}N_{\mathrm{HI}})`$ as a function of $`N_{\mathrm{HI}}`$.
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# 1 Introduction
## 1 Introduction
The microphysics of particle acceleration in the solar system is a fascinating realm and provides useful paradigms for particle acceleration in other places of the Universe. From the observed flux of cosmic rays it is clear that much more powerful particle accelerators than main sequence stars and their winds must be at work in our Galaxy and in other galaxies. Supernova remnants \[Ellison et al. (1997)\] and radio galaxies \[Biermann (1995)\] are among the most interesting putative sources of cosmic rays. The cosmic ray pressure is rather large and compares with that of the interstellar magnetic fields implying that cosmic rays are an important ingredient in the dynamics of the interstellar medium, and possibly also in the dynamics of the intergalactic medium (in superclusters). This imposes strong constraints on the energetics of possible cosmic ray sources from which one can infer the particle acceleration efficiency. For instance, supernova blast waves can supply the flux of cosmic rays up to the so-called knee in their spectrum at an energy of $`10^{15}`$ eV only if their particle acceleration efficiency is as large as $`13\%`$ \[Drury (1990)\] which basically invalidates simple test-particle approaches for the description of the acceleration mechanism. This has enforced the two-fluid theory for shock acceleration in which the momentum flux due to accelerated particles is self-consistently included into the dynamics of the shock wave \[Achterberg et al. (1984)\]. In this contribution it is argued that if radio galaxies are the most powerful particle accelerators in the Universe responsible for the diffuse isotropic gamma-ray background, and possibly also for the ultrahigh-energy cosmic rays up to at least $`5\times 10^{19}`$ eV (the so-called Greisen-Zatsepin-Kuzḿin cutoff), they require a particle acceleration efficiency even larger than that of supernova remnants.
## 2 Energy density of the diffuse isotropic $`\gamma `$-ray background
The energy losses of relativistic particles inevitably lead to $`\gamma `$-radiation by which the acceleration sites can be traced. The cumulative flux from all unresolved cosmic accelerators appears as an diffuse isotropic flux in an inertial frame. A diffuse high-latitude $`\gamma `$-ray background has been measured with CGRO from MeV to $`100`$ GeV photon energies which very likely arises from unresolved extragalactic sources \[Sreekumar et al. (1998)\]. The resolved sources already contribute $`15\%`$ of the diffuse isotropic flux and show the same average spectral shape as the diffuse background. The background spectrum above 30 MeV as determined with EGRET on board CGRO is given by
$$\frac{dN}{dE}=(7.32\pm 0.34)\times 10^9\left(\frac{E}{451\mathrm{MeV}}\right)^{2.1\pm 0.03}\mathrm{cm}^2\mathrm{s}^1\mathrm{MeV}^1\mathrm{sr}^1.$$
(1)
The spectrum continues into the MeV range. Most of the background flux below 10 MeV can be explained by the contribution of Supernovae Ia from the era of galaxy formation at $`z_\mathrm{f}=25`$ \[The et al. (1993)\]. The uncertainties in this thermal component of the diffuse $`\gamma `$-ray background are not very large, the known element abundances tightly constrain the total number of SNIa, and hence the total energy flux in $`\gamma `$-rays, in the Universe leaving little room for other sources contributing in the MeV range. The uncertainty in the cosmic star formation history due to possible dust enshrouding of the earliest galaxies should amount to less than a factor of three. Integrating the above non-thermal background spectrum between 10 MeV and 30 GeV (the flux at higher energies must be considered somewhat unreliable due to the strongly decreasing aperture of EGRET) one obtains the energy density
$$u_\gamma =(5.03\pm 0.64)\times 10^6\mathrm{eVcm}^3$$
(2)
or
$$\mathrm{\Omega }_\gamma h^2=(2.09\pm 0.27)\times 10^{10}$$
(3)
where $`h=H_{}/100\mathrm{km}\mathrm{s}^1\mathrm{Mpc}^1`$ parametrizes the Hubble constant.
## 3 Metagalactic gamma-ray absorption and production
Calorimetry of the particles accelerated throughout the Universe requires to know whether all $`\gamma `$-rays emitted by a source also reach the present-day observer. However, $`\gamma `$-rays of energy $`E`$ can interact with low-energy photons of energy $`ϵ`$ from the diffuse isotropic background over cosmological distance scales $`l`$ producing electron-positron pairs $`\gamma +\gamma e^++e^{}`$, if their energy exceeds the threshold energy
$$ϵ_{\mathrm{th}}=\frac{2(m_\mathrm{e}c^2)^2}{(1\mu )(1+z)^2E}1\left(\frac{1+z}{4}\right)^2\left(\frac{E}{30\mathrm{GeV}}\right)^1\mathrm{eV}(\mu =0)$$
(4)
where $`\mu `$ denotes the cosine of the scattering angle \[Gould & Shréder (1966), Coppi & Aharonian (1997)\]. The $`\gamma `$-ray attenuation $`e^\tau `$ due to pair production becomes important if the mean free path $`\lambda `$ becomes smaller than $`l`$, i.e. if the optical depth obeys $`\tau =l/\lambda 1`$. For the computation of $`\tau `$ one first needs to know the pair production cross section
$$\sigma _{\gamma \gamma }=\frac{3\sigma _\mathrm{T}}{16}(1\beta ^2)\left[2\beta (\beta ^22)+(3\beta ^4)\mathrm{ln}\left(\frac{1+\beta }{1\beta }\right)\right]$$
(5)
where $`\beta =\sqrt{11/\gamma ^2}`$ with $`\gamma ^2=ϵ/ϵ_{\mathrm{th}}`$, and where $`\sigma _\mathrm{T}`$ denotes the Thomson cross section \[Jauch & Rohrlich (1976)\]. Then one needs the geodesic radial displacement function $`dl/dz=\frac{c}{H_{}}[(1+z)\overline{E}(z)]^1`$ to compute the line integral from $`z=0`$ to some $`z=z_{}`$. For a cosmological model with $`\mathrm{\Omega }=1`$ and $`\mathrm{\Lambda }=0`$ the function $`\overline{E}(z)`$ simplifies to $`(1+z)^{3/2}`$. Hence one obtains the optical depth
$`\tau _{\gamma \gamma }(E,z_{})={\displaystyle _0^z_{}}𝑑z{\displaystyle \frac{dl}{dz}}{\displaystyle _1^{+1}}𝑑\mu {\displaystyle \frac{1\mu }{2}}{\displaystyle _{ϵ_{\mathrm{th}}}^{\mathrm{}}}𝑑ϵn_\mathrm{b}(ϵ)(1+z)^3\sigma _{\gamma \gamma }(E,ϵ,\mu ,z)`$
$`={\displaystyle \frac{c}{H_{}}}{\displaystyle _0^z_{}}𝑑z(1+z)^{1/2}{\displaystyle _0^2}𝑑x{\displaystyle \frac{x}{2}}{\displaystyle _{ϵ_{\mathrm{th}}}^{\mathrm{}}}𝑑ϵn_\mathrm{b}(ϵ)\sigma _{\gamma \gamma }(E,ϵ,x1,z)`$ (6)
adopting a non-evolving present-day background density $`n_\mathrm{b}`$, i.e. $`n_\mathrm{b}^{}(z,ϵ^{})dϵ^{}=(1+z)^3n_\mathrm{b}(ϵ)dϵ`$ where the dash indicates comoving-frame quantities. The simplifying assumption that the photon density transforms geometrically corresponds to a situation in which an initial short burst of star formation at $`z_\mathrm{f}>z_{}`$ produced most of the diffuse infrared-to-ultraviolet background radiation. Fig.1 shows the spectrum of the low-energy diffuse background used to solve Eq.(6) numerically.
Figure 2 shows the resulting $`\tau (E,z)=1`$ (omitting the subscript hereafter) curve for the microwave-to-ultraviolet diffuse background spectrum shown in Fig.1. It is obvious that $`\gamma `$-rays above $`1050`$ GeV can not reach us from beyond redshifts of $`z=z_\mathrm{f}=24`$. Higher energy $`\gamma `$-rays can reach us only from sources at lower redshifts (e.g. $`\gamma `$-rays with energies up to 10 TeV have been observed from Mrk 501 at $`z=0.033`$ in accord with Fig.2 \[Mannheim (1998)\]).
Corollary I: If the DIGB originates from unresolved sources distributed in redshift similar to galaxies, its spectrum must steepen above $`30`$ GeV due to $`\gamma `$-ray pair attenuation.
Here is has been tacitly assumed that the $`\gamma `$-rays which have turned into electron-positron pairs do not show up again. This is, in fact, not quite true, since the pairs are subject to inverse-Compton scattering off the microwave background thereby replenishing $`\gamma `$-rays. The 2.7 K background is more important as a target than the shorter wavelength background, since there is no threshold condition for Thomson scattering contrary to pair production and since 2.7 K photons greatly outnumber the latter. The inverse-Compton scattered microwave photons turn into $`\gamma `$-rays of energy
$$E_{\mathrm{ic}}10\left(\frac{1+z}{4}\right)\left(\frac{E}{30\mathrm{GeV}}\right)^2\mathrm{MeV}$$
(7)
conserving the energy of the absorbed $`\gamma `$-ray which corresponds to a constant $`E^2dN/dE`$, i.e. the expected slope of the differential spectrum is about -2 (-2.1 observed). A negligible amount of energy is lost to lower frequency synchrotron emission, if magnetic fields are present in the interagalactic medium.
Corollary II: Energy conservation in the reprocessing of $`\gamma `$-rays from higher to lower energies by pair production and subsequent inverse-Compton scattering produces an approximate $`dN/dEE^2`$ power law DIGB between $`10`$ MeV and $`30`$ GeV.
## 4 Particle acceleration efficiency
The energy density of the accelerated particles in the sources ($`u_{\mathrm{acc}}`$) must obey $`u_{\mathrm{acc}}u_\gamma `$ for consistency. Defining the radiative efficiency
$$\xi _{\mathrm{rad}}=\frac{u_\gamma }{u_{\mathrm{acc}}}=\mathrm{Min}[1,\frac{t_\mathrm{a}}{t_\mathrm{c}}]$$
(8)
where $`t_\mathrm{a}`$ and $`t_\mathrm{c}`$ denote the adiabatic and energy loss time scales, respectively, the energy requirement for the sources is minimized for $`\xi _{\mathrm{rad}}=1`$. In extragalactic radio sources, $`\xi _{\mathrm{rad}}1`$ corresponds to very large electron ($`\gamma _\mathrm{e}10^5`$) or proton Lorentz factors ($`\gamma _\mathrm{p}10^{10}`$) consistent with the non-thermal $`\gamma `$-ray spectra extending above 10 GeV.
The resolved extragalactic $`\gamma `$-ray sources belong to the so-called blazar subclass of radio-loud active galactic nuclei (AGN) in which one observes jets emerging from accreting supermassive black holes preferentially at small angles to their relativistic velocity vector. If the entire class of radio-loud AGN is responsible for the DIGB, the particle acceleration efficiency is constrained by
$$\xi _{\mathrm{acc}}=\frac{u_{\mathrm{acc}}}{u_\mathrm{j}}=\frac{u_\gamma }{\xi _{\mathrm{rad}}u_\mathrm{j}}$$
(9)
where $`u_\mathrm{j}`$ denotes the total (kinetic + magnetic + randomized relativistic particle) energy density in extragalactic jets. Since $`u_\gamma `$ is known (Eq.(1)), a determination of $`u_\mathrm{j}`$ is necessary to further constrain $`\xi _{\mathrm{acc}}`$. This energy density can be inferred from the relative strengths of the various bumps in the overall diffuse background spectrum (Fig.3) in the following way.
## 5 Origins of the diffuse isotropic background radiation
The microwave bump is the well-known signature of the big bang at the time of decoupling with its energy density given by the Stefan-Boltzmann law $`u_{3\mathrm{K}}=\sigma T^4`$. The bump in the far-infrared is due to star formation in early galaxies (at $`z=z_\mathrm{f}`$), since part of the stellar light, which is visible as the bump at visible wavelengths, is reprocessed by dust obscuring the star-forming regions. The energy density of the two bumps can be inferred from the present-day heavy element abundances. Heavy elements have a mass fraction $`Z=0.03`$ and were produced in early bursts of star formation by nucleosynthesis with radiative efficiency $`ϵ=0.007`$ yielding
$$u_{\mathrm{ns}}\frac{\rho _{}Zϵc^2}{1+z_\mathrm{f}}6\times 10^3\left(\frac{\mathrm{\Omega }_{}h^2}{0.01}\right)\left(\frac{1+z_\mathrm{f}}{4}\right)^1\mathrm{eV}\mathrm{cm}^3$$
(10)
About half of the energy is contained in either bump. It has been shown recently by a number of groups that probably all galaxies (except dwarfs) contain supermassive black holes in their centers which are actively accreting over a fraction of $`t_{\mathrm{agn}}/t_{}10^2`$ of their lifetime implying that the electromagnetic radiation released by the accreting black holes amounts to
$$u_{\mathrm{accr}}\frac{ϵ_{\mathrm{accr}}M_{\mathrm{bh}}}{ZϵM_{}}\frac{t_{\mathrm{agn}}}{t_{}}u_{\mathrm{ns}}1.4\times 10^4\mathrm{eV}\mathrm{cm}^3$$
(11)
adopting the accretion efficiency $`ϵ_{\mathrm{accr}}=0.1`$ and the black hole mass fraction $`M_{\mathrm{bh}}/M_{}=0.005`$ \[Rees & Silk (1998)\]. Most of the accretion power emerges in the ultraviolet where the diffuse background is unobservable owing to photoelectric absorption by the neutral component of the interstellar medium. However, a fraction of $`u_\mathrm{x}/u_{\mathrm{bh}}20\%`$ (from the average quasar spectral energy distribution \[Sanders et al. (1989)\]) shows up in hard X-rays due to coronal emission from the accretion disk to produce the diffuse isotropic X-ray background bump with $`u_\mathrm{x}2.8\times 10^5\mathrm{eV}\mathrm{cm}^3`$ \[Gruber (1992)\].
Jets with non-thermal ($`\gamma `$-ray) emission show up only in the radio-loud fraction $`\xi _{\mathrm{rl}}20\%`$ of all AGN and their kinetic power roughly equals the accretion power \[Rawlings & Saunders (1991)\]. Hence one obtains for the energy density in extragalactic jets
$$u_\mathrm{j}=\left(\frac{\xi _{\mathrm{rl}}}{0.2}\right)u_{\mathrm{accr}}\left(\frac{\xi _{\mathrm{rl}}}{0.2}\right)2.8\times 10^5\mathrm{eV}\mathrm{cm}^3$$
(12)
Substituting this value into Eq.(9) one obtains a limit for the acceleration efficiency
$$\xi _{\mathrm{acc}}0.18\left(\frac{\xi _{\mathrm{rad}}}{1.0}\right)^1\left(\frac{\xi _{\mathrm{rl}}}{0.2}\right)^1$$
(13)
which has to be compared with the 13% efficiency required for supernova remnants. Since a radiative efficiency much lower than unity must be considered realistically (e.g., to account for the pdV work of the jets against a surrounding intracluster medium), the acceleration efficiency may have to be as large as $`30\%`$ or more. Possible escapes to this conclusion are (i) the radio-loud fraction was much larger at high redshifts than at low redshifts, (ii) there is another contributor to the gamma-ray background, such as the decay of topological defects \[Sigl (1996)\], or (iii) the flux of the DIGB is actually lower due to the foreground from a galactic halo \[Dixson et al. (1998)\].
## 6 A remark on cosmic rays from radio galaxies
An extragalactic cosmic ray component seems to be evident from the observed change in slope of the local spectrum above $`3\times 10^{18}`$ eV. The slope of this extragalactic component is much steeper than $`E^2`$ consistent with the steepening due to energy losses for a cosmologically distributed and evolving source population, such as extragalactic radio sources \[Rachen et al. (1993)\]. The total energy density in the loss-steepened extragalactic cosmic ray component and the DIGB are of the same order of magnitude, as required if the $`\gamma `$-rays trace neutral pions and the cosmic rays trace neutrons produced in photo-meson production events.
## 7 Summary
The particle acceleration efficiency of the jets in radio galaxies must be larger than $`18\%`$ if the $`\gamma `$-rays from these jets are to explain the DIGB.
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# Bloch Walls and Macroscopic String States in Bethe’s solution of the Heisenberg Ferromagnetic Linear Chain
## Abstract
We present a calculation of the lowest excited states of the Heisenberg ferromagnet in 1-d for any wave vector. These turn out to be string solutions of Bethe’s equations with a macroscopic number of particles in them. These are identified as generalized quantum Bloch wall states, and a simple physical picture provided for the same.
Introduction The question of elementariness of excitations in low dimensional magnetic systems is receiving much attention currently. It was perhaps first addressed in the seminal work of Bethe in 1931. In addition to providing the celebrated Ansatz named after him, Bethe asked if Bloch’s magnons are the “most elementary” excitations in 1-d. He came to the conclusion that they were not, and instead found that the bound states of spin reversals were. After the original paper of Bethe, the ferromagnet has received comparatively less attention than its antipode, namely the antiferromagnet. One source of revival of interest in the ferromagnet is in connection with stochastic dynamical systems, albeit with a complex Aharonov Bohm magnetic flux. Another notable recent exception is a work by Sutherland, who shows that the excited states of the ferromagnet contain a singlet state at momentum $`\pi `$, with an excitation energy (EE) that is very low, of $`O(1/N)`$, where $`N`$ is the length of the ring.
At the semiclassical level, domain wall arguments lead one to expect in dimensions $`d`$ ( with volume $`=N^d`$) the Bloch wall excitations to be of $`O(N^{d2})`$ , and hence to be amongst contenders for the lowest EE in $`d=1`$. Such “large deviation” excitations carry spin as well as momentum as we show below. These configurations of spins will be discussed within the context of Bethe’s Ansatz (BA) for the $`s=1/2`$ ferromagnet presently.
In this work we ask ( and answer) the following question: For a given value of the total momentum or total spin of the Bethe ferromagnet, what is the lowest excited state? The nontriviality of the question arises from the fact that within the famous Bethe formula for the boundstate of $`n`$ magnons, $`\omega _{Bethe}(q)=J\frac{2}{n}(1\mathrm{cos}(q))`$, the lower limit on the total momentum $`q`$ depends implicitly upon $`n`$. Its dependence has not been fully explicated earlier, at least as far as we could find in the literature. In this work, we use a combination of exact diagonalization and analytic methods to attack the problem. A new and essential tool that we develop is the argument of continuity of certain regular root solutions with respect to the density, regarded as a continuous variable, leading to a differential equation formulation of Bethe’s equations (BE) that bypasses the knowledge ( or otherwise ) of the quantum numbers.
Our findings are readily stated: The lowest excitations for any momentum $`q`$ arise from special string solutions of BE. These special string solutions involve a macroscopic number of particles in a given string, and hence we call them macroscopic strings. Sutherland’s solution at momentum $`\pi `$ is a particular case of these. We have found the lowest solution for every value of the total momentum, these correspond to a definite value of total spin as well. Such states are of the type that one would expect from Bethe’s formula for $`n`$ magnon bound states, with $`n`$ of the size of the lattice. The formula of Bethe cannot, however, be used for such large bound states since we show that there are significant corrections to the traditional assumption of a uniformly spaced vertical Bethe string in the complex plane: the curvature and nonuniformity of spacing produces essential differences. Our states can be represented by a new formula: $`\omega _{BW}(q)=\frac{2\pi }{N}Jq(1q/(2\pi ))`$ for $`0q\pi `$. We find that the solution for a given $`q`$ corresponds to a particular value of total spin, $`S_{tot}=N(1/2q/2\pi )`$. If we write $`n`$ the spin deviation from the saturated ferromagnet ( $`S^z=N/2n`$) in terms of the density $`d`$ as $`n=dN`$, the spectrum can be written as $`\omega _{BW}=\frac{4\pi ^2}{N}Jd(1d)`$ with $`q=2\pi d`$ and $`S_{tot}=N(1/2d)`$. We show finally that these states correspond to generalized quantum Bloch walls: thus the states with lowest EE at any wave vector of $`O(1)`$ are Bloch walls.
Before presenting our calculations, we also note that the bound states of Bethe for $`s=1/2`$ have been identified recently with solitonic excitations of the non linear classical, i.e. large $`s`$, Landau Lifschitz equations $`\dot{\stackrel{}{S}}(x)=\stackrel{}{S}(x)\times ^2/x^2\stackrel{}{S}(x)`$. Several explicit solutions of these are known . We note that Bloch walls form a certain class of exact solutionsof the nonlinear classical equations , namely non linear spin waves.
Calculations: We consider the ferromagnetic Heisenberg Hamiltonian: $`=2J_{l=1,N}[s^2\stackrel{}{S}_l\stackrel{}{S}_{l+1}]`$, where for most part we consider $`s=1/2`$ and $`\stackrel{}{S}_l=\frac{1}{2}\stackrel{}{\sigma }_l`$, in terms of the usual Pauli spin operators.
Our first result is from observations on wavefunctions of the energy eigenstates obtained from exact numerical diagonalization of chains upto length $`N=16`$ in a momentum resolved basis. We observed that in all cases, the lowest state at momentum $`q=2\pi n/N`$ has total angular momentum $`S=N/2n`$. This implies that this state can be obtained in the $`n`$ particle sector where it is a “maximal”, i.e. highest weight state. The corresponding Bethe wavefunction has all pseudo momenta nonvanishing, and the momentum and particle density related as $`q=2\pi d`$.
We next write the BE for the Heisenberg chain in the Orbach parametrization as:
$`Nf(\alpha _l)=2\pi I_l+{\displaystyle \underset{ml}{}}f((\alpha _l\alpha _m)/2)l=1,2\mathrm{}n`$ (1)
where $`f(x)=\frac{1}{i}\mathrm{log}(\frac{x+i}{xi})=2\text{ArcCot}(x)`$, $`\alpha _l=\mathrm{cot}(k_l/2)`$ and $`\{I_l\}`$ are the Bethe integers, and $`k_l`$ the Bethe pseudo momenta. We take the branch cut of $`f(x)`$ to be on the imaginary axis running from $`i`$ to $`+i`$. The energy and total momentum are given respectively by: $`ϵ=J_{l=1,n}\frac{4}{1+\alpha _l^2}`$ and $`q=\frac{2\pi }{N}_{l=1,n}I_l`$. Let us note that we are interested in the lowest energy states for a given $`q`$, requiring a knowledge of the integers $`I_l`$. These integers, as shown by Bethe, differ by 2 for scattering states, and by either one or zero for bound states in general. Eliminating the integer sets with zeros in them, a very plausible state is one with $`I_l=1`$ for $`1ln`$, and indeed we found from numerical studies of BE for small N and small n ( with $`n<<N`$), that this was indeed so: the resulting state is invariably the lowest energy state for small $`q`$. Emboldened by this exercise we found the following exact solution analytically in the limit of a thermodynamic $`n`$ as well as $`N`$, but at low density i.e. $`d=n/N<<1`$ . Since we are interested in excitation energies that are vanishing in the thermodynamic limit, the corresponding variables $`\alpha _l`$ scale with system size and it is convenient to introduce new scaled variables $`z_l=\alpha _l/n`$. In terms of these, the lhs of Eq(1) becomes $`2/(z_ld)`$ on using the large $`x`$ expansion of $`f(x)`$ and ignoring terms of size $`O(1/N^2)`$. On the rhs we cannot make the expansion in general since, at high densities, a core is formed where the separations $`n(z_lz_m)/2`$ between pairs of particles can become arbitrarily close to the branch points of $`f(x)`$. In the low density limit we find it is possible to obtain a perturbative solution, using the crucial observation that the typical interparticle separation $`1/\sqrt{d}`$, and hence we can use the smallness of $`d`$ as an expansion parameter in a perturbative sense. Setting $`I_l=1`$ we find the approximate equation:
$`{\displaystyle \frac{1}{z_l}}=\pi d+{\displaystyle \frac{2d}{n}}{\displaystyle \underset{ml}{}}{\displaystyle \frac{1}{z_lz_m}}.`$ (2)
Note that Eq(2) has corrections from the expansion of the phase shift, that is typically of $`O(d)`$ smaller than the least term retained. The solution of the system Eq(2) can actually be found exactly, but we save it for a future publication. We find at low densitites the following result:
$`z_l={\displaystyle \frac{1}{\pi d}}+{\displaystyle \frac{i\sqrt{2}}{\pi \sqrt{d}}}x_l{\displaystyle \frac{2}{3\pi }}(x_l^2+1{\displaystyle \frac{1}{n}})+O(\sqrt{d}),`$ (3)
where $`x_l`$ satisfy $`H_n(\sqrt{n}x_l)=0`$, $`H_n`$ being the $`n`$th order Hermite polynomial. In the limit of large $`n`$ the $`x_j`$ form a continuum stretching from $`\sqrt{2}`$ to $`\sqrt{2}`$ with the familiar semicircular density of states $`\rho (x)=\frac{1}{\pi }\sqrt{2x^2}`$. This solution can be used to obtain the energy to order $`d^2`$. The energy $`Nϵ=4J/(nd)_i1/z_i^2`$ for low $`d`$ can be found as $`Nϵ=4J\pi ^2[d+d^2\{3<(\beta _i^{(1)})^2>2<\beta _i^{(2)}>\}+O(d^3)]`$, where the averages are normalized sums over the indicated variables. Using the explicit expression Eq(3) and converting the sums to integrals over the semicircular density of states we get finally the low density formula: $`Nϵ=4J\pi ^2d(1d)+O(d^3)`$. Below we will argue that there are no corrections to the above formula beyond the first term: it is exact!. Thus provisionally we write
$`Nϵ=4J\pi ^2d(1d)=2J\pi q(1{\displaystyle \frac{q}{2\pi }}).`$ (4)
We note that the low density Equation(2) must be abandoned once the minimum separation $`n(z_iz_j)`$ hits the value $`2i`$, this happens at $`2=|n(z_0z_i)|=\frac{\sqrt{2}}{\sqrt{d}\pi }n|(x_0x_1)`$. However $`n|(x_0x_1)|=1/\rho (0)=\pi /\sqrt{2}`$, thus $`d1/4`$. Indeed we found for small systems that $`d1/4`$ cannot be treated easily numerically: the new difficulty is that the quantum numbers are no longer simple as we discuss below. For $`d1/4`$ the low density result Eq(3) and the full solution of Eq(2) are extremely close.
We now discuss the techniques used for solving Eq(1) numerically at larger densities. There are two main problems. One is that the integers jump around in a complicated way that is not known beforehand and it is clearly not feasible to try all combinations. The second problem is the formation of the core, i.e. successive roots that are placed very close to a separation $`2i`$, which causes singularities in the equation and results in numerical inaccuracies. Our strategy is to start from low densities where we know the roots, and change $`d`$ slowly and study the evolution of the roots. For this we first convert Eq(1) into a set of first order ordinary differential equations (ODE) with $`\mathrm{log}(d)`$ as a time like flow-parameter. Taking the derivative of Eq(1) with respect to $`d`$ and defining new variables $`t=\mathrm{log}d`$ and $`f_l=z_le^t`$ we obtain:
$`{\displaystyle \underset{m}{}}`$ $`A_{lm}`$ $`\dot{f_m}=\pi f_l^2\dot{I_l}{\displaystyle \underset{m}{}}A_{lm}(f_lf_m),`$ (5)
$`\mathrm{where}`$ $`A_{lm}`$ $`={\displaystyle \frac{2e^tf_l^2}{n}}{\displaystyle \frac{1}{[4e^{2t}/n^2+(f_lf_m)^2]}},ml`$ (6)
$`A_{ll}`$ $`=1{\displaystyle \underset{ml}{}}A_{lm}`$ (7)
and the derivatives are with respect to the “time” variable $`t`$. The time derivatives of the integers are delta functions and hence drop out of the equations at almost all times. Also since the flow of the roots themselves is smooth, we expect the delta-function singularities to be precisely cancelled by other terms in the equation. Hence in evolving the above ODE we can drop the first term on the right hand side of the equation. This immediately solves the problem of our lack of knowledge of the integers since they do not occur anywhere else in the differential equations. From these solutions we can recover the integers and use them in the root-finder to get more accurate solutions. We find that till densities around $`d=0.45`$ the solutions obtained from the ODE are very accurate. At higher densities, significant numerical errors show up because of the singularities associated with the core. In these cases we correct our ODE solutions by fixing the core by hand and using the root-finder to self-consistently solve for the roots outside the core. We plot in Fig. 1 the solutions, for a system of $`16`$ particles, at different densities. The table below shows the integer sets at four different densities (since they occur in pairs we show only half of them):
| d=0.2 | 1 | 1 | 1 | 1 | 1 | 1 | 1 | 1 |
| --- | --- | --- | --- | --- | --- | --- | --- | --- |
| d=0.3 | 1 | 1 | 1 | 0 | 1 | 1 | 2 | 1 |
| d=0.4 | 1 | 0 | 1 | 0 | 1 | 1 | 2 | 2 |
| d=0.5 | 0 | 0 | 0 | 0 | 0 | 0 | 0 | 8 |
We now discuss the energies that we obtain from the BA solutions. For $`N16`$, we have verified that all the solutions obtained from the numerical solutions of the BE using the above scheme, match with those obtained from exact numerical diagonalization. With the BE we can go to much larger system sizes. We find that the gap vanishes at every finite $`q`$, with system size dependence $`1/N`$. In Fig. 2 we plot the system size dependence of the gap at two densities, namely at quarter and half fillings. The latter case corresponds to the $`q=\pi `$ state considered by Sutherland and we verify his result $`N\delta E=J\pi ^2`$. At $`d=1/4`$ $`N\delta E`$ seems to asymptote to the value $`3J\pi ^2/4`$. Remarkably both these asymptotic values of $`N\delta E`$ at densities $`1/2`$ and $`1/4`$ can be obtained from the low density formula in Eq(4). In Fig. 2 we also plot the energy-wavevector curve, obtained from the solution of the $`16`$-particle problem and compare it with Eq(4). Note that the discrepancies at large $`q`$ are finite size effects and would vanish in the $`N\mathrm{}`$ limit.
Variational results Having found the excited states, we now turn to the explicit connection with Bloch wall states. We now show, remarkably enough, that the expression Eq(4) for the gap can be obtained from a simple variational calculation. We work with arbitrary spin $`s`$ of the particles. A neat way to generate Bloch walls is via a unitary rotation operator $`Q=\mathrm{exp}i\frac{2\pi }{N}_mmS_m^z`$ acting upon an appropriate state, $`|0_n>(\widehat{S}_0^{})^n|ferro>`$, where $`\widehat{S}_q^{}_j\mathrm{exp}(iq.r_j)S_j^{}`$ is a spin wave creation operator carrying momentum $`q`$ and $`|ferro>`$ is the state with all spins up. Using $`[H,\widehat{S}_0^{}]=0`$, we see that $`|0_n>`$ is the ground state in the $`n`$-particle sector with zero total momentum. Thus finally we write the variational Bloch Wall state $`|BW>Q|0_n>`$. Using the quasi commutator $`Q\widehat{S}_q^{}=\widehat{S}_{q+2\pi /N}^{}Q`$, we find that $`|BW>=(\widehat{S}_{2\pi /N}^{})^n|ferro>`$, and has total momentum $`q=2\pi d`$. In the semiclassical limit, $`s>>1`$, the above state is readily visualized as classical spins that are tipped from the z axis and rotate along a cone. The variational calculation of the excitation energy $`\delta E`$, i.e. $`<0_n|Q^{}Q|0_n>`$ can be done easily by transforming the rotation onto the spin operators. Writing $`n=dN`$, with $`0d2s`$, we find $`\delta E=\frac{4J\pi ^2}{N}d(2sd)`$. At $`s=1/2`$ this is also the result of the calculation in Eq(4).
At this point we admit that we were surprised, as the reader might well be, that the results of an elaborate bound state calculation with a macroscopic number of complex roots agrees with the result of a simple looking variational wavefunction that resembles a Bose condensate of spin waves. This phenomenon is presumably a consequence of the shallow nature of the bound state. We note that the variational states satisfy $`\sqrt{<^2><>^2}/<>O(1/\sqrt{N})`$ which shows that in the $`N\mathrm{}`$ limit these become exact eigenstates, thereby providing independent evidence for the exactness of the main result of our work.
For large $`s`$, the energy as well as momentum of these states agrees with the semiclassical estimates using a Poisson bracket structure to construct a semiclassical momentum operator. The variational results for all values of $`s`$ thus collapse onto the same formula. The conjecture of implies just this kind of a result, but for the solitons, i.e. for bound states with small number of spin deviations. Needless to say, we believe that our variational results are exact (in the thermodynamic limit) for all spin, since we have established them at $`s=1/2`$ and also for very large $`s`$.
We note that the Bloch wall states $`|BW>`$ carry a total spin that is easy to calculate using simple extension of the above calculation: $`<S_{tot}^2>=<S_z>^2=N^2(sd)^2`$. Finally we note that these variational results for Bloch walls are generalizable to higher dimensions.
Concluding remarks We finally note in summary that the Bethe formula for $`n`$ magnon bound states: $`\omega _{Bethe}=J\frac{2}{n}(1\mathrm{cos}(q))`$ is (a) exact for $`q>>2\pi n/N`$, (b) invalid for $`q<2\pi n/N`$ and (c) has significant corrections when $`q2\pi n/N`$ . The evidence for (a) is in Bethe’s paper itself, the corrections to it are exponential in N, (b) is numerical. Evidence for (c) has been presented in this paper, where we find instead: $`\omega _{BW}=2\pi J/N|q|(1\frac{|q|}{2\pi })`$.
<sup>1</sup> Also at the Poornaprajna Institute, Bangalore. dabhi@rri.ernet.in <sup>2</sup> Also at the JNCASR, Bangalore. bss@physics.iisc.ernet.in
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# Embedded Stellar Clusters in the W3/W4/W5 Molecular Cloud Complex
## 1 Introduction
The equilibrium of molecular clouds is frequently perturbed by passages through spiral arms, cloud-cloud collisions, shocks from OB stellar winds and ionization fronts, molecular outflows, and other energetic forces. Depending on the local conditions, these events can either compress or disperse the molecular cloud, and consequently, induce or halt any future star formation (Elmegreen, 1992). These interactions occur more frequently among clouds located in the disk of the Milky Way, and therefore potentially, the average star formation characteristics of clouds in the disk can differ from those that are far removed from the Galactic Plane. Since most of the molecular material is confined to the disk of the Galaxy (Clemens, Sanders, & Scoville, 1988; Dame et al., 1987), studying the stellar content of Galactic Plane molecular clouds is necessary to establish the conditions under which most star form.
The biggest challenge in studying molecular clouds in the disk of the Galaxy is that they are generally rather distant from the sun at several kiloparsecs and are often blended together in projection, especially in the inner Galaxy (see, e.g., Lee, Snell, & Dickman 1990). Thus detecting and uniquely associating star formation sites within individual molecular clouds can often be problematic. To minimize these difficulties, the Five College Radio Astronomy Observatory (FCRAO) has recently completed a <sup>12</sup>CO survey of the Outer Galaxy (Heyer et al., 1998). This survey encompasses 330 deg<sup>2</sup> of the second Galactic quadrant at subarcminute resolution and sampling, and represents the most detailed examination to date of the molecular interstellar medium. Most lines of sight in the survey contain a single molecular cloud, and moreover, this region of the Galaxy includes the closest approach of the Perseus Spiral Arm to the sun at a distance of $``$ 2 kpc (Georgelin & Georgelin, 1976). In these respects, this is one of the best regions in the Milky Way to investigate the stellar properties of Galactic Plane molecular clouds.
In this paper, we investigate the stellar characteristics of molecular clouds within the Perseus Spiral Arm, and more specifically, clouds in the vicinity of the W3/W4/W5 H ii regions, using the FCRAO Outer Galaxy <sup>12</sup>CO Survey, the IRAS Point Source Catalog, published radio continuum observations, and new near-infrared and molecular line data. The W3/W4/W5 chain of H ii regions is ionized by members of the Cas OB6 association and extends over 150 pc along the Perseus Arm. The winds and ionizing flux from the massive stars have clearly impacted the interstellar medium by creating a galactic chimney out of the atomic hydrogen gas (Normandeau, Taylor, & Dewdney, 1996), shaping molecular clouds into cometary globules with parsec sized tails (Heyer et al., 1996), and possibly inducing a second generation of OB star formation (Lada et al., 1978; Thronson, Campbell, & Hoffman, 1980). Despite the vigorous star formation activity in the past, a substantial mass of molecular gas remains (Heyer & Terebey, 1998; Digel et al., 1996; Lada et al., 1978). The W3 molecular cloud alone has $``$ 10<sup>5</sup> M of molecular material spread over a $``$ 60 pc region (Deane et al., 2000; Lada et al., 1978) and is one of the most massive molecular clouds in the outer Galaxy (Heyer et al., 1998). The luminous star forming sites W3 Main, W3(OH), W3 North, and AFGL 333 (Thronson, Campbell, & Hoffman, 1980), a ridge of dense molecular gas (Tieftrunk et al., 1998), and the presence of embedded clusters throughout the W3/W4/W5 cloud complex (Tieftrunk et al., 1998; Deharveng et al., 1997; Megeath et al., 1996; Hodapp, 1994; Carpenter et al., 1993) all attest to the continued star formation activity in this region. In addition to the massive W3 Giant Molecular Cloud (GMC), a number of small clouds with similar <sup>12</sup>CO(1–0) velocities as the W3/W4/W5 region are found scattered throughout the area (Heyer et al., 1998).
To investigate the most recent generation of star formation throughout the W3/W4/W5 region, we have selected a sample of likely embedded star forming sites using the IRAS point source catalog. In conjunction with new <sup>13</sup>CO(1–0) and $`K^{}`$ band near-infrared observations, we investigate the spatial distribution of star forming regions, their associated molecular cloud properties, and the incidence of stellar clusters, and use these data to establish where most of the stars are now forming in the W3/W4/W5 region. These results are presented as follows. Section 2 discusses the criteria used to identify star forming regions from the IRAS Point Source Catalog and describe the observations and data reduction procedures for the new molecular line and near-infrared surveys. In Section 3, we characterize the stellar content associated with the IRAS sources based upon the far-infrared luminosity, published radio continuum observations, and the incidence of any stellar clusters detected in the $`K^{}`$ band mosaics. The implication of these results on star formation in the W3/W4/W5 region is discussed in Section 4, and our conclusions are summarized in Section 5.
## 2 Observations
The region analyzed for this study encompasses the area between galactic longitudes $`\mathrm{}=130\mathrm{°}`$ to $`139\mathrm{°}`$ and latitudes $`b=2.2\mathrm{°}`$ to $`+4.5\mathrm{°}`$. Figures 1 and 2 show this region of the Galaxy as observed in <sup>12</sup>CO(1–0) from the FCRAO Outer Galaxy Survey (Heyer et al., 1998) and in $`\lambda `$ 21 cm radio continuum emission from the DRAO Galactic Plane Survey (Normandeau, Taylor, & Dewdney, 1997). The W3, W4, and W5 H ii regions and the associated molecular clouds are evident in these figures. A list of candidate embedded star forming regions within these clouds were selected using the IRAS point source catalog, and follow-up observations of these sources were conducted to ascertain their stellar content. The selection criteria for the IRAS point sources and the ancillary observations are described below.
### 2.1 IRAS Point Source Selection
Guided by the far-infrared colors of known embedded star forming regions (Carpenter, Snell, & Schloerb, 1991; Kenyon et al., 1990; Beichman et al., 1986) and the properties of the IRAS sources observed by Wouterloot et al. (1990), the following five criteria were used to identify embedded star forming regions in the IRAS Point Source Catalog, Version 2.1. (1) The source has a high or moderate quality flux measurement at both 25 µm and 60 µm. (2) The flux density at 25 µm is $`0.4`$ Jy. (3) The flux density at 60 µm is $`1.0`$ Jy. (4) The ratio of the 60 µm to the 25 µm flux density is $`>`$ 1.0. (5) The ratio of the 100 µm to the 60 µm flux density is $``$ 4.0. The two flux density ratio criteria were designed to isolate sources with rising spectral energy distributions, but to eliminate the reddest objects that are often indicative of infrared cirrus. Any upper limits to the 100 µm flux density were used as appropriate in evaluating the ratios.
Of the sources meeting the above criteria, 34 were identified with a molecular cloud that has a <sup>12</sup>CO velocity between -57 km s<sup>-1</sup> $`v_{\mathrm{lsr}}32`$ km s<sup>-1</sup> and are likely in the Perseus spiral arm (Heyer et al., 1998). These 34 sources are listed in Table 1, along with the galactic and equatorial coordinates, the velocity of the <sup>12</sup>CO(1–0) emission coincident with the IRAS source, and any source identifications. All but two of these sources, as noted in Table 1, were imaged at $`K^{}`$ band and mapped in <sup>13</sup>CO as described below. The spatial distribution of the IRAS point sources is shown by the open circles in Figure 1, where the circle diameter is proportional to the far-infrared luminosity (see Section 3.2). The source list does not include W3 Main, which is confused at 60 µm in the IRAS survey and did not meet the selection criteria. Near-infrared observations of this source have been presented by Megeath et al. (1996; see also Tieftrunk et al. 1998 and Hodapp 1994). Also, two of the IRAS sources have previously been associated with galaxies (Weinberger, 1980). However, the strong <sup>12</sup>CO emission and the apparent stellar cluster (see Section 4.1) associated with these two sources suggests that the far-infrared emission originates from embedded stars and not an extragalactic object.
### 2.2 $`K^{}`$ Band Imaging
Near-infrared mosaics in a $`K^{}`$ band filter (Wainscoat & Cowie, 1992) were obtained for 32 of the 34 IRAS sources in Table 1. Time constraints prevented us from observing IRAS 02081+6225 and 02204+6128. The images were obtained over a two night period in 1996 October using QUIRC at the University of Hawaii 2.24 m telescope on Mauna Kea through thin cirrus with seeing conditions of $``$ 0.5-0.6″. QUIRC contains a 1024 $`\times `$ 1024 HgCdTe array and was used at the $`f`$/10 focus to provide a plate scale of 0.186″ pixel<sup>-1</sup> and an instantaneous field of view of $`3.2^{}\times 3.2^{}`$. For each IRAS source, a 5′ $`\times `$ 5′ mosaic aligned in the equatorial coordinate system was obtained that consists of 12 dithered frames with an exposure time of 30 seconds per frame. Sky frames were constructed by median filtering images free of extended nebulosity. The sky-subtracted frames were then corrected by a flat field image derived from a series of exposures of the dome interior with and without illumination from incandescent lights. In constructing the mosaics, three frames were coadded per pixel position in the mosaic in order to maintain a constant noise level across the final image. Coadded pixels near the edge of the mosaic that have only one or two observations were discarded. Astrometry for 23 of the 32 mosaics (see Table 1) were established using images from the 2 Micron All Sky Survey (2MASS). 2MASS images for the 9 remaining mosaics were not available at the time of this study, and we assumed that the center of these mosaics corresponds to the IRAS point source position. Based on the results from registering the 23 mosaics with available 2MASS images, we expect that the astrometry for these 9 mosaics to be accurate to be $``$ 30″.
Stars were identified in the mosaics using DAOFIND in IRAF. The noise for each mosaic was measured empirically, and a 5$`\sigma `$ detection threshold was used to create an initial source list. All mosaics were then visually inspected to remove a few saturated stars and any obvious non-stellar objects (e.g. ghosts, nebulosity knots), and to add any stars that were not identified by DAOFIND. Photometry was performed using the point-spread fitting task DAOPHOT in IRAF. The point spread function (PSF) was determined for each mosaic using several bright, isolated stars in the image. After fitting the PSF to each star in the point source list and subtracting the fit from the mosaic, the resulting image was examined for any additional point sources that were initially missed due to source confusion. These stars were added to the detection list and were also measured using DAOPHOT. Finally, objects identified by DAOPHOT as being unusually extended based upon the “sharp” statistic (e.g. galaxies, nebulosity) were removed form the point source list. The number of extended objects removed by this criteria amounted to less than 10% of the total number of sources identified. The photometry was calibrated by observing standard stars in the UKIRT faint standards list (Casali & Hawarden, 1992) and assuming a $`K^{}`$ band extinction coefficient of 0.08 mag/airmass. The RMS scatter over the two nights in the photometric zero points derived from the standard observations is 3%.
The differential completeness limit of the survey was established by adding artificial stars of known magnitude to one of the mosaics and determining the fraction of the stars that could be recovered at the 5$`\sigma `$ detection threshold. When adding artificial stars, care was taken not to add objects near a known star or within nebulous regions. The completeness limit in these more confused regions will obviously occur at a brighter magnitude. Approximately 90% of the stars with a $`K^{}`$ magnitude of 17.5<sup>m</sup> were recoverable in this automated procedure. The star counts analyzed in the paper are therefore for objects with $`K^{}`$ magnitudes between 11.5<sup>m</sup> (the saturation limit) and 17.5<sup>m</sup>.
### 2.3 <sup>13</sup>CO Mapping
A region of size ($`\mathrm{\Delta }\mathrm{}\times \mathrm{\Delta }b`$) = ($`6.2^{}\times 4.8^{}`$) toward the 32 IRAS point sources was mapped in <sup>13</sup>CO(1–0) (110.201370 GHz) using the SEQUOIA receiver array on the 14 m telescope operated by the Five College Radio Astronomy Observatory (FCRAO) in the spring of 1998. At the time of the observations, SEQUOIA had 12 pixel elements. The full-width-at-half-maximum (FWHM) beam size of the FCRAO telescope at the observed frequency is 47<sup>′′</sup>, and the maps were sampled every 22<sup>′′</sup>, or approximately the Nyquist sampling interval. The backends for each pixel in the SEQUOIA array consisted of an autocorrelator spectrometer configured to achieve a velocity resolution of 0.064 km s<sup>-1</sup> over a 54 km s<sup>-1</sup> velocity interval. The data were obtained in frequency switching mode with an offset of 4 MHz between the nominal and reference frequency. The <sup>13</sup>CO data presented here have been corrected by the main beam efficiency, previously measured to be $`\eta _\mathrm{B}`$ = 0.45. The RMS noise is typically $`\mathrm{\Delta }`$$`T_{\mathrm{MB}}`$ = 0.7-1.1 K per channel.
## 3 Results
### 3.1 Images
Molecular line and $`K^{}`$ band images of the 32 IRAS point sources observed for this study are shown in Figure 3. Four images are shown for each IRAS source: (1) a <sup>12</sup>CO(1–0) integrated intensity map from the FCRAO Outer Galaxy Survey over a $`30^{}\times 30^{}`$ area centered on the IRAS point source position (far left panels), (2) the <sup>13</sup>CO integrated intensity map over a $`6.2^{}\times 4.8^{}`$ region, (3) the $`K^{}`$ band mosaic, and (4) the $`K^{}`$ stellar surface density map. The <sup>12</sup>CO images are shown over a larger extent than the <sup>13</sup>CO maps and $`K^{}`$ band mosaics in order to place the IRAS source in context of the large scale molecular cloud in the region. Figure 3 indicates that most of the <sup>13</sup>CO maps peak near the IRAS point source positions and suggests that the IRAS sources are indeed related with the molecular gas. In the remainder of this section, we characterize the stellar population associated with these IRAS sources as inferred from the far-infrared luminosity, the presence of an embedded massive star as indicated by radio continuum emission, and the identification of stellar clusters from $`K^{}`$ band star counts.
### 3.2 Far-Infrared Luminosities
The stellar content associated with the IRAS sources can be constrained to first order by assuming that dust absorbs a substantial fraction of the stellar bolometric luminosity and re-emits the radiation in the far-infrared. Since many of the IRAS sources are associated with a cluster of stars (see Section 3.4), the far-infrared emission actually sets a limit on the most massive star that may be forming in these regions. The luminosity emitted in the 12 µm, 25 µm, 60 µm bands (L<sub>FIR</sub>) was computed by summing the observed flux densities in the individual IRAS band passes using the formula
$`\mathrm{L}_{\mathrm{FIR}}`$ $`=`$ $`4\pi \mathrm{D}^2{\displaystyle \underset{\mathrm{i}}{}}(\mathrm{S}_{\nu \mathrm{i}}\mathrm{\Delta }\nu _\mathrm{i})`$ (1)
$`=`$ $`0.30\left({\displaystyle \frac{\mathrm{D}}{\mathrm{kpc}}}\right)^2{\displaystyle \underset{\mathrm{i}}{}}\left({\displaystyle \frac{\mathrm{S}_{\nu \mathrm{i}}}{\mathrm{Jy}}}\right)\left({\displaystyle \frac{\mathrm{\Delta }\nu _\mathrm{i}}{10^{12}\mathrm{Hz}}}\right)\mathrm{L}_{},`$
where D is the distance to the source in kiloparsecs (assumed to be 2.35 kpc; Massey, Johnson, & DeGioia-Eastwood 1995), $`\mathrm{\Delta }\nu `$ the IRAS band width, and S<sub>ν</sub> the observed flux density. The sum does not extend over the 100 µm band since many sources are confused at 100 µm and this band could not be applied consistently for the entire sample. Thus the computed far-infrared luminosities (see column 2 in Table 2) will underestimate the actual bolometric and far-infrared luminosities. For IRAS sources in the W3/W4/W5 region that do have high quality 100 µm detections, we found that the 12 µm-60 µm luminosity underestimates the total IRAS far-infrared luminosity by $``$ 30% on average.
A histogram of the derived far-infrared luminosities for the 34 IRAS sources is shown in Figure 4. The observed luminosities range between 9 L and 46,000 L with the peak of the distribution at $``$ 100 L. The brightest sources have luminosities similar to that of early B type zero age main sequence (ZAMS) stars (Panagia, 1973) assuming that most of the far-infrared luminosity originates from a single object. The lowest luminosities in the histogram are a result of the selection criteria. The flux density criteria alone used to select the IRAS sources implies a detection limit of 8 L. Further, since the 60 µm flux density is on average 11 times larger than the 25 µm flux density for sources in our sample, the detection limit set by the spectral energy distribution and flux density limits is 23 L. The decline in the number of sources with luminosities fainter than $``$ 100 L then is likely a result of incompleteness in the IRAS point source catalog. This 100 L limit corresponds to a 1 Myr, 3 M pre-main-sequence object (Palla & Stahler, 1993), or a late B ZAMS star.
### 3.3 Radio Continuum Emission
The stellar content associated with the IRAS sources can be further constrained by using radio continuum observations to estimate the spectral type of the most massive star. The list of radio continuum sources in the W3/W4/W5 region were taken primarily from the $`\lambda `$ 20 cm NRAO/VLA Sky Survey (NVSS; Condon et al. 1998), but also published targeted observations (Kurtz, Churchwell, & Wood, 1994; McCutcheon et al., 1991; Carpenter, Snell, & Schloerb, 1991). Sources in the NVSS catalog were deemed associated with an IRAS source if the radio and far-infrared coordinates agreed to within 30″. While no spectral information is available from the NVSS catalog to confirm that these objects are actually compact H ii regions, the random probability of a false association between an extragalactic radio continuum source and the IRAS point source is only $``$ 0.01 for the adopted 30″ matching radius (Condon et al., 1998). The observed radio continuum flux was used to estimate the number of ionizing photons (see, e.g., Carpenter, Snell, & Schloerb 1991) and consequently the spectral type of the ionizing star (Panagia, 1973). The inferred spectral types and references for the radio continuum observations are provided in columns 3 and 4 of Table 2. Ten of the 32 IRAS point sources have radio continuum detections, and the inferred spectral types range from B2 ZAMS to O7 ZAMS. If the stellar parameters from Vacca, Garmany, & Shull (1996) are used instead of Panagia (1973) to infer the spectral type, all but one of the detected sources will have a spectral type later than B0.5 ZAMS.
Two of the IRAS sources (IRAS 01546+6319 and 02511+6023) associated with a B type star as based on their radio continuum flux have a 12 µm-60 µm luminosity that is 2 orders of magnitude less than that expected for such a massive star (Panagia, 1973). Either the association between these radio continuum and IRAS sources is incorrect and the radio continuum source is an extragalactic object, or a substantial fraction of the bolometric luminosity is radiated at shorter wavelengths. The latter situation may occur if a star forming region is relatively evolved and the circumstellar dust no longer completely absorbs and re-emit the stellar radiation, although no obvious bright star is present in the $`K^{}`$ mosaics or on the Palomar Observatory Sky Survey prints to indicate that this may be the case. Nonetheless, we assume in the remainder of this paper these two IRAS sources are indeed massive star forming regions. Our general conclusions will not change if this assumption is incorrect.
### 3.4 Identification of Stellar Clusters
As another means to characterize the stellar content associated with the IRAS Point Sources, we searched for stellar clusters in the $`K^{}`$ images using the procedure adopted by Carpenter, Snell, & Schloerb (1995; see also Carpenter et al. 1997). Briefly, histograms of the stellar field star density were generated for each mosaic using 20$`{}_{}{}^{\prime \prime }\times `$20<sup>′′</sup> counting bins sampled every 10<sup>′′</sup>. The observed frequency distribution of counts at low stellar surface densities in these histograms usually resembles a Poisson distribution, which is identified with field stars and embedded stars randomly distributed across the $`K^{}`$ band mosaic. By fitting a Poisson distribution to these lower surface density bins, the mean stellar surface density of randomly distributed stars can be determined. Stellar surface density bins that significantly exceed this mean surface density are identified as possible clusters. Contour maps of the stellar surface density are shown in Figure 3 for each of the mosaics. The lowest contour in each map begin at 2$`\sigma `$ for a Poisson distribution above the mean stellar surface density with 3$`\sigma `$ contour intervals.
A cluster was identified in Figure 3 if the total number of stars within the 2$`\sigma `$ contour represents a 5$`\sigma `$ enhancement with respect to the expected stellar background level. Further, we require that the identified cluster be near the center of the mosaic and the IRAS point source position. In a few instances, apparent clusters were identified near the edge of the mosaic. In nearly all of these cases, the extinction through the cloud was large enough and variable that a Poisson distribution is a poor representation of the star counts, and these clusters most likely represent regions within the image where the extinction becomes low. Any such “cluster” is almost certainly a projection of unrelated field stars and was excluded from the final cluster list.
Table 2 summarizes the properties of the identified clusters, including the effective cluster radius ($`R_{\mathrm{eff}}`$), the number of stars observed within the 2$`\sigma `$ boundary (N<sub>s</sub>), and the number of stars inferred for the cluster after subtracting off the expected field star population (N<sub>cluster</sub>). The effective radius is defined as $`\sqrt{A/\pi }`$, where $`A`$ is the area within the 2$`\sigma `$ contour. Of the 32 IRAS sources, 19 have identifiable clusters, with the number of cluster members ranging from 20 to 240 stars. Our results agree with previous studies in that clusters are identified around IRAS 02232+6138 (Tieftrunk et al., 1998), IRAS 02575+6017 (Deharveng et al., 1997; Hodapp, 1994; Carpenter et al., 1993), and IRAS 02593+6016 (Carpenter et al., 1993). However, Tieftrunk et al. (1998) visually noted a second cluster near IRAS 02232+6138 that is within the field of view of our mosaic. This grouping of stars does not meet the surface density criteria adopted here to be identified as a cluster.
The cluster membership listed in Table 2 are lower limits to the actual stellar population in these regions, mainly due to the finite sensitivity of the observations. To compare these clusters with other star forming regions, we computed the fraction of stars in nearby embedded clusters that would be detectable at the distance of W3 considering both differences in sensitivity and resolution. In the UKIRT survey of the MonR2 cluster by Carpenter et al. (1997), 246 of 378 stars (65%) would be detectable with our QUIRC observations at the distance of W3/W4/W5. Similarly, we could detect $``$ 410 (53%) of the $``$ 780 stars in the inner $`5^{}\times 5^{}`$ of the Orion Nebula Cluster (Hillenbrand & Carpenter, 2000), and $``$ 75% of the 94 stars brighter than $`m_k`$ = 14.5<sup>m</sup> in the NGC 1333 cluster (Lada, Alves, & Lada, 1996). Assuming that the extinction toward the W3/W4/W5 clusters is not substantially different than toward these comparison regions, the richer W3/W4/W5 clusters ($``$ 200 stars) are comparable to MonR2, but not as rich as the Orion Nebula Cluster. Several of the W3/W4/W5 clusters though have significantly fewer stars than these comparison regions.
### 3.5 Molecular Cloud Properties
The properties of the molecular clouds associated with the IRAS sources as derived from <sup>13</sup>CO(1–0) emission are summarized in Table 3. The cloud sizes are defined as the circular radius needed to produce the area within the FWHM integrated intensity contour. The cloud masses were computed from the integrated <sup>13</sup>CO(1–0) intensity by assuming that the emission is optically thin and that the <sup>13</sup>CO/H<sub>2</sub> abundance is $`1.5\times 10^6`$ (Bachiller & Cernicharo, 1986). Statistical equilibrium calculations indicate that for H<sub>2</sub> volume densities between 300 cm<sup>-3</sup> and 10<sup>4</sup> cm<sup>-3</sup> and kinetic temperatures between 10 K and 20 K, the fraction of the <sup>13</sup>CO molecules in the J=1 rotational state varies between 0.41 to 0.54, and 0.48 was adopted as a typical value. A factor of 1.36 was included in the calculations to include the mass contribution from helium and other elements. The <sup>13</sup>CO integrated intensity used to calculate the masses includes the region within the FWHM intensity contour, but was also extrapolated to include emission outside this contour level by assuming a gaussian intensity distribution. Lower limits to the cloud sizes and masses are reported for sources in which the <sup>13</sup>CO maps did not completely encompass the half power contour level. The derived cloud masses range from 27 M to $`>`$ 3200 M (see Table 3). Many of these clouds, and in particular the more massive objects, are parts of the extended W3 and W5 GMCs. Several of the clouds are small and distant from any GMC (see Figure 1) and are nearly completely mapped with these <sup>13</sup>CO observations.
## 4 Discussion
### 4.1 Stellar Clusters in W3/W4/W5
From the results presented in Section 3 we can begin to assess the properties of the regions that are forming most of the cluster stellar population. The 19 identified clusters contain a total of 1595 stars within the 2 $`\sigma `$ boundary after subtracting off the expected field star contamination. (If the 1 $`\sigma `$ boundary is used to define the clusters, the total membership increases by 33%.) Figure 5 shows the normalized cumulative distribution of the total number of stars in all 19 clusters as a function of the number of stars in an individual cluster. As Figure 5 shows, 52% of the cluster members are found in just the 5 richest clusters. Even if we assume that the 13 IRAS sources without identified clusters have 20 stars each (the smallest population cluster identified in this survey), the 5 richest clusters would still contain 45% of the total cluster membership. These results do not change significantly if we include the W3 Main cluster, which has $``$ 87 stars to a comparable sensitivity limit as our survey (Megeath et al., 1996), or the other clusters identified along the W3 ridge (Tieftrunk et al., 1998).
We can also investigate the fraction of the cluster population associated with embedded OB stars. The 5 rich clusters described above that contain the majority of the cluster population are each associated with a massive star, and in total, clusters around OB stars contain 61% of the cluster population. Not all of the embedded OB stars though are associated with clusters. IRAS 02230+6202, 02511+6023, and 02531+6032 do not show a significant enhancement in the stellar surface density despite having a radio continuum detection. Based on the field star initial mass function, one would expect an OB type star to form along with a few hundred lower mass objects (Miller & Scalo, 1979). It would be remarkable then if these regions are truly forming a single massive star, and would imply that the stellar mass function in these regions is strongly skewed toward high mass objects. With high spatial resolution data at only one wavelength, however, we cannot yet rule out the possibility that the lack of a cluster in these regions is merely due to a large amount of extinction, and additionally in the case IRAS 02230+6202, a bright nebular background (see Fig. 3), that prevents us from detecting the underlying cluster. Deeper observations in the near- and mid-infrared should be able to establish more definitively if these OB stars are truly forming in isolation.
In principle the above results could change if a large number of low luminosity IRAS sources exist that did not meet our selection criteria but contribute substantially to the cluster population. Based on the peak of the histogram in Figure 4, we assume that the luminosity completeness limit for the IRAS point sources is $``$ 200 L, and a power law function was fitted to the histogram bins more luminous than this limit. Extrapolation of this power fit suggests that there could be $``$ 40 additional sources between 10 L and 200 L that are not already in our sample. Even if we assume these sources have 20 stars each (the smallest cluster size detected here), the IRAS sources associated with OB will still contain more than a third of the total cluster population. These suggest indicate that the formation of dense stellar clusters surrounding OB stars represent a significant component of the cluster population in the W3/W4/W5 region.
Our results for the W3/W4/W5 region can be compared with other near-infrared surveys of molecular clouds. The regions that have been most extensively studied on a global scale are the Orion A and Orion B molecular clouds, and in both of these clouds, dense stellar clusters contribute significantly to the total young stellar population. In Orion B, approximately 96% of the stellar population is found in just 4 clusters (Lada et al., 1991; Li, Evans, & Lada, 1997). Similarly, in Orion A, at least 60% of the total stellar population is found within the Orion Nebula Cluster alone (Allen & Hillenbrand, 2000; Meyer & Lada, 2000), while the rest is distributed more uniformly throughout the molecular cloud and in several small clusters (Strom, Strom, & Merrill, 1993). Near-infrared surveys of other molecular clouds are not as extensive as those in Orion and are biased toward known star forming regions, but nevertheless, observations of NGC 2264 (Lada, Young, & Greene, 1993; Piche, 1993) and NGC 1333 (Lada, Alves, & Lada, 1996) also suggest that at least half of the stars within these regions form within clusters as opposed a more uniformly distributed population (see also the review by Clarke, Bonnell, & Hillenbrand 2000). While our survey is also not sensitive to any pervasive distributed stellar population, our results for the W3/W4/W5 region are similar to Orion in that a significant fraction of the cluster population is confined to just a few rich clusters.
### 4.2 Global Star Formation Characteristics
Of the 32 IRAS sources in our sample that were mosaicked at $`K^{}`$ band, 21 are located in projection against the W3/W4/W5 H ii region/molecular cloud complex (see Figure 2). The remaining 11 sources, or 34% of the sample, are separated from the W3/W4/W5 complex by as much as 100 pc. The molecular clouds associated with these 11 sources have radii of $``$ 0.5 pc with masses ranging from 27 M to 410 M, with an average mass of 130 M. This contrasts with the massive star forming sites found in W3 that have cloud masses upwards of a few thousand solar masses or more.
Despite their relatively small mass, the 11 isolated clouds are not devoid of star formation nor is their stellar population limited to low mass stars. Eight of the 11 sources are forming a stellar cluster (see Table 2), and 3 sources (IRAS 01546+6319, 02044+6031, and 02561+6147) are associated with an early B type as indicated by the radio continuum emission. The fraction of the total cluster population found within these 11 clouds is 629/1595 = 39% despite the fact they have only $``$ 11% of the total molecular mass as derived from the <sup>13</sup>CO observations. These percentages are likely upper limits since it is more difficult to identify IRAS point sources (and consequently clusters) near W3 than in more isolated regions. Indeed, the W3 Main cluster, two clusters a few arcminutes from W3(OH), and a cluster around BD +61 411 (Tieftrunk et al., 1998; Megeath et al., 1996) are found along the W3 ridge but were not identified with the IRAS point source selection criteria adopted in this paper. From visual inspection of the $`K`$ band image presented in Tieftrunk et al. (1998), we estimate that these clusters in total may contain a few hundred stars. Even after accounting for these stars, the isolated clouds contribute a non-negligible fraction to the total cluster population currently forming in the W3/W4/W5 region.
The presence of massive stars and stellar clusters in these low mass clouds is somewhat surprising in that such characteristics are generally associated with massive GMCs. An interesting question is then how these isolated, star forming clouds originated. While the observations presented here cannot answer this question directly, we can speculate on their origin based upon the distribution of clouds in the vicinity of the W3/W4/W5 H ii regions. The correspondence between the <sup>12</sup>CO and radio continuum emission shows that the H ii regions have re-shaped the spatial distribution of molecular gas in some instances (see Figure 2). For example, the molecular gas near ($`\mathrm{},b`$) $``$ (135°,0°) wraps around the southern edge of the W4 H ii region, and the cometary globule near ($`\mathrm{},b`$) $``$ (134.8°,1.3°) possibly formed from the interaction between the molecular gas and radiation pressure from OB stars in the W4 H ii region (Heyer et al., 1996). The apparent effect of these interactions is that the small, dense globules of gas (which were either pre-existing or were formed by radiation pressure) remain for longer time scales than the more diffuse molecular material which is ionized or dispersed. Such globules are prominent in other OB associations and evolved H ii regions as well, such as Orion OB 1 (Ogura & Sugitani, 1998) and IC 1396 (Patel et al., 1998). It is tempting to speculate then that the isolated clouds tens of parsecs away from W3/W4/W5 are the last remnants of a GMC that has been dispersed by OB stars. In fact, the Per OB1 association occupies the area south of the W3/W4/W5 between $`\mathrm{}`$ = 130° to 138°and $`b=`$5° to $`1`$° where several of the low mass clouds are found and is at nearly the same distance as the W3/W4/W5 H ii regions (Garmany & Stencel, 1992).
### 4.3 Massive Star Formation in W3
The W3 region has long been proposed as a classic example of induced or “triggered” star formation in molecular clouds. In this scenario, shocks arising from ionization fronts compress the ambient molecular material such that a gravitational instability develops in the post-shocked gas, leading to gravitational collapse and the formation of a new generation of stars. In the specific case of the W3 molecular cloud, it has been suggested that the expansion of the W4 H ii region triggered the formation of W3 Main, W3(OH), W3 N, and AFGL 333, (Lada et al., 1978; Thronson, Campbell, & Hoffman, 1980). Alternatively, it has been proposed that diffuse H ii regions appear adjacent to embedded star forming sites because the associated dense molecular material has effectively slowed the expansion of the H ii regions, and thus no cause and effect relationship exists between the H ii regions and the newly formed stars. These two possibilities can be investigated with the data obtained for this study. If triggering is an important manner in which massive stars form, then embedded star forming regions throughout the W3/W4/W5 region should preferentially be located along the interfaces between ionization fronts and molecular clouds. Alternatively, if the expanding H ii regions have had minimal impact on the massive star formation activity, then massive star forming regions should be located randomly throughout the W3 molecular cloud. Thronson, Campbell, & Hoffman (1980) cited similar arguments in suggesting that massive star formation has indeed been triggered in the W3 region, and we re-investigate this issue here using the more extensive and sensitive observations provided by IRAS and recent radio continuum surveys.
Of the molecular clouds in the W3/W4/W5 region, only the W3 molecular cloud extends tens of parsecs beyond the H ii regions to provide a clear distinction between IRAS sources located near and distant from the ionization fronts as seen in Figures 1 and 2. (Most of the IRAS sources in the W5 molecular cloud are located along the W5 H ii region, consistent with the triggered star formation hypothesis. However, unlike the W3 region, the small extent of the W5 molecular cloud does not provide a “control” field unaffected by the ionization front.) The W3 molecular cloud extends for $``$ 60 pc and contains $`\mathrm{¿}\mathrm{}`$ 10<sup>5</sup> M of molecular material (Deane et al., 2000; Lada et al., 1978). Of the IRAS sources in our sample, the three most luminous sources are found along the interface between the W4 H ii region and the W3 molecular cloud. This interface also includes the W3 Main (IRAS 02219+6152) massive star forming region (Megeath et al., 1996), which is not in our sample since this object is confused at 60 µm in the IRAS survey. By contrast, in the interior of the W3 cloud (i.e. west of the W4 ionization front), no IRAS sources are found that meet our selection criteria and have a far-infrared luminosity in excess of 500 L.
Since W3 Main was not picked out by our IRAS selection criteria, we examine the possibility that massive star forming regions elsewhere in the W3 molecular cloud may also have been missed. Using the IRAS point source catalog, all sources were examined that have far-infrared luminosities in excess of 200 L but have low quality flux densities at 25 µm and/or 60 µm. Where appropriate, the reported upper limits were used in estimating the luminosity. All sources more luminous than 200 L were considered since two of the IRAS point sources in our sample are associated with a B type star despite having a low far-infrared luminosity (see Section 3.3). Using this search criteria, 12 additional sources were picked out in the W3 region. Eight of these sources are near W3 Main, W3(OH), or AFGL 333 and are confused at 60 µm and 100 µm, and one source is found midway between AFGL 333 and W3(OH). The other 4 IRAS sources are found interior to the cloud and are at least 30′ (20 parsecs) from the W4 ionization front. Each of these four sources have far-infrared luminosities less than 430 L, and none are associated with a NVSS radio continuum point source (Condon et al., 1998).
Recently, Ballantyne, Kerton, & Martin (2000) identified a massive O8 star that is ionizing an extended H ii region near the southwestern edge of the W3 molecular cloud at ($`\mathrm{},b)`$ $``$ (133.425°,0.055°). The source is associated with extended far-infrared emission and is not listed as an IRAS 60µm point source, and possibly may represent a more evolved star forming region than the point sources analyzed here. This source is distant from the W4 ionization front and may indeed be an example of a massive star formation region that has formed spontaneously (Ballantyne, Kerton, & Martin, 2000). Nonetheless, for the W3 molecular cloud as a whole, we conclude that the massive star forming regions are located preferentially (but not exclusively) along the W4 ionization front and are not randomly distributed throughout the cloud. This is quite remarkable considering that $``$ 75% of the molecular mass in the W3 molecular cloud (based on the observed <sup>12</sup>CO(1–0) integrated intensity) is contained within the extended molecular gas component. These observations are thus consistent with the conjecture that massive star formation along the eastern edge of the W3 molecular cloud has been triggered by the W4 ionization front as previously proposed by Lada et al. (1978) and Thronson, Campbell, & Hoffman (1980).
## 5 Summary
We have analyzed the star formation properties of molecular clouds in the vicinity of the W3/W4/W5 H ii regions using the FCRAO Outer Galaxy <sup>12</sup>CO(1–0) Survey, the IRAS point source catalog, published radio continuum observations, $`K^{}`$ band imaging, and <sup>13</sup>CO(1–0) observations. We use these data to identify 34 IRAS point sources in the Perseus spiral arm that have far-infrared colors characteristic of embedded star forming regions. The stellar population associated with 32 of these IRAS sources are investigated based upon far-infrared and radio continuum emission and the identification of stellar clusters in the $`K^{}`$ band mosaics. Our main conclusions are:
1. Nineteen stellar clusters are identified that contain a total of 1595 stars with $`K^{}`$ magnitudes between 11.5<sup>m</sup> and 17.5<sup>m</sup>. Approximately half of the total cluster population is found within the 5 richest clusters. Further, IRAS sources associated with embedded OB stars as traced by radio continuum emission contain 61% of the cluster population. Thus clusters around OB stars contribute substantially to the stellar population currently forming in the W3/W4/W5 region.
2. Eleven of the 32 IRAS point sources are located as far as 100 pc from the W3/W4/W5 H ii region/molecular cloud complex. Despite the low average cloud mass (130 M) associated with these 11 sources, three of these objects are associated with an OB star and 8 are forming a stellar cluster. These 11 clouds contain $``$ 39% of the total stellar population identified in the IRAS sources, and thus contribute significantly to the current star formation production. We speculate that these clouds are the remnant fragments of a once large molecular cloud that has since been dispersed.
3. The conjecture that the W4 H ii region as triggered star formation along the eastern edge of the W3 molecular cloud is re-investigated by analyzing the spatial distribution of star forming regions. We find that 4 of the 5 massive star forming sites in the W3 GMC, as traced by luminous IRAS point sources and/or radio continuum emission, are found along the interface between the W4 H ii region and the W3 molecular cloud. The western portion of the W3 cloud, despite having $``$ 75% of the molecular mass, contains only one known OB star forming region (Ballantyne, Kerton, & Martin, 2000). These characteristics are consistent with the long-held notion that most of the massive star formation activity in the W3 molecular cloud has been induced by the expansion of the W4 H ii region (Lada et al., 1978; Thronson, Campbell, & Hoffman, 1980).
JMC acknowledges support from the James Clerk Maxwell Telescope Fellowship while part of this work was carried out, and current support from the Owens Valley Radio Observatory and Long Term Space Astrophysics Grant NAG5-8217. The Five College Radio Astronomy Observatory is operated with support from NSF grant 94–20159. This publication makes use of data products from the Two Micron All Sky Survey, which is a joint project of the University of Massachusetts and the Infrared Processing and Analysis Center, funded by the National Aeronautics and Space Administration and the National Science Foundation.
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# 1 Introduction
## 1 Introduction
A new class of lattice Dirac operators $`D`$ have been recently proposed on the basis of the algebraic relation
$$\gamma _5(\gamma _5D)+(\gamma _5D)\gamma _5=2a^{2k+1}(\gamma _5D)^{2k+2}$$
(1.1)
where $`k`$ stands for a non-negative integer, and $`k=0`$ corresponds to the ordinary Ginsparg-Wilson relation for which an explicit example of the operator free of species doubling has been given by Neuberger. It has been shown in that we can in fact construct the lattice Dirac operator, which is free of species doublers, for all values of $`k`$. Here $`\gamma _5`$ is a hermitian chiral Dirac matrix and $`\gamma _5D`$ is also hermitian.
When one defines
$$\mathrm{\Gamma }_5\gamma _5(a\gamma _5D)^{2k+1}$$
(1.2)
the relation (1.1) is written as
$$\mathrm{\Gamma }_5(\gamma _5D)+(\gamma _5D)\mathrm{\Gamma }_5=0.$$
(1.3)
The index relation on the lattice is generally written as
$$Tr\mathrm{\Gamma }_5=n_+n_{},$$
(1.4)
which is confirmed by
$`Tr\mathrm{\Gamma }_5`$ $``$ $`{\displaystyle \underset{\lambda _n}{}}\varphi _n^{}\mathrm{\Gamma }_5\varphi _n`$ (1.5)
$`=`$ $`{\displaystyle \underset{\lambda _n=0}{}}\varphi _n^{}\mathrm{\Gamma }_5\varphi _n+{\displaystyle \underset{\lambda _n0}{}}\varphi _n^{}\mathrm{\Gamma }_5\varphi _n`$
$`=`$ $`{\displaystyle \underset{\lambda _n=0}{}}\varphi _n^{}\mathrm{\Gamma }_5\varphi _n`$
$`=`$ $`{\displaystyle \underset{\lambda _n=0}{}}\varphi _n^{}[\gamma _5(a\gamma _5D)^{2k+1}]\varphi _n`$
$`=`$ $`n_+n_{}=index`$
where $`n_\pm `$ stand for the number of normalizable zero modes in
$$\gamma _5D\varphi _n=0$$
(1.6)
for the hermitian operator $`\gamma _5D`$ with simultaneous eigenvalues $`\gamma _5\varphi _n=\pm \varphi _n`$. We also used the relation following from (1.3)
$$\gamma _5D\mathrm{\Gamma }_5\varphi _n=\lambda _n\mathrm{\Gamma }_5\varphi _n$$
(1.7)
if
$$\gamma _5D\varphi _n=\lambda _n\varphi _n,$$
(1.8)
which suggests that either $`\mathrm{\Gamma }_5\varphi _n`$ for $`\lambda _n0`$ is orthogonal to $`\varphi _n`$ or else $`\mathrm{\Gamma }_5\varphi _n=0`$. The positive definite inner product is defined by summing over all the lattice points
$$\varphi _n^{}\varphi _n=(\varphi _n,\varphi _n)\underset{x}{}a^4\varphi _n^{}(x)\varphi _n(x)$$
(1.9)
but the coordinate $`x`$ is often omitted in writing $`\varphi _n`$.
The Euclidean path integral for a fermion is defined by
$$𝒟\overline{\psi }𝒟\psi \mathrm{exp}[\overline{\psi }D\psi ]$$
(1.10)
where
$$\overline{\psi }D\psi \underset{x,y}{}\overline{\psi }(x)D(x,y)\psi (y)$$
(1.11)
and the summation runs over all the points on the lattice. The relation (1.3) is re-written as
$$\gamma _5\mathrm{\Gamma }_5\gamma _5D+D\mathrm{\Gamma }_5=0$$
(1.12)
and thus the Euclidean action is invariant under the global “chiral” transformation
$`\overline{\psi }(x)\overline{\psi }^{}(x)=\overline{\psi }(x)+i{\displaystyle \underset{z}{}}\overline{\psi }(z)ϵ\gamma _5\mathrm{\Gamma }_5(z,x)\gamma _5`$
$`\psi (y)\psi ^{}(y)=\psi (y)+i{\displaystyle \underset{w}{}}ϵ\mathrm{\Gamma }_5(y,w)\psi (w)`$ (1.13)
with an infinitesimal constant parameter $`ϵ`$. Under this transformation, one obtains a Jacobian factor
$$𝒟\overline{\psi }^{}𝒟\psi ^{}=J𝒟\overline{\psi }𝒟\psi $$
(1.14)
with
$$J=\mathrm{exp}[2iTrϵ\mathrm{\Gamma }_5]=\mathrm{exp}[2iϵ(n_+n_{})]$$
(1.15)
where we used the index relation (1.5). This derivation may be regarded as a lattice counter part of the continuum path integral.
In Ref. it was shown by using the method in , which is a lattice extension of the method in , that the index $`n_+n_{}`$ appearing in the Jacobian factor is related to the Pontryagin number for any operator in (1.1) if the operator $`\gamma _5D`$ satisfies suitable conditions. In this paper, we evaluate $`Tr\mathrm{\Gamma }_5`$ in a more explicit and elementary manner on the basis of explicit formulas for $`\gamma _5D`$ in the continuum limit<sup>1</sup><sup>1</sup>1The continuum limit in this paper stands for the so-called “naive”continuum limit with $`a0`$, and the lattice size is gradually extended to infinity for any finite $`a`$ in the process of taking the limit $`a0`$.. We show that these operators $`\gamma _5D`$ for all $`k`$ in fact reproduce the correct chiral anomaly and consequently correct Pontryagin number.
## 2 A brief summary of the model and notation
The operator $`\mathrm{\Gamma }_5`$ appearing in the index relation (1.5) has an explicit expression
$$\mathrm{\Gamma }_5=\gamma _5H_{(2k+1)}$$
(2.1)
with
$$H_{(2k+1)}(\gamma _5aD)^{2k+1}=\frac{1}{2}\gamma _5[1+D_W^{(2k+1)}\frac{1}{\sqrt{(D_W^{(2k+1)})^{}D_W^{(2k+1)}}}].$$
(2.2)
The operator $`D_W^{(2k+1)}`$ is in turn expressed as a generalization of the ordinary Wilson Dirac operator as
$$D_W^{(2k+1)}=i(\overline{)}C)^{2k+1}+(B)^{2k+1}(\frac{m_0}{a})^{2k+1}.$$
(2.3)
See Appendix for further details of the general solution to (1.1).
The ordinary Wilson Dirac operator $`D_W`$, which corresponds to $`D_W^{(1)}`$, is given by
$`D_W(x,y)`$ $``$ $`i\gamma ^\mu C_\mu (x,y)+B(x,y){\displaystyle \frac{1}{a}}m_0\delta _{x,y},`$
$`C_\mu (x,y)`$ $`=`$ $`{\displaystyle \frac{1}{2a}}[\delta _{x+\widehat{\mu }a,y}U_\mu (y)\delta _{x,y+\widehat{\mu }a}U_\mu ^{}(x)],`$
$`B(x,y)`$ $`=`$ $`{\displaystyle \frac{r}{2a}}{\displaystyle \underset{\mu }{}}[2\delta _{x,y}\delta _{y+\widehat{\mu }a,x}U_\mu ^{}(x)\delta _{y,x+\widehat{\mu }a}U_\mu (y)],`$
$`U_\mu (y)`$ $`=`$ $`\mathrm{exp}[iagA_\mu (y)],`$ (2.4)
where we added a constant mass term to $`D_W`$. Our matrix convention is that $`\gamma ^\mu `$ are anti-hermitian, $`(\gamma ^\mu )^{}=\gamma ^\mu `$, and thus $`\overline{)}C\gamma ^\mu C_\mu (n,m)`$ is hermitian
$$\overline{)}C^{}=\overline{)}C.$$
(2.5)
Since the operators $`\overline{)}C`$ and $`B`$ form the basis for any fermion operator on the lattice, we summarize the basic properties of $`\overline{)}C`$ and $`B`$.
### 2.1 Operators $`\overline{)}C`$ and $`B`$ and Brillouin zone
For a square lattice, for which we work in this paper, one can explicitly show that the simplest lattice fermion action
$$S=\overline{\psi }i\overline{)}C\psi $$
(2.6)
is invariant under the transformation
$$\psi ^{}=𝒯\psi ,\overline{\psi }^{}=\overline{\psi }𝒯^1$$
(2.7)
where $`𝒯`$ stands for any one of the following 16 operators
$$1,T_1T_2,T_1T_3,T_1T_4,T_2T_3,T_2T_4,T_3T_4,T_1T_2T_3T_4,$$
(2.8)
and
$$T_1,T_2,T_3,T_4,T_1T_2T_3,T_2T_3T_4,T_3T_4T_1,T_4T_1T_2.$$
(2.9)
The operators $`T_\mu `$ are defined by
$$T_\mu \gamma _\mu \gamma _5\mathrm{exp}(i\pi x^\mu /a)$$
(2.10)
and satisfy the relation
$$T_\mu T_\nu +T_\nu T_\mu =2\delta _{\mu \nu }$$
(2.11)
with $`T_\mu ^{}=T_\mu =T_\mu ^1`$ for anti-hermitian $`\gamma _\mu `$. We denote the 16 operators by $`𝒯_n,n=015`$, in the following with $`𝒯_0=1`$. By recalling that the operator $`T_\mu `$ adds the momentum $`\pi /a`$ to the fermion momentum $`k_\mu `$, we cover the entire Brillouin zone
$$\frac{\pi }{2a}k_\mu <\frac{3\pi }{2a}$$
(2.12)
by the operation (2.7) starting with the free fermion defined in
$$\frac{\pi }{2a}k_\mu <\frac{\pi }{2a}.$$
(2.13)
The operators in (2.8) commute with $`\gamma _5`$, whereas those in (2.9) anti-commute with $`\gamma _5`$ and thus change the sign of chiral charge, reproducing the 15 species doublers for (2.6) with correct chiral charge assignment; $`_{n=0}^{15}(1)^n\gamma _5=0`$.
One may define the near continuum configurations by the momentum $`k_\mu `$ carried by the fermion
$$\frac{\pi }{2a}ϵk_\mu \frac{\pi }{2a}ϵ$$
(2.14)
or
$$\frac{\pi }{2}ϵak_\mu \frac{\pi }{2}ϵ$$
(2.15)
for sufficiently small $`a`$ and $`ϵ`$ combined with the operation $`𝒯_n`$ in (2.8) and (2.9). To identify each species doubler clearly in the near continuum configurations, we also keep $`r/a`$ and $`m_0/a`$ finite for $`a`$ small , and the gauge fields are assumed to be sufficiently smooth. For these configurations, we can approximate the operator $`D_W`$ by
$$D_W=i\overline{)}D+M_n+O(ϵ^2)+O(a(gA_\mu )^2)$$
(2.16)
for each species doubler, where the mass parameters $`M_n`$ stand for $`M_0=\frac{m_0}{a}`$ and one of
$`{\displaystyle \frac{2r}{a}}{\displaystyle \frac{m_0}{a}},(4,1);{\displaystyle \frac{4r}{a}}{\displaystyle \frac{m_0}{a}},(6,1)`$
$`{\displaystyle \frac{6r}{a}}{\displaystyle \frac{m_0}{a}},(4,1);{\displaystyle \frac{8r}{a}}{\displaystyle \frac{m_0}{a}},(1,1)`$ (2.17)
for $`n=115`$. Here we denoted (multiplicity, chiral charge) in the bracket for species doublers. In (2.16) we used the relation valid for the configurations (2.15), for example,
$`D_We^{ikx}`$ $``$ $`{\displaystyle \underset{y}{}}D_W(x,y)e^{iky}`$ (2.18)
$`=`$ $`[{\displaystyle \underset{\mu }{}}\gamma ^\mu {\displaystyle \frac{\mathrm{sin}ak_\mu }{a}}+{\displaystyle \frac{r}{a}}{\displaystyle \underset{\mu }{}}(1\mathrm{cos}ak_\mu ){\displaystyle \frac{m_0}{a}}]e^{ikx}`$
$`=`$ $`[\gamma ^\mu k_\mu (1+O(ϵ^2))+{\displaystyle \frac{r}{a}}O(ϵ^2){\displaystyle \frac{m_0}{a}}]e^{ikx}`$
for vanishing gauge fields.
For the near continuum configurations, we thus have from (2.3)
$$D_W^{(2k+1)}=i(\overline{)}D)^{2k+1}+M_n^{(2k+1)}+O(ϵ^2)$$
(2.19)
where the mass parameters $`M_n^{(2k+1)}`$ stand for
$$M_0^{(2k+1)}(\frac{m_0}{a})^{2k+1}$$
(2.20)
and one of
$`({\displaystyle \frac{2r}{a}})^{2k+1}({\displaystyle \frac{m_0}{a}})^{2k+1},(4,1);({\displaystyle \frac{4r}{a}})^{2k+1}({\displaystyle \frac{m_0}{a}})^{2k+1},(6,1)`$
$`({\displaystyle \frac{6r}{a}})^{2k+1}({\displaystyle \frac{m_0}{a}})^{2k+1},(4,1);({\displaystyle \frac{8r}{a}})^{2k+1}({\displaystyle \frac{m_0}{a}})^{2k+1},(1,1)`$ (2.21)
for $`n=115`$, in the same notation as in (2.17).
To avoid the appearance of species doublers in $`\gamma _5D`$, we choose $`M_0^{(2k+1)}<0`$ and all other mass parameters $`M_n^{(2k+1)}>0,n0,`$ namely
$$0<m_0<2r.$$
(2.22)
The choice
$$2m_0^{2k+1}=1$$
(2.23)
normalizes properly the Dirac operator $`H_{(2k+1)}`$ in (2.2)
$$H_{(2k+1)}(i\gamma _5a\overline{)}D)^{2k+1}+\gamma _5(i\gamma _5a\overline{)}D)^{2(2k+1)}$$
(2.24)
in the near continuum configurations for all $`|M_n|`$ large.
## 3 Evaluation of the lattice Jacobian
For an operator $`O(x,y)`$ defined on the lattice, one may define
$$O_{mn}\underset{x,y}{}\varphi _m^{}(x)O(x,y)\varphi _n(y),$$
(3.1)
and the trace
$`TrO`$ $`=`$ $`{\displaystyle \underset{n}{}}O_{nn}`$ (3.2)
$`=`$ $`{\displaystyle \underset{n}{}}{\displaystyle \underset{x,y}{}}\varphi _n^{}(x)O(x,y)\varphi _n(y)`$
$`=`$ $`{\displaystyle \underset{x}{}}({\displaystyle \underset{n,y}{}}\varphi _n^{}(x)O(x,y)\varphi _n(y)).`$
The local version of the trace (or anomaly) is then defined by
$$trO(x,x)\underset{n,y}{}\varphi _n^{}(x)O(x,y)\varphi _n(y).$$
(3.3)
For the operator of our interest, we have
$`tr\mathrm{\Gamma }_5(x)`$
$`=tr[\gamma _5(\gamma _5aD)^{2k+1}]`$
$`=tr(\gamma _5aD)^{2k+1}`$
$`=tr{\displaystyle \frac{1}{2}}\gamma _5[1+D_W^{(2k+1)}{\displaystyle \frac{1}{\sqrt{(D_W^{(2k+1)})^{}D_W^{(2k+1)}}}}]`$
$`=tr{\displaystyle \frac{1}{2}}\gamma _5[D_W^{(2k+1)}{\displaystyle \frac{1}{\sqrt{(D_W^{(2k+1)})^{}D_W^{(2k+1)}}}}]`$ (3.4)
$`={\displaystyle \frac{1}{2}}{\displaystyle \underset{n=0}{\overset{15}{}}}tr{\displaystyle _{\frac{\pi }{2a}}^{\frac{\pi }{2a}}}{\displaystyle \frac{d^4k}{(2\pi )^4}}e^{ikx}𝒯_n^1\gamma _5D_W^{(2k+1)}{\displaystyle \frac{1}{\sqrt{(D_W^{(2k+1)})^{}D_W^{(2k+1)}}}}𝒯_ne^{ikx}`$
where we used the plane wave basis defined in the domain (2.13) combined with the operation $`𝒯_n`$. In this calculation, we repeatedly used the relation
$$tr\gamma _5=0$$
(3.5)
which is expected to be valid in lattice theory. We also used a short hand notation
$$Oe^{ikx}=\underset{y}{}O(x,y)e^{iky}.$$
(3.6)
There are various ways to evaluate the above trace (3.4). We evaluate the trace (3.4) by following the procedure used for the overlap Dirac operator in Refs.. Some of the basic papers of the lattice anomaly calculation are found in -. In this section we simplify the expression of the Jacobian, and its explicit evaluation is presented in the next section.
### 3.1 General analysis of the trace
Our starting formula is (by using the momentum domain (2.12))
$$\frac{1}{2}(\frac{1}{a})^4_{\frac{\pi }{2}}^{\frac{3\pi }{2}}\frac{d^4p}{(2\pi )^4}tr\gamma _5\stackrel{~}{D}_W^{(2k+1)}(p)\frac{1}{\sqrt{(\stackrel{~}{D}_W^{(2k+1)}(p))^{}\stackrel{~}{D}_W^{(2k+1)}(p)}}$$
(3.7)
with
$`\stackrel{~}{D}_W^{(2k+1)}(p)`$ $``$ $`(a^{2k+1})D_W^{(2k+1)}(p)`$ (3.8)
$`=`$ $`i[i{\displaystyle \underset{\mu }{}}\gamma ^\mu \mathrm{sin}p_\mu +a\stackrel{~}{\overline{)}C}]^{2k+1}+[r{\displaystyle \underset{\mu }{}}(1\mathrm{cos}p_\mu )+a\stackrel{~}{B}]^{2k+1}`$
$`(m_0)^{2k+1}`$
and we defined the integration variable
$$p_\mu =ak_\mu .$$
(3.9)
We used the definitions
$$e^{ikx}a\overline{)}Ce^{ikx}h(x)[i\underset{\mu }{}\gamma ^\mu \mathrm{sin}ak_\mu +a\stackrel{~}{\overline{)}C}]h(x)$$
(3.10)
and
$$e^{ikx}aBe^{ikx}h(x)[r\underset{\mu }{}(1\mathrm{cos}ak_\mu )+a\stackrel{~}{B}]h(x).$$
(3.11)
for a sufficiently smooth function $`h(x)`$. In the following we often omit writing $`h(x)`$. Consequently,
$`D_W^{(2k+1)}(k_\mu )`$ $``$ $`e^{ikx}D_W^{(2k+1)}e^{ikx}`$ (3.12)
$`=`$ $`i[i{\displaystyle \underset{\mu }{}}\gamma ^\mu {\displaystyle \frac{\mathrm{sin}ak_\mu }{a}}+\stackrel{~}{\overline{)}C}]^{2k+1}+[{\displaystyle \frac{r}{a}}{\displaystyle \underset{\mu }{}}(1\mathrm{cos}ak_\mu )+\stackrel{~}{B}]^{2k+1}`$
$`({\displaystyle \frac{m_0}{a}})^{2k+1}`$
and $`D_W^{(2k+1)}(p)`$ is defined by setting $`k_\mu =p_\mu /a`$ in $`D_W^{(2k+1)}(k_\mu )`$.
In the continuum limit $`a0`$ with $`p_\mu =ak_\mu `$ kept fixed, the operator $`\stackrel{~}{\overline{)}C}`$ approaches
$$\stackrel{~}{\overline{)}C}=\underset{\mu }{}\gamma ^\mu (\mathrm{cos}ak_\mu _\mu +ig\mathrm{cos}ak_\mu A_\mu )+O(a)=\underset{\mu }{}\gamma ^\mu \mathrm{cos}p_\mu D_\mu +O(a)$$
(3.13)
and the leading term of $`\stackrel{~}{B}`$ is known to be
$$\stackrel{~}{B}=ir\underset{\mu }{}\mathrm{sin}ak_\mu D_\mu +O(a)=ir\underset{\mu }{}\mathrm{sin}p_\mu D_\mu +O(a)$$
(3.14)
with the covariant derivative defined by
$$D_\mu _\mu +igA_\mu .$$
(3.15)
Note that the “conventional naive continuum limit” is defined by $`a0`$ with $`k_\mu `$ kept fixed, instead of $`p_\mu =ak_\mu `$ being kept fixed as in the above limit.
In the denominator of (3.7), one has a factor
$`(\stackrel{~}{D}_W^{(2k+1)}(p))^{}\stackrel{~}{D}_W^{(2k+1)}(p)`$
$`=[(i{\displaystyle \underset{\mu }{}}\gamma ^\mu \mathrm{sin}p_\mu +a\stackrel{~}{\overline{)}C})^2]^{2k+1}`$
$`+\{[r{\displaystyle \underset{\mu }{}}(1\mathrm{cos}p_\mu )+a\stackrel{~}{B}]^{2k+1}(m_0)^{2k+1}\}^2`$
$`i[[i{\displaystyle \underset{\mu }{}}\gamma ^\mu \mathrm{sin}p_\mu +a\stackrel{~}{\overline{)}C}]^{2k+1},[r{\displaystyle \underset{\mu }{}}(1\mathrm{cos}p_\mu )+a\stackrel{~}{B}]^{2k+1}]`$
$`=\{{\displaystyle \underset{\mu }{}}(\mathrm{sin}p_\mu ai\stackrel{~}{C}_\mu )^2+{\displaystyle \frac{a^2}{4}}[\gamma ^\mu ,\gamma ^\nu ][\stackrel{~}{C}_\mu ,\stackrel{~}{C}_\nu ]\}^{2k+1}`$
$`+\{[r{\displaystyle \underset{\mu }{}}(1\mathrm{cos}p_\mu )+a\stackrel{~}{B}]^{2k+1}(m_0)^{2k+1}\}^2`$
$`i[[i{\displaystyle \underset{\mu }{}}\gamma ^\mu \mathrm{sin}p_\mu +a\stackrel{~}{\overline{)}C}]^{2k+1},[r{\displaystyle \underset{\mu }{}}(1\mathrm{cos}p_\mu )+a\stackrel{~}{B}]^{2k+1}].`$ (3.16)
Note that the first two terms in this last expression commute with $`\gamma _5`$, while the last term anti-commutes with $`\gamma _5`$. The last term of (3.16) is the interference term: From the structure of the commutator, one can confirm that it consists of terms with a factor
$$ia^2\gamma ^\mu [\stackrel{~}{C}_\mu ,\stackrel{~}{B}]=ia^2gr\underset{\mu ,\nu }{}\gamma ^\mu \mathrm{cos}p_\mu \mathrm{sin}p_\nu F_{\mu \nu }+O(a^3)$$
(3.17)
and the $`2k`$ factors of $`[i_\mu \gamma ^\mu \mathrm{sin}p_\mu +a\stackrel{~}{\overline{)}C}]`$ and the $`2k`$ factors of $`[r_\mu (1\mathrm{cos}p_\mu )+a\stackrel{~}{B}]`$. We also note that
$$[\stackrel{~}{C}_\mu ,\stackrel{~}{C}_\nu ]=ig\mathrm{cos}p_\mu \mathrm{cos}p_\nu F_{\mu \nu }+O(a).$$
(3.18)
To simplify various expressions in the following, we define the variables
$$c_\mu =\mathrm{cos}ak_\mu =\mathrm{cos}p_\mu ,s_\mu =\mathrm{sin}ak_\mu =\mathrm{sin}p_\mu $$
(3.19)
and
$$\overline{)}s=\underset{\mu }{}\gamma ^\mu \mathrm{sin}p_\mu ,s^2=\underset{\mu }{}(s_\mu )^2.$$
(3.20)
### 3.2 Contribution of mass terms
We now examine the integrand of (3.7) with only the “mass terms” in the numerator retained
$$\frac{1}{a^4}tr\gamma _5\{[r\underset{\mu }{}(1c_\mu )+a\stackrel{~}{B}]^{2k+1}(m_0)^{2k+1}\}\frac{1}{\sqrt{(\stackrel{~}{D}_W^{(2k+1)}(p))^{}\stackrel{~}{D}_W^{(2k+1)}(p)}}.$$
(3.21)
The numarator contains no $`\gamma ^\mu `$’s. We expand the denominator in powers of the interference term, which contains an odd number of $`\gamma ^\mu `$’s. By remembering that the rest of the denominator factor contains an even number of $`\gamma ^\mu `$’s and thus commute with $`\gamma _5`$, the odd powers in the interference term in this expansion anti-commute with $`\gamma _5`$ and thus vanish after taking the trace. Only the even powers in the interference term could survive the trace operation with $`\gamma _5`$. Since the interference term is of order $`O(a^2)`$, only the zeroth order term and the second order term in the interference term are important in the limit $`a0`$.
We can see that the second order term in the interference vanishes. Since the second order term is already of order $`O(a^4)`$ because of (3.17), one can set $`a=0`$ in all the remaining $`2k`$ powers of $`[i\overline{)}s+a\stackrel{~}{\overline{)}C}]`$ and the $`2k`$ powers of $`[r_\mu (1c_\mu )+a\stackrel{~}{B}]`$. Namely these terms are replaced by $`i\overline{)}s`$ and $`r_\mu (1c_\mu )`$, respectively. Since these factors commute with each other, the interference term consists of a sum of terms of the structure
$$(i\overline{)}s)^lia^2gr\underset{\mu ,\nu }{}\gamma ^\mu c_\mu s_\nu (i\overline{)}s)^m[r\underset{\mu }{}(1c_\mu )]^{2k}$$
(3.22)
with $`l+m=2k`$, if one uses (3.17): Thess terms are linear in $`\gamma ^\mu `$ if one uses the relations
$`(i\overline{)}s)^2=s^2,`$
$`(i\overline{)}s)\gamma ^\nu (i\overline{)}s)=2(\overline{)}s)s_\nu s^2\gamma ^\nu .`$ (3.23)
We thus have only two $`\gamma ^\mu `$’s with order $`O(a^4)`$ in the numerator of (3.21), which vanishes after the trace with $`\gamma _5`$. Note that the terms which contain $`\gamma ^\mu `$’s in the denominator is of order $`O(a^2)`$ as in (3.16), and thus these terms cannot be used to supply extra $`\gamma ^\mu `$’s.
We have thus established that only the zeroth order term in the interference survives in (3.21), namely, one can set the interference term to $`0`$ in the denominator. In this case, from the expression in (3.16) we see that the $`\gamma ^\mu `$ factors appear only in the combination
$$\frac{a^2}{4}[\gamma ^\mu ,\gamma ^\nu ][\stackrel{~}{C}_\mu ,\stackrel{~}{C}_\nu ].$$
(3.24)
The second power in this factor is just sufficient to survive the trace operation and cancel $`1/a^4`$ in front of the integral. This means that we can set $`a=0`$ everywhere except in the prefactor in (3.24). The expression (3.21) is then replaced by
$$\frac{1}{a^4}tr\gamma _5\{[r\underset{\mu }{}(1c_\mu )]^{2k+1}(m_0)^{2k+1}\}\frac{1}{\sqrt{F_{(k)}}}$$
(3.25)
where
$`F_{(k)}`$ $``$ $`\{s^2+ig{\displaystyle \frac{a^2}{4}}[\gamma ^\mu ,\gamma ^\nu ]c_\mu c_\nu F_{\mu \nu }\}^{2k+1}`$ (3.26)
$`+\{[r{\displaystyle \underset{\mu }{}}(1c_\mu )]^{2k+1}(m_0)^{2k+1}\}^2.`$
### 3.3 Contribution of kinetic term
Similarly, we analyze the integrand of (3.7) with the “kinetic term” in the numerator
$$\frac{1}{a^4}tr\gamma _5\{i[i\overline{)}s+a\stackrel{~}{\overline{)}C}]^{2k+1}\}\frac{1}{\sqrt{(\stackrel{~}{D}_W^{(2k+1)}(p))^{}\stackrel{~}{D}_W^{(2k+1)}(p)}}.$$
(3.27)
Since the numerator is now odd in powers of $`\gamma ^\mu `$, only the odd powers of the interference term could survive. The third power of the interference is $`O(a^6)`$, and only the first power in the interference need to be analyzed.
We first rewrite the numerator factor as
$`[i\overline{)}s+a\stackrel{~}{\overline{)}C}]^{2k+1}`$ (3.28)
$`=\{{\displaystyle \underset{\mu }{}}(s_\mu ai\stackrel{~}{C}_\mu )^2+{\displaystyle \frac{a^2}{4}}[\gamma ^\mu ,\gamma ^\nu ][\stackrel{~}{C}_\mu ,\stackrel{~}{C}_\nu ]\}^k[i\overline{)}s+a\stackrel{~}{\overline{)}C}]`$
which shows that we have only one $`\gamma ^\mu `$ which is not multiplied by $`a`$. As we have already explained, the interference term in the denominator
$$i[[i\overline{)}s+a\stackrel{~}{\overline{)}C}]^{2k+1},[r\underset{\mu }{}(1c_\mu )+a\stackrel{~}{B}]^{2k+1}]$$
(3.29)
is written as a sum of terms with a single commutator
$$ia^2\gamma ^\mu [\stackrel{~}{C}_\mu ,\stackrel{~}{B}]=ia^2gr\underset{\mu ,\nu }{}\gamma ^\mu c_\mu s_\nu F_{\mu \nu }+O(a^3)$$
(3.30)
multiplied by the $`2k`$ factors of $`[i\overline{)}s+a\stackrel{~}{\overline{)}C}]`$ and the $`2k`$ factors of $`[r_\mu (1c_\mu )+a\stackrel{~}{B}]`$. In such a term, if one exchanges the order of $`[i\overline{)}s+a\stackrel{~}{\overline{)}C}]`$ and $`[r_\mu (1c_\mu )+a\stackrel{~}{B}]`$, one generates another commutator as in (3.30). We then have a factor $`a^4`$ and thus we can set $`a=0`$ in all other terms in the integrand. We recognize that such a term contains $`\gamma ^\mu `$ only in the combination with $`\overline{)}s`$ and the two factors of the above commutator (3.30). From this combination together with $`\overline{)}s`$ in the numerator (3.28), we cannot form a non-vanishing contraction with the antisymmetric $`ϵ^{\mu \nu \alpha \beta }`$ tensor.
This means that we can write the interference term as a sum of terms of the structure
$`[i\overline{)}s+a\stackrel{~}{\overline{)}C}]^l(i)a^2gr{\displaystyle \underset{\mu ,\nu }{}}\gamma ^\mu c_\mu s_\nu F_{\mu \nu }[i\overline{)}s+a\stackrel{~}{\overline{)}C}]^m`$
$`\times [r{\displaystyle \underset{\mu }{}}(1c_\mu )+a\stackrel{~}{B}]^{2k}`$ (3.31)
where $`l+m=2k`$. If both of $`l`$ and $`m`$ are even, we can use
$$[i\overline{)}s+a\stackrel{~}{\overline{)}C}]^2=\underset{\mu }{}(s_\mu ai\stackrel{~}{C}_\mu )^2+\frac{a^2}{4}[\gamma ^\mu ,\gamma ^\nu ][\stackrel{~}{C}_\mu ,\stackrel{~}{C}_\nu ]$$
(3.32)
and $`\gamma ^\mu `$ always appears in the combination $`a\gamma ^\mu `$ except for the numerator term (3.28) which contains $`\overline{)}s`$, and the commutator term (3.30), which contains $`a^2\gamma _\mu `$. Such a combination could give rise to a non-vanishing result. On the other hand, if both of $`l`$ and $`m`$ are odd, one has to deal with a left-over term
$`[i\overline{)}s+a\stackrel{~}{\overline{)}C}](i)a^2gr{\displaystyle \underset{\mu ,\nu }{}}\gamma ^\mu c_\mu s_\nu F_{\mu \nu }[i\overline{)}s+a\stackrel{~}{\overline{)}C}]`$
$`=[i\overline{)}s+a\stackrel{~}{\overline{)}C}](2)(i)a^2gr{\displaystyle \underset{\mu ,\nu }{}}c_\mu s_\nu F_{\mu \nu }[is_\mu +a\stackrel{~}{C}_\mu ]`$
$`+[i\overline{)}s+a\stackrel{~}{\overline{)}C}](i)a^3gr{\displaystyle \underset{\mu ,\nu ,\alpha }{}}\gamma ^\alpha \gamma ^\mu c_\alpha c_\mu s_\nu D_\alpha F_{\mu \nu }`$
$`\{{\displaystyle \underset{\mu }{}}(s_\mu ai\stackrel{~}{C}_\mu )^2+{\displaystyle \frac{a^2}{4}}[\gamma ^\mu ,\gamma ^\nu ][\stackrel{~}{C}_\mu ,\stackrel{~}{C}_\nu ]\}(i)a^2gr{\displaystyle \underset{\mu ,\nu }{}}\gamma ^\mu c_\mu s_\nu F_{\mu \nu }.`$
where we used $`\gamma ^\mu \gamma ^\alpha +\gamma ^\alpha \gamma ^\mu =2\eta ^{\mu \alpha }`$. The first term of this equation contains the factor $`a^2[i\overline{)}s+a\stackrel{~}{\overline{)}C}]`$, which should be replaced by $`a^2i\overline{)}s`$ and should be combined with the factor $`i\overline{)}s`$ in the numerator factor (together with $`a^2[\gamma ^\mu ,\gamma ^\nu ][\stackrel{~}{C}_\mu ,\stackrel{~}{C}_\nu ]`$ from other factors) to obtain a possible non-zero result. But such a term cannot make a non-vanishing contraction with $`ϵ^{\mu \nu \alpha \beta }`$. The second term is of orfer $`O(a^3)`$ and contains 3 $`\gamma ^\mu `$’s. If this term is combined with $`[i\overline{)}s+a\stackrel{~}{\overline{)}C}]`$ in the numerator, it becomes of order $`O(a^5)`$ due to (3.32). If it is combined with the drivative operator such as $`a\stackrel{~}{C}_\mu `$ in (3.16), when commuting with other denominator factors, it becomes $`O(a^4)`$; such a term contains $`\overline{)}s^2=s^2`$ if combined with $`\overline{)}s`$ in the numerator and vanishes after trace with $`\gamma _5`$. Thus we can set the first and second terms to $`0`$ in the above equation (3.33).
Only the last term in (3.33) can survive: Among those surviving terms, the even $`l`$ terms and odd $`l`$ terms cancel pairwise except one term in the interference term.
By this way, we can write the total interference term as
$`(2k+1)ia^2gr{\displaystyle \underset{\mu ,\nu }{}}\gamma ^\mu c_\mu s_\nu F_{\mu \nu }\{s^2+ig{\displaystyle \frac{a^2}{4}}[\gamma ^\mu ,\gamma ^\nu ]c_\mu c_\nu F_{\mu \nu }]\}^k`$
$`\times [r{\displaystyle \underset{\mu }{}}(1c_\mu )]^{2k}`$ (3.34)
where the factor $`2k+1`$ comes from the $`2k+1`$ powers of $`[r_\mu (1c_\mu )+a\stackrel{~}{B}]`$. We also set $`a=0`$ in all the terms without $`\gamma ^\mu `$ since this does not influence the surviving terms. We can also set $`a=0`$ in the numerator factor except for the combination $`a^2[\gamma ^\mu ,\gamma ^\nu ]`$. Note that the order of $`\gamma ^\mu `$ and $`[\gamma ^\mu ,\gamma ^\nu ]`$ can be changed freely in the expansion of the denominator in powers of $`a^2`$, since the surviving terms are contracted with $`\gamma _5`$ to give rise to $`ϵ^{\mu \nu \alpha \beta }`$.
To summarize this tedious analysis, we can write the integrand with the “kinetic” term (3.27) as
$`{\displaystyle \frac{1}{a^4}}tr\gamma _5\{(2k+1)r(ig{\displaystyle \frac{a^2}{4}})[\gamma ^\mu ,\gamma ^\nu ]c_\mu c_\nu F_{\mu \nu }({\displaystyle \underset{\alpha }{}}{\displaystyle \frac{s_\alpha ^2}{4c_\alpha }})`$
$`\times \{s^2+ig{\displaystyle \frac{a^2}{4}}[\gamma ^\mu ,\gamma ^\nu ]c_\mu c_\nu F_{\mu \nu }]\}^{2k}[r{\displaystyle \underset{\mu }{}}(1c_\mu )]^{2k}\}{\displaystyle \frac{1}{F_{(k)}^{3/2}}}`$ (3.35)
where $`F_{(k)}`$ is defined in (3.26). In writing this final expression we used the following sequence of rewriting
$`\overline{)}s{\displaystyle \underset{\mu ,\nu }{}}\gamma ^\mu c_\mu s_\nu F_{\mu \nu }`$ $`=`$ $`{\displaystyle \underset{\alpha }{}}{\displaystyle \underset{\mu ,\nu }{}}\gamma ^\alpha \gamma ^\mu c_\mu s_\nu s_\alpha F_{\mu \nu }`$ (3.36)
$`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{\alpha }{}}{\displaystyle \underset{\mu ,\nu }{}}[\gamma ^\alpha ,\gamma ^\mu ]c_\mu s_\nu s_\nu \delta _{\nu ,\alpha }F_{\mu \nu }`$
$`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{\mu ,\nu }{}}[\gamma ^\nu ,\gamma ^\mu ]c_\mu c_\nu (s_\nu s_\nu /c_\nu )F_{\mu \nu }`$
$`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{\mu ,\nu }{}}[\gamma ^\mu ,\gamma ^\nu ]c_\mu c_\nu F_{\mu \nu }({\displaystyle \underset{\alpha }{}}s_\alpha ^2/4c_\alpha )`$
Namely, only the term with two $`\gamma ^\mu `$’s contributes to the final result, and the odd term in $`s_\nu `$ after integration over the momentum vanishes. In the last step, we used the lattice hypercubic symmetry by taking into account the contraction with the $`ϵ^{\mu \nu \alpha \beta }`$ symbol later.
## 4 Formula for the chiral anomaly and parameter independence
The basic formula for the chiral anomaly is given by (3.7), (3.25) and (3.35). The next step is to expand the integrand in the powers of
$`(ig\frac{a^2}{4})[\gamma ^\mu ,\gamma ^\nu ]c_\mu c_\nu F_{\mu \nu }`$ and retain only the terms which contain the second power of this factor. We then combine the expansion with the formula
$$tr\gamma _5\{(ig\frac{1}{4})[\gamma ^\mu ,\gamma ^\nu ]c_\mu c_\nu F_{\mu \nu }\}^2=(\underset{\alpha =1}{\overset{4}{}}c_\alpha )g^2trϵ^{\mu \nu \alpha \beta }F_{\mu \nu }F_{\alpha \beta }$$
(4.1)
where $`ϵ^{1234}=1`$. We thus write only the coefficients of the factor
$$g^2trϵ^{\mu \nu \alpha \beta }F_{\mu \nu }F_{\alpha \beta }$$
(4.2)
in the following.
The contribution from the “mass terms” (3.25) is given by
$$\frac{1}{16}_{\frac{\pi }{2}}^{\frac{3\pi }{2}}\frac{d^4p}{(2\pi )^4}(\underset{\alpha =1}{\overset{4}{}}c_\alpha )\frac{(2k+1)M_{(2k+1)}[3(2k+1)(s^2)^{4k}4k(s^2)^{2k1}H]}{H^{5/2}}$$
(4.3)
where
$`H(s^2)^{2k+1}+M_{(2k+1)}^2`$
$`M_{(2k+1)}[r{\displaystyle \underset{\mu }{}}(1c_\mu )]^{2k+1}m_0^{2k+1}.`$ (4.4)
The contribution from the “kinetic term” (3.35) is written as
$`{\displaystyle \frac{1}{16}}{\displaystyle _{\frac{\pi }{2}}^{\frac{3\pi }{2}}}{\displaystyle \frac{d^4p}{(2\pi )^4}}({\displaystyle \underset{\alpha =1}{\overset{4}{}}}c_\alpha )(2k+1)\{r({\displaystyle \underset{\beta }{}}{\displaystyle \frac{s_\beta ^2}{c_\beta }})[r{\displaystyle \underset{\mu }{}}(1c_\mu )]^{2k}\}`$
$`\times \{4k(s^2)^{2k1}H3(2k+1)(s^2)^{4k}\}{\displaystyle \frac{1}{H^{5/2}}}.`$ (4.5)
Thus the total contribution is given by
$`I_{2k+1}`$ $`=`$ $`{\displaystyle \frac{2k+1}{16}}{\displaystyle _{\frac{\pi }{2}}^{\frac{3\pi }{2}}}{\displaystyle \frac{d^4p}{(2\pi )^4}}({\displaystyle \underset{\alpha =1}{\overset{4}{}}}c_\alpha )\{M_{(2k+1)}r({\displaystyle \underset{\beta }{}}{\displaystyle \frac{s_\beta ^2}{c_\beta }})[r{\displaystyle \underset{\mu }{}}(1c_\mu )]^{2k}\}`$ (4.6)
$`\times \{3(2k+1)(s^2)^{4k}4k(s^2)^{2k1}H\}{\displaystyle \frac{1}{H^{5/2}}}.`$
### 4.1 Parameter independence
To analyze the parameter independence of the coefficient of the chiral anomaly, we follow the procedure in Refs.. We first rewrite the integral in the domain $`\frac{\pi }{2}p_\mu \frac{3\pi }{2}`$ to the integral in the domain $`\frac{\pi }{2}p_\mu \frac{\pi }{2}`$ by using the variables $`s_\mu =\mathrm{sin}p_\mu `$ as
$$_{\frac{\pi }{2}}^{\frac{3\pi }{2}}\frac{d^4p}{(2\pi )^4}(\underset{\alpha =1}{\overset{4}{}}c_\alpha )=\underset{ϵ_\mu =\pm }{}(\underset{\mu }{}ϵ_\mu )_1^1\frac{d^4s}{(2\pi )^4}$$
(4.7)
where the symbol $`ϵ_\mu `$ takes care of the 16 (would-be) species doublers
$$ϵ_\mu =(\pm ,\pm ,\pm ,\pm ).$$
(4.8)
The following formula is also valid
$$c_\mu =ϵ_\mu (1s_\mu ^2)^{1/2}.$$
(4.9)
In this new notation, we have
$`H(s)(s^2)^{2k+1}+M_{(2k+1)}^2`$
$`M_{(2k+1)}[r{\displaystyle \underset{\mu }{}}(1ϵ_\mu (1s_\mu ^2)^{1/2})]^{2k+1}m_0^{2k+1}`$ (4.10)
and we evaluate
$`(2k+1)H(s)+{\displaystyle \frac{1}{5}}H(s)^{7/2}{\displaystyle \underset{\mu }{}}s_\mu {\displaystyle \frac{}{s_\mu }}H(s)^{5/2}`$
$`=(2k+1)H(s)+{\displaystyle \frac{1}{3}}H(s)^{5/2}{\displaystyle \underset{\mu }{}}s_\mu {\displaystyle \frac{}{s_\mu }}H(s)^{3/2}`$
$`=(2k+1)M_{(2k+1)}\{[r{\displaystyle \underset{\mu }{}}(1ϵ_\mu (1s_\mu ^2)^{1/2})]^{2k+1}+m_0^{2k+1}`$
$`+[r{\displaystyle \underset{\mu }{}}(1ϵ_\mu (1s_\mu ^2)^{1/2})]^{2k}r{\displaystyle \underset{\mu }{}}ϵ_\mu s_\mu ^2(1s_\mu ^2)^{1/2}\}`$
$`=(2k+1)M_{(2k+1)}`$ (4.11)
$`\times \{M_{(2k+1)}[r{\displaystyle \underset{\mu }{}}(1ϵ_\mu (1s_\mu ^2)^{1/2})]^{2k}r{\displaystyle \underset{\mu }{}}s_\mu ^2ϵ_\mu (1s_\mu ^2)^{1/2}\}`$
which is a generalization of the identity discussed in . By using these relations one can prove
$`{\displaystyle \frac{1}{2k+1}}[{\displaystyle \underset{\mu }{}}s_\mu {\displaystyle \frac{}{s_\mu }}+4]({\displaystyle \frac{(s^2)^{4k}}{H(s)^{5/2}}})`$
$`={\displaystyle \frac{(s^2)^{4k}}{H(s)^{5/2}}}`$
$`+\{M_{(2k+1)}[r{\displaystyle \underset{\mu }{}}(1ϵ_\mu (1s_\mu ^2)^{1/2})]^{2k}r{\displaystyle \underset{\mu }{}}s_\mu ^2ϵ_\mu (1s_\mu ^2)^{1/2}\}`$
$`\times 5M_{(2k+1)}({\displaystyle \frac{(s^2)^{4k}}{H(s)^{7/2}}})`$ (4.12)
and
$`{\displaystyle \frac{1}{2k+1}}[{\displaystyle \underset{\mu }{}}s_\mu {\displaystyle \frac{}{s_\mu }}+4]({\displaystyle \frac{(s^2)^{2k1}}{H(s)^{3/2}}})`$
$`={\displaystyle \frac{(s^2)^{2k1}}{H(s)^{3/2}}}`$
$`+\{M_{(2k+1)}[r{\displaystyle \underset{\mu }{}}(1ϵ_\mu (1s_\mu ^2)^{1/2})]^{2k}r{\displaystyle \underset{\mu }{}}s_\mu ^2ϵ_\mu (1s_\mu ^2)^{1/2}\}`$
$`\times 3M_{(2k+1)}({\displaystyle \frac{(s^2)^{2k1}}{H(s)^{5/2}}}).`$ (4.13)
One can then show that
$`{\displaystyle \frac{I_{2k+1}}{m_0^{2k+1}}}={\displaystyle \frac{2k+1}{16}}{\displaystyle \underset{ϵ_\mu =\pm }{}}({\displaystyle \underset{\mu }{}}ϵ_\mu ){\displaystyle _1^1}{\displaystyle \frac{d^4s}{(2\pi )^4}}`$
$`\times \{3(2k+1)[{\displaystyle \frac{(s^2)^{4k}}{H^{5/2}}}5[M_{(2k+1)}r({\displaystyle \underset{\beta }{}}{\displaystyle \frac{s_\beta ^2}{c_\beta }})[r{\displaystyle \underset{\mu }{}}(1c_\mu )]^{2k}]M_{(2k+1)}{\displaystyle \frac{(s^2)^{4k}}{H^{7/2}}}]`$
$`4k[{\displaystyle \frac{(s^2)^{2k1}}{H^{3/2}}}3[M_{(2k+1)}r({\displaystyle \underset{\beta }{}}{\displaystyle \frac{s_\beta ^2}{c_\beta }})[r{\displaystyle \underset{\mu }{}}(1c_\mu )]^{2k}]M_{(2k+1)}{\displaystyle \frac{(s^2)^{2k1}}{H^{5/2}}}]\}`$
$`={\displaystyle \frac{2k+1}{16}}{\displaystyle \underset{ϵ_\mu =\pm }{}}({\displaystyle \underset{\mu }{}}ϵ_\mu ){\displaystyle _1^1}{\displaystyle \frac{d^4s}{(2\pi )^4}}`$
$`\times \{[3({\displaystyle \underset{\mu }{}}s_\mu {\displaystyle \frac{}{s_\mu }}+4)({\displaystyle \frac{(s^2)^{4k}}{H(s)^{5/2}}})]`$
$`+{\displaystyle \frac{4k}{2k+1}}[({\displaystyle \underset{\mu }{}}s_\mu {\displaystyle \frac{}{s_\mu }}+4)({\displaystyle \frac{(s^2)^{2k1}}{H(s)^{3/2}}})]\}`$ (4.14)
Similarly, we can show the relations by using (4.11)
$`{\displaystyle \frac{1}{2k+1}}[{\displaystyle \underset{\mu }{}}s_\mu {\displaystyle \frac{}{s_\mu }}+4]({\displaystyle \frac{[r_\mu (1c_\mu )]^{2k+1}(s^2)^{4k}}{H(s)^{5/2}}})`$
$`={\displaystyle \frac{\{[r_\mu (1c_\mu )]^{2k+1}[r_\mu (1c_\mu )]^{2k}r_\mu ϵ_\mu s_\mu ^2(1s_\mu ^2)^{1/2}\}(s^2)^{4k}}{H(s)^{5/2}}}`$
$`+\{M_{(2k+1)}[r{\displaystyle \underset{\mu }{}}(1c_\mu )]^{2k}r{\displaystyle \underset{\mu }{}}s_\mu ^2ϵ_\mu (1s_\mu ^2)^{1/2}\}`$
$`\times 5M_{(2k+1)}({\displaystyle \frac{[r_\mu (1c_\mu )]^{2k+1}(s^2)^{4k}}{H(s)^{7/2}}})`$ (4.15)
and
$`{\displaystyle \frac{1}{2k+1}}[{\displaystyle \underset{\mu }{}}s_\mu {\displaystyle \frac{}{s_\mu }}+4]({\displaystyle \frac{[r_\mu (1c_\mu )]^{2k+1}(s^2)^{2k1}}{H(s)^{3/2}}})`$
$`={\displaystyle \frac{\{[r_\mu (1c_\mu )]^{2k+1}[r_\mu (1c_\mu )]^{2k}r_\mu ϵ_\mu s_\mu ^2(1s_\mu ^2)^{1/2}\}(s^2)^{2k1}}{H(s)^{3/2}}}`$
$`+\{M_{(2k+1)}[r{\displaystyle \underset{\mu }{}}(1c_\mu )]^{2k}r{\displaystyle \underset{\mu }{}}s_\mu ^2ϵ_\mu (1s_\mu ^2)^{1/2}\}`$
$`\times 3M_{(2k+1)}({\displaystyle \frac{[r_\mu (1c_\mu )]^{2k+1}(s^2)^{2k1}}{H(s)^{5/2}}})`$ (4.16)
where $`c_\mu =ϵ_\mu (1s_\mu ^2)^{1/2}`$ is understood. Using these relations, we have
$`r^{2k+1}{\displaystyle \frac{I_{2k+1}}{r^{2k+1}}}={\displaystyle \frac{2k+1}{16}}{\displaystyle \underset{ϵ_\mu =\pm }{}}({\displaystyle \underset{\mu }{}}ϵ_\mu ){\displaystyle _1^1}{\displaystyle \frac{d^4s}{(2\pi )^4}}`$
$`\times \{3(2k+1)\{{\displaystyle \frac{\{[r_\mu (1c_\mu )]^{2k+1}[r_\mu (1c_\mu )]^{2k}r_\mu ϵ_\mu s_\mu ^2(1s_\mu ^2)^{1/2}\}(s^2)^{4k}}{H(s)^{5/2}}}`$
$`[M_{(2k+1)}[r{\displaystyle \underset{\mu }{}}(1c_\mu )]^{2k}r{\displaystyle \underset{\mu }{}}s_\mu ^2ϵ_\mu (1s_\mu ^2)^{1/2}]`$
$`\times 5M_{(2k+1)}({\displaystyle \frac{[r_\mu (1c_\mu )]^{2k+1}(s^2)^{4k}}{H(s)^{7/2}}})\}`$
$`4k\{{\displaystyle \frac{\{[r_\mu (1c_\mu )]^{2k+1}[r_\mu (1c_\mu )]^{2k}r_\mu ϵ_\mu s_\mu ^2(1s_\mu ^2)^{1/2}\}(s^2)^{2k1}}{H(s)^{3/2}}}`$
$`[M_{(2k+1)}[r{\displaystyle \underset{\mu }{}}(1c_\mu )]^{2k}r{\displaystyle \underset{\mu }{}}s_\mu ^2ϵ_\mu (1s_\mu ^2)^{1/2}]`$
$`\times 3M_{(2k+1)}({\displaystyle \frac{[r_\mu (1c_\mu )]^{2k+1}(s^2)^{2k1}}{H(s)^{5/2}}})\}\}`$
$`={\displaystyle \frac{2k+1}{16}}{\displaystyle \underset{ϵ_\mu =\pm }{}}({\displaystyle \underset{\mu }{}}ϵ_\mu ){\displaystyle _1^1}{\displaystyle \frac{d^4s}{(2\pi )^4}}`$
$`\times \{[3({\displaystyle \underset{\mu }{}}s_\mu {\displaystyle \frac{}{s_\mu }}+4)({\displaystyle \frac{[r_\mu (1c_\mu )]^{2k+1}(s^2)^{4k}}{H(s)^{5/2}}})]`$
$`+{\displaystyle \frac{4k}{2k+1}}[({\displaystyle \underset{\mu }{}}s_\mu {\displaystyle \frac{}{s_\mu }}+4)({\displaystyle \frac{[r_\mu (1c_\mu )]^{2k+1}(s^2)^{2k1}}{H(s)^{3/2}}})]\}.`$ (4.17)
The integrand becomes singular only if the following two relations simultaneously hold
$$s^2=0,[r\underset{\mu }{}(1ϵ_\mu (1s_\mu ^2)^{1/2})]^{2k+1}m_0^{2k+1}=0$$
(4.18)
namely, only when $`m_0/r=0,2,4,6,8`$. We are working in the physical region
$$0<m_0<2r$$
(4.19)
and thus the above integrals are regular, and we have from (4.14) and (4.17)
$$\frac{I_{2k+1}}{m_0^{2k+1}}=\frac{I_{2k+1}}{r^{2k+1}}=0$$
(4.20)
after partial integration. It can be confirmed that boundary terms at $`s_\mu =\pm 1`$ give vanishing contributions after a summation over $`_{ϵ_\mu =\pm }`$. This shows that the coefficient of the anomaly is stable under a smooth variation of the parameters $`r`$ and $`m_0`$, which is expected for a topological quantity such as the chiral anomaly.
### 4.2 Explicit evaluation of the chiral anomaly
Since the coefficient of the anomaly is independent of the parameters $`r`$ and $`m_0`$, we evaluate the anomaly in the limit where both of $`r`$ and $`m_0`$ go to $`0`$. To be precise we introduce an auxiliary parameter $`a`$, and take a limit $`a0`$ with both of
$$\frac{r}{a},\frac{m_0}{a}$$
(4.21)
kept fixed in the physical region (4.19). The parameter $`a`$ plays the role of an effective lattice spacing, though our formulas (4.3) and (4.5) are derived in the limit of the lattice spacing $`a=0`$.
There are various ways to evaluate the coefficient of the anomaly , and one of these methods is given in in the analysis of the overlap operator with $`k=0`$. We present a calculation which reveals a close connnection with the naive continuum limit. We first observe that the contribution of the “kinetic” term (4.5), which arises from the interference term in the denominator, vanishes in the above limit (4.21) <sup>2</sup><sup>2</sup>2 This property has been used in the treatment of the overlap operator in Refs. and .. To show this we examine
$`{\displaystyle \frac{1}{16}}{\displaystyle \underset{ϵ_\mu =\pm }{}}({\displaystyle \underset{\mu }{}}ϵ_\mu ){\displaystyle _1^1}{\displaystyle \frac{d^4s}{(2\pi )^4}}\{r({\displaystyle \underset{\beta }{}}{\displaystyle \frac{s_\beta ^2}{c_\beta }})[r{\displaystyle \underset{\mu }{}}(1c_\mu )]^{2k}\}`$
$`\times \{4k(s^2)^{2k1}H3(2k+1)(s^2)^{4k}\}{\displaystyle \frac{1}{H^{5/2}}}.`$ (4.22)
where $`c_\mu =ϵ_\mu (1s_\mu ^2)^{1/2}`$ is understood. We define the integration domain
$$ϵs_\mu ϵ$$
(4.23)
for all $`\mu `$ with a sufficiently small but finite $`ϵ`$. Since $`s^2>0`$ and the denominator of the integrand is regular for the integration domain outside the above domain, the integral outside the domain (4.23) vanishes in the limit $`a0`$. Note that the denominator of $`_\beta \frac{s_\beta ^2}{c_\beta }`$ does not cause any divergence in the integral (4.22). In fact one can even take $`ϵ0`$ in such a manner that
$$a/ϵ^l0$$
(4.24)
for a suitable fixed positive integer $`l`$. This is because the integral outside the domain (4.23) vanishes at least linearly in $`a`$, and thus one can let $`ϵ0`$ simultaneously with the above constraint, where the denominator $`ϵ^l`$ stands for the possible infrared singularity in this calculational procedure.
We thus examine the remaining integral
$`{\displaystyle \frac{1}{16}}{\displaystyle \underset{ϵ_\mu =\pm }{}}({\displaystyle \underset{\mu }{}}ϵ_\mu ){\displaystyle _ϵ^ϵ}{\displaystyle \frac{d^4s}{(2\pi )^4}}\{r({\displaystyle \underset{\beta }{}}{\displaystyle \frac{s_\beta ^2}{c_\beta }})[r{\displaystyle \underset{\mu }{}}(1c_\mu )]^{2k}\}`$
$`\times \{4k(s^2)^{2k1}H3(2k+1)(s^2)^{4k}\}{\displaystyle \frac{1}{H^{5/2}}}.`$ (4.25)
Since $`|c_\mu |1`$ in the above integral, we can ignore the variation of $`c_\mu `$. If one rescales the integration variable $`s_\mu =as_\mu ^{}`$ and defines
$$r^{}=r/a,m_0^{}=m_0/a$$
(4.26)
the above integral is written as
$`{\displaystyle \frac{1}{16}}{\displaystyle \underset{ϵ_\mu =\pm }{}}({\displaystyle \underset{\mu }{}}ϵ_\mu ){\displaystyle _{ϵ/a}^{ϵ/a}}{\displaystyle \frac{d^4s^{}}{(2\pi )^4}}\{r^{}({\displaystyle \underset{\beta }{}}{\displaystyle \frac{s_\beta ^2}{c_\beta }})[r^{}{\displaystyle \underset{\mu }{}}(1c_\mu )]^{2k}\}`$
$`\times \{4k((s^{})^2)^{2k1}H3(2k+1)((s^{})^2)^{4k}\}{\displaystyle \frac{1}{H^{5/2}}}`$ (4.27)
where $`H(s^{})`$ is parametrized by $`r^{}`$ and $`m_0^{}`$, which are kept fixed in the limit $`a0`$. Note that the factor $`(_\beta \frac{s_\beta ^2}{c_\beta })`$ in the numerator, which is written in terms of the original variables, is $`O(ϵ^2)`$. The above integral is convergent in this limit $`a0`$ and of order $`O(ϵ^2)`$, and thus it can be made arbitrarily small. We can even make it vanish precisely by taking the limit (4.24). We can thus ignore the contribution of the “kinetic” term, which arises from the interference term in the denominator, in the above limit (4.21) .
We now come to the main contribution of the “mass terms” in (4.3). It turns out to be more convenient to go back to (3.25), which gives rise to (4.3). If one uses the notation of (3.4) instead of (3.7), (3.25) is written as
$$\underset{n=0}{\overset{15}{}}(\frac{1}{2})(1)^n\frac{1}{a^4}_{\pi /2}^{\pi /2}\frac{d^4p}{(2\pi )^4}tr\gamma _5\{[r\underset{\mu }{}(1\pm c_\mu )]^{2k+1}(m_0)^{2k+1}\}\frac{1}{\sqrt{F_{(k)}(n,p_\mu )}}$$
(4.28)
where
$`F_{(k)}(n,p_\mu )`$ $``$ $`\{(i\overline{)}s+{\displaystyle \underset{\mu }{}}a\gamma ^\mu c_\mu D_\mu )^2\}^{2k+1}`$ (4.29)
$`+\{[r{\displaystyle \underset{\mu }{}}(1\pm c_\mu )]^{2k+1}(m_0)^{2k+1}\}^2.`$
The summation runs over the 16 would-be species doublers and the factor $`r_\mu (1\pm c_\mu )`$ arises from each momentum domain. For later convenience, we modified the denominator factor in (3.25) by replacing $`s^2+ig\frac{a^2}{4}[\gamma ^\mu ,\gamma ^\nu ]c_\mu c_\nu F_{\mu \nu }`$ with $`(i\overline{)}s+_\mu a\gamma ^\mu c_\mu D_\mu )^2`$ , but this does not change the result as was explained in detail in the passage from (3.21) to (3.25) in Section 3.2.
In the present integral, we can also define the domain
$$ϵp_\mu ϵ$$
(4.30)
for arbitrarily small but finite $`ϵ`$. One can again confirm that the integral outside this domain vanishes at least linearly in $`a`$ for $`a0`$. This is because we retain the second order term in $`a^2[\gamma ^\mu ,\gamma ^\nu ]F_{\mu \nu }`$ to survive the trace with $`\gamma _5`$, and this cancels the factor $`1/a^4`$ in front of the integral. The resulting integral is finite outside the above domain, and it vanishes at least linearly in $`a`$ in the limit (4.21). We can also let $`ϵ0`$ as in (4.24).
We thus examine (4.28) inside the domain (4.30)
$`{\displaystyle \underset{n=0}{\overset{15}{}}}({\displaystyle \frac{1}{2}})(1)^n{\displaystyle \frac{1}{a^4}}{\displaystyle _ϵ^ϵ}{\displaystyle \frac{d^4p}{(2\pi )^4}}tr\gamma _5{\displaystyle \frac{[r_\mu (1\pm c_\mu )]^{2k+1}(m_0)^{2k+1}}{\sqrt{F_{(k)}(n,p_\mu )}}}`$
$`={\displaystyle \underset{n=0}{\overset{15}{}}}({\displaystyle \frac{1}{2}})(1)^n{\displaystyle _{ϵ/a}^{ϵ/a}}{\displaystyle \frac{d^4k}{(2\pi )^4}}tr\gamma _5{\displaystyle \frac{[r_\mu (1\pm \mathrm{cos}ak_\mu )]^{2k+1}(m_0)^{2k+1}}{\sqrt{F_{(k)}(n,ak_\mu )}}}`$ (4.31)
For sufficiently small $`ϵ`$, we have
$`[r{\displaystyle \underset{\mu }{}}(1\pm \mathrm{cos}ak_\mu )]^{2k+1}/a^{2k+1}(m_0)^{2k+1}/a^{2k+1}=M_n^{(2k+1)}+O(ϵ^2)`$
$`F_{(k)}(n,ak_\mu )/a^{2(2k+1)}=(i\overline{)}k+\overline{)}D)^{2(2k+1)}+(M_n^{(2k+1)})^2+O(ϵ^2)`$ (4.32)
where the mass parameters $`M_n^{(2k+1)}`$ are defined in (2.20) and (2.21). The above integral in the limit $`a0`$ with (4.24) approaches
$`{\displaystyle \frac{1}{2}}{\displaystyle \underset{n=0}{\overset{15}{}}}(1)^ntr{\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle \frac{d^4k}{(2\pi )^4}}\gamma _5{\displaystyle \frac{M_n^{(2k+1)}}{\sqrt{([i\overline{)}k+\overline{)}D]^2)^{2k+1}+(M_n^{(2k+1)})^2}}}`$ (4.33)
$`=`$ $`{\displaystyle \frac{1}{2}}tr{\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle \frac{d^4k}{(2\pi )^4}}\gamma _5{\displaystyle \frac{1}{\sqrt{((i\overline{)}k+\overline{)}D)^2/\widehat{M}_0^2)^{2k+1}+1}}}`$
$`{\displaystyle \frac{1}{2}}{\displaystyle \underset{n=1}{\overset{15}{}}}(1)^ntr{\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle \frac{d^4k}{(2\pi )^4}}\gamma _5{\displaystyle \frac{1}{\sqrt{((i\overline{)}k+\overline{)}D)^2/\widehat{M}_n^2)^{2k+1}+1}}}`$
by recalling $`\widehat{M}_0<0`$, and $`\widehat{M}_n^2[(M_n^{(2k+1)})^2]^{1/(2k+1)}`$. The sum of integrals in (4.33) gives rise to the anomaly for all $`\widehat{M}_n^2\mathrm{}`$ in the final stage:
$$\underset{M\mathrm{}}{lim}tr_{\mathrm{}}^{\mathrm{}}\frac{d^4k}{(2\pi )^4}e^{ikx}\gamma _5f(\frac{\overline{)}D^2}{M^2})e^{ikx}=\frac{g^2}{32\pi ^2}trϵ^{\mu \nu \alpha \beta }F_{\mu \nu }F_{\alpha \beta }.$$
(4.34)
Here we defined
$$f(x)=\frac{1}{\sqrt{x^{2k+1}+1}}$$
(4.35)
which satisfies
$`f(0)=1,f(\mathrm{})=0,`$
$`f^{}(x)x|_{x=0}=f^{}(x)x|_{x=\mathrm{}}=0.`$ (4.36)
The left-hand side of (4.34) is known to be independent of the choice of $`f(x)`$ which satisfies the mild condition (4.36).
By combining (1.15), (3.4) and (4.34), we recover the Atiyah-Singer index theorem ( in continuum $`R^4`$ space)
$$n_+n_{}=𝑑x\frac{g^2}{32\pi ^2}trϵ^{\mu \nu \alpha \beta }F_{\mu \nu }F_{\alpha \beta }$$
(4.37)
in the (naive) continuum limit for any operator in (1.1), if one follows the construction of $`\gamma _5D`$ in .
## 5 Discussion
We have shown in an explicit and elementary manner that a new class of lattice Dirac operators proposed in satisfy the correct anomaly relation. We have in particular shown that the anomaly coefficient is independent of a small variation in the parameters $`r`$ and $`m_0`$, which characterize these Dirac operators. This is in agreement with the more general but formal analysis in .
The Dirac operator in is defined in a somewhat indirect manner as
$$detH=(detH_{(2k+1)})^{1/(2k+1)}$$
(5.1)
as in (A.9) in Appendix. This definition is sufficient for the non-perturbative analysis, and as we have shown in this paper, it is also sufficient to evaluate the chiral anomaly explicitly. However, the perturbative treatment of these general class of Dirac operators is not well understood yet. It would therefore be interesting to extend the analyses in , for example, to the present general class of operators.
As for the locality issue of the present class of operators, the fact that the anomaly calculation in a naive continuum limit makes sense suggests that there is a certain range of gauge field configurations which make these operators local. The fact that we take a 2k+1th root of $`H_{(2k+1)}`$ in (5.1) by itself may not spoil much of the locality, since the eigenfunctions in (A.17) are defined in terms of $`H_{(2k+1)}`$ and thus they may reflect the locality properties of $`H_{(2k+1)}`$, which may not differ qualitatively from those of the overlap operator. In any case, a direct analysis of this locality issue is left as an important problem.
Our construction of the Dirac operator (5.1) on the basis of the defining algebraic relation (1.1) suggests that we obtain better chiral properties if one increases the parameter $`k`$. An intuitive argument for this expectation is that the right-hand side $`2a^{2k+1}(\gamma _5D)^{2k+2}`$ of the algebra (1.1), which breaks chiral symmetry, becomes more irrelevant for larger $`k`$ in the sense of renormalization group. At the same time, however, our construction requires a larger lattice for larger $`k`$ since the basic operator appearing in our construction spreads over far apart lattice points for large $`k`$. To maintain the locality of the Dirac operator, we need to take a smaller lattice spacing for larger $`k`$.
As for the chiral fermions, the present calculation of anomaly is readily extended to the evaluation of the fermion number anomaly of chiral theory and also to the so-called covariant form of non-Abelian anomalies in the continuum limit. But the construction of chiral fermion theory at finite lattice spacing is a challenging un-solved problem not only in our general Dirac operators but also in the original overlap operator.
T-W. Chiu has recently informed us that a numerical study of some of basic properties of the operator $`\gamma _5D`$ with $`k=1`$, such as index theorem, chiral anomaly and the propagator, is in progress.
One of us (KF) thanks Ting-Wai Chiu for numerous helpful discussions from the very beginning of the present investigation.
## Appendix A Basic construction of general Dirac operators
We start with (1.1) written in the form
$$H\gamma _5+\gamma _5H=2H^{2k+2}$$
(A.1)
or equivalently
$$\mathrm{\Gamma }_5H+H\mathrm{\Gamma }_5=0$$
(A.2)
where $`H=a\gamma _5D`$ and $`\mathrm{\Gamma }_5=\gamma _5H^{2k+1}`$. This algebraic relation implies that
$$\gamma _5H^2=[\gamma _5H+H\gamma _5]HH[\gamma _5H+H\gamma _5]+H^2\gamma _5=H^2\gamma _5.$$
(A.3)
Namely, the algebraic relation (A.1) is equivalent to the two relations
$`H^{2k+1}\gamma _5+\gamma _5H^{2k+1}=2H^{2(2k+1)},`$
$`\gamma _5H^2H^2\gamma _5=0.`$ (A.4)
If one defines $`H_{(2k+1)}H^{2k+1}`$, the first relation of (A.4) becomes
$$H_{(2k+1)}\gamma _5+\gamma _5H_{(2k+1)}=2H_{(2k+1)}^2$$
(A.5)
with $`\mathrm{\Gamma }_5=\gamma _5H_{(2k+1)}`$, which is identical to the conventional Ginsparg-Wilson relation with $`k=0`$ in (1.1). We utilize this property to construct a solution to (A.1).
The physical condition for the operator $`H`$ in (A.1) in the near continuum configuration is
$$H\gamma _5ai\overline{)}D+\gamma _5(\gamma _5ai\overline{)}D)^{2k+2}$$
(A.6)
where the first term stands for the leading term in chiral symmetric terms, and the second term stands for the leading term in chiral symmetry breaking terms. Thus $`H_{(2k+1)}`$ should satisfy
$`H_{(2k+1)}`$ $``$ $`[\gamma _5ai\overline{)}D+\gamma _5(\gamma _5ai\overline{)}D)^{2k+2}]^{2k+1}`$ (A.7)
$``$ $`(\gamma _5ai\overline{)}D)^{2k+1}+\gamma _5(\gamma _5ai\overline{)}D)^{2(2k+1)}`$
as can be confirmed by noting $`\gamma _5\overline{)}D+\overline{)}D\gamma _5=0`$. Here only the leading terms in chiral symmetric and chiral symmetry breaking terms, respectively, are written.
One can thus construct a solution for $`H_{(2k+1)}`$ by
$$H_{(2k+1)}=\frac{1}{2}\gamma _5[1+\gamma _5H_W^{(2k+1)}\frac{1}{\sqrt{H_W^{(2k+1)}H_W^{(2k+1)}}}]$$
(A.8)
in terms of the hermitian $`H_W^{(2k+1)}\gamma _5D_W^{(2k+1)}=(H_W^{(2k+1)})^{}`$. The operator $`D_W^{(2k+1)}`$ is defined in (2.3). The physical condition (A.7) is satisfied by (2.19), as was noted in the body of the text.
We now discuss how to reconstruct $`H`$, which satisfies (A.1), from $`H_{(2k+1)}`$ defined above. The basic idea is to define in the representation where $`H_{(2k+1)}`$ is diagonal
$$H=(H_{(2k+1)})^{1/(2k+1)}$$
(A.9)
in such a manner that $`H`$ thus obtained satisfies the second constraint in (A.4). For this purpose, we first recall the essence of the general representation of the algebra (A.1).
If one defines the eigenvalue problem
$$H_{(2k+1)}\varphi _n=(a\lambda _n)^{2k+1}\varphi _n,(\varphi _n,\varphi _n)=1$$
(A.10)
one can classify the eigenstates into the 3 classes:
(i) $`n_\pm `$ (“zero modes”),
$$H_{(2k+1)}\varphi _n=0,\gamma _5\varphi _n=\pm \varphi _n,$$
(A.11)
(ii) $`N_\pm `$ (“highest states”),
$$H_{(2k+1)}\varphi _n=\pm \varphi _n,\gamma _5\varphi _n=\pm \varphi _n,respectively,$$
(A.12)
(iii)“paired states” with $`0<|(a\lambda _n)^{2k+1}|<1`$,
$$H_{(2k+1)}\varphi _n=(a\lambda _n)^{2k+1}\varphi _n,H_{(2k+1)}(\mathrm{\Gamma }_5\varphi _n)=(a\lambda _n)^{2k+1}(\mathrm{\Gamma }_5\varphi _n).$$
(A.13)
where
$$\mathrm{\Gamma }_5=\gamma _5H_{(2k+1)}.$$
(A.14)
Note that $`\mathrm{\Gamma }_5(\mathrm{\Gamma }_5\varphi _n)\varphi _n`$ for $`0<|(a\lambda _n)^{2k+1}|<1`$.
We have a chirality sum rule
$$n_++N_+=n_{}+N_{}$$
(A.15)
where $`N_\pm `$ stand for the number of “highest states” in the classification (ii) above.
All the states $`\varphi _n`$ with $`0<|(a\lambda _n)^{2k+1}|<1`$, which appear pairwise with $`(a\lambda _n)^{2k+1}=\pm |(a\lambda _n)^{2k+1}|`$, can be normalized to satisfy the relations
$`\mathrm{\Gamma }_5\varphi _n`$ $`=`$ $`[1(a\lambda _n)^{2(2k+1)}]^{1/2}\varphi _n,`$
$`\gamma _5\varphi _n`$ $`=`$ $`(a\lambda _n)^{2k+1}\varphi _n+[1(a\lambda _n)^{2(2k+1)}]^{1/2}\varphi _n,`$ (A.16)
where $`\varphi _n`$ stands for the eigenstate with an eigenvalue opposite to that of $`\varphi _n`$.
We can define the solution $`H`$ of (A.1) operationally by
$$H\varphi _na\lambda _n\varphi _n$$
(A.17)
by using the same set of eigenfunctions and (the real $`2k+1`$th roots of) eigenvalues
$$\{\varphi _n\},\{a\lambda _n\}$$
(A.18)
as for $`H_{(2k+1)}`$ in (A.10). Note that the operator $`\mathrm{\Gamma }_5=\gamma _5H_{(2k+1)}=\gamma _5H^{2k+1}`$, which reverses the signature of eigenvalues of “paired states” and defines the index, is identical to (A.2) and (A.5).
We can confirm the second constraint in (A.4) and the defining algebraic relation (A.2) for any $`\varphi _n`$ in (A.17) by using (A.16),
$`[H^2\gamma _5\gamma _5H^2]\varphi _n=0`$
$`[\mathrm{\Gamma }_5H+H\mathrm{\Gamma }_5]\varphi _n=0.`$ (A.19)
The general representation of the algebra (A.1) is obtained from the standard representation, which is defined by $`H`$ and $`\{\varphi _n\}`$ in (A.17), and $`\gamma _5`$ in (A.16), by applying a suitable unitary transformation.
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# PION ELECTROPRODUCTION ON THE NUCLEON AND THE GENERALIZED GDH SUM RULE
## 1 Introduction
The spin structure of the nucleon in the resonance region is of particular interest to understand the rapid transition from resonance dominated coherent processes to incoherent processes of deep inelastic scattering (DIS) off the constituents. By scattering polarized lepton beams off polarized targets, it has become possible to determine the spin structure functions $`g_1`$ and $`g_2`$. The results of the first experiments at CERN $`^\mathrm{?}`$ and SLAC $`^\mathrm{?}`$ sparked considerable interest in the community, because the first moment of $`g_1`$, $`\mathrm{\Gamma }_1=_0^1g_1(x)𝑑x`$, was found to be substantially smaller than expected from the quark model, in particular from the Ellis-Jaffe sum rule $`^\mathrm{?}`$.
Here we present the results of the recently developed unitary isobar model (MAID) $`^\mathrm{?}`$ for the spin asymmetries, structure functions and relevant sum rules in the resonance region. This model describes the presently available data for single-pion photo- and electroproduction up to a total $`cm`$ energy $`W_{\mathrm{max}}=2`$ GeV and for $`Q^2`$ 4 (GeV/c)<sup>2</sup>. It is based on effective Lagrangians for Born terms and vector meson exchange (background) and resonance contributions modeled by Breit-Wigner functions. All major resonances below $`W=1700`$ MeV are included. The respective multipoles are constructed in a gauge-invariant and unitary way for each partial wave. The eta production is included in a similar way $`^\mathrm{?}`$, while the contribution of more-pion and higher channels is modeled by comparison with the total cross sections and simple phenomenological assumptions.
## 2 Formalism
The differential cross section for exclusive electroproduction of mesons from polarized targets using polarized electrons, e.g. $`\stackrel{}{p}(\stackrel{}{e},e^{}\pi ^0)p`$ can be parametrized in terms of 18 response functions $`^\mathrm{?}`$, a total of 36 is possible if in addition also the recoil polarization is observed. Due to the azimuthal symmetry most of them vanish by integration over the angle $`\varphi `$ and only 5 total cross sections remain. The differential cross section for the electron is then given by
$$\frac{d\sigma }{d\mathrm{\Omega }dE^{}}=\mathrm{\Gamma }\sigma (\nu ,Q^2),$$
(1)
$$\sigma =\sigma _T+ϵ\sigma _L+P_y\sqrt{2ϵ(1+ϵ)}\sigma _{LT}+hP_x\sqrt{2ϵ(1ϵ)}\sigma _{LT^{}}+hP_z\sqrt{1ϵ^2}\sigma _{TT^{}},$$
(2)
where $`\mathrm{\Gamma }`$ is the flux of the virtual photon field and the $`\sigma _i,`$ $`i=L`$, $`T`$, $`LT`$, $`LT^{}`$, $`TT^{}`$, are functions of the $`lab`$ energy of the virtual photon $`\nu `$ and the squared four-momentum transferred $`Q^2`$. These response functions can be separated by varying the transverse polarization $`ϵ`$ of the virtual photon as well as the polarizations of the electron ($`h`$) and proton ($`P_z`$ parallel, $`P_x`$ perpendicular to the virtual photon, in the scattering plane and $`P_y`$ perpendicular to the scattering plane). In particular, $`\sigma _T`$ and $`\sigma _{TT^{}}`$ can be expressed in terms of the total cross sections for excitation of hadronic states with spin projections $`3/2`$ and $`1/2`$: $`\sigma _T=(\sigma _{3/2}+\sigma _{1/2})/2`$ and $`\sigma _{TT^{}}=(\sigma _{3/2}\sigma _{1/2})/2.`$
Here we use Hand’s notation with the equivalent photon $`cm`$ energy $`K=(W^2m^2)/(2W)`$ for the virtual photon flux. Correspondingly, the phase space factors of the cross sections are given by $`q/K`$, where $`q`$ is the pion momentum in the $`cm`$.
In inclusive electron scattering $`\stackrel{}{e}+\stackrel{}{N}X`$, only 4 cross sections $`\sigma _T`$, $`\sigma _L`$, $`\sigma _{LT^{}}`$ and $`\sigma _{TT^{}}`$ appear, the fifth cross section, $`\sigma _{LT}`$, vanishes due to unitarity when all open channels are summed up. The individual channels, however, give finite contributions.
The Gerasimov-Drell-Hearn (GDH) sum rule is only derived for real photons. It is based on unitarity and low-energy theorems and the assumption of the convergence of an unsubtracted dispersion relation,
$`I_{GDH}`$ $`=`$ $`{\displaystyle \frac{m^2}{8\pi ^2\alpha }}{\displaystyle _{\nu _0}^{\mathrm{}}}\left(\sigma _{1/2}\sigma _{3/2}\right){\displaystyle \frac{d\nu }{\nu }}={\displaystyle \frac{\kappa ^2}{4}}.`$ (3)
This sum rule is often presented without the leading factor in front of the integral, the numerical conversion is $`8\pi ^2\alpha /m^2=254.8\mu b`$.
It can be generalized in various ways. Three forms often used in the literature are summing up only contributions from $`\sigma _{TT^{}}`$ with no longitudinal terms,
$`I_{GDH}^{(a)}(Q^2)`$ $`=`$ $`{\displaystyle \frac{m^2}{8\pi ^2\alpha }}{\displaystyle _{\nu _0}^{\mathrm{}}}(1x)\left(\sigma _{1/2}\sigma _{3/2}\right){\displaystyle \frac{d\nu }{\nu }},`$ (4)
$`I_{GDH}^{(b)}(Q^2)`$ $`=`$ $`{\displaystyle \frac{m^2}{8\pi ^2\alpha }}{\displaystyle _{\nu _0}^{\mathrm{}}}{\displaystyle \frac{K}{\stackrel{}{k}_\gamma }}\left(\sigma _{1/2}\sigma _{3/2}\right){\displaystyle \frac{d\nu }{\nu }},`$ (5)
$`I_{GDH}^{(c)}(Q^2)`$ $`=`$ $`{\displaystyle \frac{m^2}{8\pi ^2\alpha }}{\displaystyle _{\nu _0}^{\mathrm{}}}\left(\sigma _{1/2}\sigma _{3/2}\right){\displaystyle \frac{d\nu }{\nu }}.`$ (6)
The factor $`K/\stackrel{}{k}_\gamma `$ can also be expressed as $`(1x)/\sqrt{1+\gamma ^2}`$. The relations between the $`\sigma _i`$ and the quark structure functions $`g_1`$ and $`g_2`$ can be read off the following equations, which define further possible generalizations of the GDH integral $`^\mathrm{?}`$ and the Burkhardt-Cottingham (BC) sum rule $`^\mathrm{?}`$, which in addition also include longitudinal-transverse interference terms,
$`I_1(Q^2)`$ $`=`$ $`{\displaystyle \frac{2m^2}{Q^2}}{\displaystyle _0^{x_0}}g_1(x,Q^2)𝑑x`$ (7)
$`=`$ $`{\displaystyle \frac{m^2}{8\pi ^2\alpha }}{\displaystyle _{\nu _0}^{\mathrm{}}}{\displaystyle \frac{1x}{1+\gamma ^2}}\left(\sigma _{1/2}\sigma _{3/2}2\gamma \sigma _{LT^{}}\right){\displaystyle \frac{d\nu }{\nu }},`$
$`I_2(Q^2)`$ $`=`$ $`{\displaystyle \frac{2m^2}{Q^2}}{\displaystyle _0^{x_0}}g_2(x,Q^2)𝑑x`$ (8)
$`=`$ $`{\displaystyle \frac{m^2}{8\pi ^2\alpha }}{\displaystyle _{\nu _0}^{\mathrm{}}}{\displaystyle \frac{1x}{1+\gamma ^2}}\left(\sigma _{3/2}\sigma _{1/2}{\displaystyle \frac{2}{\gamma }}\sigma _{LT^{}}\right){\displaystyle \frac{d\nu }{\nu }},`$
$`I_3(Q^2)`$ $`=`$ $`{\displaystyle \frac{2m^2}{Q^2}}{\displaystyle _0^{x_0}}(g_1(x,Q^2)+g_2(x,Q^2))𝑑x`$ (9)
$`=`$ $`{\displaystyle \frac{m^2}{4\pi ^2\alpha }}{\displaystyle _{\nu _0}^{\mathrm{}}}{\displaystyle \frac{1x}{Q}}\sigma _{LT^{}}𝑑\nu =I_1+I_2,`$
where $`\gamma =Q/\nu `$ and $`x=Q^2/2m\nu `$ the Bjorken scaling variable, with $`x_0`$ ($`\nu _0`$) referring to the inelastic threshold of one-pion production. Since $`\sigma _{LT^{}}=𝒪(Q)`$, the real photon limit of the integral $`I_1`$ is given by the GDH sum rule $`I_1(0)=I_{GDH}(0)=\kappa _N^2/4,`$ with $`\kappa _N`$ the anomalous magnetic moment of the nucleon. At large $`Q^2`$ the structure functions should depend only on $`x,`$ i.e. $`I_12m\mathrm{\Gamma }_1/Q^2`$ with $`\mathrm{\Gamma }_1=g_1(x)dx=`$ const. In the case of the proton, all experiments for $`Q^2>1`$GeV<sup>2</sup> yield $`\mathrm{\Gamma }_1>0.`$ Therefore, a strong variation of $`I_1(Q^2)`$ with a zero-crossing at $`Q^2<1`$ GeV<sup>2</sup> is required in order to reconcile the GDH sum rule with the measurements in the DIS region. The $`I_2`$ integral of Eq. (8) is constrained by the BC sum rule, which requires that the inelastic contribution for $`0<x<x_0`$ equals the negative value of the elastic contribution, i.e.
$$I_2(Q^2)=\frac{2m^2}{Q^2}_0^{x_0}g_2(x,Q^2)𝑑x=\frac{1}{4}\frac{G_M(Q^2)G_E(Q^2)}{1+Q^2/4m^2}G_M(Q^2),$$
(10)
where $`G_M`$ and $`G_E`$ are the magnetic and electric Sachs form factors respectively. At large $`Q^2`$ the integral vanishes as $`Q^{10}`$, while at the real photon limit $`I_2(0)=\kappa _N^2/4+e_N\kappa _N/4`$, the two terms on the right hand side corresponding to the contributions of $`\sigma _{TT^{}}`$ and $`\sigma _{LT^{}}`$ respectively. Finally, Eq. (9) defines an integral $`I_3(Q^2)`$ as the sum of $`I_1(Q^2)`$ and $`I_2(Q^2)`$ and is given by the unweighted integral over the longitudinal-transverse interference cross section $`\sigma _{LT^{}}`$. At the real photon point this integral is given by the GDH and BC sum rules, $`I_3(0)=e_N\kappa _N/4`$. In particular this vanishes for the neutron target.
## 3 Unitary Isobar Model
Our calculation for the response functions $`\sigma _i`$ is based on the Unitary Isobar Model (UIM) for one-pion photo- and electroproduction of Ref. $`^\mathrm{?}`$, accessible in the internet as the MAID program. The model is constructed with effective Lagrangians for Born terms, vector meson exchange in the $`t`$ channel (background), and the dominant resonances up to the third resonance region are modeled using Breit-Wigner functions with energy-dependent widths. For each partial wave the multipoles satisfy gauge invariance and unitarity. As in any realistic model a special effort is needed to describe the $`s`$-channel multipoles $`S_{11}`$ and $`S_{31}`$. Even close at threshold these multipoles pick up sizeable imaginary parts that cannot be explained by nucleon resonances. In fact the $`S_{11}(1535)`$, $`S_{11}(1650)`$ and the $`S_{31}(1620)`$ play only a minor role for the complex phase of the $`E_{0+}`$ multipoles even at higher energies. The main effect arises from pion rescattering. This we can take into account by $`K`$-matrix unitarization. Furthermore we introduce a gauge invariant contact term proportional to the anomalous magnetic moment of the nucleon $`\kappa _N`$,
$$j_\kappa ^\mu =\frac{ieg}{2m}\kappa _NF(q_0^2)\frac{\sigma ^{\mu \nu }k_\nu }{2m}\gamma _5.$$
(11)
The form factor $`F(q_0^2)`$, $`q_0`$ being the asymptotic pion momentum, vanishes at threshold, consistent with chiral symmetry, but gives rise to a cancellation of unphysically high momentum components in the Born terms at high energies.
Due to unitarity each partial wave has to fulfill Watson’s theorem,
$`t_{\gamma ,\pi }^\alpha `$ $`=`$ $`t_{\gamma ,\pi }^\alpha (background)+t_{\gamma ,\pi }^\alpha (resonances)`$
$`=`$ $`\pm t_{\gamma ,\pi }^\alpha e^{i\delta _{\pi N}^\alpha }.`$
In an isobar model this condition has to be constructed explicitly. In Maid98 the background is real (except for the S-waves) and a phase is added to the resonance. In Maid2000 both background and resonance contributions are unitarized separately for all partial waves up to $`l=3`$ in the following way
$$t_{\gamma ,\pi }^\alpha =t_{\gamma ,\pi }^\alpha (Born+\omega ,\rho )(1+it_{\pi N}^I)+t_{\gamma ,\pi }^\alpha (resonances)e^{i\mathrm{\Psi }^\alpha }.$$
(13)
In Fig. 1 we show the $`M_{1+}^{(3/2)}`$ multipole and the $`M1`$ transition form factor of the $`\mathrm{\Delta }`$ resonance on the left and the $`R_{EM}=E_{1+}^{3/2}/M_{1+}^{3/2}`$ and $`R_{SM}=S_{1+}^{3/2}/M_{1+}^{3/2}`$ ratios on the right side. The experimental data agree very well with our empirical fit for $`G_M^{}`$. For the small ratios the experimental information is not yet reliable enough to justify anything else than a constant value. In our model we use $`R_{EM}=2.2\%`$ and $`R_{SM}=6.5\%`$. Recent analyses of experimental data on $`p(e,e^{}\pi ^0)p`$ from JLab $`^\mathrm{?}`$ show a trend to $`E/M`$ ratios very close to zero and increasing negative values for $`S/M`$ at large $`Q^2`$.
The UIM is able to describe the single-pion electroproduction channel quite well. However, at higher energies the contributions from other channels become increasingly important. In the structure functions $`\sigma _T`$ and $`\sigma _{TT^{}}`$ we account for the $`\eta `$ and the multi-pion production contributions extracting the necessary information from the existing data for the total cross section $`^\mathrm{?}`$. In Fig. 2 we show the individual channels for the total helicity dependent cross sections $`\sigma _{1/2}`$, $`\sigma _{3/2}`$, $`\sigma _T`$ and $`\mathrm{\Delta }\sigma =2\sigma _{TT^{}}`$ at $`Q^2=0`$. Due to the non-regularized Born terms in the $`1\pi `$ channels the cross sections start to rise again at energies $`W>1.8`$ GeV. However, because of the energy weighting, the effect is negligible for the integrals.
## 4 Integrals
### 4.1 Results for Real Photons
In Tab. 1 we show our results for the GDH integral and the forward spin polarizability over the lab energy range of 200-450 MeV together with the latest Mainz data $`^\mathrm{?}`$. In comparison with the results of the dispersion theoretical partial wave analysis HDT $`^\mathrm{?}`$ and the SAID solution SM99K $`^\mathrm{?}`$ our MAID results agree very well with the experiment. The additional information on the individual channels, however, offers interesting insights in the different calculations, especially the $`\pi ^+`$ result of SAID for $`\gamma _0`$ is a factor 3 larger than MAID and about a factor 2 above the data. The reason is the enhanced sensitivity of background contributions in the $`\pi ^+`$ channel, especially the S-wave near threshold.
Tab. 2 shows the GDH integral over the full energy range up to $`W_{max}=2`$ GeV. The preliminary experimental result is obtained only from measurements at Mainz and covers the energy range from 200 to 800 MeV photon $`lab`$ energy. However, our theoretical calculations indicate a very close cancellation between the low energy contribution from threshold up to 200 MeV (30 $`\mu b`$) and of energies above 800 MeV (-34$`\mu b`$). For our detailed comparison we included recent calculations of reggeized $`\pi \pi `$ photoproduction by Holvoet and Vanderhaeghen $`^\mathrm{?}`$ that include $`\gamma ,\pi \mathrm{\Delta }`$ Born terms and additional $`D_{13}(1520)`$ excitations. This $`2\pi `$ contribution to the GDH integral is about twice as large as compared to our simple phenomenological multi-pion parametrization used for finite $`Q^2`$.
Our calculation that also include very recent Regge-type calculations for two-pion photoproduction $`^\mathrm{?}`$ shows a very good agreement with the sum rule for the proton target but also exhibits a big deviation for the neutron target. However, on the neutron target the photoproduction information is rather limited above the $`\mathrm{\Delta }`$ region and it has to be further investigated if the high-energy region is perhaps more important than for the proton target. Furthermore, for both nucleon targets high energy contributions beyond the two-pion production have to be studied that make a very big contribution in deep inelastic scattering at finite $`Q^2>1GeV^2`$.
### 4.2 Results for Virtual Photons
In Fig. 3 and Fig. 4 we give our predictions for the integrals $`I_{GDH}(Q^2)`$, $`I_1(Q^2)`$ and $`I_2(Q^2)`$ in the resonance region, i.e. integrated up to $`W_{\mathrm{max}}=2`$ GeV for the proton and neutron targets. A comparison of the 3 different forms, defined in Eqs. 4, 5, 6 show significantly different slopes at $`Q^2=0`$ and quite different zero positions, where the GDH integral crosses from negative values observed for real photons to positive values known from deep inelastic scattering. In the case of the integral $`I_1,`$ our model is able to generate the expected drastic change in the helicity structure at low $`Q^2`$.
We find a zero-crossing at $`Q^2=0.75`$ (GeV/c)<sup>2</sup> if we include only the one-pion contribution. This value is lowered to 0.52 (GeV/c)<sup>2</sup> and 0.45 (GeV/c)<sup>2</sup> when we include the $`\eta `$ and the multi-pion contributions respectively. The SLAC analysis of the proton yields $`I_1=0.1\pm 0.06`$ at $`Q^2=0.5`$ (GeV/c)<sup>2</sup>, while our result at this point is only slightly positive. For the neutron our calculation is fully consistent within the SLAC analysis at $`Q^2=0.5`$ (GeV/c)<sup>2</sup> in contrast to the large discrepancy observed at $`Q^2=0`$, see Tab. 2.
Comparing with the generalizations of the GDH sum rule in Fig. 3, it can be seen that the slope at $`Q^2=0`$ and the existence of a minimum for small $`Q^2`$ depends on the inclusion of the longitudinal contributions, i.e. the minimum disappears when $`\sigma _{LT^{}}`$ is added. Concerning the integral $`I_2`$, our full result is in good agreement with the prediction of the BC sum rule. The deviation is within $`10\%`$ and should be attributed to contributions beyond $`W_{\mathrm{max}}=2`$ GeV and the uncertainties in our calculation for $`\sigma _{LT^{}}`$. As seen in Eq. (9) the integral $`I_3`$ depends only on this $`\sigma _{LT^{}}`$ contribution. From the sum rule result a value of $`e_N\kappa _N/4`$ is expected at $`Q^2=0`$, i.e. 0.45 for the proton and zero for the neutron target. While our value arising entirely from the $`1\pi `$ channel (0.59) gets relatively close to the sum rule result for the proton, in the neutron case this sum rule is heavily violated (0.78). So far it is not clear where such a large negative contribution should arise for the neutron target in order to cancel the $`1\pi `$ contribution. Either it is due to the high-energy tail that may converge rather slowly for the unweighted integral $`I_3`$, or the multi-pion channels could contribute in such a way, while the eta channel is very unlikely. On the other hand the convergence of the BC sum rule cannot be given for granted. In fact Ioffe et al. $`^\mathrm{?}`$ have argued that the BC sum rule is valid only in the scaling region, while it is violated by higher twist terms at low $`Q^2`$. In any case a careful study of the multi-pion contribution for both proton and neutron targets will be very helpful, in particular one can expect longitudinal contributions from the non-resonant background.
## 5 Summary
In summary, we have applied our recently developed unitary isobar model for pion electroproduction to calculate generalized GDH integrals and the BC sum rule for both proton and neutron targets. Our results indicate that both the experimental analysis and the theoretical models have to be quite accurate in order to fully describe the helicity structure of the cross section in the resonance region.
While our results agree quite well for the GDH and BC sum rules for the proton, we find substantial deviations for the neutron target, in particular the sum rule $`I_3(0)I_1(0)+I_2(0)=0`$ is heavily violated by the contribution from the single-pion channel which is even larger than in the case of the proton. Concerning the theoretical description, the treatment of the multi-pion channels has to be improved with more refined models. On the experimental side, the upcoming results from measurements with real $`^\mathrm{?}`$ and virtual photons $`^\mathrm{?}`$ hold the promise to provide new precision data in the resonance region.
## Acknowledgments
We would like to thank M. Vanderhaeghen for the contributions to the $`\pi \pi `$ channels and J. Arends for the information on the experimental data analysis. This work was supported by the Deutsche Forschungsgemeinschaft (SFB 443).
## References
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# Quasi exactly solvable matrix Schrödinger operators.
## 1 Introduction
In a recent paper a classification of $`2\times 2`$ matrix quasi exactly solvable (QES) Schrödinger operators in one spatial dimension is attempted. This problem was first addressed in and further developped in and . Here we consider a suitable class of finite dimensional vector spaces of polynomials in a real variable and we construct two families of operators preserving sub-classes of these vector spaces. The first family is related to one of the cases treated in ; the second, which is not considered in , generalizes an equation considered in ,. The two corresponding QES equations respectively constitute “coupled channel” generalisations of the anharmonic and Lamé QES scalar equations.
Following the basic idea of QES operators , we consider the finite dimensional vector space of couples of polynomials of given degree $`n`$ and $`m`$ in a real variable $`x`$. We slightly generalize this vector space by setting
$$𝒱=P\left(\begin{array}{c}𝒫(n)\hfill \\ 𝒫(m)\hfill \end{array}\right)$$
(1)
where $`𝒫(n)`$ denotes the set of real polynomials of degree at most $`n`$ in $`x`$ while $`P`$ is a fixed invertible $`2\times 2`$ matrix operator; $`P`$ can be interpreted as a change of basis in the vector space $`𝒫(n)𝒫(m)`$. With such an interpretation it is reasonnable to choose $`P`$ of the form
$$P=\left(\begin{array}{cc}1& P_{12}\\ 0& 1\end{array}\right)(\mathrm{resp}.P=\left(\begin{array}{cc}1& 0\\ P_{21}& 1\end{array}\right))$$
(2)
for $`nm`$ (resp. $`mn`$). In this paper, we limit ourself to scalar operators $`P_{12}`$ (or $`P_{21})`$ of the form
$$P_{12}=\kappa _0\frac{}{x}+\kappa _1+\kappa _2x\frac{}{x}+\kappa _3x$$
(3)
where $`\kappa _j`$ are constants. The vector space defined in Eq. (6) of can be set in the form (1) with $`\kappa _0=1`$, $`\kappa _{1,2,3}=0`$.
## 2 Polynomial potential
We consider an operator of the form
$$H(y)=\frac{d^2}{dy^2}1\mathrm{I}_2+M_6(y)$$
(4)
where $`M_6(y)`$ is a $`2\times 2`$ hermitian matrix whose entries are even polynomials of degree at most six in $`y`$. This operator is a natural generalisation of the famous QES anharmonic oscillator to $`2\times 2`$ matrix operator. After the standard “gauge transformation” of $`H(y)`$ with a factor
$$\varphi (y)=y^ϵ\mathrm{exp}\{\frac{p_2}{2}y^4+p_1y^2\}$$
(5)
and the change of variable $`x=y^2`$, the equivalent operator $`\widehat{H}(x)`$ can be computed :
$$\widehat{H}(x)=\varphi ^1(x)H(y)\varphi (x)_{y=\sqrt{x}}$$
(6)
Then we pose the problem : what is the most general choice of $`M_6`$ such that $`\widehat{H}(x)`$ preserves a vector space of the form (1),(2),(3)?.
The following solutions are obtained after straighforward calculations (we exclude the case where $`M_6(y)`$ is diagonal since it corresponds to a direct sum of two scalar QES operators) :
$$ϵ=0\mathrm{or}1,n=m2,\kappa _1=\kappa _2=\kappa _3=0.$$
(7)
The corresponding potentials $`M_6`$ has the form
$`M_6(y)`$ $`=`$ $`\{4p_2^2y^6+8p_1p_2y^4+(4p_1^28mp_2+2(12ϵ)p_2)y^2\}1\mathrm{I}_2`$ (8)
$`+`$ $`(8p_2y^2+4p_1)\sigma _38mp_2\kappa _0\sigma _1`$
where $`\sigma _1,\sigma _3`$ are the Pauli matrices, $`p_2,p_1,\kappa _0`$ are free real parameters and $`m`$ is an integer. In particular, the non diagonal term is parametrized by an arbitrary constant which cannot be suppressed because of the $`y`$-dependent term proportional to $`\sigma _3`$. If the parameter $`ϵ`$ is choosen as an arbitrary real number, then the potential $`M_6`$ has a supplementary term of the form $`ϵ(ϵ1)/y^2`$.
When the parameters of the case 1 of Ref. are choosen so that the potential is a polynomial matrix (i.e. $`\alpha _2=\alpha _0=0`$, $`\alpha _1=1`$, $`\beta _0=1/2\mathrm{or}3/2`$ in Eq. (34) of ) the potential reduces to the matrix $`M_6(y)`$ above. The way of obtaining this result here is slightly different because the method starts from the natural vector space $`𝒫(m)𝒫(n)`$. The more elaborated QES operators obtained in can also be produced by our technique but this is not aimed in this note.
Let us also point out that the “gauge factor $`U`$” considered in is limited to be a function of the variable $`x`$ and, therefore, is not supposed to contain any derivative operator like our operator $`P`$ (see (3)). This explains that the QES polynomial potential (8) was not found in . With the restriction that $`U`$ is a function of $`x`$ only, these authors correctly reach the conclusion that Hamiltonian preserving a space like $`𝒫(n)𝒫(m)`$ with $`|nm|>1`$ are essentially diagonal, in contrast with the present operator related to the case $`nm=2`$.
## 3 Application to N-body hamiltonians
By using the idea of the previous section, a matrix version of the QES many-body problem of Ref. can be constructed. Let us consider the Calogero Hamiltonian (we note it $`H_{cal}`$) supplemented by a matrix-valued potential $`V^{}`$ :
$$H=H_{cal}+V^{}=\frac{1}{2}\underset{j=1}{\overset{N}{}}[\frac{^2}{x_j^2}+x_j^2]+\underset{j<i}{}\frac{\nu (\nu 1)}{(x_jx_i)^2}+V^{}$$
(9)
Along with we assume $`V^{}`$ to depend only on the variable $`\tau `$
$$\tau \underset{j<i}{\overset{N}{}}(x_jY)(x_iY),Y\underset{j=1}{\overset{N}{}}x_j,$$
(10)
and we look for eigenfunctions of the hamiltonian (9) of the form
$$\mathrm{\Psi }(x)=\psi _0(x)\tau ^ϵ\mathrm{exp}\{\frac{p_2}{2}\tau ^4+p_1\tau ^2\}\varphi (\tau )$$
(11)
where $`\psi _0`$ denotes the ground state of the standard Calogero system :
$$\psi _0(x)=[\underset{i<j}{}|x_ix_j|)]^\nu \mathrm{exp}(X^2/2),X^2_{j=1}^Nx_j^2$$
(12)
while $`\varphi (\tau )`$ represents a couple of polynomials in $`\tau `$.
After a standard algebra, the operator acting on $`\varphi (\tau )`$ can be isolated :
$$h\tau \frac{^2}{\tau ^2}+(4\tau +2b)\frac{}{\tau }+V^{}$$
(13)
and it can be shown that this operator preserves the space (1) (again with $`n=m2`$, $`\kappa _1=\kappa _2=\kappa _3=0`$) provided $`V^{}`$ is of the form
$$V^{}(\tau )=p_2^2\tau ^3+2p_2(1p_1)\tau ^2+(a2p_2\sigma _3)\tau +(1p_1)\sigma _3+\frac{\gamma }{\tau }+2m\kappa _0\sigma _1$$
(14)
with the definitions
$$ap_1(2p_1)+p_2(2m+3ϵ1+b),b\frac{1}{2}(1+\nu N)(N1),\gamma 2ϵ(ϵ1+b).$$
(15)
As a consequence, (9),(14) constitutes a QES matrix extension (labelled by the parameters $`p_1,p_2,ϵ`$) of the exactly solvable Calogero hamiltonian.
## 4 Lamé type potential.
As a second example, we consider the family of operators
$$H(z)=\frac{d^2}{dz^2}+\left[\begin{array}{cc}Ak^2\mathrm{sn}^2+\delta (1+k^2)/2& 2\theta k\mathrm{cn}\mathrm{dn}\\ 2\theta k\mathrm{cn}\mathrm{dn}& Ck^2\mathrm{sn}^2\delta (1+k^2)/2\end{array}\right]$$
(16)
where $`A,C,\delta ,\theta `$ are constants while $`\mathrm{sn},\mathrm{cn},\mathrm{dn}`$ respectively abbreviate the Jacobi elliptic functions of argument $`z`$ and modulus $`k`$
$$\mathrm{sn}(z,k),\mathrm{cn}(z,k),\mathrm{dn}(z,k).$$
(17)
These functions are periodic with period $`4K(k),4K(k),2K(k)`$ respectively ($`K(k)`$ is the complete elliptic integral of the first type). The above hamiltonian is therefore to be considered on the Hilbert space of periodic functions on $`[0,4K(k)]`$. For completeness, we mention the properties of the Jacobi functions which are needed in the calculations
$$\mathrm{cn}^2+\mathrm{sn}^2=1,\mathrm{dn}^2+k^2\mathrm{sn}^2=1$$
(18)
$$\frac{d}{dz}\mathrm{sn}=\mathrm{cn}\mathrm{dn},\frac{d}{dz}\mathrm{cn}=\mathrm{sn}\mathrm{dn},\frac{d}{dz}\mathrm{dn}=k^2\mathrm{sn}\mathrm{cn}$$
(19)
The relevant change of variable which eliminates the transcendental functions sn, cn, dn from (16) in favor of algebraic expressions is (for $`k`$ fixed)
$$x=\mathrm{sn}^2(z,k)$$
(20)
In particular the second derivative term in (16) becomes
$$\frac{d^2}{dz^2}=4x(1x)(1k^2x)\frac{d^2}{d^2x}+2(3k^2x^22(1+k^2)x+1)\frac{d}{dx}$$
(21)
Several possibilities of extracting prefactors then lead to equivalent forms of (16), say $`\widehat{H}(x)`$, which are matrix operators build with the derivative $`d/dx`$ and polynomial coefficients in $`x`$. The requirement that $`\widehat{H}(x)`$ preserves a space of the form (1) leads to two possible sets of values for $`A,C,\theta `$ (we do not consider the case $`\theta =0`$ since it corresponds to two decoupled scalar Lamé equations).
Case 1
$`A=4m^2+6m+3\delta `$
$`C=4m^2+6m+3+\delta `$
$`\theta =\frac{1}{2}[(4m+3)^2\delta ^2]^{\frac{1}{2}}`$
The parameter $`\delta `$ remains free, and also $`k`$ which fixes the period of the potential. Four invariant spaces are available. In order to present them we conveniently define
$$R_1=\frac{4m+3\delta }{4m+3+\delta },$$
(22)
We have then
$`𝒱_1`$ $`=`$ $`\left(\begin{array}{cc}1& 0\\ 0& \mathrm{cn}\mathrm{dn}\end{array}\right)\left(\begin{array}{cc}1& \kappa x\\ 0& 1\end{array}\right)\left(\begin{array}{c}𝒫(m)\\ 𝒫(m)\end{array}\right),\kappa ^2=k^2R_1`$ (29)
$`𝒱_2`$ $`=`$ $`\left(\begin{array}{cc}\mathrm{cn}\mathrm{dn}& 0\\ 0& 1\end{array}\right)\left(\begin{array}{cc}1& 0\\ \kappa x& 1\end{array}\right)\left(\begin{array}{c}𝒫(m)\\ 𝒫(m)\end{array}\right),\kappa ^2=k^2/R_1`$ (36)
$`𝒱_3`$ $`=`$ $`\left(\begin{array}{cc}\mathrm{sn}\mathrm{cn}& 0\\ 0& \mathrm{sn}\mathrm{dn}\end{array}\right)\left(\begin{array}{cc}1& \kappa \\ 0& 1\end{array}\right)\left(\begin{array}{c}𝒫(m1)\\ 𝒫(m)\end{array}\right),\kappa ^2=k^2R_1`$ (43)
$`𝒱_4`$ $`=`$ $`\left(\begin{array}{cc}\mathrm{sn}\mathrm{dn}& 0\\ 0& \mathrm{sn}\mathrm{cn}\end{array}\right)\left(\begin{array}{cc}1& \kappa \\ 0& 1\end{array}\right)\left(\begin{array}{c}𝒫(m1)\\ 𝒫(m)\end{array}\right),\kappa ^2=R_1/k^2`$ (50)
Case 2
$`A=4m^2+2m+1\delta `$
$`C=4m^2+2m+1+\delta `$
$`\theta =\frac{1}{2}[(4m+1)^2\delta ^2]^{\frac{1}{2}}`$
The associated invariant vector spaces read, defining $`R_2=(4m+1\delta )/(4m+1+\delta )`$,
$`𝒱_5`$ $`=`$ $`\left(\begin{array}{cc}\mathrm{cn}& 0\\ 0& \mathrm{dn}\end{array}\right)\left(\begin{array}{cc}1& \kappa \\ 0& 1\end{array}\right)\left(\begin{array}{c}𝒫(m1)\\ 𝒫(m)\end{array}\right),\kappa ^2=k^2R_2`$ (57)
$`𝒱_6`$ $`=`$ $`\left(\begin{array}{cc}\mathrm{dn}& 0\\ 0& \mathrm{cn}\end{array}\right)\left(\begin{array}{cc}1& \kappa \\ 0& 1\end{array}\right)\left(\begin{array}{c}𝒫(m1)\\ 𝒫(m)\end{array}\right),\kappa ^2=R_2/k^2`$ (64)
$`𝒱_7`$ $`=`$ $`\left(\begin{array}{cc}\mathrm{sn}& 0\\ 0& \mathrm{sn}\mathrm{cn}\mathrm{dn}\end{array}\right)\left(\begin{array}{cc}1& \kappa x\\ 0& 1\end{array}\right)\left(\begin{array}{c}𝒫(m1)\\ 𝒫(m1)\end{array}\right),\kappa ^2=k^2R_2`$ (71)
$`𝒱_8`$ $`=`$ $`\left(\begin{array}{cc}\mathrm{sn}\mathrm{cn}\mathrm{dn}& 0\\ 0& \mathrm{sn}\end{array}\right)\left(\begin{array}{cc}1& 0\\ \kappa x& 1\end{array}\right)\left(\begin{array}{c}𝒫(m1)\\ 𝒫(m1)\end{array}\right),\kappa ^2=k^2/R_2`$ (78)
The operator (16) was studied in , for $`\delta =1`$. For this particular value of $`\delta `$, the corresponding eigenvalue equation $`H\psi =\omega ^2\psi `$ determines the normal modes of the sphaleron classical solution in the Abelian Higgs model in 1+1 dimension. It therefore plays a crucial role in the understanding of the instabilities of the sphaleron in this model.
The above results demonstrate that the remarkable algebraic properties of the Lamé equation also hold for the operator (16), irrespectively of the value of $`\delta `$. The associated eigenvalue equation therefore constitutes a (one parameter) $`2\times 2`$ matrix equation analog of the scalar Lamé equation.
## 5 Generalization
The kind of operators presented in Sect. 3 can be generalized to matrix potentials of the form
$$H(z)=\frac{d^2}{dz^2}+\left[\begin{array}{cc}V_1(\mathrm{sn}^2)& \theta \mathrm{sn}^{\alpha _1}\mathrm{cn}^{\alpha _2}\mathrm{dn}^{\alpha _3}\\ \theta \mathrm{sn}^{\alpha _1}\mathrm{cn}^{\alpha _2}\mathrm{dn}^{\alpha _3}& V_2(\mathrm{sn}^2)\end{array}\right]$$
(79)
where $`V_1,V_2`$ are polynomials, $`\theta `$ is a constant and $`\alpha _j`$ are non-negative integers.
The similarity transformation
$$\widehat{H}(x)=U^1(z)H(z)U(z),U(z)=\mathrm{diag}(\mathrm{sn}^{\beta _1}\mathrm{cn}^{\beta _2}\mathrm{dn}^{\beta _3},\mathrm{sn}^{\gamma _1}\mathrm{cn}^{\gamma _2}\mathrm{dn}^{\gamma _3})$$
(80)
sets the operator (79) into a form with polynomial coefficients in the variable $`x=\mathrm{sn}^2`$ provided
* $`\beta _j,\gamma _j=0\mathrm{or}1`$ , $`j=1,2,3`$
* $`\alpha _j\pm (\beta _j\gamma _j)=`$ non-negative even integer , $`j=1,2,3`$.
After making a choice of $`\alpha _j,\beta _j,\gamma _j`$ satisfying the above conditions, the possible forms of $`V_1,V_2`$ and of $`P,m,n`$ in Eq. (1) have to be determined in order for $`H(z)`$ to be QES.
Taking $`k=0`$, the standard trigonometric functions are recovered :
$$\mathrm{sn}(\mathrm{z},0)=\mathrm{sin}\mathrm{z},\mathrm{cn}(\mathrm{z},0)=\mathrm{cos}\mathrm{z},\mathrm{dn}(\mathrm{z},0)=1.$$
(81)
The periodic potential below, which is exactly solvable , furnishes a particular example of this type
$$V(z)÷\left[\begin{array}{cc}\mathrm{cos}^2(z)& \mathrm{cos}(z)\mathrm{sin}(z)\\ \mathrm{cos}(z)\mathrm{sin}(z)& \mathrm{sin}^2(z)\end{array}\right].$$
(82)
It determines the normal modes about some static solutions of the Goldstone model in 1+1 dimensions .
## 6 Concluding remarks
The examples of operators presented above give evidences of the difficulty to classify the coupled-channel (or matrix) QES Schrodinger equations. The way of constructing the QES potential $`M_6`$ in Sect. 2 further provides a clear link between the approaches and to this mathematical problem; we hope that this note will motivate further investigations of it.
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# Massive random matrix ensembles at 𝜷 = 1 & 4 : QCD in three dimensions
## I introduction
Spontaneous breaking of global symmetries has long been a subject of an extensive study for quantum field theories in various dimensions. Besides well-known cases such as 2D Gross-Neveu and Schwinger models and 4D QCD that exhibit breakdown of discrete or continuous chiral symmetry, it has been suspected that 3D QED and QCD may also undergo spontaneous breaking of flavor symmetry . For such odd dimensional field theories, there is no comprehensive theorem that predicts the surviving part of the flavor symmetry as in the case of even dimensions , as the notion of fermion chirality is absent. However, if there are even $`N_f=2n`$ number of massless 2-component complex fermions $`\psi _j`$, one can group them appropriately into 4-component complex fermions $`\mathrm{\Psi }_j`$ :
$$=\underset{j=1}{\overset{2n}{}}\overline{\psi }_jD_\mu \sigma _\mu \psi _j=\underset{j=1}{\overset{n}{}}\overline{\mathrm{\Psi }}_jD_\mu \gamma _\mu \mathrm{\Psi }_j,$$
(1.1)
where $`\gamma _\mu =\sigma _\mu \sigma _3(\mu =1,2,3)`$ in Euclidean 3D space. Then one can define a quasi-chirality according to two Hermitian matrices, e.g. $`\gamma _4=𝟙\sigma _\mathrm{𝟙}`$ and $`\gamma _5=𝟙\sigma _\mathrm{𝟚}`$, that anticommute with the Dirac operator $`D_\mu \gamma _\mu `$ and with each other. Generators of the flavor $`U(2n)`$ group are rearranged into that of $`U(n)`$ group times $`\{1,\gamma _4,\gamma _5,\gamma _S=𝟙\sigma _\mathrm{𝟛}\}`$ in the 4-component notation. In order to predict spontaneous flavor symmetry breaking pattern along the same line as in 4D QCD, one introduces a small symmetry-breaking mass term that is parity-invariant, i.e., $`_{j=1}^nm_j\overline{\mathrm{\Psi }}_j\mathrm{\Psi }_j`$ but not $`_{j=1}^nm_j\overline{\mathrm{\Psi }}_j\gamma _S\mathrm{\Psi }_j`$. In the 2-component notation it leads to include masses $`\{m\}=(m_1,\mathrm{},m_n,m_1,\mathrm{},m_n)`$. As the fermion determinant
$$\underset{j=1}{\overset{n}{}}det(iD_\mu \gamma _\mu +im_j)=\underset{j=1}{\overset{n}{}}det((D_\mu \sigma _\mu )^2+m_j^2)$$
(1.2)
is positive definite under this restriction, one can appeal to the Vafa-Witten theorem and predict that if the flavor symmetry is spontaneously broken, the absolute values of the fermion condensate $`\overline{\psi }_j\psi _j`$ are equal for all $`j=1,\mathrm{},2n`$ and their signs are the same as those of respective masses. That is, the continuous part of the global symmetry group is broken according to
$$U(2n)U(n)\times U(n)$$
(1.3)
by the order parameter
$$\mathrm{\Sigma }=\frac{1}{2n}\underset{j=1}{\overset{2n}{}}|\overline{\psi }_j\psi _j|,$$
(1.4)
while the discrete $`𝒁_2`$ group (a product of the parity and the exchange of fields $`\psi _j\psi _{n+j}`$) remains unbroken. The formation of the quark condensate was indeed observed in Monte-Carlo simulations on a lattice . This peculiar pattern of flavor symmetry breaking can also be predicted for 3D large-$`N_c`$ QCD by the Coleman-Witten argument , as remarked in Ref.. One can repeat the above argument valid for complex (fundamental representation of $`SU(N_c3)`$ gauge group) fermions towards the cases with even flavors of pseudoreal (fundamental representation of $`SU(2)`$ gauge group) and of real (adjoint representation of $`SU(N_c)`$ gauge group) fermions. We assign Dyson indices $`\beta =2,1,4`$, respectively, to these three cases, according to the anti-unitary symmetries of the associated Dirac operators . Then the continuous parts of the global symmetry groups are predicted to be broken down as
$`Sp(2n)`$ $``$ $`Sp(n)\times Sp(n)(\beta =1),`$ (1.6)
$`SO(2n)`$ $``$ $`SO(n)\times SO(n)(\beta =4).`$ (1.7)
These symmetry breaking patterns determine the forms of the low-energy effective Lagrangian of associated Nambu-Goldstone bosons:
$$_{\mathrm{eff}}=\frac{F^2}{4}\mathrm{tr}_\mu U_\mu U^{}\mathrm{\Sigma }\mathrm{tr}𝒎U\mathrm{\Gamma }U^{}+\mathrm{},$$
(1.8)
where $`U(𝐱)`$ takes its value in the coset manifolds
$$_F=\mathrm{AIII}_{n,n},\mathrm{CII}_{n,n},\mathrm{BDI}_{n,n},$$
(1.9)
respectively for $`\beta =2,1,4`$. The mass and chiral matrices in the above are defined as
$`𝒎=\{\begin{array}{cc}\mathrm{diag}(m_1,\mathrm{},m_n)\sigma _3\hfill & (\beta =2,4)\hfill \\ \mathrm{diag}(m_1,\mathrm{},m_n)\sigma _3i\sigma _2\hfill & (\beta =1)\hfill \end{array},`$ (1.12)
$`\mathrm{\Gamma }=\{\begin{array}{cc}𝟙_𝕟\sigma _\mathrm{𝟛}\hfill & (\beta =2,4)\hfill \\ 𝟙_𝕟\sigma _\mathrm{𝟛}𝕚\sigma _\mathrm{𝟚}\hfill & (\beta =1)\hfill \end{array}.`$ (1.15)
Given non-linear $`\sigma `$ models (NL$`\sigma `$Ms) in 4D, Leutwyler and Smilga proposed to extract out of them nonperturbative, exact information on the spectra of Dirac operators by imposing a constraint on the parameters: the linear dimension $`L`$ of the system be much smaller than the Compton length of Nambu-Goldstone bosons $`F/\sqrt{m\mathrm{\Sigma }}`$. In this ‘ergodic’ regime where the zero momentum mode of $`U(𝐱)`$ dominates, the effective finite-volume partition function simplifies into a finite-dimensional integral and its small-$`m`$ expansion yields a sequence of spectral sum rules. Verbaarschot and collaborators (see Ref. for a exhaustive list of references) made an important observation along this line, that the 0D finite-volume partition functions could as well be derived from models much simpler than QCD, random matrix ensembles (RMEs). In our 3D QCD context , it means that in the limit
$$L\mathrm{},m_i0,\mu _iL^3\mathrm{\Sigma }m_i\text{ : fixed},$$
(1.16)
the finite-volume partition functions
$`𝒵(\{\mu \})`$ $`=`$ $`{\displaystyle __F}𝑑U\mathrm{exp}\left(\mathrm{tr}𝝁U\mathrm{\Gamma }U^{}\right),`$ (1.17)
$`𝝁`$ $``$ $`L^3\mathrm{\Sigma }𝒎,`$ (1.18)
have alternative representation in terms of large-$`N`$ RMEs
$$Z(\{m\})=_𝒟𝑑H\mathrm{e}^{\beta \mathrm{tr}V(H^2)}\underset{i=1}{\overset{n}{}}det\left(H^2+m_j^2\right),$$
(1.19)
where the integral domains $`𝒟`$ are sets of $`N\times N`$ complex hermitian ($`u(N)=T(\mathrm{A}_N)`$), real symmetric ($`o(N)=T(\mathrm{AI}_N)`$), and quaternion selfdual ($`sp(N)=T(\mathrm{AII}_N)`$) matrices $`H`$ for $`\beta =2,1,4`$, respectively. The determinant in the $`\beta =4`$ case is understood as a quaternion determinant (Tdet). Their proofs consist of the ‘color-flavor’ (or Hubbard-Stratonovich) transformation that converts the integration variables into matrices with small ($`n\times n`$) dimensions, and the saddle point method under which
$$N\mathrm{},m_i0,\mu _i\pi \overline{\rho }(0)m_i\text{ : fixed}.$$
(1.20)
Here $`\overline{\rho }(0)`$ stands for the large-$`N`$ spectral density of the random matrix $`H`$:
$$\overline{\rho }(x)=\underset{N\mathrm{}}{lim}\mathrm{tr}\delta (xH),$$
(1.21)
at the spectral origin. The RMEs (1.19) are motivated by the microscopic theories (3D Euclidean QCD) on a lattice, with a crude simplification of replacing matrix elements of the Hermitian Dirac operator $`i/D=(i_\mu +A_\mu )\sigma _\mu `$ by random numbers $`H_{jk}`$ distributed according to the weight $`\mathrm{e}^{\beta \mathrm{tr}V(H^2)}`$. Under this correspondence, the microscopic limit (1.20) is equivalent to Leutwyler-Smilga limit (1.16), since the size $`N`$ of the matrix $`H`$ is interpreted as the number of sites $`L^3`$ of the lattice on which QCD is discretized, and the Dirac spectral density at zero virtuality $`\overline{\rho }(0)`$ is related to the quark condensate by the Banks-Casher relation $`\mathrm{\Sigma }=\pi \overline{\rho }(0)/L^3`$ .
On the other hand, the situation is far more subtle in the case with an odd number ($`N_f=2n+1`$) of 2-component fermion flavors . In the massless case, the fermion determinant $`det^{2n+1}(iD_\mu \sigma _\mu )`$ is not positive definite, and its phase is gauge dependent. This dependence can be compensated by the Chern-Simons term that is yielded by a gauge-invariant regularization, but this anomalous term explicitly breaks the parity . Therefore, even in a presence of small masses $`\{m\}=(m_1,\mathrm{},m_n,m_1,\mathrm{},m_n,0)`$ that respects the $`𝒁_2`$ invariance (combining parity and flavor exchange) classically, one cannot appeal to the previous argument to derive low-energy effective Lagrangians. With this situation in mind, we nevertheless adopt a pattern of spontaneous flavor symmetry breaking proposed by Verbaarschot and Zahed for $`\beta =2`$ and its generalizations to $`\beta =1`$ and 4:
$`U(2n+1)`$ $``$ $`U(n)\times U(n+1)(\beta =2),`$ (1.23)
$`Sp(2n+1)`$ $``$ $`Sp(n)\times Sp(n+1)(\beta =1),`$ (1.24)
$`SO(2n+1)`$ $``$ $`SO(n)\times SO(n+1)(\beta =4),`$ (1.25)
leading to NL$`\sigma `$Ms of Nambu-Goldstone fields over the coset manifolds
$$_F=\mathrm{AIII}_{n,n+1},\mathrm{CII}_{n,n+1},\mathrm{BDI}_{n,n+1},$$
(1.26)
respectively. Then, by the same token as in the case of even $`N_f`$, one can write down corresponding RMEs :
$$Z(\{m\})=_𝒟𝑑H\mathrm{e}^{\beta \mathrm{tr}V(H^2)}detH\underset{i=1}{\overset{n}{}}det\left(H^2+m_j^2\right),$$
(1.27)
which are equivalent, in the limit (1.20), to the ‘finite-volume partition functions’
$`𝒵(\{\mu \})`$ $`=`$ $`{\displaystyle __F}𝑑U\{\begin{array}{cc}\mathrm{cosh}\left(\mathrm{tr}𝝁U\mathrm{\Gamma }U^{}\right)\hfill & (N:\mathrm{even})\hfill \\ \mathrm{sinh}\left(\mathrm{tr}𝝁U\mathrm{\Gamma }U^{}\right)\hfill & (N:\mathrm{odd})\hfill \end{array},`$ (1.30)
$`\mathrm{\Gamma }`$ $`=`$ $`\{\begin{array}{cc}\mathrm{diag}(𝟙_𝕟,𝟙_{𝕟+\mathrm{𝟙}})\hfill & (\beta =2,4)\hfill \\ \mathrm{diag}(𝟙_𝕟,𝟙_{𝕟+\mathrm{𝟙}})𝕚\sigma _\mathrm{𝟚}\hfill & (\beta =1)\hfill \end{array}.`$ (1.33)
As the Chern-Simons term cannot be incorporated within these RMEs, their physical relevance is unclear . An immediate problem is that if the rank $`N`$ of the matrix $`H`$ is odd, the partition function (1.27) (or (1.30)) is zero, and the (unnormalized) correlation functions are odd under a simultaneous change of signs of the arguments, which are unacceptable as physical observables. The above relationships between RMEs and NL$`\sigma `$Ms for 3D QCD consist, together with its counterpart for 4D QCD ($`_F=\mathrm{A}_n,\mathrm{AII}_n,\mathrm{AI}_n`$ and $`𝒟=T(\mathrm{AIII}_{N,N^{}}),T(\mathrm{BDI}_{N,N^{}}),T(\mathrm{CII}_{N,N^{}})`$ for $`\beta =2,1,4`$, respectively), a part of Zirnbauer’s complete classification scheme of RMEs in terms of Riemannian symmetric spaces .
These RMEs are technically suited for the computation of correlations of eigenvalues $`\{x\}`$ of the Dirac operator $`i/DH`$ in the microscopic asymptotic limit where the energy eigenvalues are scaled as the quark masses are in Eq.(1.20),
$$N\mathrm{},x0,\lambda =\pi \overline{\rho }(0)x:\text{fixed}.$$
(1.34)
For the chiral RMEs describing 4D QCD in the low-energy ergodic regime, such Dirac spectral correlators have been computed previously for the massless and recently for the massive cases with all three values of $`\beta `$. On the other hand, for the (non-chiral) RMEs describing 3D QCD, they have been analytically computed solely for the unitary ($`\beta =2`$) ensemble, in the massless case , as well as in the massive case . For other values of $`\beta `$, a numerical work based on a finite-$`N`$ formula for the correlation functions recently appeared only in the massless case . The subject of this Article is to complete this program by analytically computing the partition and correlation functions for the orthogonal ($`\beta =1`$) and symplectic ($`\beta =4`$) ensembles in a presence of finite scaled mass parameters. We employ a slightly modified version of the method used in our previous articles .
We finally remark on the universality issue. It was noticed by Şener and Verbaarschot (see also Ref.) and proved by Widom that the diagonal element $`S(x,y)`$ of the quaternion kernel for an orthogonal or symplectic ensemble, and accordingly all spectral correlation functions thereof, can be constructed from the scalar kernel $`K(x,y)`$ for a unitary ensemble sharing the same weight function:
$`S^T`$ $`=`$ $`\left(I(IK)\epsilon KD\right)^1K(\beta =1),`$ (1.36)
$`S`$ $`=`$ $`\left(I(IK)DK\epsilon \right)^1K(\beta =4),`$ (1.37)
where $`I`$, $`D`$, $`\epsilon `$, $`S`$, $`K`$ stand for integral operators with convolution kernels $`\delta (xy)`$, $`\delta ^{}(xy)`$, $`\frac{1}{2}\mathrm{sgn}(xy)`$, $`S(x,y)`$, $`K(x,y)`$, respectively, <sup>T</sup> and <sup>-1</sup> stand for transpose and inverse operators. Since the scalar kernel in the asymptotic limit (1.34), (1.20) is insensitive to the details of the potential $`V(x^2)`$ either in the absence or in the presence of finite and nonzero $`\mu `$’s , the universality of correlation functions for orthogonal and symplectic ensembles are automatically guaranteed. Therefore it suffices for us to concentrate onto Gaussian ensembles, $`V(x^2)=x^2/2`$. This choice leads to Wigner’s semi-circle law
$$\overline{\rho }(x)=\frac{1}{\pi }\sqrt{2Nx^2}.$$
(1.38)
## II orthogonal ensemble
For $`\beta =1`$, we treat the following three cases separately:
$`𝐀`$ $`:`$ $`\{m\}=(m_1,\mathrm{},m_n,m_1,\mathrm{},m_n),`$
$`𝐁`$ $`:`$ $`\{m\}=(m_1,\mathrm{},m_n,m_1,\mathrm{},m_n,0),\mathrm{even}N,`$
$`𝐂`$ $`:`$ $`\{m\}=(m_1,\mathrm{},m_n,m_1,\mathrm{},m_n,0),\mathrm{odd}N.`$
### A even $`𝐍_𝐟`$
We first consider the case with $`N_f2n`$ flavors and $`\{m\}=(m_1,\mathrm{},m_n,m_1,\mathrm{},m_n)`$.
We express the partition function (1.19) of the RME in terms of eigenvalues $`\{x_j\}`$ of $`H`$ (up to a constant independent of $`m`$):
$$Z(\{m\})=\frac{1}{N!}_{\mathrm{}}^{\mathrm{}}\mathrm{}_{\mathrm{}}^{\mathrm{}}\underset{j=1}{\overset{N}{}}dx_j\underset{j=1}{\overset{N}{}}\left(\mathrm{e}^{x_j^2/2}\underset{k=1}{\overset{n}{}}(x_j^2+m_k^2)\right)\underset{j>k}{\overset{N}{}}|x_jx_k|.$$
(2.1)
The $`p`$-level correlation function of the matrix $`H`$ is defined as
$`\rho (x_1,\mathrm{},x_p;\{m\})`$ $`=`$ $`{\displaystyle \underset{j=1}{\overset{p}{}}}\mathrm{tr}\delta (x_jH)`$ (2.2)
$`=`$ $`{\displaystyle \frac{\mathrm{\Xi }_p(x_1,\mathrm{},x_p;\{m\})}{\mathrm{\Xi }_0(\{m\})}},`$ (2.3)
$`\mathrm{\Xi }_p(x_1,\mathrm{},x_p;\{m\})`$ $`=`$ $`{\displaystyle \frac{1}{(Np)!}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}\mathrm{}{\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle \underset{j=p+1}{\overset{N}{}}}dx_j{\displaystyle \underset{j=1}{\overset{N}{}}}\left(\mathrm{e}^{x_j^2/2}{\displaystyle \underset{k=1}{\overset{n}{}}}(x_j^2+m_k^2)\right){\displaystyle \underset{j>k}{\overset{N}{}}}|x_jx_k|`$ (2.4)
$`(\mathrm{\Xi }_0=Z)`$. We define new variables $`z_j`$ as
$`z_{2j1}=im_j(j=1,\mathrm{},n),`$ (2.5)
$`z_{2j}=im_j(j=1,\mathrm{},n),`$ (2.6)
$`z_{2n+j}=x_j(j=1,\mathrm{},p).`$ (2.7)
Then the multiple integral (2.4) is expressed as
$`\mathrm{\Xi }_p(z_1,\mathrm{},z_{2n+p})`$ $`=`$ $`{\displaystyle \frac{1}{_{j=1}^{2n}\sqrt{w(z_j)}_{j>k}^{2n}(z_jz_k)}}`$ (2.9)
$`\times {\displaystyle \frac{1}{(Np)!}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}\mathrm{}{\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle \underset{j=2n+p+1}{\overset{2n+N}{}}}dz_j{\displaystyle \underset{j=1}{\overset{2n+N}{}}}\sqrt{w(z_j)}{\displaystyle \underset{j>k}{\overset{2n+N}{}}}(z_jz_k){\displaystyle \underset{j>k>2n}{\overset{2n+N}{}}}\mathrm{sgn}(z_jz_k),`$
where $`w(z)=\mathrm{e}^{z^2}`$. Eq.(2.9) resembles an $`(2n+p)`$-level correlation function of the conventional Gaussian ensemble with $`2n+N`$ levels. However, conventionally the levels $`z_1,\mathrm{},z_{2n+p}`$ are all real, while in the present case some of them ($`z_1,\mathrm{},z_{2n}`$) are pure imaginary. We carefully incorporate this fact into the following evaluation.
Let us denote the integrand in Eq.(2.9) as
$$p(z_1,\mathrm{},z_{2n+N})=\underset{j=1}{\overset{2n+N}{}}\sqrt{w(z_j)}\underset{j>k}{\overset{2n+N}{}}(z_jz_k)\underset{j>k>2n}{\overset{2n+N}{}}\mathrm{sgn}(z_jz_k).$$
(2.10)
For $`N`$ even, an identity
$$\underset{j>k}{\overset{2n+N}{}}\mathrm{sgn}(z_jz_k)=\mathrm{Pf}[\mathrm{sgn}(z_kz_j)]_{j,k=1,\mathrm{},2n+N}$$
(2.11)
holds for real $`z_1,\mathrm{},z_{2n+N}`$. By taking the limit $`z_1<z_2<\mathrm{}<z_{2n}\mathrm{}`$, we find another identity
$$\underset{j>k>2n}{\overset{2n+N}{}}\mathrm{sgn}(z_jz_k)=\mathrm{Pf}[F_{jk}]_{j,k=1,\mathrm{},2n+N},$$
(2.12)
where
$$F_{jk}=\{\begin{array}{cc}\mathrm{sgn}(kj)\hfill & (j,k=1,\mathrm{},2n)\hfill \\ 1\hfill & (j=1,\mathrm{},2n;k=2n+1,\mathrm{},2n+N)\hfill \\ 1\hfill & (j=2n+1,\mathrm{},2n+N;k=1,\mathrm{},2n)\hfill \\ \mathrm{sgn}(z_kz_j)\hfill & (j,k=2n+1,\mathrm{},2n+N)\hfill \end{array}.$$
(2.13)
Substitution of Eq.(2.12) into Eq.(2.10) yields
$$p(z_1,\mathrm{},z_{2n+N})=\underset{j=1}{\overset{2n+N}{}}\sqrt{w(z_j)}\underset{j>k}{\overset{2n+N}{}}(z_jz_k)\mathrm{Pf}[F_{jk}]_{j,k=1,\mathrm{},2n+N}.$$
(2.14)
For $`N`$ odd, we similarly obtain
$$p(z_1,\mathrm{},z_{2n+N})=\underset{j=1}{\overset{2n+N}{}}\sqrt{w(z_j)}\underset{j>k}{\overset{2n+N}{}}(z_jz_k)\mathrm{Pf}\left[\begin{array}{cc}[F_{jk}]_{j,k=1,\mathrm{},2n+N}\hfill & [g_j]_{j=1,\mathrm{},2n+N}\hfill \\ \left[g_k\right]_{k=1,\mathrm{},2n+N}\hfill & 0\hfill \end{array}\right],$$
(2.15)
with $`g_j=g_k=1`$ ($`j,k=1,\mathrm{},2n+N`$). The Pfaffians in the above can be represented as quaternion determinants . In doing so, we introduce monic skew-orthogonal polynomials $`R_j(z)=z^j+\mathrm{}`$, which satisfy the skew-orthogonality relation:
$`R_{2j},R_{2k+1}_R=R_{2k+1},R_{2j}_R=r_j\delta _{jk},`$ (2.16)
$`R_{2j},R_{2k}_R=R_{2j+1},R_{2k+1}_R=0,`$ (2.17)
where
$$f,g_R=_{\mathrm{}}^{\mathrm{}}dz\sqrt{w(z)}g(z)_{\mathrm{}}^zdz^{}\sqrt{w(z^{})}f(z^{})(fg).$$
(2.18)
Explicit forms for the skew-orthogonal polynomials and their norms associated with the Gaussian weight $`w(z)`$ are known :
$`R_{2j}(z)`$ $`=`$ $`{\displaystyle \frac{1}{2^{2j}}}H_{2j}(z),`$ (2.19)
$`R_{2j+1}(z)`$ $`=`$ $`{\displaystyle \frac{1}{2^{2j+1}}}\left(H_{2j+1}(z)H_{2j}^{}(z)\right),`$ (2.20)
$`r_j`$ $`=`$ $`2^{2j+1}(2j)!\sqrt{\pi },`$ (2.21)
in terms of the Hermite polynomials
$$H_j(z)=(1)^j\mathrm{e}^{z^2}\frac{d^j}{dz^j}\mathrm{e}^{z^2}.$$
(2.22)
Now we present the following theorems:
Theorem 1
For even $`N`$, we can rewrite $`p(z_1,\mathrm{},z_{2n+N})`$ as
$$p(z_1,\mathrm{},z_{2n+N})=(\underset{j=0}{\overset{n+N/21}{}}r_j)\mathrm{Tdet}[f_{jk}(z_j,z_k)]_{j,k=1,\mathrm{},2n+N}.$$
(2.23)
The quaternion elements $`f_{jk}(z_j,z_k)`$ are represented as
$$f_{jk}(z_j,z_k)=\left[\begin{array}{cc}S(z_j,z_k)& I(z_j,z_k)\\ D(z_j,z_k)& S(z_k,z_j)\end{array}\right].$$
(2.24)
The functions $`S(z_j,z_k)`$, $`D(z_j,z_k)`$ and $`I(z_j,z_k)`$ are given by
$`S(z_j,z_k)`$ $`=`$ $`{\displaystyle \underset{\mathrm{}=0}{\overset{n+N/21}{}}}{\displaystyle \frac{\mathrm{\Phi }_2\mathrm{}(z_j)\mathrm{\Psi }_{2\mathrm{}+1}(z_k)\mathrm{\Phi }_{2\mathrm{}+1}(z_j)\mathrm{\Psi }_2\mathrm{}(z_k)}{r_{\mathrm{}}}},`$ (2.25)
$`D(z_j,z_k)`$ $`=`$ $`{\displaystyle \underset{\mathrm{}=0}{\overset{n+N/21}{}}}{\displaystyle \frac{\mathrm{\Psi }_2\mathrm{}(z_j)\mathrm{\Psi }_{2\mathrm{}+1}(z_k)\mathrm{\Psi }_{2\mathrm{}+1}(z_j)\mathrm{\Psi }_2\mathrm{}(z_k)}{r_{\mathrm{}}}},`$ (2.26)
$`I(z_j,z_k)`$ $`=`$ $`{\displaystyle \underset{\mathrm{}=0}{\overset{n+N/21}{}}}{\displaystyle \frac{\mathrm{\Phi }_2\mathrm{}(z_j)\mathrm{\Phi }_{2\mathrm{}+1}(z_k)\mathrm{\Phi }_{2\mathrm{}+1}(z_j)\mathrm{\Phi }_2\mathrm{}(z_k)}{r_{\mathrm{}}}}+F_{jk},`$ (2.27)
where
$`\mathrm{\Psi }_j(z)=\sqrt{w(z)}R_j(z),`$ (2.28)
$`\mathrm{\Phi }_j(z_k)=\{\begin{array}{ccc}_{\mathrm{}}^{\mathrm{}}𝑑z\sqrt{w(z)}R_j(z)\mathrm{sgn}(z_kz)\hfill & (k=2n+1,\mathrm{},2n+N)\hfill & \\ _{\mathrm{}}^{\mathrm{}}𝑑z\sqrt{w(z)}R_j(z)\hfill & (k:\mathrm{otherwise})\hfill & \end{array}.`$ (2.31)
Theorem 2
For odd $`N`$, we have
$$p(z_1,\mathrm{},z_{2n+N})=\left(\underset{j=0}{\overset{n+[N/2]1}{}}r_j\right)s_{2n+N1}\mathrm{Tdet}[f_{jk}^{\mathrm{odd}}(z_j,z_k)]_{j,k=1,\mathrm{},2n+N}.$$
(2.32)
The quaternion elements are represented as
$$f_{jk}^{\mathrm{odd}}(z_j,z_k)=\left[\begin{array}{cc}S^{\mathrm{odd}}(z_j,z_k)& I^{\mathrm{odd}}(z_j,z_k)\\ D^{\mathrm{odd}}(z_j,z_k)& S^{\mathrm{odd}}(z_k,z_j)\end{array}\right].$$
(2.33)
and
$$s_j=_{\mathrm{}}^{\mathrm{}}\mathrm{\Psi }_j(z)𝑑z.$$
(2.34)
The functions $`S^{\mathrm{odd}}`$, $`D^{\mathrm{odd}}`$ and $`I^{\mathrm{odd}}`$ are given in terms of $`S`$, $`D`$ and $`I`$ defined in Eq.(2.26), according to
$`S^{\mathrm{odd}}(z_j,z_k)`$ $`=`$ $`S(z_j,z_k)|_{}+{\displaystyle \frac{\mathrm{\Psi }_{2n+N1}(z_k)}{s_{2n+N1}}},`$ (2.35)
$`D^{\mathrm{odd}}(z_j,z_k)`$ $`=`$ $`D(z_j,z_k)|_{},`$ (2.36)
$`I^{\mathrm{odd}}(z_j,z_k)`$ $`=`$ $`I(z_j,z_k)|_{}+{\displaystyle \frac{\mathrm{\Phi }_{2n+N1}(z_j)\mathrm{\Phi }_{2n+N1}(z_k)}{s_{2n+N1}}}.`$ (2.37)
Here $``$ stands for a set of substitutions
$$R_j(z)R_j(z)\frac{s_j}{s_{2n+N1}}R_{2n+N1}(z)(j=0,\mathrm{},2n+N2),$$
(2.38)
associated with a change in the upper limit of the sum
$$n+\frac{N}{2}1n+\left[\frac{N}{2}\right]1.$$
(2.39)
Theorem 3
Let the quaternion elements $`q_{jk}`$ of a selfdual $`n\times n`$ matrix $`Q_n`$ depend on $`n`$ real or complex variables $`z_1,\mathrm{},z_n`$ as
$$q_{jk}=f_{jk}(z_j,z_k).$$
(2.40)
We assume that $`f_{jk}(z_j,z_k)`$ satisfies the following conditions.
$`{\displaystyle f_{nn}(z_n,z_n)𝑑\mu (z_n)}=c_n,`$ (2.42)
$`{\displaystyle f_{jn}(z_j,z_n)f_{nk}(z_n,z_k)𝑑\mu (z_n)}=f_{jk}(z_j,z_k)+\lambda f_{jk}(z_j,z_k)f_{jk}(z_j,z_k)\lambda .`$ (2.43)
Here $`d\mu (z)`$ is a suitable measure, $`c_n`$ is a constant scalar, and $`\lambda `$ is a constant quaternion. Then we have
$$\mathrm{Tdet}Q_n𝑑\mu (z_n)=(c_nn+1)\mathrm{Tdet}Q_{n1},$$
(2.44)
where $`Q_{n1}`$ is the $`(n1)\times (n1)`$ matrix obtained by removing the row and the column which contain $`z_n`$.
It is straightforward to show that the quaternion elements $`f_{jk}(z_j,z_k)`$ and $`f_{jk}^{\mathrm{odd}}(z_j,z_k)`$ in Theorem 1 and Theorem 2, respectively, both satisfy the conditions imposed on $`f_{jk}(z_j,z_k)`$ in Theorem 3. This means that we can inductively write
$$\mathrm{\Xi }_p(z_1,\mathrm{},z_{2n+p})=\frac{_{j=0}^{p+[N/2]1}r_j}{_{j=1}^{2n}\sqrt{w(z_j)}_{j>k}^{2n}(z_jz_k)}\times \{\begin{array}{ccc}& \mathrm{Tdet}[f_{jk}(z_j,z_k)]_{j,k=1,\mathrm{},2n+p}\hfill & (N:\text{even})\hfill \\ s_{2n+N1}\hfill & \mathrm{Tdet}[f_{jk}^{\mathrm{odd}}(z_j,z_k)]_{j,k=1,\mathrm{},2n+p}\hfill & (N:\text{odd})\hfill \end{array}.$$
(2.45)
Since the final result in the asymptotic limit $`N\mathrm{}`$ should be insensitive to the parity of $`N`$, we consider only even $`N`$ henceforth. Then the $`p`$-level correlation function (2.3) is written as
$$\rho (x_1,\mathrm{},x_p;\{m\})=\frac{\mathrm{\Xi }_p(z_1,\mathrm{},z_{2n+p})}{\mathrm{\Xi }_0(z_1,\mathrm{},z_{2n})}=\frac{\mathrm{Tdet}[f_{jk}(z_j,z_k)]_{j,k=1,\mathrm{},2n+p}}{\mathrm{Tdet}[f_{jk}(z_j,z_k)]_{j,k=1,\mathrm{},2n}}.$$
(2.46)
We introduce a set of notations
$`S_{jk}^{II}=S(z_j,z_k)(j,k=1,\mathrm{},2n),`$ (2.47)
$`S_{jk}^{IR}=S(z_j,z_{2n+k})(j=1,\mathrm{},2n;k=1,\mathrm{},p),`$ (2.48)
$`S_{jk}^{RI}=S(z_{2n+j},z_k)(j=1,\mathrm{},p;k=1,\mathrm{},2n),`$ (2.49)
$`S_{jk}^{RR}=S(z_{2n+j},z_{2n+k})(j,k=1,\mathrm{},p)`$ (2.50)
and similarly for $`D`$ and $`I`$. Using Dyson’s identity (A.16) we can rewrite the correlation functions as
$`\rho (x_1,\mathrm{},x_p;\{m\})`$ $`=`$ $`(1)^{p(p1)/2}{\displaystyle \frac{\mathrm{Pf}\left[\begin{array}{cccc}I^{II}& S^{II}& I^{IR}& S^{IR}\\ (S^{II})^T& D^{II}& (S^{RI})^T& D^{IR}\\ I^{RI}& S^{RI}& I^{RR}& S^{RR}\\ (S^{IR})^T& D^{RI}& (S^{RR})^T& D^{RR}\end{array}\right]}{\mathrm{Pf}\left[\begin{array}{cc}I^{II}& S^{II}\\ (S^{II})^T& D^{II}\end{array}\right]}}`$ (2.57)
$`=`$ $`(1)^{p(p1)/2}{\displaystyle \frac{\mathrm{Pf}\left[\begin{array}{ccc}D^{II}& (S^{RI})^T& (D^{RI})^T\\ S^{RI}& I^{RR}& S^{RR}\\ D^{RI}& (S^{RR})^T& D^{RR}\end{array}\right]}{\mathrm{Pf}\left[D^{II}\right]}}.`$ (2.61)
In the last line we have exploited a Pfaffian identity that holds for antisymmetric matrices $`A`$, $`B`$ of even ranks and a row vector $`v`$:
$$\mathrm{Pf}\left[\begin{array}{cc}A& \begin{array}{c}v\\ \mathrm{}\\ v\end{array}\\ & \\ v^T\mathrm{}v^T& B\end{array}\right]=\mathrm{Pf}[A]\mathrm{Pf}[B].$$
(2.62)
Now we proceed to evaluate the component functions of the quaternion kernel $`f_{jk}(z_j,z_k)`$ in the asymptotic limit (1.20) and (1.34), where unfolded microscopic variables
$`\sqrt{2N}z_{2j1}`$ $``$ $`\zeta _{2j1}=i\mu _j(j=1,\mathrm{},n),`$ (2.63)
$`\sqrt{2N}z_{2j}`$ $``$ $`\zeta _{2j}=i\mu _j(j=1,\mathrm{},n),`$ (2.64)
$`\sqrt{2N}z_{2n+j}`$ $``$ $`\lambda _j(j=1,\mathrm{},p)`$ (2.65)
are kept fixed. We note that all elements of the sub-matrices that appear in the second line of Eq.(2.61) are expressed as (derivatives or integrals of) an analytic function
$`\overline{S}(z,z^{})`$ $`=`$ $`{\displaystyle \underset{j=0}{\overset{n+N/21}{}}}{\displaystyle \frac{\overline{\mathrm{\Phi }}_{2j}(z)\mathrm{\Psi }_{2j+1}(z^{})\overline{\mathrm{\Phi }}_{2j+1}(z)\mathrm{\Psi }_{2j}(z^{})}{r_j}}`$ (2.66)
$`=`$ $`{\displaystyle \frac{\mathrm{e}^{(z^2+z^{}{}_{}{}^{2})/2}}{2^{2n+N}\sqrt{\pi }\mathrm{\Gamma }(2n+N)}}{\displaystyle \frac{H_{2n+N}(z)H_{2n+N1}(z^{})H_{2n+N1}(z)H_{2n+N}(z^{})}{zz^{}}}`$ (2.68)
$`+{\displaystyle \frac{\mathrm{e}^{z^{}{}_{}{}^{2}/2}}{2^{2n+N}\sqrt{\pi }\mathrm{\Gamma }(2n+N)}}H_{2n+N1}(z^{}){\displaystyle _0^z}\mathrm{e}^{u^2/2}H_{2n+N}(u)𝑑u,`$
where
$$\overline{\mathrm{\Phi }}_j(z)=\left\{_{\mathrm{}}^z_z^{\mathrm{}}\right\}\mathrm{\Psi }_j(u)du,$$
(2.69)
with real and/or imaginary arguments:
$`D_{jk}^{II}`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \frac{}{z_j}}\overline{S}(z_j,z_k),`$ (2.71)
$`S_{jk}^{RI}`$ $`=`$ $`\overline{S}(x_j,z_k),`$ (2.72)
$`D_{jk}^{RI}`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \frac{}{x_j}}\overline{S}(x_j,z_k),`$ (2.73)
$`S_{jk}^{RR}`$ $`=`$ $`\overline{S}(x_j,x_k),`$ (2.74)
$`D_{jk}^{RR}`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \frac{}{x_j}}\overline{S}(x_j,x_k),`$ (2.75)
$`I_{jk}^{RR}`$ $`=`$ $`2{\displaystyle _{x_j}^{x_k}}\overline{S}(x_j,x)𝑑x\mathrm{sgn}(x_jx_k).`$ (2.76)
In the second line of Eq.(2.68), we have singled out the unitary scalar kernel and applied to it the Christoffel-Darboux formula. Substituting asymptotic formulas for the Hermite polynomials
$`H_{2k}(z)`$ $``$ $`{\displaystyle \frac{(1)^k2^{2k}k!}{\sqrt{\pi k}}}\mathrm{cos}(2\sqrt{k}z),`$ (2.77)
$`H_{2k+1}(z)`$ $``$ $`{\displaystyle \frac{(1)^k2^{2k+1}k!}{\sqrt{\pi }}}\mathrm{sin}(2\sqrt{k}z),`$ (2.78)
valid under $`k\mathrm{}`$, $`z0`$, $`\sqrt{k}z:`$ fixed, one can show that $`\overline{S}(z,z^{})`$ approaches the sine kernel :
$$\frac{1}{\sqrt{2N}}\overline{S}(\frac{\zeta }{\sqrt{2N}},\frac{\zeta ^{}}{\sqrt{2N}})\frac{\mathrm{sin}(\zeta \zeta ^{})}{\pi (\zeta \zeta ^{})}K(\zeta \zeta ^{}).$$
(2.79)
The Pfaffian elements in the asymptotic limit, after taking into account an unfolding by the factor $`\sqrt{2N}`$, are then expressed in terms of $`K(\zeta )`$:
$`𝐃_{jk}^{II}`$ $``$ $`{\displaystyle \frac{1}{2N}}D(z_j,z_k){\displaystyle \frac{1}{2}}K^{}(\zeta _j\zeta _k),`$ (2.81)
$`𝐒_{jk}^{RI}`$ $``$ $`{\displaystyle \frac{1}{\sqrt{2N}}}S(x_j,z_k)K(\lambda _j\zeta _k),`$ (2.82)
$`𝐃_{jk}^{RI}`$ $``$ $`{\displaystyle \frac{1}{2N}}D(x_j,z_k){\displaystyle \frac{1}{2}}K^{}(\lambda _j\zeta _k),`$ (2.83)
$`𝐒_{jk}^{RR}`$ $``$ $`{\displaystyle \frac{1}{\sqrt{2N}}}S(x_j,x_k)K(\lambda _j\lambda _k),`$ (2.84)
$`𝐃_{jk}^{RR}`$ $``$ $`{\displaystyle \frac{1}{2N}}D(x_j,x_k){\displaystyle \frac{1}{2}}K^{}(\lambda _j\lambda _k),`$ (2.85)
$`𝐈_{jk}^{RR}`$ $``$ $`I(x_j,x_k)2{\displaystyle _0^{\lambda _j\lambda _k}}K(\lambda )𝑑\lambda \mathrm{sgn}(\lambda _j\lambda _k).`$ (2.86)
These matrix elements constitute the finite-volume partition function
$`𝒵(\{\mu \})`$ $``$ $`\mathrm{\Xi }_0(\{{\displaystyle \frac{\mu }{\sqrt{2N}}}\})`$ (2.87)
$`=`$ $`\mathrm{const}.{\displaystyle \frac{\mathrm{Pf}[𝐃^{II}]}{_{j=1}^n\mu _j_{j>k}^n(\mu _j^2\mu _k^2)^2}},`$ (2.88)
and the scaled spectral correlation functions
$`\rho _s(\lambda _1,\mathrm{},\lambda _p;\{\mu \})`$ $``$ $`\left({\displaystyle \frac{1}{\sqrt{2N}}}\right)^p\rho ({\displaystyle \frac{\lambda _1}{\sqrt{2N}}},\mathrm{},{\displaystyle \frac{\lambda _p}{\sqrt{2N}}};\{{\displaystyle \frac{\mu }{\sqrt{2N}}}\})`$ (2.89)
$`=`$ $`(1)^{p(p1)/2}{\displaystyle \frac{\mathrm{Pf}\left[\begin{array}{ccc}𝐃^{II}& (𝐒^{RI})^T& (𝐃^{RI})^T\\ 𝐒^{RI}& 𝐈^{RR}& 𝐒^{RR}\\ 𝐃^{RI}& (𝐒^{RR})^T& 𝐃^{RR}\end{array}\right]}{\mathrm{Pf}\left[𝐃^{II}\right]}}.`$ (2.93)
In the quenched limit $`\mu _1,\mathrm{},\mu _n\mathrm{}`$ when the ratio of two Pfaffians is replaced by a minor $`\mathrm{Pf}\left[\genfrac{}{}{0pt}{}{𝐈^{RR}𝐒^{RR}}{(𝐒^{RR})^T𝐃^{RR}}\right]`$, the correlation functions approach those of the Gaussian orthogonal ensemble . By the same token, it satisfies a sequence
$$\rho _s(\{\lambda \};\mu _1,\mathrm{},\mu _n,\mu _1,\mathrm{},\mu _n)\stackrel{\mu _n\mathrm{}}{}\rho _s(\{\lambda \};\mu _1,\mathrm{},\mu _{n1},\mu _1,\mathrm{},\mu _{n1})\stackrel{\mu _{n1}\mathrm{}}{}\mathrm{},$$
(2.94)
as each of the masses is decoupled by going to infinity. To illustrate this decoupling, we exhibit in FIG.1 a plot of the spectral density $`\rho _s(\lambda ;\mu ,\mu )`$ ($`p=1,n=1`$).
### B odd $`𝐍_𝐟`$, even $`𝐍`$
Next we consider the case with $`N_f2n+1`$ flavors, $`\{m\}=(m_1,\mathrm{},m_n,m_1,\mathrm{},m_n,0)`$, and with even $`N`$. We express the partition function (1.27) of the RME in terms of eigenvalues $`\{x_j\}`$ of $`H`$:
$`Z(\{m\})={\displaystyle \frac{1}{N!}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}\mathrm{}{\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle \underset{j=1}{\overset{N}{}}}dx_j{\displaystyle \underset{j=1}{\overset{N}{}}}\left(\mathrm{e}^{x_j^2/2}x_j{\displaystyle \underset{k=1}{\overset{n}{}}}(x_j^2+m_k^2)\right){\displaystyle \underset{j>k}{\overset{N}{}}}|x_jx_k|.`$ (2.95)
The $`p`$-level correlation function of the matrix $`H`$ is defined as
$`\rho (x_1,\mathrm{},x_p;\{m\})`$ $`=`$ $`{\displaystyle \underset{j=1}{\overset{p}{}}}\mathrm{tr}\delta (x_jH)`$ (2.96)
$`=`$ $`{\displaystyle \frac{\mathrm{\Xi }_p(x_1,\mathrm{},x_p;\{m\})}{\mathrm{\Xi }_0(\{m\})}},`$ (2.97)
$`\mathrm{\Xi }_p(x_1,\mathrm{},x_p;\{m\})`$ $`=`$ $`{\displaystyle \frac{1}{(Np)!}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}\mathrm{}{\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle \underset{j=p+1}{\overset{N}{}}}dx_j{\displaystyle \underset{j=1}{\overset{N}{}}}\left(\mathrm{e}^{x_j^2/2}x_j{\displaystyle \underset{k=1}{\overset{n}{}}}(x_j^2+m_k^2)\right){\displaystyle \underset{j>k}{\overset{N}{}}}|x_jx_k|`$ (2.98)
$`(\mathrm{\Xi }_0=Z)`$. We define new variables $`z_j`$ as
$`z_0=0,`$ (2.99)
$`z_{2j1}=im_j(j=1,\mathrm{},n),`$ (2.100)
$`z_{2j}=im_j(j=1,\mathrm{},n),`$ (2.101)
$`z_{2n+j}=x_j(j=1,\mathrm{},p).`$ (2.102)
Then the multiple integral (2.98) is expressed as
$`\mathrm{\Xi }_p(z_0,\mathrm{},z_{2n+p})`$ $`=`$ $`{\displaystyle \frac{1}{_{j=0}^{2n}\sqrt{w(z_j)}_{j>k0}^{2n}(z_jz_k)}}`$ (2.104)
$`\times {\displaystyle \frac{1}{(Np)!}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}\mathrm{}{\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle \underset{j=2n+p+1}{\overset{2n+N}{}}}dz_j{\displaystyle \underset{j=0}{\overset{2n+N}{}}}\sqrt{w(z_j)}{\displaystyle \underset{j>k0}{\overset{2n+N}{}}}(z_jz_k){\displaystyle \underset{j>k>2n}{\overset{2n+N}{}}}\mathrm{sgn}(z_jz_k),`$
where $`w(z)=\mathrm{e}^{z^2}`$.
Let us denote the integrand in Eq.(2.104) as
$$p(z_0,\mathrm{},z_{2n+N})=\underset{j=0}{\overset{2n+N}{}}\sqrt{w(z_j)}\underset{j>k0}{\overset{2n+N}{}}(z_jz_k)\underset{j>k>2n}{\overset{2n+N}{}}\mathrm{sgn}(z_jz_k).$$
(2.105)
An identity
$$\underset{j>k0}{\overset{2n+N+1}{}}\mathrm{sgn}(z_jz_k)=\mathrm{Pf}[\mathrm{sgn}(z_kz_j)]_{j,k=0,\mathrm{},2n+N+1}$$
(2.106)
holds for real $`z_0,\mathrm{},z_{2n+N+1}`$. By taking the limit $`z_0<z_1<\mathrm{}<z_{2n}\mathrm{}`$ and $`z_{2n+N+1}+\mathrm{}`$, we find another identity
$$\underset{j>k>2n}{\overset{2n+N+1}{}}\mathrm{sgn}(z_jz_k)=\mathrm{Pf}\left[\begin{array}{cc}[F_{jk}]_{j,k=0,\mathrm{},2n+N}\hfill & [g_j]_{j=0,\mathrm{},2n+N}\hfill \\ \left[g_k\right]_{k=0,\mathrm{},2n+N}\hfill & 0\hfill \end{array}\right],$$
(2.107)
where
$$F_{jk}=\{\begin{array}{cc}\mathrm{sgn}(kj)\hfill & (j,k=0,\mathrm{},2n)\hfill \\ 1\hfill & (j=0,\mathrm{},2n;k=2n+1,\mathrm{},2n+N)\hfill \\ 1\hfill & (j=2n+1,\mathrm{},2n+N;k=0,\mathrm{},2n)\hfill \\ \mathrm{sgn}(z_kz_j)\hfill & (j,k=2n+1,\mathrm{},2n+N)\hfill \end{array}$$
(2.108)
and $`g_j=g_k=1`$ ($`j,k=0,\mathrm{},2n+N`$). Substitution of Eq.(2.107) into Eq.(2.105) yields
$$p(z_0,\mathrm{},z_{2n+N})=\underset{j=0}{\overset{2n+N}{}}\sqrt{w(z_j)}\underset{j>k0}{\overset{2n+N}{}}(z_jz_k)\mathrm{Pf}\left[\begin{array}{cc}[F_{jk}]_{j,k=0,\mathrm{},2n+N}\hfill & [g_j]_{j=0,\mathrm{},2n+N}\hfill \\ \left[g_k\right]_{k=0,\mathrm{},2n+N}\hfill & 0\hfill \end{array}\right].$$
(2.109)
The Pfaffian in the above can be represented as a quaternion determinant, due to the following theorem (which is essentially Theorem 2):
Theorem 2’
For even $`N`$, we can rewrite $`p(z_0,\mathrm{},z_{2n+N})`$ as
$$p(z_0,\mathrm{},z_{2n+N})=\left(\underset{j=0}{\overset{n+N/21}{}}r_j\right)s_{2n+N}\mathrm{Tdet}[f_{jk}^{\mathrm{even}}(z_j,z_k)]_{j,k=0,\mathrm{},2n+N}.$$
(2.110)
The quaternion elements are represented as
$$f_{jk}^{\mathrm{even}}(z_j,z_k)=\left[\begin{array}{cc}S^{\mathrm{even}}(z_j,z_k)& I^{\mathrm{even}}(z_j,z_k)\\ D^{\mathrm{even}}(z_j,z_k)& S^{\mathrm{even}}(z_k,z_j)\end{array}\right].$$
(2.111)
and $`s_j`$ is defined in Eq.(2.34). The functions $`S^{\mathrm{even}}`$, $`D^{\mathrm{even}}`$ and $`I^{\mathrm{even}}`$ are given in terms of $`S`$, $`D`$ and $`I`$ defined in Eq.(2.26), according to
$`S^{\mathrm{even}}(z_j,z_k)`$ $`=`$ $`S(z_j,z_k)|_\mathrm{\#}+{\displaystyle \frac{\mathrm{\Psi }_{2n+N}(z_k)}{s_{2n+N}}},`$ (2.112)
$`D^{\mathrm{even}}(z_j,z_k)`$ $`=`$ $`D(z_j,z_k)|_\mathrm{\#},`$ (2.113)
$`I^{\mathrm{even}}(z_j,z_k)`$ $`=`$ $`I(z_j,z_k)|_\mathrm{\#}+{\displaystyle \frac{\mathrm{\Phi }_{2n+N}(z_j)\mathrm{\Phi }_{2n+N}(z_k)}{s_{2n+N}}}.`$ (2.114)
Here $`\mathrm{\#}`$ stands for a set of substitutions
$$R_j(z)R_j(z)\frac{s_j}{s_{2n+N}}R_{2n+N}(z)(j=0,\mathrm{},2n+N1).$$
(2.115)
It is straightforward to show that the quaternion element $`f_{jk}^{\mathrm{even}}(z_j,z_k)`$ in Theorem 2’ satisfies the conditions imposed on $`f_{jk}(z_j,z_k)`$ in Theorem 3. This means that we can inductively write
$$\mathrm{\Xi }_p(z_0,\mathrm{},z_{2n+p})=\frac{(_{j=0}^{p+N/21}r_j)s_{2n+N}}{_{j=0}^{2n}\sqrt{w(z_j)}_{j>k0}^{2n}(z_jz_k)}\mathrm{Tdet}[f_{jk}^{\mathrm{even}}(z_j,z_k)]_{j,k=0,\mathrm{},2n+p}.$$
(2.116)
Then the $`p`$-level correlation function (2.97) is written as
$$\rho (x_1,\mathrm{},x_p;\{m\})=\frac{\mathrm{\Xi }_p(z_0,\mathrm{},z_{2n+p})}{\mathrm{\Xi }_0(z_0,\mathrm{},z_{2n})}=\frac{\mathrm{Tdet}[f_{jk}^{\mathrm{even}}(z_j,z_k)]_{j,k=0,\mathrm{},2n+p}}{\mathrm{Tdet}[f_{jk}^{\mathrm{even}}(z_j,z_k)]_{j,k=0,\mathrm{},2n}}.$$
(2.117)
We introduce a set of notations
$`S_{jk}^{II}=S^{\mathrm{even}}(z_j,z_k)(j,k=0,\mathrm{},2n),`$ (2.118)
$`S_{jk}^{IR}=S^{\mathrm{even}}(z_j,z_{2n+k})(j=0,\mathrm{},2n;k=1,\mathrm{},p),`$ (2.119)
$`S_{jk}^{RI}=S^{\mathrm{even}}(z_{2n+j},z_k)(j=1,\mathrm{},p;k=0,\mathrm{},2n),`$ (2.120)
$`S_{jk}^{RR}=S^{\mathrm{even}}(z_{2n+j},z_{2n+k})(j,k=1,\mathrm{},p)`$ (2.121)
and similarly for $`D`$ and $`I`$. A simplification
$`S_{jk}^{II}=S^{\mathrm{even}}(\mathrm{},z_k)={\displaystyle \frac{\mathrm{\Psi }_{2n+N}(z_k)}{s_{2n+N}}}S_k^I,`$ (2.122)
$`S_{jk}^{IR}=S^{\mathrm{even}}(\mathrm{},z_{2n+k})={\displaystyle \frac{\mathrm{\Psi }_{2n+N}(z_{2n+k})}{s_{2n+N}}}S_k^R,`$ (2.123)
$`I_{jk}^{IR}=I_{kj}^{RI}=I^{\mathrm{even}}(\mathrm{},z_{2n+k})={\displaystyle \frac{\mathrm{\Phi }_{2n+N}(z_{2n+k})}{s_{2n+N}}}I_k^R,`$ (2.124)
results from the definitions (2.26) and (2.113). Using Dyson’s identity (A.16) we can rewrite the correlation functions as
$`\rho (x_1,\mathrm{},x_p;\{m\})`$ $`=`$ $`(1)^{p(p1)/2}{\displaystyle \frac{\mathrm{Pf}\left[\begin{array}{cccc}I^{II}& S^{II}& I^{IR}& S^{IR}\\ (S^{II})^T& D^{II}& (S^{RI})^T& D^{IR}\\ I^{RI}& S^{RI}& I^{RR}& S^{RR}\\ (S^{IR})^T& D^{RI}& (S^{RR})^T& D^{RR}\end{array}\right]}{\mathrm{Pf}\left[\begin{array}{cc}I^{II}& S^{II}\\ (S^{II})^T& D^{II}\end{array}\right]}}`$ (2.131)
$`=`$ $`(1)^{p(p1)/2}{\displaystyle \frac{\mathrm{Pf}\left[\begin{array}{cccc}0& S^I& I^R& S^R\\ (S^I)^T& D^{II}& (S^{RI})^T& (D^{RI})^T\\ (I^R)^T& S^{RI}& I^{RR}& S^{RR}\\ (S^R)^T& D^{RI}& (S^{RR})^T& D^{RR}\end{array}\right]}{\mathrm{Pf}\left[\begin{array}{cc}0& S^I\\ (S^I)^T& D^{II}\end{array}\right]}}.`$ (2.138)
In the last line we have exploited a Pfaffian identity that holds for antisymmetric matrices $`A`$, $`B`$ of odd ranks and a row vector $`v`$:
$$\mathrm{Pf}\left[\begin{array}{cc}A& \begin{array}{c}v\\ \mathrm{}\\ v\end{array}\\ & \\ v^T\mathrm{}v^T& B\end{array}\right]=\mathrm{Pf}\left[\begin{array}{cc}A& \begin{array}{c}1\\ \mathrm{}\\ 1\end{array}\\ & \\ 1\mathrm{}1& 0\end{array}\right]\mathrm{Pf}\left[\begin{array}{cc}0& v\\ v^T& B\end{array}\right].$$
(2.139)
Likewise previous Subsection, we proceed to evaluate the component functions of the quaternion kernel $`f_{jk}^{\mathrm{even}}(z_j,z_k)`$ in the asymptitic limit where unfolded microscopic variables
$`\sqrt{2N}z_0`$ $``$ $`\zeta _0=0,`$ (2.140)
$`\sqrt{2N}z_{2j1}`$ $``$ $`\zeta _{2j1}=i\mu _j(j=1,\mathrm{},n),`$ (2.141)
$`\sqrt{2N}z_{2j}`$ $``$ $`\zeta _{2j}=i\mu _j(j=1,\mathrm{},n),`$ (2.142)
$`\sqrt{2N}z_{2n+j}`$ $``$ $`\lambda _j(j=1,\mathrm{},p),`$ (2.143)
are kept fixed. We note that all elements of the sub-matrices that appear in the second line of Eq.(2.138) are expressed as (derivatives or integrals of) analytic functions $`\overline{S}(z,z^{})`$ and $`\overline{\mathrm{\Phi }}_j(z)`$ defined in Eqs.(2.68) and (2.69), and $`\mathrm{\Psi }_j(z)`$:
$`D_{jk}^{II}`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \frac{}{z_j}}\overline{S}(z_j,z_k),`$ (2.145)
$`S_{jk}^{RI}`$ $`=`$ $`\overline{S}(x_j,z_k)+{\displaystyle \frac{\overline{\mathrm{\Phi }}_{2\alpha +N}(x_j)}{s_{2\alpha +N}}}\overline{S}(\mathrm{},z_k){\displaystyle \frac{\mathrm{\Psi }_{2\alpha +N}(z_k)}{s_{2\alpha +N}}}\left(2{\displaystyle _{\mathrm{}}^{x_j}}\overline{S}(\mathrm{},x)𝑑x1\right),`$ (2.146)
$`D_{jk}^{RI}`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \frac{}{x_j}}\overline{S}(x_j,z_k),`$ (2.147)
$`S_{jk}^{RR}`$ $`=`$ $`\overline{S}(x_j,x_k)+{\displaystyle \frac{\overline{\mathrm{\Phi }}_{2\alpha +N}(x_j)}{s_{2\alpha +N}}}\overline{S}(\mathrm{},x_k){\displaystyle \frac{\mathrm{\Psi }_{2\alpha +N}(x_k)}{s_{2\alpha +N}}}\left(2{\displaystyle _{\mathrm{}}^{x_j}}\overline{S}(\mathrm{},x)𝑑x1\right),`$ (2.148)
$`D_{jk}^{RR}`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \frac{}{x_j}}\overline{S}(x_j,x_k),`$ (2.149)
$`I_{jk}^{RR}`$ $`=`$ $`2{\displaystyle _{x_j}^{x_k}}\overline{S}(x_j,x)𝑑x\mathrm{sgn}(x_jx_k)`$ (2.151)
$`+{\displaystyle \frac{\overline{\mathrm{\Phi }}_{2\alpha +N}(x_k)}{s_{2\alpha +N}}}\left(2{\displaystyle _{\mathrm{}}^{x_j}}\overline{S}(\mathrm{},x)𝑑x1\right){\displaystyle \frac{\overline{\mathrm{\Phi }}_{2\alpha +N}(x_j)}{s_{2\alpha +N}}}\left(2{\displaystyle _{\mathrm{}}^{x_k}}\overline{S}(\mathrm{},x)𝑑x1\right).`$
Utilizing an identity
$$_0^{\mathrm{}}\mathrm{e}^{u^2/2}H_{2k}(u)𝑑u=2^{2k1/2}\mathrm{\Gamma }(k+\frac{1}{2}),$$
(2.152)
an asymptotic formula
$`{\displaystyle _0^{\mathrm{}}}\mathrm{e}^{u^2/2}H_{2k1}(u)𝑑u`$ $`=`$ $`{\displaystyle \frac{2^{2k1}\mathrm{\Gamma }(k)}{\sqrt{\pi }}}{\displaystyle \underset{\mathrm{}=0}{\overset{k1}{}}}(1)^{\mathrm{}}{\displaystyle \frac{\mathrm{\Gamma }(\mathrm{}+1/2)}{\mathrm{\Gamma }(\mathrm{}+1)}}`$ (2.153)
$`=`$ $`{\displaystyle \frac{2^{2k1}\mathrm{\Gamma }(k)}{\sqrt{\pi }}}\left(\sqrt{\pi }{}_{2}{}^{}F_{1}^{}(1,{\displaystyle \frac{1}{2}};1;1){\displaystyle \underset{\mathrm{}=k}{\overset{\mathrm{}}{}}}(1)^{\mathrm{}}{\displaystyle \frac{\mathrm{\Gamma }(\mathrm{}+1/2)}{\mathrm{\Gamma }(\mathrm{}+1)}}\right)`$ (2.154)
$``$ $`{\displaystyle \frac{2^{2k1}\mathrm{\Gamma }(k)}{\sqrt{\pi }}}\left(\sqrt{{\displaystyle \frac{\pi }{2}}}{\displaystyle \frac{(1)^k}{2\sqrt{k}}}\right)(k1),`$ (2.155)
and Eq.(2.78), we obtain Eqs.(II A) (with $`S`$, $`D`$, $`I`$ in the second places replaced by $`S^{\mathrm{even}}`$, $`D^{\mathrm{even}}`$, $`I^{\mathrm{even}}`$) and
$`𝐒_j^IS^I(z_j){\displaystyle \frac{(1)^{n+N/2}}{\sqrt{2\pi }}}\mathrm{cos}\zeta _j,`$ (2.157)
$`𝐒_j^RS^R(x_j){\displaystyle \frac{(1)^{n+N/2}}{\sqrt{2\pi }}}\mathrm{cos}\lambda _j,`$ (2.158)
$`𝐈_j^R\sqrt{2N}I^R(x_j)(1)^{n+N/2+1}\sqrt{{\displaystyle \frac{2}{\pi }}}\mathrm{sin}\lambda _j.`$ (2.159)
These matrix elements constitute the finite-volume partition function
$`𝒵(\{\mu \})`$ $``$ $`\mathrm{\Xi }_0(\{{\displaystyle \frac{\mu }{\sqrt{2N}}}\})`$ (2.160)
$`=`$ $`\mathrm{const}.{\displaystyle \frac{\mathrm{Pf}\left[\begin{array}{cc}0& 𝐒^I\\ (𝐒^I)^T& 𝐃^{II}\end{array}\right]}{_{j=1}^n\mu _j^3_{j>k}^n(\mu _j^2\mu _k^2)^2}},`$ (2.163)
and the scaled spectral correlation functions
$`\rho _s(\lambda _1,\mathrm{},\lambda _p;\{\mu \})`$ $``$ $`\left({\displaystyle \frac{1}{\sqrt{2N}}}\right)^p\rho ({\displaystyle \frac{\lambda _1}{\sqrt{2N}}},\mathrm{},{\displaystyle \frac{\lambda _p}{\sqrt{2N}}};\{{\displaystyle \frac{\mu }{\sqrt{2N}}}\})`$ (2.164)
$`=`$ $`(1)^{p(p1)/2}{\displaystyle \frac{\mathrm{Pf}\left[\begin{array}{cccc}0& 𝐒^I& 𝐈^R& 𝐒^R\\ (𝐒^I)^T& 𝐃^{II}& (𝐒^{RI})^T& (𝐃^{RI})^T\\ (𝐈^R)^T& 𝐒^{RI}& 𝐈^{RR}& 𝐒^{RR}\\ (𝐒^R)^T& 𝐃^{RI}& (𝐒^{RR})^T& 𝐃^{RR}\end{array}\right]}{\mathrm{Pf}\left[\begin{array}{cc}0& 𝐒^I\\ (𝐒^I)^T& 𝐃^{II}\end{array}\right]}}.`$ (2.171)
It satisfies a sequence
$$\rho _s(\{\lambda \};\mu _1,\mathrm{},\mu _n,\mu _1,\mathrm{},\mu _n,0)\stackrel{\mu _n\mathrm{}}{}\rho _s(\{\lambda \};\mu _1,\mathrm{},\mu _{n1},\mu _1,\mathrm{},\mu _{n1},0)\stackrel{\mu _{n1}\mathrm{}}{}\mathrm{},$$
(2.173)
as each of the masses is decoupled by going to infinity. To illustrate this decoupling, we exhibit in FIG.2 a plot of the spectral density $`\rho _s(\lambda ;\mu ,\mu ,0)`$ ($`p=1,n=1`$).
### C odd $`𝐍_𝐟`$, odd $`𝐍`$
We finally consider the case with $`N_f2n+1`$ flavors, $`\{m\}=(m_1,\mathrm{},m_n,m_1,\mathrm{},m_n,0)`$, and with odd $`N`$. As mentioned in Introduction, this case is pathological because
$$\mathrm{\Xi }_p(x_1,\mathrm{},x_p;\{m\})=\mathrm{\Xi }_p(x_1,\mathrm{},x_p;\{m\})$$
(2.174)
and namely $`\mathrm{\Xi }_0(\{m\})=Z(\{m\})=0`$. The quantities computed below should not be considered as correlation functions of a RME, but merely as multiple integrals defined by Eqs.(2.95)$``$(2.98).
As in previous Subsection, we define variables $`z_j`$ by Eq.(2.101), and denote the integrand in $`\mathrm{\Xi }_p`$ as
$$p(z_0,\mathrm{},z_{2n+N})=\underset{j=0}{\overset{2n+N}{}}\sqrt{w(z_j)}\underset{j>k0}{\overset{2n+N}{}}(z_jz_k)\underset{j>k>2n}{\overset{2n+N}{}}\mathrm{sgn}(z_jz_k).$$
(2.175)
An identity
$$\underset{j>k0}{\overset{2n+N}{}}\mathrm{sgn}(z_jz_k)=\mathrm{Pf}[\mathrm{sgn}(z_kz_j)]_{j,k=0,\mathrm{},2n+N}$$
(2.176)
holds for real $`z_0,z_1,\mathrm{},z_{2n+N}`$. By taking the limit $`z_0<z_1<\mathrm{}<z_{2n}\mathrm{}`$, we find another identity
$$\underset{j>k>2n}{\overset{2n+N}{}}\mathrm{sgn}(z_jz_k)=\mathrm{Pf}[F_{jk}]_{j,k=0,\mathrm{},2n+N},$$
(2.177)
where $`F_{jk}`$ is defined in Eq.(2.108). Substitution of Eq.(2.177) into Eq.(2.175) yields
$$p(z_0,\mathrm{},z_{2n+N})=\underset{j=0}{\overset{2n+N}{}}\sqrt{w(z_j)}\underset{j>k0}{\overset{2n+N}{}}(z_jz_k)\mathrm{Pf}[F_{jk}]_{j,k=0,\mathrm{},2n+N}.$$
(2.178)
The Pfaffian in the above can be represented as a quaternion determinant, due to the following theorem (which is essentially Theorem 1)
Theorem 1’
For odd $`N`$, we can rewrite $`p(z_0,\mathrm{},z_{2n+N})`$ as
$$p(z_0,\mathrm{},z_{2n+N})=\left(\underset{j=0}{\overset{n+(N+1)/21}{}}r_j\right)\mathrm{Tdet}[f_{jk}(z_j,z_k)]_{j,k=0,\mathrm{},2n+N}.$$
(2.179)
The quaternion elements $`f_{jk}(z_j,z_k)`$ are represented as
$$f_{jk}(z_j,z_k)=\left[\begin{array}{cc}S(z_j,z_k)& I(z_j,z_k)\\ D(z_j,z_k)& S(z_k,z_j)\end{array}\right].$$
(2.180)
The functions $`S(z_j,z_k)`$, $`D(z_j,z_k)`$ and $`I(z_j,z_k)`$ are given by
$`S(z_j,z_k)`$ $`=`$ $`{\displaystyle \underset{\mathrm{}=0}{\overset{n+(N+1)/21}{}}}{\displaystyle \frac{\mathrm{\Phi }_2\mathrm{}(z_j)\mathrm{\Psi }_{2\mathrm{}+1}(z_k)\mathrm{\Phi }_{2\mathrm{}+1}(z_j)\mathrm{\Psi }_2\mathrm{}(z_k)}{r_{\mathrm{}}}},`$ (2.181)
$`D(z_j,z_k)`$ $`=`$ $`{\displaystyle \underset{\mathrm{}=0}{\overset{n+(N+1)/21}{}}}{\displaystyle \frac{\mathrm{\Psi }_2\mathrm{}(z_j)\mathrm{\Psi }_{2\mathrm{}+1}(z_k)\mathrm{\Psi }_{2\mathrm{}+1}(z_j)\mathrm{\Psi }_2\mathrm{}(z_k)}{r_{\mathrm{}}}},`$ (2.182)
$`I(z_j,z_k)`$ $`=`$ $`{\displaystyle \underset{\mathrm{}=0}{\overset{n+(N+1)/21}{}}}{\displaystyle \frac{\mathrm{\Phi }_2\mathrm{}(z_j)\mathrm{\Phi }_{2\mathrm{}+1}(z_k)\mathrm{\Phi }_{2\mathrm{}+1}(z_j)\mathrm{\Phi }_2\mathrm{}(z_k)}{r_{\mathrm{}}}}+F_{jk}.`$ (2.183)
As before, the above quaternion elements satisfy the conditions imposed on $`f_{jk}(z_j,z_k)`$ in Theorem 3, so that we can inductively write
$$\mathrm{\Xi }_p(z_0,\mathrm{},z_{2n+p})=\frac{_{j=0}^{p+(N+1)/21}r_j}{_{j=0}^{2n}\sqrt{w(z_j)}_{j>k0}^{2n}(z_jz_k)}\mathrm{Tdet}[f_{jk}(z_j,z_k)]_{j,k=0,\mathrm{},2n+p}.$$
(2.184)
In order to circumvent the vanishing of the partition function, we regard $`z_0,z_1,\mathrm{},z_{2n}`$ as generic variables, and define the ‘$`p`$-level correlation function’
$$\rho (x_1,\mathrm{},x_p;\{m\})\frac{\mathrm{\Xi }_p(z_0,\mathrm{},z_{2n+p})}{\mathrm{\Xi }_0(z_0,\mathrm{},z_{2n})}=\frac{\mathrm{Tdet}[f_{jk}(z_j,z_k)]_{j,k=0,\mathrm{},2n+p}}{\mathrm{Tdet}[f_{jk}(z_j,z_k)]_{j,k=0,\mathrm{},2n}}.$$
(2.185)
We introduce a set of notations
$`S_{jk}^{II}=S(z_j,z_k)(j,k=0,\mathrm{},2n),`$ (2.186)
$`S_{jk}^{IR}=S(z_j,z_{2n+k})(j=0,\mathrm{},2n;k=1,\mathrm{},p),`$ (2.187)
$`S_{jk}^{RI}=S(z_{2n+j},z_k)(j=1,\mathrm{},p;k=0,\mathrm{},2n),`$ (2.188)
$`S_{jk}^{RR}=S(z_{2n+j},z_{2n+k})(j,k=1,\mathrm{},p)`$ (2.189)
and similarly for $`D`$ and $`I`$. A simplification
$`S_{jk}^{II}=S^{II}(\mathrm{},z_k)S_k^I,`$ (2.190)
$`I_{jk}^{IR}=I_{kj}^{RI}=I^{IR}(\mathrm{},z_{2n+k})I_k^R,`$ (2.191)
$`S_{jk}^{IR}=S^{IR}(\mathrm{},z_{2n+k})S_k^R,`$ (2.192)
results again from the definition (2.182). Using Dyson’s identity (A.16) we can rewrite the ‘correlation functions’ as Eq.(2.138).
Likewise two previous Subsections, we proceed to evaluate the component functions of the quaternion kernel $`f_{jk}(z_j,z_k)`$ in the asymptitic limit where unfolded microscopic variables (2.142) are kept fixed. We note that all elements (other than $`S^I`$, $`S^R`$, and $`I^R`$) of the sub-matrices that appear in the second line of Eq.(2.138) are expressed as (derivatives or integrals of) an analytic function $`\overline{S}(z,z^{})`$ defined in Eq.(2.68) (with $`N`$ replaced by $`N+1`$), with real and/or imaginary arguments:
$`D_{jk}^{II}`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \frac{}{z_j}}\overline{S}(z_j,z_k),`$ (2.194)
$`S_{jk}^{RI}`$ $`=`$ $`\overline{S}(x_j,z_k),`$ (2.195)
$`D_{jk}^{RI}`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \frac{}{x_j}}\overline{S}(x_j,z_k),`$ (2.196)
$`S_{jk}^{RR}`$ $`=`$ $`\overline{S}(x_j,x_k),`$ (2.197)
$`D_{jk}^{RR}`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \frac{}{x_j}}\overline{S}(x_j,x_k),`$ (2.198)
$`I_{jk}^{RR}`$ $`=`$ $`2{\displaystyle _{x_j}^{x_k}}\overline{S}(x_j,x)𝑑x\mathrm{sgn}(x_jx_k).`$ (2.199)
We obtain Eqs.(II A) and
$`𝐒_j^IS^I(z_j){\displaystyle \frac{(1)^{n+(N+1)/2}}{\sqrt{2\pi }}}\mathrm{sin}\zeta _j,`$ (2.201)
$`𝐒_j^RS^R(x_j){\displaystyle \frac{(1)^{n+(N+1)/2}}{\sqrt{2\pi }}}\mathrm{sin}\lambda _j,`$ (2.202)
$`𝐈_j^R\sqrt{2N}I^R(x_j)(1)^{n+(N+1)/2}\sqrt{{\displaystyle \frac{2}{\pi }}}\mathrm{cos}\lambda _j.`$ (2.203)
Note the phase shifts of the trigonometric functions between Eq.(II C) and its even-$`N`$ counterpart (II B). These matrix elements constitute the ‘finite-volume partition function’
$`𝒵(\{\mu \})`$ $``$ $`\mathrm{\Xi }_0(\{{\displaystyle \frac{\mu }{\sqrt{2N}}}\})`$ (2.204)
$`=`$ $`\mathrm{const}.{\displaystyle \frac{\mathrm{Pf}\left[\begin{array}{cc}0& 𝐒^I\\ (𝐒^I)^T& 𝐃^{II}\end{array}\right]}{_{j=1}^n\mu _j^3_{j>k}^n(\mu _j^2\mu _k^2)^2}},`$ (2.207)
and the ‘scaled correlation functions’
$`\rho _s(\lambda _1,\mathrm{},\lambda _p;\{\mu \})`$ $``$ $`\left({\displaystyle \frac{1}{\sqrt{2N}}}\right)^p\rho ({\displaystyle \frac{\lambda _1}{\sqrt{2N}}},\mathrm{},{\displaystyle \frac{\lambda _p}{\sqrt{2N}}};\{{\displaystyle \frac{\mu }{\sqrt{2N}}}\})`$ (2.208)
$`=`$ $`(1)^{p(p1)/2}{\displaystyle \frac{\mathrm{Pf}\left[\begin{array}{cccc}0& 𝐒^I& 𝐈^R& 𝐒^R\\ (𝐒^I)^T& 𝐃^{II}& (𝐒^{RI})^T& (𝐃^{RI})^T\\ (𝐈^R)^T& 𝐒^{RI}& 𝐈^{RR}& 𝐒^{RR}\\ (𝐒^R)^T& 𝐃^{RI}& (𝐒^{RR})^T& 𝐃^{RR}\end{array}\right]}{\mathrm{Pf}\left[\begin{array}{cc}0& 𝐒^I\\ (𝐒^I)^T& 𝐃^{II}\end{array}\right]}}.`$ (2.215)
We again remind the reader that under the identification (2.142), the above ‘finite-volume partition function’ vanishes and the ‘correlation function’ diverges as its Pfaffian denominator vanishes.
## III symplectic ensemble
For $`\beta =4`$, we concentrate on even flavor cases for a technical reason, and treat the following two cases separately:
$`𝐀:\{m\}=(m_1,m_1,\mathrm{},m_\alpha ,m_\alpha ,m_1,m_1,\mathrm{},m_\alpha ,m_\alpha ),`$
$`𝐁:\{m\}=(m_1,m_1,\mathrm{},m_\alpha ,m_\alpha ,m_1,m_1,\mathrm{},m_\alpha ,m_\alpha ,0,0).`$
The $`\beta =4`$ case with an odd number of fermions will not be treated in this Article.
### A $`𝐍_𝐟`$ = 0 mod 4
We first consider the case with $`N_f4\alpha `$ flavors and $`\{m\}`$$`=`$$`(m_1`$,$`m_1`$,$`\mathrm{}`$,$`m_\alpha `$,$`m_\alpha `$,$`m_1`$,$`m_1`$,$`\mathrm{}`$,$`m_\alpha `$,$`m_\alpha )`$. We express the partition function (1.19) of the RME in terms of eigenvalues $`\{x_j\}`$ of $`H`$:
$`Z(\{m\})={\displaystyle \frac{1}{N!}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}\mathrm{}{\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle \underset{j=1}{\overset{N}{}}}dx_j{\displaystyle \underset{j=1}{\overset{N}{}}}\left(\mathrm{e}^{2x_j^2}{\displaystyle \underset{k=1}{\overset{\alpha }{}}}(x_j^2+m_k^2)\right){\displaystyle \underset{j>k}{\overset{N}{}}}(x_jx_k)^4.`$ (3.1)
The $`p`$-level correlation function of the matrix $`H`$ is defined as
$`\rho (x_1,\mathrm{},x_p;\{m\})`$ $`=`$ $`{\displaystyle \underset{j=1}{\overset{p}{}}}\mathrm{tr}\delta (x_jH)`$ (3.2)
$`=`$ $`{\displaystyle \frac{\mathrm{\Xi }_p(x_1,\mathrm{},x_p;\{m\})}{\mathrm{\Xi }_0(\{m\})}},`$ (3.3)
$`\mathrm{\Xi }_p(x_1,\mathrm{},x_p;\{m\})`$ $`=`$ $`{\displaystyle \frac{1}{(Np)!}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}\mathrm{}{\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle \underset{j=p+1}{\overset{N}{}}}dx_j{\displaystyle \underset{j=1}{\overset{N}{}}}\left(\mathrm{e}^{2x_j^2}{\displaystyle \underset{k=1}{\overset{\alpha }{}}}(x_j^2+m_k^2)\right){\displaystyle \underset{j>k}{\overset{N}{}}}(x_jx_k)^4`$ (3.4)
$`(\mathrm{\Xi }_0=Z)`$. We define new variables $`z_j`$ as
$$\begin{array}{c}z_{2j1}=im_j\hfill \\ z_{2j}=im_j\hfill \end{array}\}(j=1,\mathrm{},\alpha ).$$
(3.5)
Then the multiple integral (3.4) is expressed as
$`\mathrm{\Xi }_p(x_1,\mathrm{},x_p;\{z\})={\displaystyle \frac{1}{_{j=1}^{2\alpha }w(z_j)_{j>k}^{2\alpha }(z_jz_k)}}`$ (3.6)
$`\times {\displaystyle \frac{1}{(Np)!}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}\mathrm{}{\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle \underset{j=2\alpha +p+1}{\overset{2\alpha +N}{}}}dx_j{\displaystyle \underset{j=1}{\overset{2\alpha }{}}}w(z_j){\displaystyle \underset{j=1}{\overset{N}{}}}w(x_j)^2{\displaystyle \underset{j>k}{\overset{2\alpha }{}}}(z_jz_k){\displaystyle \underset{j=1}{\overset{N}{}}}{\displaystyle \underset{k=1}{\overset{2\alpha }{}}}(x_jz_k)^2{\displaystyle \underset{j>k}{\overset{N}{}}}(x_jx_k)^4,`$ (3.7)
where $`w(z)=\mathrm{e}^{z^2}.`$
Let us denote the integrand in Eq.(3.7) as
$$p(z_1,\mathrm{},z_{2\alpha };x_1,\mathrm{},x_N)=\underset{j=1}{\overset{2\alpha }{}}w(z_j)\underset{j=1}{\overset{N}{}}w(x_j)^2\underset{j>k}{\overset{2\alpha }{}}(z_jz_k)\underset{j=1}{\overset{N}{}}\underset{k=1}{\overset{2\alpha }{}}(x_jz_k)^2\underset{j>k}{\overset{N}{}}(x_jx_k)^4.$$
(3.8)
The above expressions can be represented as a quaternion determinant. In doing so, we introduce monic skew-orthogonal polynomials $`Q_j(z)=z^j+\mathrm{}`$ that satisfy
$`Q_{2j},Q_{2k+1}_Q=Q_{2k+1},Q_{2j}_Q=q_j\delta _{jk},`$ (3.9)
$`Q_{2j},Q_{2k}_Q=Q_{2j+1},Q_{2k+1}_Q=0,`$ (3.10)
where
$$f,g_Q=_{\mathrm{}}^{\mathrm{}}𝑑zw(z)^2\left(f(z)g^{}(z)f^{}(z)g(z)\right).$$
(3.11)
Explicit forms for the skew-orthogonal polynomials and their norms associated with the Gaussian weight $`w(z)`$ are known in terms of the Hermite polynomials :
$`Q_{2j}(z)`$ $`=`$ $`{\displaystyle \frac{1}{2^{3j+1/2}}}\mathrm{e}^{z^2}{\displaystyle _{\mathrm{}}^z}\mathrm{e}^{z^{}^2}H_{2j+1}(\sqrt{2}z^{})𝑑z^{},`$ (3.12)
$`Q_{2j+1}(z)`$ $`=`$ $`{\displaystyle \frac{1}{2^{3j+3/2}}}H_{2j+1}(\sqrt{2}z),`$ (3.13)
$`q_j`$ $`=`$ $`2^{4j1/2}\sqrt{\pi }(2j+1)!.`$ (3.14)
It can be readily seen that (cf. Eq.(26) of Ref.) $`p(z_1,\mathrm{},z_{2\alpha };x_1,\mathrm{},x_N)`$ is represented as a Pfaffian:
$`p(z_1,\mathrm{},z_{2\alpha };x_1,\mathrm{},x_N)=det\left[\begin{array}{cccc}\mathrm{\Psi }_0(z_1)& \mathrm{\Psi }_1(z_1)& \mathrm{}& \mathrm{\Psi }_{2N+2\alpha 1}(z_1)\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ \mathrm{\Psi }_0(z_{2\alpha })& \mathrm{\Psi }_1(z_{2\alpha })& \mathrm{}& \mathrm{\Psi }_{2N+2\alpha 1}(z_{2\alpha })\\ \mathrm{\Psi }_0^{}(x_1)& \mathrm{\Psi }_1^{}(x_1)& \mathrm{}& \mathrm{\Psi }_{2N+2\alpha 1}^{}(x_1)\\ \mathrm{\Psi }_0(x_1)& \mathrm{\Psi }_1(x_1)& \mathrm{}& \mathrm{\Psi }_{2N+2\alpha 1}(x_1)\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ \mathrm{\Psi }_0^{}(x_N)& \mathrm{\Psi }_1^{}(x_N)& \mathrm{}& \mathrm{\Psi }_{2N+2\alpha 1}^{}(x_N)\\ \mathrm{\Psi }_0(x_N)& \mathrm{\Psi }_1(x_N)& \mathrm{}& \mathrm{\Psi }_{2N+2\alpha 1}(x_N)\end{array}\right]`$ (3.23)
$`=\left({\displaystyle \underset{j=0}{\overset{\alpha +N1}{}}}q_j\right)\mathrm{Pf}\left[\begin{array}{cc}\left[\begin{array}{cc}I(z_{2j},z_{2k})& I(z_{2j},z_{2k1})\\ I(z_{2j1},z_{2k})& I(z_{2j1},z_{2k1})\end{array}\right]_{j,k=1,\mathrm{},\alpha }\hfill & \left[\begin{array}{cc}S(z_{2j},x_k)& I(z_{2j},x_k)\\ S(z_{2j1},x_k)& I(z_{2j1},x_k)\end{array}\right]_{j=1,\mathrm{},\alpha ;k=1,\mathrm{},N}\hfill \\ \left[\begin{array}{cc}S(z_{2k},x_j)& S(z_{2k1},x_j)\\ I(x_j,z_{2k})& I(x_j,z_{2k1})\end{array}\right]_{j=1,\mathrm{},N;k=1,\mathrm{},\alpha }\hfill & \left[\begin{array}{cc}D(x_j,x_k)& S(x_k,x_j)\\ S(x_j,x_k)& I(x_j,x_k)\end{array}\right]_{j,k=1,\mathrm{},N}\hfill \end{array}\right].`$ (3.34)
Here
$$\mathrm{\Psi }_j(z)=w(z)Q_j(z),$$
(3.36)
and the functions $`S(x,y)`$, $`D(x,y)`$, $`I(x,y)`$ are given by
$`S(x,y)={\displaystyle \underset{j=0}{\overset{\alpha +N1}{}}}{\displaystyle \frac{\mathrm{\Psi }_{2j}(x)\mathrm{\Psi }_{2j+1}^{}(y)\mathrm{\Psi }_{2j+1}(x)\mathrm{\Psi }_{2j}^{}(y)}{q_j}},`$ (3.37)
$`D(x,y)={\displaystyle \underset{j=0}{\overset{\alpha +N1}{}}}{\displaystyle \frac{\mathrm{\Psi }_{2j}^{}(x)\mathrm{\Psi }_{2j+1}^{}(y)\mathrm{\Psi }_{2j+1}^{}(x)\mathrm{\Psi }_{2j}^{}(y)}{q_j}}={\displaystyle \frac{}{x}}S(x,y),`$ (3.38)
$`I(x,y)={\displaystyle \underset{j=0}{\overset{\alpha +N1}{}}}{\displaystyle \frac{\mathrm{\Psi }_{2j}(x)\mathrm{\Psi }_{2j+1}(y)\mathrm{\Psi }_{2j+1}(x)\mathrm{\Psi }_{2j}(y)}{q_j}}={\displaystyle _x^y}S(x,z)𝑑z.`$ (3.39)
Using Dyson’s identity (A.16), we can convert the Pfaffian into a quaternion determinant:
$$p(z_1,\mathrm{},z_{2\alpha };x_1,\mathrm{},x_N)=\left(\underset{j=0}{\overset{\alpha +N1}{}}q_j\right)\mathrm{Tdet}\left[\begin{array}{cc}\left[g(z_{2j1},z_{2j};z_{2k1},z_{2k})\right]_{j,k=1,\mathrm{},\alpha }\hfill & \left[h(z_{2j1},z_{2j};x_k)\right]_{j=1,\mathrm{},\alpha ;k=1,\mathrm{},N}\hfill \\ [\widehat{h}(z_{2k1},z_{2k};x_j)]_{j=1,\mathrm{},N;k=1,\mathrm{},\alpha }\hfill & \left[f(x_j,x_k)\right]_{j,k=1,\mathrm{},N}\hfill \end{array}\right],$$
(3.40)
where
$`g(z_{2j1},z_{2j};z_{2k1},z_{2k})=\left[\begin{array}{cc}I(z_{2j1},z_{2k})\hfill & I(z_{2j1},z_{2k1})\hfill \\ I(z_{2j},z_{2k})\hfill & I(z_{2j},z_{2k1})\hfill \end{array}\right],`$ (3.43)
$`h(z_{2j1},z_{2j};x_k)=\left[\begin{array}{cc}S(z_{2j1},x_k)\hfill & I(z_{2j1},x_k)\hfill \\ S(z_{2j},x_k)\hfill & I(z_{2j},x_k)\hfill \end{array}\right],`$ (3.46)
$`f(x_j,x_k)=\left[\begin{array}{cc}S(x_j,x_k)\hfill & I(x_j,x_k)\hfill \\ D(x_j,x_k)\hfill & S(x_k,x_j)\hfill \end{array}\right],`$ (3.49)
and $`\widehat{h}`$ stands for the dual of $`h`$ (see Appendix A). The skew-orthogonality relations (3.9) lead to
$`{\displaystyle _{\mathrm{}}^{\mathrm{}}}f(x^{},x)f(x,x^{\prime \prime })𝑑x=f(x^{},x^{\prime \prime }),`$ (3.50)
$`{\displaystyle _{\mathrm{}}^{\mathrm{}}}h(z,z^{};x)f(x,x^{})𝑑x=h(z,z^{};x^{}),`$ (3.51)
$`{\displaystyle _{\mathrm{}}^{\mathrm{}}}h(z,z^{},x)\widehat{h}(w,w^{};x)𝑑x=g(z,z^{};w,w^{}),`$ (3.52)
which mean that the condition of Theorem 3 is satisfied. Therefore we can inductively write
$`\mathrm{\Xi }_p(x_1,\mathrm{},x_p;z_1,\mathrm{},z_{2\alpha })=`$ (3.53)
$`{\displaystyle \frac{_{j=0}^{\alpha +N1}q_j}{_{j=1}^{2\alpha }w(z_j)_{j>k}^{2\alpha }(z_jz_k)}}\mathrm{Tdet}\left[\begin{array}{cc}[g(z_{2j1},z_{2j};z_{2k1},z_{2k})]_{j,k=1,\mathrm{},\alpha }\hfill & [h(z_{2j1},z_{2j};x_k)]_{j=1,\mathrm{},\alpha ;k=1,\mathrm{},p}\hfill \\ [\widehat{h}(z_{2k1},z_{2k};x_j)]_{j=1,\mathrm{},p;k=1,\mathrm{},\alpha }\hfill & [f(x_j,x_k)]_{j,k=1,\mathrm{},p}\hfill \end{array}\right].`$ (3.56)
Then the $`p`$-level correlation function (3.3) is written as
$`\rho (x_1,\mathrm{},x_p;\{m\})`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Xi }_p(x_1,\mathrm{},x_p;z_1,\mathrm{},z_{2n})}{\mathrm{\Xi }_0(z_1,\mathrm{},z_{2n})}}`$ (3.57)
$`=`$ $`{\displaystyle \frac{\mathrm{Tdet}\left[\begin{array}{cc}[g(z_{2j1},z_{2j};z_{2k1},z_{2k})]_{j,k=1,\mathrm{},\alpha }\hfill & [h(z_{2j1},z_{2j};x_k)]_{j=1,\mathrm{},\alpha ;k=1,\mathrm{},p}\hfill \\ [\widehat{h}(z_{2k1},z_{2k};x_j)]_{j=1,\mathrm{},p;k=1,\mathrm{},\alpha }\hfill & [f(x_j,x_k)]_{j,k=1,\mathrm{},p}\hfill \end{array}\right]}{\mathrm{Tdet}[g(z_{2j1},z_{2j};z_{2k1},z_{2k})]_{j,k=1,\mathrm{},\alpha }}}.`$ (3.60)
We introduce a set of notations
$`S_{jk}^{II}=S(z_j,z_k)(j,k=1,\mathrm{},2\alpha ),`$ (3.61)
$`S_{jk}^{IR}=S(z_j,x_k)(j=1,\mathrm{},2\alpha ;k=1,\mathrm{},p),`$ (3.62)
$`S_{jk}^{RR}=S(x_j,x_k)(j,k=1,\mathrm{},p),`$ (3.63)
and similarly for $`D`$ and $`I`$. Using Dyson’s identity (A.16) again, we can convert the quaternion determinant back to a Pfaffian, so that the correlation functions read
$$\rho (x_1,\mathrm{},x_p;\{m\})=(1)^{p(p1)/2}\frac{\mathrm{Pf}\left[\begin{array}{ccc}I^{II}& I^{IR}& S^{IR}\\ (I^{IR})^T& I^{RR}& S^{RR}\\ (S^{IR})^T& (S^{RR})^T& D^{RR}\end{array}\right]}{\mathrm{Pf}[I^{II}]}.$$
(3.64)
Now we proceed to evaluate the component functions of the quaternion kernel, which are expressed as (derivatives or integrals of) the analytic function $`S(x,y)`$ with real and/or imaginary arguments, in the asymptotic limit where unfolded microscopic variables
$`\sqrt{2N}z_{2j1}`$ $``$ $`\zeta _{2j1}=i\mu _j(j=1,\mathrm{},\alpha ),`$ (3.65)
$`\sqrt{2N}z_{2j}`$ $``$ $`\zeta _{2j}=i\mu _j(j=1,\mathrm{},\alpha ),`$ (3.66)
$`\sqrt{2N}x_j`$ $``$ $`\lambda _j(j=1,\mathrm{},p),`$ (3.67)
are kept fixed. We note that $`S(x,y)`$ has a compact expression
$`S(x,y)`$ $`=`$ $`{\displaystyle \frac{\mathrm{e}^{x^2y^2}}{2^{2\alpha +2N+1}\sqrt{\pi }\mathrm{\Gamma }(2\alpha +2N)}}{\displaystyle \frac{H_{2\alpha +2N}(\sqrt{2}x)H_{2\alpha +2N1}(\sqrt{2}y)H_{2\alpha +2N1}(\sqrt{2}x)H_{2\alpha +2N}(\sqrt{2}y)}{xy}}`$ (3.68)
$`+`$ $`{\displaystyle \frac{\mathrm{e}^{y^2}}{2^{2\alpha +2N}\sqrt{\pi }\mathrm{\Gamma }(2\alpha +2N)}}H_{2\alpha +2N}(\sqrt{2}y){\displaystyle _{\mathrm{}}^x}\mathrm{e}^{u^2}H_{2\alpha +2N1}(\sqrt{2}u)𝑑u.`$ (3.69)
In the second line of Eq.(3.69), we have singled out the unitary scalar kernel and applied to it the Christoffel-Darboux formula. Substituting asymptotic formulas for the Hermite polynomials (2.78), one can show that $`S(x,y)`$ approaches the sine kernel :
$$\frac{1}{\sqrt{2N}}S(\frac{\zeta }{\sqrt{2N}},\frac{\zeta ^{}}{\sqrt{2N}})\frac{\mathrm{sin}2(\zeta \zeta ^{})}{2\pi (\zeta \zeta ^{})}=K(2(\zeta \zeta ^{})).$$
(3.70)
The Pfaffian elements in the the asymptotic limit, after taking into account an unfolding by the factor $`\sqrt{2N}`$, are then expressed in terms of $`K(\zeta )`$:
$`𝐈_{jk}^{II}`$ $``$ $`I(z_j,z_k){\displaystyle _0^{\zeta _j\zeta _k}}K(2\zeta )𝑑\zeta ,`$ (3.72)
$`𝐒_{jk}^{IR}`$ $``$ $`{\displaystyle \frac{1}{\sqrt{2N}}}S(z_j,x_k)K(2(\zeta _j\lambda _k)),`$ (3.73)
$`𝐈_{jk}^{IR}`$ $``$ $`I(z_j,x_k){\displaystyle _0^{\zeta _j\lambda _k}}K(2\zeta )𝑑\zeta ,`$ (3.74)
$`𝐒_{jk}^{RR}`$ $``$ $`{\displaystyle \frac{1}{\sqrt{2N}}}S(x_j,x_k)K(2(\lambda _j\lambda _k)),`$ (3.75)
$`𝐃_{jk}^{RR}`$ $``$ $`{\displaystyle \frac{1}{2N}}D(x_j,x_k)2K^{}(2(\lambda _j\lambda _k)),`$ (3.76)
$`𝐈_{jk}^{RR}`$ $``$ $`I(x_j,x_k){\displaystyle _0^{\lambda _j\lambda _k}}K(2\lambda )𝑑\lambda .`$ (3.77)
These matrix elements constitute the finite-volume partition function
$`𝒵(\{m\})`$ $``$ $`\mathrm{\Xi }_0(\{{\displaystyle \frac{\mu }{\sqrt{2N}}}\})`$ (3.78)
$`=`$ $`\mathrm{const}.{\displaystyle \frac{\mathrm{Pf}[𝐈^{II}]}{_{j=1}^\alpha \mu _j_{j>k}^\alpha (\mu _j^2\mu _k^2)^2}},`$ (3.79)
and the scaled spectral correlation functions
$`\rho _s(\lambda _1,\mathrm{},\lambda _p;\{\mu \})`$ $``$ $`\left({\displaystyle \frac{1}{\sqrt{2N}}}\right)^p\rho ({\displaystyle \frac{\lambda _1}{\sqrt{2N}}},\mathrm{},{\displaystyle \frac{\lambda _p}{\sqrt{2N}}};\{{\displaystyle \frac{\mu }{\sqrt{2N}}}\})`$ (3.80)
$`=`$ $`(1)^{p(p1)/2}{\displaystyle \frac{\mathrm{Pf}\left[\begin{array}{ccc}𝐈^{II}& 𝐈^{IR}& 𝐒^{IR}\\ (𝐈^{IR})^T& 𝐈^{RR}& 𝐒^{RR}\\ (𝐒^{IR})^T& (𝐒^{RR})^T& 𝐃^{RR}\end{array}\right]}{\mathrm{Pf}\left[𝐈^{II}\right]}}.`$ (3.84)
In the quenched limit $`\mu _1,\mathrm{},\mu _\alpha \mathrm{}`$ when the ratio of two Pfaffians is replaced by a minor $`\mathrm{Pf}\left[\genfrac{}{}{0pt}{}{𝐈^{RR}𝐒^{RR}}{(𝐒^{RR})^T𝐃^{RR}}\right]`$, the correlation functions approach those of the Gaussian symplectic ensemble . By the same token, it satisfies a sequence
$`\rho _s(\{\lambda \};\mu _1,\mu _1,\mathrm{},\mu _n,\mu _n,\mu _1,\mu _1,\mathrm{},\mu _n,\mu _n)\stackrel{\mu _n\mathrm{}}{}`$ (3.85)
$`\rho _s(\{\lambda \};\mu _1,\mu _1,\mathrm{},\mu _{n1},\mu _{n1},\mu _1,\mu _1,\mathrm{},\mu _{n1},\mu _{n1})\stackrel{\mu _{n1}\mathrm{}}{}\mathrm{},`$ (3.86)
as each of the masses is decoupled by going to infinity. To illustrate this decoupling, we exhibit in FIG.3 a plot of the spectral density $`\rho _s(\lambda ;\mu ,\mu ,\mu ,\mu )`$ ($`p=1,\alpha =1`$).
### B $`𝐍_𝐟`$ = 2 mod 4
Next we consider the case with $`N_f4\alpha +2`$ flavors and $`\{m\}=(m_1,m_1,\mathrm{},m_\alpha ,m_\alpha ,m_1,m_1,\mathrm{}`$, $`m_\alpha ,m_\alpha ,0,0)`$. We express the partition function (1.19) of the RME in terms of eigenvalues $`\{x_j\}`$ of $`H`$
$`Z(\{m\})={\displaystyle \frac{1}{N!}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}\mathrm{}{\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle \underset{j=1}{\overset{N}{}}}dx_j{\displaystyle \underset{j=1}{\overset{N}{}}}\left(\mathrm{e}^{2x_j^2}x_j^2{\displaystyle \underset{k=1}{\overset{\alpha }{}}}(x_j^2+m_k^2)\right){\displaystyle \underset{j>k}{\overset{N}{}}}(x_jx_k)^4.`$ (3.87)
The $`p`$-level correlation function of the matrix $`H`$ is defined as
$`\rho (x_1,\mathrm{},x_p;\{m\})`$ $`=`$ $`{\displaystyle \underset{j=1}{\overset{p}{}}}\mathrm{tr}\delta (x_jH)`$ (3.88)
$`=`$ $`{\displaystyle \frac{\mathrm{\Xi }_p(x_1,\mathrm{},x_p;\{m\})}{\mathrm{\Xi }_0(\{m\})}},`$ (3.89)
$`\mathrm{\Xi }_p(x_1,\mathrm{},x_p;\{m\})`$ $`=`$ $`{\displaystyle \frac{1}{(Np)!}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}\mathrm{}{\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle \underset{j=p+1}{\overset{N}{}}}dx_j{\displaystyle \underset{j=1}{\overset{N}{}}}\left(\mathrm{e}^{2x_j^2}x_j^2{\displaystyle \underset{k=1}{\overset{\alpha }{}}}(x_j^2+m_k^2)\right){\displaystyle \underset{j>k}{\overset{N}{}}}(x_jx_k)^4`$ (3.90)
$`(\mathrm{\Xi }_0=Z)`$. We define new variables $`z_j`$ as
$`z_0=0,`$ (3.91)
$`\begin{array}{c}z_{2j1}=im_j\hfill \\ z_{2j}=im_j\hfill \end{array}\}(j=1,\mathrm{},\alpha ).`$ (3.94)
Then the multiple integral (3.90) is expressed as
$`\mathrm{\Xi }_p(x_1,\mathrm{},x_p;\{z\})={\displaystyle \frac{1}{_{j=0}^{2\alpha }w(z_j)_{j>k0}^{2\alpha }(z_jz_k)}}`$ (3.95)
$`\times {\displaystyle \frac{1}{(Np)!}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}\mathrm{}{\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle \underset{j=2\alpha +p+1}{\overset{2\alpha +N}{}}}dx_j{\displaystyle \underset{j=0}{\overset{2\alpha }{}}}w(z_j){\displaystyle \underset{j=1}{\overset{N}{}}}w(x_j)^2{\displaystyle \underset{j>k0}{\overset{2\alpha }{}}}(z_jz_k){\displaystyle \underset{j=1}{\overset{N}{}}}{\displaystyle \underset{k=0}{\overset{2\alpha }{}}}(x_jz_k)^2{\displaystyle \underset{j>k}{\overset{N}{}}}(x_jx_k)^4,`$ (3.96)
where $`w(z)=\mathrm{e}^{z^2}`$.
Let us denote the integrand in Eq.(3.96) as
$$p(z_0,\mathrm{},z_{2\alpha };x_1,\mathrm{},x_N)=\underset{j=0}{\overset{2\alpha }{}}w(z_j)\underset{j=1}{\overset{N}{}}w(x_j)^2\underset{j>k0}{\overset{2\alpha }{}}(z_jz_k)\underset{j=1}{\overset{N}{}}\underset{k=0}{\overset{2\alpha }{}}(x_jz_k)^2\underset{j>k}{\overset{N}{}}(x_jx_k)^4.$$
(3.97)
It can be readily seen that $`p(z_0,\mathrm{},z_{2\alpha };x_1,\mathrm{},x_N)`$ is represented as a Pfaffian:
$`p(z_0,\mathrm{},z_{2\alpha };x_1,\mathrm{},x_N)=det\left|\begin{array}{cccc}\mathrm{\Psi }_0(z_0)& \mathrm{\Psi }_1(z_0)& \mathrm{}& \mathrm{\Psi }_{2N+2\alpha }(z_0)\\ \mathrm{\Psi }_0(z_1)& \mathrm{\Psi }_1(z_1)& \mathrm{}& \mathrm{\Psi }_{2N+2\alpha }(z_1)\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ \mathrm{\Psi }_0(z_{2\alpha })& \mathrm{\Psi }_1(z_{2\alpha })& \mathrm{}& \mathrm{\Psi }_{2N+2\alpha }(z_{2\alpha })\\ \mathrm{\Psi }_0^{}(x_1)& \mathrm{\Psi }_1^{}(x_1)& \mathrm{}& \mathrm{\Psi }_{2N+2\alpha }^{}(x_1)\\ \mathrm{\Psi }_0(x_1)& \mathrm{\Psi }_1(x_1)& \mathrm{}& \mathrm{\Psi }_{2N+2\alpha }(x_1)\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ \mathrm{\Psi }_0^{}(x_N)& \mathrm{\Psi }_1^{}(x_N)& \mathrm{}& \mathrm{\Psi }_{2N+2\alpha }^{}(x_N)\\ \mathrm{\Psi }_0(x_N)& \mathrm{\Psi }_1(x_N)& \mathrm{}& \mathrm{\Psi }_{2N+2\alpha }(x_N)\end{array}\right|`$ (3.107)
$`=`$ $`\left({\displaystyle \underset{j=0}{\overset{\alpha +N1}{}}}q_j\right)\mathrm{Pf}\left[\begin{array}{ccc}[G_{jk}]_{j,k=1,\mathrm{},\alpha }& [A_j]_{j=1,\mathrm{},\alpha }& [H_{jk}]_{j=1,\mathrm{},\alpha ;k=1,\mathrm{},N}\\ [\widehat{A}_k]_{k=1,\mathrm{},\alpha }& \mathrm{\Omega }& [B_k]_{k=1,\mathrm{},N}\\ [\widehat{H}_{kj}]_{j=1,\mathrm{},N;k=1,\mathrm{},\alpha }& [\widehat{B}_j]_{j=1,\mathrm{},N}& [F_{jk}]_{j,k=1,\mathrm{},N}\end{array}\right],`$ (3.111)
where
$`G_{jk}`$ $`=`$ $`\left[\begin{array}{cc}\stackrel{~}{I}(z_{2j},z_{2k})& \stackrel{~}{I}(z_{2j},z_{2k1})\\ \stackrel{~}{I}(z_{2j1},z_{2k})& \stackrel{~}{I}(z_{2j1},z_{2k1})\end{array}\right],`$ (3.114)
$`H_{jk}`$ $`=`$ $`\left[\begin{array}{cc}\stackrel{~}{S}(z_{2j},x_k)& \stackrel{~}{I}(z_{2j},x_k)\\ \stackrel{~}{S}(z_{2j1},x_k)& \stackrel{~}{I}(z_{2j1},x_k)\end{array}\right],`$ (3.117)
$`F_{jk}`$ $`=`$ $`\left[\begin{array}{cc}\stackrel{~}{D}(x_j,x_k)& \stackrel{~}{S}(x_k,x_j)\\ \stackrel{~}{S}(x_j,x_k)& \stackrel{~}{I}(x_j,x_k)\end{array}\right],`$ (3.120)
$`A_j`$ $`=`$ $`\left[\begin{array}{cc}\mathrm{\Psi }_{2N+2\alpha }(z_{2j})& \stackrel{~}{I}(z_{2j},z_0)\\ \mathrm{\Psi }_{2N+2\alpha }(z_{2j1})& \stackrel{~}{I}(z_{2j1},z_0)\end{array}\right],`$ (3.123)
$`B_k`$ $`=`$ $`\left[\begin{array}{cc}\mathrm{\Psi }_{2N+2\alpha }^{}(x_k)& \mathrm{\Psi }_{2N+2\alpha }(x_k)\\ \stackrel{~}{S}(z_0,x_k)& \stackrel{~}{I}(z_0,x_k)\end{array}\right],`$ (3.126)
$`\mathrm{\Omega }`$ $`=`$ $`\left[\begin{array}{cc}0& \mathrm{\Psi }_{2N+2\alpha }(z_0)\\ \mathrm{\Psi }_{2N+2\alpha }(z_0)& 0\end{array}\right].`$ (3.129)
The functions $`\stackrel{~}{S}`$, $`\stackrel{~}{D}`$, $`\stackrel{~}{I}`$ are given in terms of $`S`$, $`D`$ and $`I`$ defined in Eq.(3.38) according to
$`\stackrel{~}{S}(x,y)`$ $`=`$ $`S(x,y)|_{}+{\displaystyle \frac{\mathrm{\Psi }_{2\alpha +2N}(x)}{s_{2\alpha +2N}}},`$ (3.130)
$`\stackrel{~}{D}(x,y)`$ $`=`$ $`D(x,y)|_{},`$ (3.131)
$`\stackrel{~}{I}(x,y)`$ $`=`$ $`I(x,y)|_{}+{\displaystyle \frac{\mathrm{\Psi }_{2\alpha +2N}^{}(x)\mathrm{\Psi }_{2\alpha +2N}^{}(y)}{s_{2\alpha +2N}}},`$ (3.132)
and
$$s_j=_{\mathrm{}}^{\mathrm{}}\mathrm{\Psi }_j(x)𝑑x.$$
(3.133)
Here $``$ stands for a set of substitutions
$$Q_j(z)Q_j(z)\frac{s_j}{s_{2\alpha +2N}}Q_{2\alpha +2N}(z)(j=0,\mathrm{},2\alpha +2N1).$$
(3.134)
Using Dyson’s identity (A.16), we find
$`p(z_0,\mathrm{},z_{2\alpha };x_1,\mathrm{},x_N)`$ (3.135)
$`=`$ $`\left({\displaystyle \underset{j=0}{\overset{\alpha +N1}{}}}q_j\right)\mathrm{Tdet}\left[\begin{array}{ccc}[g(z_{2j1},z_{2j};z_{2k1},z_{2k})]_{j,k=1,\mathrm{},\alpha }& [a(z_{2j1},z_{2j})]_{j=1,\mathrm{},\alpha }& [h(z_{2j1},z_{2j};x_k)]_{j=1,\mathrm{},\alpha ;k=1,\mathrm{},N}\\ [\widehat{a}(z_{2k1},z_{2k})]_{k=1,\mathrm{},\alpha }& \omega & [b(x_k)]_{k=1,\mathrm{},N}\\ [\widehat{h}(z_{2k1},z_{2k};x_j)]_{j=1,\mathrm{},N;k=1,\mathrm{},\alpha }& [\widehat{b}(x_j)]_{j=1,\mathrm{},N}& [f(x_j,x_k)]_{j,k=1,\mathrm{},N}\end{array}\right],`$ (3.139)
where
$`g(z_{2j1},z_{2j};z_{2k1},z_{2k})`$ $`=`$ $`\left[\begin{array}{cc}\stackrel{~}{I}(z_{2j1},z_{2k})& \stackrel{~}{I}(z_{2j1},z_{2k1})\\ \stackrel{~}{I}(z_{2j},z_{2k})& \stackrel{~}{I}(z_{2j},z_{2k1})\end{array}\right],`$ (3.143)
$`h(z_{2j1},z_{2j};x_k)`$ $`=`$ $`\left[\begin{array}{cc}\stackrel{~}{S}(z_{2j1},x_k)& \stackrel{~}{I}(z_{2j1},x_k)\\ \stackrel{~}{S}(z_{2j},x_k)& \stackrel{~}{I}(z_{2j},x_k)\end{array}\right],`$ (3.146)
$`f(x_j,x_k)`$ $`=`$ $`\left[\begin{array}{cc}\stackrel{~}{S}(x_j,x_k)& \stackrel{~}{I}(x_j,x_k)\\ \stackrel{~}{D}(x_j,x_k)& \stackrel{~}{S}(x_k,x_j)\end{array}\right],`$ (3.149)
$`a(z_{2j1},z_{2j})`$ $`=`$ $`\left[\begin{array}{cc}\mathrm{\Psi }_{2N+2\alpha }(z_{2j1})& \stackrel{~}{I}(z_{2j1},z_0)\\ \mathrm{\Psi }_{2N+2\alpha }(z_{2j})& \stackrel{~}{I}(z_{2j},z_0)\end{array}\right],`$ (3.152)
$`b(x_k)`$ $`=`$ $`\left[\begin{array}{cc}\stackrel{~}{S}(z_0,x_k)& \stackrel{~}{I}(z_0,x_k)\\ \mathrm{\Psi }_{2N+2\alpha }^{}(x_k)& \mathrm{\Psi }_{2N+2\alpha }(x_k)\end{array}\right],`$ (3.155)
$`\omega `$ $`=`$ $`\left[\begin{array}{cc}\mathrm{\Psi }_{2N+2\alpha }(z_0)& 0\\ 0& \mathrm{\Psi }_{2N+2\alpha }(z_0)\end{array}\right],`$ (3.158)
and $`\widehat{a}`$, $`\widehat{b}`$, $`\widehat{h}`$ stand for the duals of $`a`$, $`b`$, $`h`$. The skew-orthogonality relations (3.9) lead to
$`{\displaystyle _{\mathrm{}}^{\mathrm{}}}f(x^{},x)f(x,x^{\prime \prime })𝑑x=f(x^{},x^{\prime \prime }),`$ (3.159)
$`{\displaystyle _{\mathrm{}}^{\mathrm{}}}b(x)f(x,x^{})𝑑x=b(x^{}),`$ (3.160)
$`{\displaystyle _{\mathrm{}}^{\mathrm{}}}h(z,z^{};x)f(x,x^{})𝑑x=h(z,z^{};x^{}),`$ (3.161)
$`{\displaystyle _{\mathrm{}}^{\mathrm{}}}b(x)\widehat{b}(x)𝑑x=\omega ,`$ (3.162)
$`{\displaystyle _{\mathrm{}}^{\mathrm{}}}h(z,z^{};x)\widehat{b}(x)𝑑x=a(z,z^{}),`$ (3.163)
$`{\displaystyle _{\mathrm{}}^{\mathrm{}}}h(z,z^{},x)\widehat{h}(w,w^{};x)𝑑x=g(z,z^{};w,w^{}),`$ (3.164)
which mean that the condition of Theorem 3 is satisfied. Therefore we can inductively write
$`\mathrm{\Xi }_p(x_1,\mathrm{},x_p;z_0,\mathrm{},z_{2\alpha })={\displaystyle \frac{_{j=0}^{\alpha +N1}q_j}{_{j=0}^{2\alpha }w(z_j)_{j>k0}^{2\alpha }(z_jz_k)}}`$ (3.165)
$`\times \mathrm{Tdet}\left[\begin{array}{ccc}[g(z_{2j1},z_{2j};z_{2k1},z_{2k})]_{j,k=1,\mathrm{},\alpha }& [a(z_{2j1},z_{2j})]_{j=1,\mathrm{},\alpha }& [h(z_{2j1},z_{2j};x_k)]_{j=1,\mathrm{},\alpha ;k=1,\mathrm{},p}\\ [\widehat{a}(z_{2k1},z_{2k})]_{k=1,\mathrm{},\alpha }& \omega & [b(x_k)]_{k=1,\mathrm{},p}\\ [\widehat{h}(z_{2k1},z_{2k};x_j)]_{j=1,\mathrm{},p;k=1,\mathrm{},\alpha }& [\widehat{b}(x_j)]_{j=1,\mathrm{},p}& [f(x_j,x_k)]_{j,k=1,\mathrm{},p}\end{array}\right].`$ (3.169)
Then the $`p`$-level correlation function (3.89) is written as
$`\rho (x_1,\mathrm{},x_p;\{m\})={\displaystyle \frac{\mathrm{\Xi }_p(x_1,\mathrm{},x_p;z_0,\mathrm{},z_{2n})}{\mathrm{\Xi }_0(z_0,\mathrm{},z_{2n})}}`$ (3.170)
$`={\displaystyle \frac{\mathrm{Tdet}\left[\begin{array}{ccc}[g(z_{2j1},z_{2j};z_{2k1},z_{2k})]_{j,k=1,\mathrm{},\alpha }& [a(z_{2j1},z_{2j})]_{j=1,\mathrm{},\alpha }& [h(z_{2j1},z_{2j};x_k)]_{j=1,\mathrm{},\alpha ;k=1,\mathrm{},p}\\ [\widehat{a}(z_{2k1},z_{2k})]_{k=1,\mathrm{},\alpha }& \omega & [b(x_k)]_{k=1,\mathrm{},p}\\ [\widehat{h}(z_{2k1},z_{2k};x_j)]_{j=1,\mathrm{},p;k=1,\mathrm{},\alpha }& [\widehat{b}(x_j)]_{j=1,\mathrm{},p}& [f(x_j,x_k)]_{j,k=1,\mathrm{},p}\end{array}\right]}{\mathrm{Tdet}\left[\begin{array}{cc}[g(z_{2j1},z_{2j};z_{2k1},z_{2k})]_{j,k=1,\mathrm{},\alpha }& [a(z_{2j1},z_{2j})]_{j=1,\mathrm{},\alpha }\\ [\widehat{a}(z_{2k1},z_{2k})]_{k=1,\mathrm{},\alpha }& \omega \end{array}\right]}}.`$ (3.176)
Now we make replacement of the elements of the quaternion kernel back to those defined in Eq.(3.38):
$`\stackrel{~}{S}(x,y)S(x,y),\stackrel{~}{D}(x,y)D(x,y),\stackrel{~}{I}(x,y)I(x,y),`$ (3.177)
which does not change the values of the quaternion determinants in Eqs.(3.169) and (3.176). We introduce a set of notations
$`S_{jk}^{II}=S(z_j,z_k)(j,k=0,\mathrm{},2\alpha ),`$ (3.178)
$`S_{jk}^{IR}=S(z_j,x_k)(j=0,\mathrm{},2\alpha ;k=1,\mathrm{},p),`$ (3.179)
$`S_{jk}^{RR}=S(x_j,x_k)(j,k=1,\mathrm{},p),`$ (3.180)
and similarly for $`D`$ and $`I`$, and
$`Q_j^I`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Psi }_{2N+2\alpha }(z_j)}{s_{2N+2\alpha }}}(j=0,\mathrm{},2\alpha ),`$ (3.181)
$`Q_j^R`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Psi }_{2N+2\alpha }(x_j)}{s_{2N+2\alpha }}}(j=1,\mathrm{},p),`$ (3.182)
$`P_j^R`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Psi }_{2N+2\alpha }^{}(x_j)}{s_{2N+2\alpha }}}(j=1,\mathrm{},p).`$ (3.183)
Using Dyson’s identity (A.16), we can convert the quaternion determinant back to a Pfaffian, so that the correlation functions read
$$\rho (x_1,\mathrm{},x_p;\{m\})=(1)^{p(p1)/2}\frac{\mathrm{Pf}\left[\begin{array}{cccc}I^{II}& Q^I& I^{IR}& S^{IR}\\ (Q^I)^T& 0& (Q^R)^T& (P^R)^T\\ (I^{IR})^T& Q^R& I^{RR}& S^{RR}\\ (S^{IR})^T& P^R& (S^{RR})^T& D^{RR}\end{array}\right]}{\mathrm{Pf}\left[\begin{array}{cc}I^{II}& Q^I\\ (Q^I)^T& 0\end{array}\right]}.$$
(3.184)
Now we proceed to evaluate the component functions in the asymptitic limit where unfolded microscopic variables
$`\sqrt{2N}z_0`$ $``$ $`\zeta _0=0,`$ (3.185)
$`\sqrt{2N}z_{2j1}`$ $``$ $`\zeta _{2j1}=i\mu _j(j=1,\mathrm{},\alpha ),`$ (3.186)
$`\sqrt{2N}z_{2j}`$ $``$ $`\zeta _{2j}=i\mu _j(j=1,\mathrm{},\alpha ),`$ (3.187)
$`\sqrt{2N}x_j`$ $``$ $`\lambda _j(j=1,\mathrm{},p),`$ (3.188)
are kept fixed. We obtain Eq.(III A) and
$`𝐐_j^I{\displaystyle \frac{1}{\sqrt{2N}}}Q_j^I{\displaystyle \frac{1}{2}},`$ (3.190)
$`𝐐_j^R{\displaystyle \frac{1}{\sqrt{2N}}}Q_j^R{\displaystyle \frac{1}{2}},`$ (3.191)
$`𝐏_j^RP_j^R0.`$ (3.192)
These matrix elements constitute the finite-volume partition function
$`𝒵(\{\mu \})`$ $``$ $`\mathrm{\Xi }_0(\{{\displaystyle \frac{\mu }{\sqrt{2N}}}\})`$ (3.193)
$`=`$ $`\mathrm{const}.{\displaystyle \frac{\mathrm{Pf}\left[\begin{array}{cc}𝐈^{II}& 𝐐^I\\ (𝐐^I)^T& 0\end{array}\right]}{_{j=1}^\alpha \mu _j^3_{j>k}^\alpha (\mu _j^2\mu _k^2)^2}},`$ (3.196)
and the scaled spectral correlation functions
$`\rho _s(\lambda _1,\mathrm{},\lambda _p;\{\mu \})`$ $``$ $`\left({\displaystyle \frac{1}{\sqrt{2N}}}\right)^p\rho ({\displaystyle \frac{\lambda _1}{\sqrt{2N}}},\mathrm{},{\displaystyle \frac{\lambda _p}{\sqrt{2N}}};\{{\displaystyle \frac{\mu }{\sqrt{2N}}}\})`$ (3.197)
$`=`$ $`(1)^{p(p1)/2}{\displaystyle \frac{\mathrm{Pf}\left[\begin{array}{cccc}𝐈^{II}& 𝐐^I& 𝐈^{IR}& 𝐒^{IR}\\ (𝐐^I)^T& 0& (𝐐^R)^T& (𝐏^R)^T\\ (𝐈^{IR})^T& 𝐐^R& 𝐈^{RR}& 𝐒^{RR}\\ (𝐒^{IR})^T& 𝐏^R& (𝐒^{RR})^T& 𝐃^{RR}\end{array}\right]}{\mathrm{Pf}\left[\begin{array}{cc}𝐈^{II}& 𝐐^I\\ (𝐐^I)^T& 0\end{array}\right]}}.`$ (3.204)
It satisfies a sequence
$`\rho _s(\{\lambda \};\mu _1,\mu _1,\mathrm{},\mu _n,\mu _n,\mu _1,\mu _1,\mathrm{},\mu _n,\mu _n,0,0)\stackrel{\mu _n\mathrm{}}{}`$ (3.205)
$`\rho _s(\{\lambda \};\mu _1,\mu _1,\mathrm{},\mu _{n1},\mu _{n1},\mu _1,\mu _1,\mathrm{},\mu _{n1},\mu _{n1},0,0)\stackrel{\mu _{n1}\mathrm{}}{}\mathrm{},`$ (3.206)
as each of the masses is decoupled by going to infinity. To illustrate this decoupling, we exhibit in FIG.4 a plot of the spectral density $`\rho _s(\lambda ;\mu ,\mu ,\mu ,\mu ,0,0)`$ ($`p=1,\alpha =1`$).
###### Acknowledgements.
This work was supported in part (SMN) by JSPS Research Fellowships for Young Scientists, and by Grant-in-Aid No. 411044 from the Ministry of Education, Science, and Culture, Japan.
## A quaternion determinant
A quaternion is defined as a linear combination of four basic units $`\{1,e_1,e_2,e_3\}`$:
$$q=q_0+𝐪𝐞=q_0+q_1e_1+q_2e_2+q_3e_3.$$
(A.1)
Here the coefficients $`q_0,q_1,q_2,`$ and $`q_3`$ are real or complex numbers. The first part $`q_0`$ is called the scalar part of $`q`$. The quaternion basic units satisfy the multiplication laws
$`11=1,1e_j=e_j1=e_j(j=1,2,3),`$ (A.2)
$`e_1^2=e_2^2=e_3^2=e_1e_2e_3=1.`$ (A.3)
The multiplication is associative and in general not commutative. The dual $`\widehat{q}`$ of a quaternion $`q`$ is defined as
$$\widehat{q}=q_0𝐪𝐞.$$
(A.4)
For a selfdual $`N\times N`$ matrix $`Q`$ with quaternion elements $`q_{jk}`$ has a dual matrix $`\widehat{Q}=[\widehat{q}_{kj}]`$. The quaternion units can be represented as $`2\times 2`$ matrices
$`1\left[\begin{array}{cc}1& 0\\ 0& 1\end{array}\right],e_1\left[\begin{array}{cc}0& 1\\ 1& 0\end{array}\right],`$ (A.9)
$`e_2\left[\begin{array}{cc}0& i\\ i& 0\end{array}\right],e_3\left[\begin{array}{cc}i& 0\\ 0& i\end{array}\right].`$ (A.14)
We define a quaternion determinant Tdet of a selfdual $`Q`$ (i.e., $`Q=\widehat{Q}`$) as
$$\mathrm{Tdet}Q=\underset{PS_N}{}(1)^N\mathrm{}\underset{1}{\overset{\mathrm{}}{}}(q_{ab}q_{bc}\mathrm{}q_{da})_0,$$
(A.15)
where $`P`$ denotes any permutation of the indices $`(1,\mathrm{},N)`$ consisting of $`\mathrm{}`$ exclusive cycles of the form $`(abc\mathrm{}da)`$ and $`(1)^N\mathrm{}`$ is the parity of $`P`$. The subscript $`0`$ means that the scalar part of the product is taken over each cycle. Note that a quaternion determinant of a selfdual quaternion matrix is always a scalar. The quaternion determinant of $`Q`$ can as well be represented by its $`2N\times 2N`$ complex matrix representation $`C(Q)`$ :
$$\mathrm{Tdet}Q=\mathrm{Pf}[JC(Q)],J=𝟙_{}\left[\begin{array}{cc}\mathrm{𝟘}& \mathrm{𝟙}\\ \mathrm{𝟙}& \mathrm{𝟘}\end{array}\right].$$
(A.16)
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# Metastable states in glassy systems
## I Introduction
Slowly relaxing systems such as glasses or compacting granular media can be viewed as having fast, local, quasi-equilibrium dynamics, plus a slow, nonequilibrium drift. These two superposed motions can take different forms: ‘cage’ vibrations plus structural rearrangements in glasses; bulk fluctuations plus domain-wall motion in coarsening problems; etc. At a given long time, the fast motion covers a region of phase space which one may picture as a ‘metastable state’.
Even though metastable states are a familiar and appealing concept, it turns out that defining them in an unambiguous manner in non-mean field models is quite subtle. This in turn has as a consequence that such standard ideas in glass theory as ‘configurational entropy’ (related to the number of metastable states) are not only hard to calculate, but are indeed, with the exception of some fortunate cases, approximate as concepts.
In this paper we show how these questions can be put on a well defined basis using a formalism that does not rely on specifically mean-field concepts.
First of all: given that one can simulate and in certain cases calculate analytically the complete history of a sample starting from a quench, why should one have any need to introduce the apparently unneeded notion of ‘metastable state’ ? Indeed, these states only come into play when one wishes to make arguments such as: ‘phase space contains such and such a distribution of states, which will be accessed with such and such a probability by a typical dynamical history. Long-time out of equilibrium observables can be directly calculated by averaging the observables over some subset of states — and further reference to dynamics may be omitted’.
This kind of ‘ergodic’ argument was pioneered by Edwards , who proposed that in compacting or slowly flowing granular systems one can obtain the correct dynamical observables by averaging the values they take over all blocked configurations of a certain volume. It later turned out that mean-field glass models relaxing at zero temperature had exactly Edwards’ ergodicity property : at long times any nonequilibrium observable is correctly given by the typical value it takes over all local energy minima of the appropriate energy density.
A first problem arises when one wishes to apply this concept at finite temperatures (or vibration, in the case of granular media). There again, the mean-field case offers a suggestion: at non zero temperature Edwards’ argument works as well, provided one substitutes ‘energy minima’ by ‘free-energy minima’ (‘states’). This construction is possible because within mean-field we have a well defined notion of free-energy landscape, whose local minima are in some (but not all) cases related to completely stable distributions. However, as discussed by Franz and Virasoro , one needs to consider ‘quasi-states’ with finite lifetimes in order to understand the situation at finite waiting times.
In finite-dimensional problems, a high-lying metastable state cannot have an infinite lifetime: there is always a finite probability of escape through the nucleation of droplet of a more favourable phase. Hence, which distribution one considers as metastable depends always on which lifetimes one is considering. For example, the concept of ‘configurational entropy’ (the logarithm of the number of states), ubiquitous in glass theory, has in finite dimensions only a meaning with a timescale attached. Moreover, even the mere definition of an Edwards distribution is not as simple, quite apart from the question of the validity of the ergodicity-like hypothesis it assumes.
In section 2 we shall review the notion of metastable state within mean-field glass models and how the knowledge of their distribution allows in certain cases to reproduce some results obtained from the full solution of the out of equilibrium dynamics. We shall also mention some limitations found even at this level of the identification ‘free-energy minimum $``$ stable state’.
In section 3 we discuss a strategy valid in any dimension for the definition and calculation of metastable states based on the evolution operator developed by Gaveau and Schulman . We discuss how one can thus recast the question of configurational entropy and Edwards’ distribution in a form relevant in finite dimensions at nonzero temperature (or stronger vibration, in the case of granular media) by considering finite-lifetime metastable states, in the spirit of the ‘quasi-states’ discussed in .
In section 4 we apply this method to a simple mean-field glass model. We show how one can rederive in this way both the number of states and the dynamics inside a state, within a framework whose applicability goes beyond mean-field.
## II A Fortunate case: mean-field models
Consider the mean-field model of ferromagnet:
$$E=\frac{1}{2N}\underset{i,j}{}S_iS_jh\underset{i}{}S_i$$
(1)
where the sum is over all spins. The spins can be Ising $`S_i=\pm 1`$ or spherical $`_iS_i^2=N`$. One can easily obtain a free energy in terms of local magnetisations $`m_i=S_i`$ (where $``$ means the average over the Gibbs measure):
$$f(\{m_i\})=\frac{1}{2N}\underset{ij}{}m_im_jh\underset{i}{}m_iTS(\{m_i\})$$
(2)
where $`T=1/\beta `$ denotes the temperature and $`S(\{m_i\})`$ is the usual entropic term:
$`S(\{m_i\})=`$ $`{\displaystyle \underset{i}{}}{\displaystyle \frac{1}{2}}\left[(1m_i)\mathrm{ln}(1m_i)+(1+m_i)\mathrm{ln}(1+m_i)\right]\text{for}S_i=\pm 1`$ (3)
$`=`$ $`{\displaystyle \frac{1}{2}}\mathrm{ln}(1q)\text{ spherical model}`$ (4)
and $`Nq=_{i=1}^Nm_i^2`$.
The states are represented by the minima of $`f`$. At $`T<T_c`$ there are two, and the deeper one dominates the Boltzmann distribution. If $`h>0`$ one of the states becomes metastable: within mean-field its lifetime is infinite, but in finite dimensions it will decay.
The TAP approach
Next, consider the glassy Hamiltonians:
$$E=\frac{1}{p!}\underset{i_1,\mathrm{},i_p}{}J_{i_1,\mathrm{},i_p}S_{i_1}\mathrm{}S_{i_p}$$
(5)
where the $`J_{i_1,\mathrm{},i_p}`$ are independent random variables with variance $`p!/2N^{p1}`$. The Ising version with $`p=2`$ is the Sherrington-Kirkpatrick model, the mean field version of spin-glasses. The models with $`p3`$ are instead systems having the behaviour resembling structural glasses .
It turns out that one can find for these models a free energy function analogous to (2) (the ‘TAP’ free energy ), in terms of local magnetisations $`m_i`$:
$$f_{TAP}(\{m_i\})=\frac{1}{p!}\underset{i_1,\mathrm{},i_p}{}J_{i_1,\mathrm{},i_p}m_{i_1}\mathrm{}m_{i_p}\frac{\beta }{4}[1q^ppq^{p1}(1q)]TS(\{m_i\})$$
(6)
where $`Nq=_im_i^2`$ and $`S(\{m_i\})`$ is given in (4) for the spherical and Ising cases, respectively.
Within the TAP approach, one signature of the glass transition is the fact that the free energy (6) has, below a critical temperature, many ($`e^{aN}`$) minima. The main difference between mean-field versions of spin and structural glasses is seen, in the TAP approach, in the way states are separated.
Given the analogies with the ferromagnet, it seems very tempting to attribute to the minima of (6) a dynamical meaning of ‘state’. For the models with $`p3`$ this has been done by starting from a configuration where the coordinates are as close as possible to having a given $`m_i^a`$, and then checking that the subsequent dynamical evolution is a stable, quasi-equilibrium situation confined to a region of phase space in such a way that $`<S_i>_{\text{time-average}}=m_i^a`$.
On the contrary, for models like the Sherrington-Kirkpatrick model the identification of all TAP-minima with stable states seems to breaks down. Let us formulate an heuristic argument to see this. Decreasing temperature, minima split in a second-order transition manner. At least a fraction of the minima are ‘born’ this way , and to get an exponential number of minima one needs that on average there is a division every $`O(1/N)`$ change in temperature. Hence, a fraction of TAP solutions are just $`O(1/N)`$ below their critical temperature, and under those circumstances barriers cannot be large enough to dynamically separate them from their ‘twins’.
FDT-temperature. Long-time out of equilibrium observables
The dynamics of model (5) following a quench below the critical temperature can be solved analytically . One finds that the system never equilibrates, and remains aging just above a threshold level of energy density $`e_{th}`$ and free energy density $`f_{th}`$ higher than the equilibrium ones. Given any two observables $`A`$ and $`B`$, one can define their dynamic correlation function:
$$C_{AB}(t,t^{})=A(t)B(t^{})$$
(7)
and the integrated response $`\chi _{AB}`$ to a field $`h_B`$ conjugate to $`B`$ acting between times $`t^{}`$ and $`t`$:
$$\chi _{AB}(t,t^{})=\frac{\delta A(t)}{\delta h_B}$$
(8)
where the averages are over the dynamical realisation.
If one makes a parametric plot of $`\chi _{AB}(t,t^{})`$ versus $`C_{AB}(t,t^{})`$ one obtains for model (5) with $`p>2`$ at long times a curve like Fig. 2. One has in addition to the straight line with gradient $`1/T`$ (as in equilibrium), another straight line of gradient, say, $`1/T_{eff}`$, associated to the slow relaxation. $`T_{eff}`$ is the same for every pair of observables $`A`$ and $`B`$, and it can be shown that it satisfies all the properties of a true temperature.
The appearance of a temperature in a system is an indication of some form of ergodicity, in this case clearly not the usual Gibbs-Boltzmann equilibrium at temperature $`T`$. Indeed, soon after the dynamical solution was obtained, Monasson and Virasoro observed that the temperature $`T_{eff}`$ could be reobtained from the TAP approach without making any reference to the dynamics: Defining the complexity (or configurational entropy) $`𝒮(f)`$ as the logarithm of the number of TAP solutions of a given free energy, one checks that
$$\frac{1}{T_{eff}}=\frac{𝒮(f)}{f}|_{f=f_{th}}$$
(9)
How this equality follows from Edwards’ assumption is discussed in .
Furthermore, one can also see that the long-time values of macroscopic observables are given by the flat average of their values taken over all TAP solutions of the dynamical energy density $`e_{th}`$. Hence, we have a strong indication for a measure à la Edwards, this time applied to TAP states.
If the system has large, but finite size, it will slowly approach equilibrium. In that case, a plot like Fig. 2 will show that the two tracts slowly tend to become parallel and $`T_{eff}(t)`$ (now a function of time) tends to $`T`$. Inspired by work of Bonilla et. al. , Nieuwenhuzen conjectured that one could extend the relation (9) for all times, using the TAP solutions of the energy level appropriate at each time.
Later on, a two-temperature $`\chi `$ versus $`C`$ plot (and hence the existence of an FDT temperature $`T_{eff}`$) was seen to occur in realistic models such as three dimensional Lennard-Jones glasses . Several of these simulations where performed at temperatures above the putative equilibrium glass temperature, so that the existence of a (slowly evolving) well-defined $`T_{eff}`$ is surely a dynamical phenomenon, unrelated to the structure of equilibrium states. If one wishes to consider this as a symptom of slowly evolving flat distribution between metastable states, one finds oneself in the embarrassing situation that it is not entirely obvious what one means by ‘metastable state’ in finite dimensions, as we now discuss.
The problem with finite dimensional and driven systems
In finite dimensions and non-zero temperature nucleation arguments suggest that a distribution with dynamical free energy density (to be defined below) higher than the equilibrium one should decay through nucleation in finite times. We are hence in a situation in which we have no absolute notion of state without making reference to a timescale (and hence to dynamics): two different distributions may be confused into a single state or be treated as two separate entities depending on whether the time to go from one to the other is smaller than or larger than the timescale considered.
If we are interested in systems driven by shear or by vibration, we have the additional problem that even in the mean-field case the distribution is not Gibbsean within a state. In a vibrated case, the notion of stability must be substituted by the notion of periodicity, so that a ‘state’ will turn out to be a structure periodic in time.
Before entering into the present approach, let us mention that a pragmatic way of dealing with these difficulties, at least at very short time-scales, is the so-called ‘inherent structure’ construction . Though it does not solve the questions of principle mentioned above , it offers a practical way around applicable to concrete problems.
## III Dynamical definition of metastable states
Let us consider a system evolving with stochastic dynamics, which for definiteness we shall consider is of the Langevin form. The probability distribution will evolve according to:
$`{\displaystyle \frac{dP(S,t)}{dt}}`$ $`=`$ $`HP(S,t)`$ (10)
$`H`$ $`=`$ $`{\displaystyle \frac{}{S_i}}\left(T{\displaystyle \frac{}{S_i}}+{\displaystyle \frac{E}{S_i}}\right)`$ (11)
where $`H`$ is the Fokker-Planck operator. The potential energy $`E`$ can be time-dependent, and furthermore one can add to $`(\frac{E}{S_i})`$ forces that do not derive from a potential.
Given a probability distribution $`P`$, one can define a dynamic free energy
$$F(t)=\left(TP(S,t)\mathrm{ln}P(S,t)+E(S)P(S,t)\right)𝑑S$$
(12)
If $`H`$ is time-independent, any stationary configuration satisfies
$$HP_{\text{stationary}}=0$$
(13)
Moreover, writing any distribution as $`P(x,t)=c_i(t)\psi _i(x)`$ where $`\psi _i(x)`$ are the right eigenvectors of $`H`$:
$$H|\psi _i=\lambda _i|\psi _i$$
(14)
the evolution equation (11) implies:
$$c_i(t)=c_i^oe^{\lambda _it}$$
(15)
We see that if $`P`$ is to vary slowly it has to be concentrated on eigenvectors with low eigenvalues $`\lambda _i`$. Indeed, each $`\lambda _i`$ is an inverse timescale .
In the following subsections we shall motivate and discuss an identification of the set of small eigenvalue $`\psi _i`$’s with metastable states.
Motivation
Consider first the system of two hard disks performing Langevin dynamics in a box (Fig. 2). Clearly, if the disks are really not interpenetrable, there are two different ergodic components, each composed of the mirror image of the configurations of the other. Symmetry implies that the spectrum of the Fokker-Planck operator is doubly degenerate. The two lowest (zero) eigenvalues correspond to two stationary distributions. One can construct an associated eigenvector as the flat distribution over all pairs of coordinates of the centers of the disks such that they do not superpose and such that the disk ‘A’ is to the right and ‘B’ to the left, and similarly a second eigenvector corresponding to having the disk ‘B’ to the right and ‘A’ to the left. These are the ‘pure states’: any linear combination of these two distributions will be an ‘impure’ state. The next higher eigenvalues are equal to the inverse of the time needed for the particles to explore their ergodic component. Note, in passing, that this hard-spheres system is the typical example in which the inherent structure construction is not meaningful while the present one has no problem.
In this example we have strictly two ergodic components, corresponding to twofold ground state degeneracy of the Fokker-Planck operator. This is indeed very general: suppose we have $`p`$ ergodic components $`𝒞_1,\mathrm{},𝒞_p`$, with typical times $`t_1,\mathrm{},t_p`$ required to explore each component. We can construct an independent eigenvector with zero eigenvalue using the stationary distribution $`P_a(x)`$ restricted to each component $`𝒞_a`$. These completely span the zero eigenvalue subspace. To show this, we calculate:
$`\text{tr}[e^{t^{}H}]`$ $`=`$ $`{\displaystyle 𝑑x𝑑yx|e^{\frac{t^{}}{2}H}|yy|e^{\frac{t^{}}{2}H}|x}={\displaystyle \underset{a=1}{\overset{p}{}}}{\displaystyle \underset{b=1}{\overset{p}{}}}{\displaystyle _{x𝒞_a}}𝑑x{\displaystyle _{y𝒞_b}}𝑑yx|e^{\frac{t^{}}{2}H}|yy|e^{\frac{t^{}}{2}H}|x`$ (16)
$`=`$ $`{\displaystyle \underset{a=1}{\overset{p}{}}}{\displaystyle _{x𝒞_a}}𝑑x{\displaystyle _{y𝒞_a}}𝑑yx|e^{\frac{t^{}}{2}H}|yy|e^{\frac{t^{}}{2}H}|x`$ (17)
If we now take $`t^{}`$ much larger than all the $`t_i`$, $`y|e^{\frac{t^{}}{2}H}|xP_a(y)`$, the equilibrium probability for $`y`$ restricted to the ergodic component $`a`$ to which $`x`$ belongs. Hence:
$$\text{tr}[e^{t^{}H}]\underset{a=1}{\overset{p}{}}\underset{b=1}{\overset{p}{}}_{x𝒞_a}𝑑x_{y𝒞_a}𝑑yP(x)P(y)p$$
(18)
This shows that the number of states ‘below the gap’ coincides with the number of ergodic components.
In the preceding examples the ergodic components are strictly separated. However, in most applications this is not the case: there is in fact a passage time between components that only becomes infinite in some limit. To understand the construction in these cases, consider a very low temperature Langevin process occurring in asymmetric and symmetric double-well potentials as in Figure 3. On the left of the figure we show the lowest levels of the spectrum for both cases, and at the top the corresponding eigenvectors. For the asymmetric case, the two lowest eigenvalues are separated by the inverse Arrhenius escape time from the highest minimum. All other eigenvalues are much higher ($`O(1)`$), and include the escape time from a maximum, etc. The eigenfunction labeled a is essentially positive and represents a “pure” state $`P_1`$, while one can make a linear combination $`P_2=(\text{a}+\text{b})`$ that will also be positive and concentrated on the metastable minimum. For the symmetric well the situation is similar, but now it is the linear combinations $`P_1=(\text{a}+\text{b})`$ and $`P_2=(\text{a}\text{b})`$ that play the role of “pure” states. Any other combination of the form $`yP_1+(1y)P_2`$ with $`0<y<1`$ will give an ‘impure’ (almost) steady state.
Note that these definitions make sense in the time-window in which we can consider the exponential Arrhenius times much larger that any other time involved ($`O(1)`$), and this will happen only in the low temperature limit.
If the temperature is non-zero, a separation of timescales can happen as a result of the thermodynamical (or other) limit. For example, it is easy to see that a mean-field ferromagnet at $`0<T<T_c`$ will have a similar spectrum, but with $`N`$ playing the role of large parameter instead of $`1/T`$. The case of finite dimensional ferromagnets is slightly more subtle: we have there a timescale for domain excitations within a state that can be as large as a power law in $`N`$ (the time it takes for a large domain to collapse), and a much longer timescale ($`e^{cN^{(d1)/d}}`$) for going from one phase to the other.
The construction of Gaveau and Schulman
In general, the low eigenvalue spectrum can correspond to more than two eigenvectors. One can now ask in general whether all these eigenvectors (or combinations of them) represent positive, stable distributions, and whether one can construct as many pure states as there are low eigenvalues. In a series of papers, Gaveau, Schulman and Lesne showed this to be the case, gave recipes for the explicit construction of the pure states, proved their unicity and exploited this construction to study metastability.
We shall be only descriptive here, we refer the reader to the references for the proofs, as well as other investigations concerning metastability. Consider for example a Fokker-Planck operator having the lowest $`p`$ eigenvalues $`\lambda _1,\mathrm{},\lambda _p`$ separated by a gap from the others, ı.e. one can find a $`t^{}`$ such that one can consider that:
$`t^{}\lambda _i`$ $`<<1`$ $`\text{for}i=1,\mathrm{},p`$ (19)
$`t^{}\lambda _i`$ $`>>1`$ $`\text{for}i=p+1,\mathrm{}`$ (20)
In the previous simple example any $`t^{}`$ such that $`t^{}`$ is positive will satisfy this as $`T0`$. The meaning of $`t^{}`$ is clear: it is a timescale much longer than the relaxation into states, but much shorter than transitions between states. Clearly, the operator $`\mathrm{exp}[t^{}H]`$ is essentially a projector onto the space ‘below the gap’ (up to terms of order $`\mathrm{exp}[t^{}\lambda _i]`$, with $`ip`$). Within the same accuracy, one can then find a basis of $`p`$ right eigenvectors $`|P_i`$ which are:
* positive: $`P_i(x)=x|P_i0`$:
* almost stationary: $`H|P_i0i=1,\mathrm{},p`$
* normalised and not zero in non-overlapping regions of space.
The last property is related to the fact that one can also find a basis of $`p`$ approximate almost stationary ($`Q_i|H0`$) left eigenvectors $`Q_i|`$, such that each $`Q_i`$ is essentially one within the support of $`P_i(x)`$, zero everywhere else and satisfy the orthogonality and normalisation conditions:
$$Q_i|P_j\delta _{ij}$$
(21)
Given any observable $`A`$, we can calculate its average within the state “$`i`$” as:
$$A_i=Q_i|A|P_i$$
(22)
One can also write approximately:
$$e^{t^{}H}\underset{i}{}|P_iQ_i|$$
(23)
As a consequence the $`|P_i`$ vectors have all the good properties to represents metastable states: they are positive normalised distributions, non zero only on different regions of the configuration space and they are stationary on time scales less than $`t^{}`$. In the proof, as in the simple example of the previous subsection, the definition is unavoidably linked to a timescale: if one considers really infinite times, before any other limit, then the distinction between states vanishes. Furthermore, it is assumed the number of states so defined remains finite in the thermodynamic limit. We believe that going to situations in which this is not the case (as we will below) is indeed not entirely innocent, but is at the heart of quite a few problems associated to the definition of “state” in glassy systems. (We have already encountered such subtleties when we discussed TAP minima in the Sherrington-Kirkpatrick model).
Driven systems.
The construction described above is not limited to purely relaxational or time-independent systems. Consider for example the case in which a system is periodically ‘vibrated’ or ‘tapped’. One can still try to look for stationary, or rather, periodic situations. One can repeat essentially the same argument by considering the evolution operator through one cycle:
$$U=𝒯_{\text{cycle}}𝑑te^{tH(t)}$$
(24)
where $`𝒯`$ denotes time-order. One has to look now for the eigenvectors of $`U`$ whose eigenvalue is close to one.
Similarly, one can also work with systems driven by constant nonconservative forces (as in a sheared fluid), and with nonlinear space-dependent friction (as in granular systems).
Pure barriers.
Given that pure states can be viewed as playing the role for finite temperature that certain energy minima play for zero temperature, one is naturally led to ask which distributions play the role of barriers (or in general saddle points) in finite temperature. For Fokker-Planck processes this can be done naturally starting from the states constructed as above. In Appendix A we show how this can be done.
We remark that a solvable example, in which saddle points play an important role in the spectrum of the evolution operators, is the Glauber evolution of the completely connected Ising model and its generalizations . It has been found that the lowest eigenvalues, i.e. the longest relaxation times, are gathered in families, each one being in correspondence with a stationary, not necessary stable, point of the static mean field free energy.
## IV Flat distribution over states at finite temperature or vibration. A definition of configurational entropy.
Motivation.
As we mentioned in Section II, the fact that one finds a two-temperature behaviour in mean-field glasses can be seen as suggesting the relevance of a measure consisting in summing over metastable states of given energy (or free energy) with equal weight. Armed with the construction for metastable states we discussed above, we shall see how this measure can be expressed in finite dimensional, or periodically driven systems.
Expressions.
Consider, in the spirit of Edwards’ distribution, an average of an observable over states, each measured with equal weight. This will be relevant for the long-time out of equilibrium dynamics under the assumption that almost all states of given energy have the same basin of attraction.
Using the same notation as in the previous section, we define the average over the measure $``$ of an observable A in the following way:
$$A_t^{}=\frac{1}{p}\underset{i=1}{\overset{p}{}}A_i$$
(25)
where the subindex $`t^{}`$ reminds us of the fact that states are now defined according to their timescale. We have:
$$\underset{i=1}{\overset{p}{}}A_i=\underset{i=1}{\overset{p}{}}Q_i|A|P_i=\underset{i=1}{\overset{p}{}}\text{tr}\left[|P_iQ_i|A\right]=\text{tr}[e^{t^{}H}A]$$
(26)
where we have used (22) and (23). Hence:
$$A_t^{}=\frac{\text{tr}[e^{t^{}H}A]}{\text{tr}[e^{t^{}H}]}$$
(27)
Note that once written this way, all reference to pure states has disappeared, except indirectly in the value chosen for $`t^{}`$.
We often need an equation like (27), but restricted to states having a certain energy, particle number, etc. In that case we generalise (27) as, for example:
$$A_t^{}(E_o)=\frac{\text{tr}[\delta (EE_o)e^{t^{}H}A]}{\text{tr}[\delta (EE_o)e^{t^{}H}]}=\frac{\text{tr}_{E_o}[e^{t^{}H}A]}{\text{tr}_{E_o}[e^{t^{}H}]}$$
(28)
where $`\text{tr}_{E_o}`$ denotes a restricted trace.
Once we make the assumption that all states of the same energy (or particle number, etc) have an equal weight for the purposes of calculating a dynamical observable, it becomes meaningful to count their number at given energy, the configurational entropy:
$$𝒮_t^{}(E_o)\mathrm{ln}\text{tr}[\delta (EE_o)e^{t^{}H}]=\mathrm{ln}\text{tr}_{E_o}[e^{t^{}H}]$$
(29)
If $`t^{}0`$ we get the microcanonical measure (and entropy, up to irrelevant constants ), and if we let $`t^{}\mathrm{}`$ we find no high-lying metastable states at all in finite dimensions. The dependence on $`t^{}`$ is hence unavoidable if one is to obtain a finite configurational entropy in that case.
Equation (29) defines the timescale-dependent configurational entropy. One also needs the average entropy within a state $`s_t^{}`$, and the corresponding average free-energy of a state $`f_t^{}=E_oTs_t^{}`$. Using the construction of (21), (22) and (23), we have:
$$s_t^{}(E_o)=\frac{_{x/E(x)=E_o}𝑑xx|e^{t^{}H}|x\mathrm{ln}[x|e^{t^{}H}|x]}{_{x^{}/E(x^{})=E_o}𝑑x^{}x^{}|e^{t^{}H}|x^{}}$$
(30)
The meaning of this equation becomes transparent in the example of the completely separated ergodic components of the previous section.
Note that also the intra-state entropy is timescale-dependent. Indeed, if we set $`t^{}0`$, we are defining as ‘states’ the configurations themselves . On the other extreme, if $`t^{}`$ is longer than the equilibration time, $`x^{}|e^{t^{}H}|x^{}`$ gives the Gibbs-measure, and the intra-state entropy becomes the usual entropy. In short: changing the timescale both changes the number of states and their nature, hence the change in configurational and intra-state entropy, respectively.
Flat measures and effective temperatures
Suppose we have a glassy system, taken through a given thermal history (an annealing protocol) to a glassy phase at time $`t`$, at which time its energy is $`E(t)=E_o`$ (and if we let the particle number or the volume change we should specify also them). The assumption of typicality of metastable states, is then:
$$A_t^{}(E_o)=\frac{\text{tr}_{E_o}[e^{t^{}H}A]}{\text{tr}_{E_o}[e^{t^{}H}]}A_{\text{history}}$$
(31)
where the average is over several realisations of the same protocol, ending at time $`t`$ with energy $`E(t)=E_o`$.
From the Langevin point of view, the left hand side of (31) corresponds to adding over all periodic trajectories starting from an energy $`E_o`$ with period $`t^{}`$. Equation (31) tells us that thermal histories give the same result as a very particular set of trajectories, in a manner analogous as when one represents chaotic dynamical systems by using only the periodic orbits.
In the zero-temperature or in the mean-field case, we can set $`t^{}=\mathrm{}`$ in (31), and thus select the states with infinite lifetime. This is Edwards’ prescription for granular media ($`T0`$) and the one we discussed above within mean-field. In finite dimensions, where metastable states eventually nucleate, we are forced to give $`t^{}`$ a finite value. The fact that the choice of $`t^{}`$ is not unique already tells us that equation (31) will be an approximation.
The central remaining question now is what is a reasonable value for $`t^{}`$. Indeed, giving a value of $`t^{}`$ determines a configurational entropy $`𝒮_t^{}(E_o)`$ and an intra-state free-energy: $`f_t^{}(E_o)=E_oTs_t^{}(E_o)`$. This in turn determines a temperature associated with the timescale as:
$$T_t^{}\left[\frac{𝒮_t^{}}{f_t^{}}\right]^1$$
(32)
(where we have eliminated the energy in favour of $`f_t^{}`$). It seems now natural to compare the different relaxation times of the correlations in a given problem with $`t^{}`$. For example, glassy systems can often be described with two timescales, a fast (‘$`\beta `$’) relaxation and a slow, (waiting time dependent if the system is aging) ‘$`\alpha `$-relaxation’ $`t_\alpha `$. If $`t^{}`$ is small ($`t^{}<<t_\beta `$), Eqn. (32) gives $`T_{eff}T`$ . If we put instead:
$$t_\beta <<t^{}t_\alpha $$
(33)
then $`T_t^{}`$ may be different (larger) than $`T`$. This is the temperature to compare with the one governing the relation between correlation and response in the regime corresponding to the $`\alpha `$-relaxation .
There are cases in which the ‘slow’ relaxations happen in several timescales, becoming more and more different as aging proceeds . Then, the definition (32) immediately yields a different temperature for every widely separated timescale, and this would agree with what one observes from fluctuation-dissipation relations .
Time dependent order parameters
We are interested in calculating two-time correlation functions:
$$<A(t)B(t^{})>\frac{1}{𝒩}\text{tr}\left[e^{(t^{}t)H}Ae^{(tt^{})H}Be^{t^{}H}\right]=\frac{1}{𝒩}\text{tr}\left[e^{(t^{}\tau )H}Ae^{\tau H}B\right]=<AB>(\tau )$$
(34)
where $`\tau =tt^{}`$, and $`𝒩`$ is the normalisation. Cyclic permutation implies:
$$<AB>(\tau )=<BA>(t^{}\tau )$$
(35)
If $`H`$ is a time-dependent Fokker-Planck operator, associated with forces deriving from a potential $`E`$, we have:
$$e^{\beta E}He^{\beta E}=H^{}$$
(36)
Using (36) in (34) we get the time-reversal property:
$$<AB>(\tau )=tr^{}[e^{(t^{}\tau )H}(e^{\beta E}Be^{\beta E})^{}e^{\tau H}(e^{\beta E}Ae^{\beta E})^{}]]=<\stackrel{~}{B}\stackrel{~}{A}>^{}(\tau )$$
(37)
where $`\stackrel{~}{A}(e^{\beta E}Ae^{\beta E})^{}`$ and $`\stackrel{~}{B}(e^{\beta E}Be^{\beta E})^{}`$
We shall need to consider the cases in which $`A`$ and $`B`$ are respectively $`S_i`$ and $`\widehat{S}_i\frac{}{S_i}`$. Let us define, for a system of $`N`$ degrees of freedom:
$`C(tt^{})`$ $`=`$ $`C(\tau )={\displaystyle \frac{1}{N}}{\displaystyle \underset{i}{}}<S_i(t)S_i(t^{})>`$ (38)
$`R(tt^{})`$ $`=`$ $`R(\tau )={\displaystyle \frac{1}{N}}{\displaystyle \underset{i}{}}<S_i(t)\widehat{S}_i(t^{})>`$ (39)
$`D(tt^{})`$ $`=`$ $`D(\tau )={\displaystyle \frac{1}{N}}{\displaystyle \underset{i}{}}<\widehat{S}_i(t)\widehat{S}_i(t^{})>`$ (40)
Because we are considering periodic trajectories, causality is violated. This means that neither $`D(\tau )`$ nor $`R(\tau )`$ for negative $`\tau `$ need to be zero. However, it is easy to show that if $`(t^{}\tau )`$ is larger than the thermalisation time (in which the system is projected to its Gibbs measure), then we recover $`D(\tau )=0`$ and $`R(\tau )R(\tau )=0`$. The existence of solutions violating causality for $`t^{}`$ large is then a symptom of large equilibration times, i.e. of glassiness.
Using the time-reversal equation (37), it is easy to derive a non-causal form of FDT, valid for any $`t^{}`$:
$$\frac{}{t^{}}C(tt^{})=T\left\{R(\tau )R(\tau )\right\}$$
(41)
The physical meaning of $`D`$, unlike that of $`C`$ and $`R`$, is unfamiliar. If we couple the states to a time-dependent, random magnetic ‘pinning fields’ $`h_i`$ such that $`\overline{h_i(t)h_j(t^{})}=F(t,t^{})\delta _{ij}`$, the fields will make a change in the number of metastable states, and it is easy to show that:
$$D(t,t^{})=\frac{1}{N}\frac{\delta }{\delta F(t,t^{})}\mathrm{ln}\left\{\text{tr}[e^{t^{}H}]\right\}$$
(42)
It is hence clear that systems with a finite number of states will have $`D=0`$ in the long-time limit.
## V The calculation
In the following we apply the theory discussed above to a simple glassy system: the spherical version of model (5). These models are thought to be mean-field versions of structural glasses, we shall not deal in this paper with models corresponding to spin-glasses for which, as mentioned above, we do not expect the present computation of states to be in correspondence with the TAP-equation based calculations in the literature.
The trace of the Fokker-Planck operator, at a fixed energy density $`E`$, can be written as a functional integral over the spin fields $`S_i(t)`$ and the response fields $`\widehat{S}_i(t)`$ with periodic boundary conditions on $`S_i(t)`$. Once the average of the trace has been performed , the action depends on the fields $`S_i(t)`$ and $`\widehat{S}_i(t)`$ through the two time functions $`C(t,t^{}),R(t,t^{}),D(t,t^{})`$ only. As a consequence one can integrate out the fields $`S_i(t)`$ and $`\widehat{S}_i(t)`$ and get an effective action on two time functions:
$`S/N`$ $`=`$ $`{\displaystyle _0^t^{}}𝑑t\left(_tR(t,t^{})+\lambda R(t,t^{})TD(t,t^{})\right)|_{t^{}=t^+}`$ (43)
$`+`$ $`{\displaystyle \frac{p}{4}}{\displaystyle _0^t^{}}𝑑t𝑑t^{}\left(D(t,t^{})C^{p1}(t,t^{})+(p1)R(t,t^{})R(t^{},t)C^{p2}(t,t^{})\right)`$ (44)
$``$ $`{\displaystyle \frac{\widehat{\lambda }}{2}}{\displaystyle _0^t^{}}𝑑t\left(C(t,t)1\right)+{\displaystyle \frac{1}{2}}\text{Tr}\mathrm{ln}\text{M}`$ (45)
where the operator $`M`$ reads:
$`M=\left(\begin{array}{cc}R(t,t^{})& C(t,t^{})\\ D(t,t^{})& R(t^{},t)\end{array}\right),`$ (48)
Since we consider times of order one with respect to N the functional integral is dominated by a saddle point contribution. We shall obtain periodic dynamic solutions which, in the glassy phase (a) break causality, (b) have non-zero action, (c) satisfy time-translational invariance, and (d) satisfy time-reversal and its consequence (41). Note that (a) and (b) are properties typical of instantons, while (c) and (d) are not. In the high temperature phase there is a periodic solution with zero action for long times corresponding essentially to the equilibrium dynamics.
The stationarity conditions on the action are equivalent to four equations on the two-time functions:
$$C^{}(\tau )=\lambda C(\tau )+2TR(\tau )+\frac{p}{2}_0^t^{}𝑑t^{\prime \prime }C^{p1}(tt^{\prime \prime })R(t^{}t^{\prime \prime })+k_0^t^{}R(tt^{\prime \prime })C^{p2}(tt^{\prime \prime })C(t^{\prime \prime }t^{})𝑑t^{\prime \prime }$$
(49)
$$R^{}(\tau )=\lambda R(\tau )+2TD(\tau )+\frac{p}{2}_0^t^{}𝑑t^{\prime \prime }C^{p1}(tt^{\prime \prime })D(t^{}t^{\prime \prime })+k_0^t^{}𝑑t^{\prime \prime }C^{p2}(tt^{\prime \prime })R(tt^{\prime \prime })R(t^{\prime \prime }t^{})+\delta (\tau )$$
(50)
$`R^{}(\tau )`$ $`=`$ $`\lambda R(\tau )+k{\displaystyle _0^t^{}}𝑑t^{\prime \prime }D(t^{}t^{\prime \prime })C^{p2}(t^{}t^{\prime \prime })C(tt^{\prime \prime })+k{\displaystyle _0^t^{}}𝑑t^{\prime \prime }C^{p2}(t^{}t^{\prime \prime })R(tt^{\prime \prime })R(t^{\prime \prime }t^{})`$ (51)
$`+`$ $`k(p2){\displaystyle _0^t^{}}𝑑t^{\prime \prime }C^{p3}(t^{}t^{\prime \prime })R(t^{}t^{\prime \prime })R(t^{\prime \prime }t^{})C(tt^{\prime \prime })\widehat{\lambda }C(tt^{})+\delta (\tau )`$ (52)
$`D^{}(\tau )`$ $`=`$ $`\lambda D(\tau )+k{\displaystyle _0^t^{}}𝑑t^{\prime \prime }D(t^{}t^{\prime \prime })R(t^{\prime \prime }t)C^{p2}(tt^{\prime \prime })+k{\displaystyle _0^t^{}}𝑑t^{\prime \prime }D(tt^{\prime \prime })C^{p2}(tt^{\prime \prime })R(t^{\prime \prime }t^{})`$ (53)
$`+`$ $`k(p2){\displaystyle _0^t^{}}𝑑t^{\prime \prime }R(tt^{\prime \prime })R(t^{\prime \prime }t)R(t^{\prime \prime }t^{})C^{p3}(tt^{\prime \prime })\widehat{\lambda }R(\tau )`$ (54)
where $`k=\frac{p(p1)}{2}`$ and we make explicitly use of the time translation invariance $`\tau =tt^{}`$. The spherical condition fix the value of $`\widehat{\lambda }`$, which can be obtained subtracting eq. (51) to eq. (50) for $`\tau =0`$:
$$\widehat{\lambda }=(p2)\left(p/2_0^\tau 𝑑t^{\prime \prime }C^{p1}(tt^{\prime \prime })D(tt^{\prime \prime })+k_0^\tau 𝑑t^{\prime \prime }R(tt^{\prime \prime })R(t^{\prime \prime }t)C^{p2}(tt^{\prime \prime })\right)2TD(0)$$
(55)
Moreover fixing the value of the energy $`E`$ gives an equation on the spherical multiplier $`\lambda `$:
$$pE=\lambda +T(R(0^+)+R(0^{}))$$
(56)
as a consequence $`E`$ and $`\lambda `$ are directly related. Using the FDT relation one can show that (49-53) reduce to a set of three independent equations on the functions: $`C(\tau )`$, $`R(\tau )+R(\tau )`$ and $`D(\tau )`$.
### A Time-reversal, non-causal solutions
Let us show that computing for very large $`t^{}`$ the trace of the Fokker-Planck operator one can recover the the number of stable states, and the dynamics within these states. The number of stable states can be obtained by a pure static computation for the p-spin spherical model ($`p>2`$) using the TAP equations .
For very large $`t^{}`$ there are two possible behaviours for the two-time functions depending on the energy (and the model) we consider:
* if at the energy value considered there are stable states then the action evaluated in the solution has a well defined limit as $`t^{}\mathrm{}`$. In this same limit, one expects that the two-time functions for finite $`\tau `$ describe the dynamics inside a stable state (as calculated previously with other methods ). A careful analysis of equations (49-53) allows one to show that the asymptotic forms of two-time quantities reads for $`\tau <<t^{}`$ :
$`C(\tau )`$ $`=`$ $`C_c(\tau )+{\displaystyle \frac{1}{t^{}}}\widehat{C}(\tau ),R(\tau )=R_c(\tau )+{\displaystyle \frac{1}{t^{}}}(\widehat{R}(\tau )+rr_c),`$ (57)
$`D(\tau )`$ $`=`$ $`{\displaystyle \frac{1}{t^{}}}D_o(\tau )+{\displaystyle \frac{1}{(t^{})^2}}\widehat{D}\left({\displaystyle \frac{\tau }{t^{}}}\right)`$ (58)
The function $`R_c(\tau )`$ is causal, and for large $`\tau `$ (but small with respect to $`t^{}`$) we have that $`R_c(\tau )0`$ and $`C_c(\tau )q`$. Together $`R_c`$ and $`C_c`$ describe the relaxation within a state. All other functions are of order one when their arguments are of order one, and tend to zero when their argument is large. The corrections of order $`1/t^{}`$ are not subleading in the computation of the trace, as this involves integrals over a large interval $`[0,t^{}]`$. The Edwards-Anderson parameter $`q`$ and $`rr_c`$ are order one constants to be determined in what follows. Note that the scaling of $`D`$ implies that $`\widehat{\lambda }`$ is of order $`1/t^{}`$. One can easily check that $`C_c(\tau )`$ and $`R_c(\tau )`$ satisfy the equations describing the equilibrium dynamics inside a state studied in :
$`C_c^{}(\tau )`$ $`=`$ $`\lambda C_c(\tau )+p/2{\displaystyle _0^t^{}}𝑑t^{\prime \prime }C_c^{p1}(tt^{\prime \prime })R_c(t^{}t^{\prime \prime })`$ (59)
$`+`$ $`k{\displaystyle _0^\mathrm{L}}R_c(tt^{\prime \prime })C_c^{p2}(tt^{\prime \prime })C_c(t^{\prime \prime }t^{})𝑑t^{\prime \prime }+p^2q^{p1}(rr_c)/2`$ (60)
$`R_c(\tau )`$ is causal and related to the correlation function through the FDT relation $`R_c(\tau )=1/TC_c^{}(\tau )\theta (\tau )`$.
* if for the energy value considered there are no stable states then the behaviour of the two-time functions is similar to the previous one except that $`D`$ has a part of order of one for finite time $`D_o`$ and a part of order or $`1/t^{}`$ for $`\tau t^{}`$:
$$D(\tau )=D_o(\tau )+\frac{1}{(t^{})}\widehat{D}\left(\frac{\tau }{t^{}}\right)$$
(61)
As a consequence $`C(\tau )`$ and $`R(\tau )`$ do not satisfy equations ‘within a state’, and one has to solve a set of three equations on $`C`$, $`R`$ and $`D`$ in which $`R`$ is not causal also for infinite times ($`t^{}`$).
In the rest of this subsection we consider an energy such that stable states exists and we compute the zero-frequency values of $`C`$, $`R`$ and $`D`$. Using the asymptotic form introduced before one finds:
$`{\displaystyle _0^t^{}}C(\tau )𝑑\tau `$ $`=`$ $`qt^{}+O(1)`$ (62)
$`{\displaystyle _0^t^{}}R_c(\tau )𝑑\tau `$ $`=`$ $`r_c={\displaystyle \frac{1q}{T}}`$ (63)
$`{\displaystyle _0^t^{}}R(\tau )𝑑\tau `$ $`=`$ $`{\displaystyle _0^t^{}}R_c(\tau )+rr_c+O(1/t^{})=r+O(1/t^{})`$ (64)
$`{\displaystyle _0^t^{}}D(\tau )𝑑\tau `$ $`=`$ $`{\displaystyle \frac{d}{t^{}}}+O(1/(t^{})^2)`$ (65)
Moreover for large $`t^{}`$ the relation (56) reduces to the usual one: $`pE=\lambda +T/2+O(1/t^{})`$. Therefore we can consider $`\lambda `$ as a fixed parameter.
Integrating eq. (49) between $`0`$ and $`t^{}`$ and taking the leading order in $`t^{}`$ we obtain:
$$q\left(\lambda +\frac{p}{2T}(1q)q^{p2}+\frac{p}{2T}(1q^{p1})+\frac{p^2}{2}(rr_c)q^{p2}\right)=0$$
(66)
Subtracting eq. (49) evaluated in $`\tau =0`$ to (66) we get the usual equation on $`\lambda `$ and $`q`$:
$$\lambda =\frac{T}{1q}+\frac{p}{2T}(1q^{p1})$$
(67)
Since the spherical multiplier is a parameter, this equation fixes the overlap. It is useful to write this equation in a way which is directly related to the static computation. Using the following notation:
$$q^{p/21}p=\lambda +\frac{p}{2T}(1q^{p1})\frac{p(p1)}{2T}q^{p2}(1q)z=(1q)q^{p/21}/T$$
(68)
where $``$ corresponds to the zero-temperature or radial energy which appear in the static computation, one can rewrite (67) as the usual static equation on $`q`$
$$1+pz+p(p1)z^2/2=0$$
(69)
From equation (66) we find the value of $`r`$:
$$r=\frac{2}{p}q^{1p/2}$$
(70)
Integrating $`(\text{50})`$ between $`0`$ and $`t^{}`$ and using eq. (66) we get the value of $`d`$:
$$q^{p1}d=2/p+4^2/p^2$$
(71)
Finally we confirm, using the relationship between $`\lambda `$ and $`q`$, that the equations on $`C_c(\tau )`$ and $`R_c(\tau )`$ (60) are indeed the same ones found in for the relaxational dynamics inside the stable state with overlap $`q`$ and energy $`E=(\lambda +T/2)/p`$.
### B The configurational entropy
To obtain the number of stable states we have to inject the solution of eqs. (49-53) into the action (43) and then take the long time limit. Using the compact notation $`Q(\tau ;t^{})=(C(\tau ;t^{}),R(\tau ;t^{}),D(\tau ;t^{}))`$ for the set of the two-time functions, we decompose the asymptotic solution as $`Q(\tau ;t^{})=Q_0(\tau ;t^{})+Q_1(\tau ;t^{})/t^{}`$, where $`Q_0`$ reads:
$$Q_0(\tau ;t^{})=(C_c(\tau ),R_c(\tau )+\frac{rr_c}{t^{}},\frac{d}{(t^{})^2})$$
(72)
and $`C_c(\tau ),R_c(\tau )`$ are the solution of the relaxational dynamics inside a stable state (60) with the values of $`r,r_c,d`$ determined above. Using this decomposition and the fact that $`Q`$ is a saddle point of (43) we find for very large $`t^{}`$:
$$S(Q)S(Q_0)\frac{1}{2}\frac{Q_1}{t^{}}\frac{\delta ^2S}{\delta Q^2}\frac{Q_1}{t^{}}+\mathrm{}t^{}>>1$$
(73)
An explicit computation shows that the second term of (73) vanishes in the long time limit provided that the corrections to $`q`$ are of an order less than $`1/t^{}`$. This seems natural to us since we are considering the relaxation dynamics inside a stable state. As a consequence the (annealed average of the) logarithm of the number of stable states coincides with $`S(Q_0)`$. When we inject $`Q_0`$ into (43) the first two lines can be easily computed and read:
$$\lambda (rr_c)+\frac{p}{4}dq^{p1}+\frac{p}{2T}(rr_c)(1q^{p1})+\frac{p(p1)}{4}(rr_c)^2q^{p2}$$
(74)
whereas the computation of the third one, which reduces only to Trln$`M/2`$, is slightly more subtle. Since the operator (48) is diagonal in Fourier space, we get:
$`{\displaystyle \frac{1}{2}}\text{Tr}\mathrm{ln}M={\displaystyle \frac{1}{2}}\mathrm{ln}(r^2qd)+{\displaystyle \frac{1}{2}}{\displaystyle \underset{\omega 0}{}}\mathrm{ln}\left(\begin{array}{cc}\widehat{R}(\omega )& \widehat{C}(\omega )\\ 0& \widehat{R}^{}(\omega )\end{array}\right),`$ (77)
where the Fourier transform of a function $`F(\tau )`$ is defined as:
$$\widehat{F}(\omega )=_0^t^{}e^{i\omega \tau }F(\tau )d\tau ,\omega =\frac{2\pi n}{t^{}}n=0,\pm 1,\pm 2,\mathrm{}$$
(78)
The function $`R_c(t,t^{})`$ is causal and the associated operator is upper triangular with diagonal elements equal to unity. Its determinant hence is one (here the Ito convention is crucial), as it should because it does not give any contribution to the action in the standard case. As a consequence we have:
$`{\displaystyle \frac{1}{2}}{\displaystyle \underset{\omega 0}{}}\mathrm{ln}\left(\begin{array}{cc}\widehat{R}(\omega )& \widehat{C}(\omega )\\ 0& \widehat{R}^{}(\omega )\end{array}\right)=\mathrm{ln}(r_c),`$ (81)
Collecting all the pieces together and using the equation on $`q,r,r_c`$ and $`d`$ obtained in section V A we find that the action $`S(Q_0)`$ reads:
$`S(Q_0)/N`$ $`=`$ $`{\displaystyle \frac{1}{2}}(1+\mathrm{ln}p\mathrm{ln}2)^2+{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{\sqrt{^2_c^2}}{_c}}\right)^2+\mathrm{ln}(\sqrt{^2_c^2})`$ (82)
$`_c`$ $`=`$ $`\sqrt{{\displaystyle \frac{2(p1)}{p}}}`$ (83)
As expected, this expression coincides with the logarithm of the number of TAP states computed by Crisanti and Sommers .
Note that this formula is correct only for $`<_c`$. For $`>_c`$ the formalism tells us that there are no stable states as follows: in this energy regime $`D`$ and $`\widehat{\lambda }`$ remain of order one even for infinite $`t^{}`$ and the action acquires a negative contribution of order $`t^{}`$. This is as it should: since for these energy values there are only metastable states with finite life-times, the longer we set $`t^{}`$ the less metastable states we find.
In general for finite dimensional glassy systems the interesting quantity will be the logarithm of the number of metastable states with finite lifetime, which can be obtained plugging the solution $`Q(\tau )`$ into the action S for a finite value of $`t^{}`$. In the following we compute this quantity at zero temperature.
### C The zero temperature case
At zero temperature the equation (50) is particularly simple. Introducing the notation $`R(\tau )=R_c(\tau )+(rr_c)/t^{}`$ we find that the equation on the Fourier transform of $`R_c`$ reads:
$$i\omega \widehat{R}_c(\omega )\lambda \widehat{R}_c(\omega )+\frac{p(p1)}{2}\widehat{R}_c(\omega )\widehat{R}_c(\omega )+1=0$$
(84)
For each frequency there are two solutions:
$$\widehat{R}_c^\pm (\omega )=\frac{1}{p(p1)}\left(\mathrm{\Lambda }\pm \sqrt{\mathrm{\Lambda }^22p(p1)}\right),\mathrm{\Lambda }=\lambda +i\omega $$
(85)
In the following we focus on the two solutions $`R_c^+(\tau )`$ and $`R_c^{}(\tau )`$ which correspond respectively to taking the Fourier transform of $`\widehat{R}_c^+(\omega )`$ and $`\widehat{R}_c^{}(\omega )`$. Using that at zero temperature $`C(\tau )=1`$ one can decompose $`S(Q)`$ in two terms such that all the dependence on $`t^{}`$ is contained only in one of them:
$`S/N(Q^\pm )`$ $`=`$ $`{\displaystyle \frac{1}{2}}(1\mathrm{ln}p\mathrm{ln}2)^2+{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{\sqrt{^2_c^2}}{_c}}\right)^2+\mathrm{ln}(\sqrt{^2_c^2})`$ (86)
$``$ $`{\displaystyle \frac{p(p1)}{4}}{\displaystyle _0^t^{}}R_c^\pm (\tau )R_c^\pm (\tau )𝑑\tau +{\displaystyle \frac{1}{2}}{\displaystyle \underset{\omega }{}}\mathrm{ln}(\widehat{R}_c^\pm (\omega )\widehat{R}_c^\pm (\omega ))`$ (87)
The computation of the second line is performed in the Appendix B. It turns out that, as in the static case , the dominant contribution is given by $`R_c^{}`$ for $`_{RSB}<<_c`$, by $`R_c^+`$ for $`_c<<_{RSB}`$ and for $`_c<<_c`$ the two saddle point contributions are the same (see Appendix B). Note that at zero temperature, $``$ is the energy density of the system. The final result is
$`S(Q)`$ $`=`$ $`\text{Re}\left({\displaystyle \frac{1}{2}}(1\mathrm{ln}p\mathrm{ln}2)^2+{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{\sqrt{^2_c^2}}{_c}}\right)^2+\mathrm{ln}(\sqrt{^2_c^2})\right)`$ (88)
$``$ $`{\displaystyle 𝑑\omega \rho _p(\omega +p)\mathrm{ln}(1\mathrm{exp}(t^{}|\omega |))}+t^{}{\displaystyle _{\mathrm{}}^0}𝑑\omega \rho _p(\omega +p)\omega `$ (89)
where $`\rho _p(x)=\sqrt{2p(p1)x^2}/(\pi p(p1))`$ is the Wigner semi-circle law. In fig. 4 we plot (88) for $`p=3`$ and different values of $`t^{}`$ as a function of the energy density $``$. For very large values of $`t^{}`$, $`S(Q)`$ converges to the logarithm of the number of stable states. Note that for a finite dimensional system we expect a similar behavior but with a vanishing curve for infinite $`t^{}`$. Finally, we remark that the formula (88) has a simple interpretation. In fact the first line coicides with the number of saddles with energy density $`E`$. Moreover, since the spectrum of the Fokker-Planck operator for an harmonic oscillator with frequency $`\omega `$ is $`E_n=(n+1/2)|\omega |\omega /2`$, the second line of (88) corresponds to the contribution due to a collection of harmonic oscillators with frequency distributed by the semicircle law centered in $`p`$. This distribution is exactly the same of the eigenvalues of the energy Hessian evaluated in saddles with energy density $``$ . As a consequence, at zero temperature, the spectrum of the Fokker-Planck operator for the p-spin spherical model coincides with the one obtained making an harmonic expansion around each saddle (also the instable ones).
## VI The two-groups Ansatz and supersymmetry breaking
Twenty years ago, when people started to search for replica symmetry breaking solutions of the Sherrigton-Kirkpatrick model, Bray and Moore proposed a two-groups Ansatz for the famous $`Q_{a,b}`$ matrix . At a first sight the results seemed a little bit strange since the $`lim_{n0}\overline{Z^n}1`$ ! ($`\overline{}`$ means the average over disorder). It turned out that in the limit of $`n`$ going to zero, the logarithm of $`\overline{Z^n}`$ equals the the logarithm of the number of TAP states, i.e. the long-time limit of the configurational entropy (29). The reason of this “coincidence” has been completely obscure until now. For instance, for many mean-field spin glass models the complexity was computed starting from TAP states and, after, it was checked that the two-groups Ansatz gave back the correct result .
Using the properties of the dynamical solutions presented in previous sections we can unveil why the Bray and Moore Ansatz allows one to calculate the long-time limit of the configurational entropy (29). In fact this Ansatz is isomorphic to zero frequency part of the dynamical calculation for $`t^{}\mathrm{}`$. This can be shown by the supersymmetric formalism for Langevin dynamics . Within this framework all the two points correlation functions between fields $`S_i(t)`$ and $`\widehat{S}_i(t)`$ can be encoded in:
$`Q(1,2)`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{N}{}}}S_i(1)S_i(2)1=(Tt_1,\overline{\theta }_1,\theta _1)`$ (90)
$`S_i(1)`$ $`=`$ $`S_i(t_1)+\overline{\theta }_1\theta _1\widehat{S}_i(t)+\overline{c}_i(t_1)\theta _1+\overline{\theta }_1c_i(t_1)`$ (91)
where $`\overline{\theta }_1,\theta _1`$ are Grassmann variables and $`\overline{c}_i(t),c_i(t)`$ are fermion fields . Using this formalism the dynamical solution (72) giving back the configurational entropy reads at large times:
$`Q(1,2)`$ $`=`$ $`Q_c(1,2)+q+(\overline{\theta }_1\theta _1+\overline{\theta }_2\theta _2){\displaystyle \frac{rr_c}{t^{}}}+\overline{\theta }_1\theta _1\overline{\theta }_2\theta _2{\displaystyle \frac{d}{(t^{})^2}}`$ (92)
$`Q_c(1,2)`$ $`=`$ $`\left(1+{\displaystyle \frac{1}{2}}(\overline{\theta }_1\overline{\theta }_2)[\theta _1+\theta _2(\theta _1\theta _2)\text{sign}(t_1t_2)]{\displaystyle \frac{1}{T}}{\displaystyle \frac{}{t_1}}\right)(C_c(t_1t_2)q)`$ (93)
The function $`Q_c(1,2)`$ is supersymmetric, whereas the last two terms in the right hand side of (92) break supersymmetry.
On the other hand, the two-groups Ansatz consists in a symmetric $`Q_{a,b}`$ matrix:
$`Q_{a,b}`$ $`=`$ $`1+{\displaystyle \frac{B}{m}}a=bm,Q_{a,b}=1{\displaystyle \frac{B}{m}}m<a=bn`$ (94)
$`Q_{a,b}`$ $`=`$ $`A+{\displaystyle \frac{B}{m}}ab;a,bm`$ (95)
$`Q_{a,b}`$ $`=`$ $`A{\displaystyle \frac{B}{m}}ab;m<a,bn`$ (96)
$`Q_{a,b}`$ $`=`$ $`A{\displaystyle \frac{C}{m^2}}am,m<bn`$ (97)
where one has to take the $`m\mathrm{}`$ and $`n0`$ limits. The functional dependence of the dynamical and replica free energy of, respectively, $`Q(1,2)`$ and $`Q_{a,b}`$ is the same . Indeed, the kinetic term in the dynamical free energy, which does not have a correspondent in the static case, is zero for the dynamical solution (92). Moreover if one puts $`B=(rr_c)T`$, $`2C=T^2d`$, $`A=q`$ the two matrices $`Q(1,2)`$ and $`Q_{a,b}`$ lead to the same results under tracing, convolution and term by term product. For instance, one can easily obtain:
$`{\displaystyle \underset{a}{}}f(Q_{a,a})`$ $`=`$ $`{\displaystyle d1f(Q(1,1))}=2B`$ (98)
$`{\displaystyle \underset{a,b}{}}f(Q_{a,b})`$ $`=`$ $`{\displaystyle d1d2f(Q(1,2))}=2(f^{}(1)f^{}(A))B+f^{\prime \prime }(A)B^2+2f^{}(A)C`$ (99)
$`\text{Tr}_{a,b}\mathrm{log}Q`$ $`=`$ $`\text{Tr}_{1,2}\mathrm{log}Q=2\mathrm{log}(1A)+\mathrm{log}((1A+B)^22AC)`$ (100)
As a consequence the computation, which make use of the two-groups Ansatz, is isomorphic to the dynamical one for $`tt^{}`$. Therefore the replica symmetry breaking scheme encoded in this Ansatz can be finally understood: it is a way to implement the dynamical computation in a replica formalism. There are, however, two important differences between the two approaches. First of all, in the dynamical computation we are not free to choose between different Ansätze the one which gives back the long-time limit of the configurational entropy (29), but we have simply to solve the equations of motion. This clearly makes the procedure inambiguous, unlike the case of the replica computation. Moreover the two approaches lead to the same results only if a dynamical solution with the correct values of $`q,d,r`$ exists. It could then happen (cf. the discussion on the configurational entropy of the SK model) that the equations on $`q,d,r`$ admit a solution but there is no dynamical solution corresponding to these values. As a consequence even if the static computation, i.e. the sum over all TAP solutions or the computation by the two-groups scheme, predicts the existence of an exponential number of stable states the more correct dynamical calculation does not.
## VII Conclusions
In this work we have shown how to put the questions related to metastable states in glasses in a manner valid for finite-dimensional systems.
We have used the construction of Gaveau and Schulman to define the metastable states. This construction requires the existence of a ‘gap’ in the lifetime, so that one can associate ‘states’ with distributions that are stable for much longer than a given time $`t^{}`$, and ‘transient processes’ those that decay much faster. There is no such gap in real glasses, so our use of this construction has to be considered partly as a definition inspired in the cases where there is.
A reasonable criterion for the relevance of any quantity will hence be that they are not too sensitive to the exact value chosen for $`t^{}`$. For example, if we consider a temperature associated with the $`\alpha `$-relaxation as $`T_t^{}`$ with $`t^{}\stackrel{<}{}t_\alpha `$, this definition is meaningful to the extent that it is stable with respect to a change in $`t^{}`$ of, say, an order of magnitude.
Next, there is the question as to whether $`T_t^{}`$ indeed reproduces the fluctuation-dissipation temperature. This and other results depend on the validity of the flatness hypotheses à la Edwards (for which positive evidence begins to appear ). In this paper we have formulated the hypothesis in a manner applicable to positive temperatures and finite dimensions (as well as to vibrated systems)— but we have not attempted to prove it. It may be, however, that writing it in the form (31), can be a good starting point for doing this. Moreover, the form (31) (and (32)) lends itself naturally to a generalisation to cases in which a system has more than two widely separated timescales and temperatures.
Finally, the computation in section V has allowed us to check the mean-field results without relying on the TAP states, themselves an intrinsically mean-field concept. The kind of solutions that dominate break causality and have positive action, but satisfy time-reversal and a non-causal form of FDT. Unlike the barrier-crossing solution of Lopatin and Ioffe , they have in this sense only some of the properties of true instantons. Moreover, the dynamical computation unveils the meaning of the two-groups Ansatz , which allows one to compute the number of stable states within a replica formalism.
Acknowledgments
We wish to thank L. Ioffe, J. L. Lebowitz and L. S. Schulman for useful suggestions.
## VIII Appendix A. Pure Barriers.
Let us show how to define ‘pure barriers’ using a supersymmetric extention of the Fokker-Planck operator.
For a system with $`N`$ degrees of freedom, introduce the $`N`$ fermion creation and annihilation operators $`a_i`$ and $`a_i^{}`$, with anticommutation relations $`[a_i,a_j^{}]_+=\delta _{ij}`$. Define the supersymmetric charges as:
$$\overline{Q}=(Tp_iiE_{,i})a_i^{};Q=p_ia_i$$
(101)
The supersymmetric operator $`H_{SUSY}`$
$$H_{SUSY}=[\overline{Q},Q]_+=H+\frac{^2E}{x_ix_j}a_j^{}a_i$$
(102)
commutes with the charges, and with the fermion number operator. In the zero-fermion subspace it coincides with the Fokker-Planck operator.
Applying the operator $`\overline{Q}`$ to all but the lowest Fokker-Planck (zero-fermion) eigenvectors, one obtains a one-fermion eigenvector. Going back to the example of low-temperature dynamics in a potential of section III, one can see that the lowest one-fermion eigenvectors correspond to distributions associated to the barriers. Indeed, one can convince oneself that in a low temperature multi-dimensional system, the lowest eigenvectors with $`k`$ fermions correspond to barriers with $`k`$ unstable directions (see Ref. , where these questions are discussed in detail, and used to derive Morse theory results).
Now, by analogy with the argument motivating the construction of pure states, it is reasonable to define ‘barrier distributions’ in general as the lowest eigenvectors of $`H_{SUSY}`$ with lowest eigenvalues in subspaces with $`k`$ fermions. This definition will make sense whenever there is a gap in the spectrum, whatever the origin of such a gap (large $`N`$, low $`T`$, etc). One can then use traces of $`e^{t^{}H_{SUSY}}`$ to calculate expectation values.
## IX Appendix B.
The aim of this section is to calculate the last line of (86). Instead of making the computation by brute force, we will use the exact results that can be obtained in the $`p=2`$ case. For $`p=2`$ the spherical model is a simple collection of harmonic oscillators with frequency distributed with a semicircle law centered around the value of the spherical multiplier. The spectrum of the Fokker-Planck operator for an harmonic oscillator with frequency $`\omega `$ is : $`E_n=(n+1/2)|\omega |\omega /2`$. Therefore for $`p=2`$ and at zero temperature, the logarithm of the trace of the evolution operator reads:
$$\overline{\mathrm{ln}(Tre^{t^{}H})}=N𝑑\omega \rho (\omega +2)\mathrm{ln}(1\mathrm{exp}(t^{}|\omega |))+N_{\mathrm{}}^0𝑑\omega \rho (\omega +2)t^{}\omega $$
(103)
where $`\rho (x)`$ is the Wigner semicircle law. One can obtain this result also by the functional formalism. In the computation of section (V C) we obtained an action which for $`p=2`$ reads:
$`S/N`$ $`=`$ $`{\displaystyle \frac{1}{2}}^2+{\displaystyle \frac{1}{2}}(\sqrt{^21})^2+\mathrm{ln}(\sqrt{^21})`$ (105)
$`{\displaystyle \frac{1}{2}}{\displaystyle \underset{n}{}}R^\pm (n)R^\pm (n)+{\displaystyle \frac{1}{2}}{\displaystyle \underset{n}{}}\mathrm{ln}(R^\pm (n)R^\pm (n))`$
where we have written $`R(t)`$ instead $`R_c(t)`$ since $`r=r_c`$ for $`p=2`$. This is directly related to the fact that there is no configurational entropy for $`p=2`$ and therefore the two-time functions relax asymptotically faster than $`1/t^{}`$. For $`<1`$ the $`p=2`$ model is a collection of stable harmonic oscillators. One can easily compute $`R_c(t)`$ starting from the result for a single oscillator and integrating over the Wigner distribution. As expected, this function coincides with $`R_c^{}(t)`$. As a consequence for $`<1`$ we expect that $`\overline{\mathrm{ln}\text{Tr}e^{t^{}H}}=NS(Q^{})`$. On the other hand for $`>1`$ the $`p=2`$ model is a collection of unstable harmonic oscillators which can mapped to the previous case changing the sign of each $`\omega `$. As a consequence for $`>1`$ we expect that $`\overline{\mathrm{ln}\text{Tr}e^{t^{}H}}=NS(Q^+)`$. In the intermediate energy regime a priori one has to consider both solutions. Since the functional computation should give back the result (103), the last line of (105) reads:
$`{\displaystyle \frac{1}{2}}{\displaystyle \underset{n}{}}R^\pm (n)R^\pm (n)+{\displaystyle \frac{1}{2}}{\displaystyle \underset{n}{}}\mathrm{ln}(R^\pm (n)R^\pm (n))=`$ (106)
$`=`$ $`{\displaystyle 𝑑\omega \rho (\omega +2)\left(t^{}\frac{\omega }{2}+\mathrm{ln}|e^{\omega t^{}/2}e^{\omega t^{}/2}|\right)}\pm i\pi {\displaystyle _{\mathrm{}}^0}𝑑\omega \rho (\omega +2)`$ (107)
if the determination of the logarithm in the first line of (105) is such that $`\mathrm{ln}1=i\pi `$. Note that for $`<1`$ all the oscillators are stable, as a consequence the last term in (106) vanishes and (105) coincides with (103). For $`>1`$ all the oscillators are unstable, as a consequence the last term in (106) equals $`i\pi `$ and cancels the $`i\pi `$ coming from the first logarithm in (105), and one obtains the same results that for $`<1`$ with an additional term $`𝑑\omega \rho (\omega +2)t^{}\omega `$. Finally, in the intermediate energy regime the first line in (105) is complex and its imaginary part cancels exactly the imaginary contribution coming from the last term in (106). In this case the the two saddle point contributions are the same, therefore to obtain (103) one can consider only one of them. However to obtain the expected value of $`R(t)`$ one has to sum on saddle points.
For $`p`$ greater than two, the equation on $`R_c`$ has the same form of (50) but for a $`p=2`$ spherical model at zero temperature with a variance of the couplings $`J^2=p(p1)/2`$. As a consequence the last line of (86) reads:
$$𝑑\omega \rho (\omega +p)\left(t^{}\frac{\omega }{2}+\mathrm{ln}|e^{\omega t^{}/2}e^{\omega t^{}/2}|\right)\pm i\pi _{\mathrm{}}^0𝑑\omega \rho (\omega +p)$$
(108)
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# 1 Some general remarks
## 1 Some general remarks
The conjectured AdS/CFT correspondence provides a rare tool for studying the strong coupling dynamics of certain gauge theories. Up to now, the correspondence has found its best application in studies of the conjectured duality between type $`IIB`$ string theory on AdS$`{}_{5}{}^{}\times `$S<sup>5</sup> describing the dynamics of $`N`$ coincident $`D3`$-branes and the large-$`N`$, large-$`g_{YM}^2N`$ limit of $`𝒩=4`$ SYM in $`d=4`$ with gauge group $`SU(N)`$. This so-called AdS<sub>5</sub>/CFT<sub>4</sub> correspondence is the prototype example for all possible dualities between various compactifications of string/$`M`$-theory and gauge field theories.
Of particular importance in establishing the AdS<sub>5</sub>/CFT<sub>4</sub> correspondence are studies of the trace and the $`SO(6)`$ $`R`$-current anomalies in $`𝒩=4`$ SYM<sub>4</sub>. Such studies have revealed, among others, the field theory interpretation of the parameter $`N`$, which in the string theory picture corresponds to the number of the coincident $`D3`$-branes, as being the dimension of the gauge group. Since the operators in the gauge theory transform under the adjoint representation of an $`SU(N)`$ gauge group, the conformal anomaly in the large-$`N`$ strong coupling regime differs from the corresponding weak-coupling value by an overall $`N^2`$ factor. Then, as the conformal anomaly is in the same supermultiplet with the $`R`$-current anomaly , the large-$`N`$ strong coupling value of the latter also differs by the same overall factor $`N^2`$ from the corresponding weak-coupling value .
Recently, there has been growing interest in studying the AdS<sub>7</sub>/CFT<sub>6</sub> correspondence as another example of the duality between string/$`M`$-theory and gauge field theory. Explicitly, this form of the correspondence is conjectured to encode the duality between $`M`$-theory compactifications on AdS$`{}_{7}{}^{}\times `$S<sup>4</sup> and the maximally supersymmetric (2,0) tensor multiplet in $`d=6`$. The former theory describes the low energy limit of $`N`$ coincident $`M5`$-branes. The latter is a mysterious, strongly coupled six-dimensional CFT without a free coupling parameter. Nevertheless, there also exists in $`d=6`$ a free (2,0) tensor multiplet which would, presumably, describe the weak-coupling regime of the above mysterious theory. In studying the AdS<sub>7</sub>/CFT<sub>6</sub> correspondence one naturally focuses on the conformal and $`R`$-symmetry anomalies of the boundary CFT<sub>6</sub> theory. Conformal anomaly studies have revealed a remarkable property : the Weyl-invariant part of the conformal anomaly in the strong-coupling regime of the (2,0) multiplet differs by an overall $`4N^3`$ factor from the corresponding weak-coupling anomaly, the latter being calculated using free-fields. The field-theoretic interpretation of the overall $`4N^3`$ factor - yet alone of the $`N`$ parameter - is far from obvious in this case and it is conceivably related to some as yet unknown realization of gauge invariance in higher dimensions.
As supersymmetry relates the trace and the $`SO(5)`$ $`R`$-current anomalies (the energy momentum and the $`R`$-current are in the same supermultiplet), it is important that studies of the trace anomaly are compatible with the well-known result for the the $`R`$-current anomaly of the (2,0) theory. In this work we present an explicit calculation of the trace anomaly of the (2,0) tensor multiplet in the presence of a background $`SO(5)`$ vector field, both using a free-field realization and also using AdS<sub>7</sub>/CFT<sub>6</sub> correspondence. In this way our calculation yields respectively the weak- and the strong-coupling results for the trace anomaly. In the next section we briefly review the structure of the trace anomaly in $`d=6`$ in the presence of background vector fields and discuss its connection to the coefficients of the two- and three-point functions of $`R`$-currents. Then we present our calculation. The free-field calculation is done using Seeley-de Witt coefficients while for the strong-coupling calculation we rely on AdS<sub>7</sub>/CFT<sub>6</sub> correspondence. We find that only one of the two possible trace anomaly structures is non-zero. We attribute the result to the maximal supersymmetry of the (2,0) tensor multiplet. The strong-coupling result differs from the weak-coupling one by an overall factor $`4N^3`$. Finally, we discuss the relevance of our result to the well-known results for the $`R`$-current anomaly of the (2,0) tensor multiplet in $`d=6`$.
## 2 Trace anomaly and $`R`$-current correlation functions in the (2,0) tensor multiplet
Let $`V_\mu (x)=V_\mu ^A(x)T^A`$, $`\mu =1,..,d`$ be a general conserved current of a $`d`$-dimensional CFT. In the case when this coincides with the $`R`$-current of the (2,0) multiplet in $`d=6`$, $`T^A`$ denote the adjoint generators of $`SO(5)`$. When the theory is coupled to a background vector field $`A_\mu ^A(x)`$, even in flat spacetime the trace of the energy momentum tensor acquires an (external) anomaly which up to total derivative terms can be written as
$`T_\mu ^\mu (x)=\alpha _VF_\mu ^{A,\nu }F_\nu ^{B,\lambda }F_\lambda ^{C,\mu }f^{ABC}+\beta _V^\mu F_{\mu \nu }^A^\lambda F_\lambda ^{A,\nu },`$ (1)
where $`F_{\mu \nu }^A(x)`$ is the standard field strength of $`A_\mu ^A(x)`$. Notice in (1) the presence of two different structures in in contrast to the $`d=4`$ case where only one structure appears. An important property of the parameters $`\alpha _V`$, $`\beta _V`$ in (1) is that they are intimately connected to the parameters appearing in the two- and three-point functions of the $`R`$-current $`V_\mu ^A(x)`$. Such a connection follows from the observation that the external trace anomaly is tied to the short distance singularities of renormalized $`n`$-point functions. In the case of interest here, assuming a coupling to the background vector field of the form $`\mathrm{d}^dx\sqrt{g}g^{\mu \nu }A_\mu ^A(x)V_\nu ^A(x)`$ we can follow and write the general renormalization group equation as
$`{\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{k!}}{\displaystyle }\mathrm{d}^6x_1\sqrt{g}g^{\mu _1\nu _1}..\mathrm{d}^6x_k\sqrt{g}g^{\mu _k\nu _k}A_{\mu _1}^{A_1}(x_1)..A_{\mu _k}^{A_k}(x_k)\mu {\displaystyle \frac{}{\mu }}V_{\nu _1}^{A_1}(x_1)..V_{\nu _k}^{A_k}(x_k)_R=`$
$`={\displaystyle \mathrm{d}^6x\sqrt{g}g^{\mu \nu }(x)T_{\mu \nu }(x)}.`$ (2)
The subscript $`R`$ in the first line of (2) denotes the renormalized $`n`$-point functions which depend on the arbitrary mass parameter $`\mu `$. Taking suitable functional derivatives of (2) with respect to $`A_\mu ^A(x)`$ we can, in principle, connect the parameters which appear in $`n`$-point functions of $`J_\mu ^A(x)`$ with the possible terms in the trace anomaly. Then, the importance of (2) is based on the fact that in a CFT the two- and three-point functions of conserved currents are determined for general $`d`$ up to a number of constant parameters. Specifically, for the case of the conserved current $`V_\mu ^A(x)`$ one has
$`V_\mu ^A(x)V_\nu ^B(y)`$ $`=`$ $`\delta ^{AB}{\displaystyle \frac{𝒞_V^{(d)}}{[(xy)^2]^{(d1)}}}I_{\mu \nu }(xy),I_{\mu \nu }(x)=\delta _{\mu \nu }2{\displaystyle \frac{x_\mu x_\nu }{x^2}},`$ (3)
$`V_\mu ^A(x)V_\nu ^B(y)V_\lambda ^C(z)`$ $`=`$ $`{\displaystyle \frac{f^{ABC}}{(xy)^{d2}(xz)^{d2}(yz)^d}}`$ (4)
$`\times I_{\nu \sigma }(xy)I_{\lambda \rho }(xz)I_{\sigma \beta }(X)t_{\mu \beta \rho }(X),`$
$`X_\mu `$ $`=`$ $`{\displaystyle \frac{(xy)_\mu }{(xy)^2}}{\displaystyle \frac{(xz)_\mu }{(xz)^2}},`$
$`t_{\mu \nu \lambda }(X)`$ $`=`$ $`𝒜^{(d)}{\displaystyle \frac{X_\mu X_\nu X_\lambda }{X^2}}+^{(d)}\left(X_\mu \delta _{\nu \lambda }+X_\nu \delta _{\mu \lambda }X_\lambda \delta _{\mu \nu }\right),`$
Therefore, had one been able to provide the explicit relation between the parameters $`𝒞_V^{(6)}`$, $`𝒜^{(6)}`$, $`^{(6)}`$ and $`\alpha _V`$, $`\beta _V`$, the calculation of the latter two would reduce to a calculation of the former three. In such a case, one could utilize the well-known AdS<sub>7</sub>/CFT<sub>6</sub> calculations for the two- and three-point functions of conserved currents to calculated the trace anomaly.
Using differential renormalization arguments, the parameter $`\beta _V`$ was evaluated in terms of $`𝒞_V^{(6)}`$ as
$`\beta _V={\displaystyle \frac{𝒞_V^{(6)}\pi ^3}{960}}.`$ (5)
The corresponding result for $`\alpha _V`$ is still missing, however on general grounds one expects it to be a linear combination of $`𝒜^{(6)}`$ and $`^{(6)}`$.
The above show that both the weak- and the strong-coupling values of the trace anomaly parameters $`\alpha _V`$ and $`\beta _V`$ for the (2,0) tensor multiplet can be obtain from the corresponding values of $`𝒞_V^{(6)}`$, $`𝒜^{(6)}`$ and $`^{(6)}`$. The weak-coupling values of the latter three parameters can be found using a free-field realization for the theory as
$`𝒞_{V,free}^{(6)}`$ $`=`$ $`{\displaystyle \frac{5}{\pi ^6}},_{free}^{(6)}={\displaystyle \frac{9}{2\pi ^9}},𝒜_{free}^{(6)}={\displaystyle \frac{3}{\pi ^9}}.`$ (6)
The strong-coupling values can be found using AdS<sub>7</sub>/CFT<sub>6</sub> correspondence which requires the consideration of maximal supergravity on AdS<sub>7</sub>.<sup>3</sup><sup>3</sup>3We consider the Euclidean version of AdS<sub>d+1</sub> space where $`\mathrm{d}\widehat{x}^i\mathrm{d}\widehat{x}_i=\frac{1}{x_0^2}(\mathrm{d}x_0\mathrm{d}x_0+\mathrm{d}x^\mu \mathrm{d}x^\mu )`$, with $`\mu =1,..,d.`$ and $`\widehat{x}_i=(x_0,x_\mu )`$. The boundary of this space is isomorphic to S<sup>d</sup> since it consists of R<sup>d</sup> at $`x_0=0`$ and a single point at $`x_0=\mathrm{}`$. The relevant part of the supergravity Lagrangian is
$`={\displaystyle \frac{1}{4g_{SG_7}^2}}{\displaystyle }\mathrm{d}^7\widehat{x}\sqrt{g}F_{ij}^A(\widehat{x})F^{A,ij}(\widehat{x}),i,j=0,..,6,`$ (7)
from where one obtains
$$𝒞_V^{(6)}=\frac{120}{\pi ^3g_{SG_7}^2},𝒜^{(6)}=\frac{72}{\pi ^6g_{SG_7}^2},^{(6)}=\frac{108}{\pi ^6g_{SG_7}^2}.$$
(8)
The important remaining piece of information is the value of $`1/g_{SG_7}^2`$ . This can be read from the general results for $`d+1=4,5,7`$, given in . One considers the equations of motion of gauged supergravity in 4, 5 and 7 dimensions correspondingly in the presence of non trivial scalar fields in the coset space $`SL(n,𝐑)/SO(n)`$. The $`SO(n)`$ group corresponds to the $`R`$-symmetry group of the boundary CFT<sub>d</sub>. The result is
$$\frac{1}{g_{SG_{d+1}}^2}=\frac{n(n2)}{4d(d1)}\frac{1}{2\kappa _{d+1}^2}=\frac{1}{2\kappa _{d+1}^2}\frac{2}{(d2)^2},$$
(9)
where we have used the relation $`n=4\frac{d1}{d2}`$ . Then, from (6) and (8) we obtain <sup>4</sup><sup>4</sup>4It is important to point out that the general result (9) is compatible with all known calculations of two- and three-point functions in $`d=3,4,6`$ (see e.g. ). In particular, for $`d=3`$ we obtain $`\frac{𝒞_V^{(3)}}{𝒞_{V,free}^{(3)}}=\frac{^{(3)}}{_{free}^{(3)}}=\frac{𝒜^{(3)}}{𝒜_{free}^{(3)}}=\frac{N^{\frac{3}{2}}4\sqrt{2}}{3\pi }`$ which shows that the ratio between the strong- and weak-coupling values for the two- and three-point functions in $`d=3`$ is $`\frac{N^{\frac{3}{2}}4\sqrt{2}}{3\pi }`$. This irrational overall factor coincides with the one found in in studies of the two- and three-point functions of the energy momentum tensor.
$`{\displaystyle \frac{𝒞_V^{(6)}}{𝒞_{V,free}^{(6)}}}`$ $`=`$ $`{\displaystyle \frac{^{(6)}}{_{free}^{(6)}}}={\displaystyle \frac{𝒜^{(6)}}{𝒜_{free}^{(6)}}}=4N^3.`$ (10)
Then, from (5) and (10) we obtain
$`\beta _V=4N^3\beta _{V,free}={\displaystyle \frac{N^3}{192\pi ^3}},`$ (11)
Our result shows that the strong-coupling value of the trace anomaly parameter $`\beta _V`$ differs from its weak-coupling value by an overall $`4N^3`$ factor.
## 3 The free-field result
In the absence of a result such as (5) for the trace anomaly parameter $`\alpha _V`$, we have to rely on some other method for calculating the trace anomaly (1) both in the weak- and also in the strong-coupling regimes of the (2,0) tensor multiplet. In each case, agreement with the result (11) would be a strong test for our calculation.
The weak-coupling values of both coefficients $`\alpha _V`$ and $`\beta _V`$ can be calculated by the method of Seeley-De Witt coefficients using a free-field realization of the (2,0) tensor multiplet. Following we can evaluate the anomaly in $`d=6`$ from the Seeley-De Witt coefficient $`b_6`$ for the general second order Laplace operator
$`\mathrm{\Delta }=^2E,`$ (12)
where $`_\mu `$ is a covariant derivative with normal bundle connection $`[_\mu ,_\nu ]=F_{\mu \nu }`$ and $`E`$ is a matrix endomorphism. The general formula for $`b_6`$ is (see and references therein)
$`b_6`$ $`=`$ $`{\displaystyle \frac{1}{(4\pi )^3}}(\alpha _6+{\displaystyle \frac{1}{6}}\alpha _2^3+\alpha _2\alpha _4),`$ (13)
$`\alpha _2`$ $`=`$ $`E,`$ (14)
$`\alpha _4`$ $`=`$ $`{\displaystyle \frac{1}{6}}^2E+{\displaystyle \frac{1}{12}}F_{\mu \nu }^2,`$ (15)
$`\alpha _6`$ $`=`$ $`{\displaystyle \frac{2}{6!}}[8(_\alpha F_{\mu \nu })^2+2(^\alpha F_{\alpha \mu })^2+12F_{\alpha \beta }^2F^{\alpha \beta }12F_\alpha {}_{}{}^{\mu }F_{\mu }^{}{}_{}{}^{\beta }F_{\beta }^{}^\alpha `$
$`+6^4E+30(_\alpha E)^2].`$
We have to evaluate $`b_6`$ for fermions and bosons taking into account the following identities
$`_\alpha F_{\mu \nu }^\alpha F^{\mu \nu }`$ $`=`$ $`2^\mu F_{\mu \alpha }^\nu F_\nu {}_{}{}^{\alpha }2F_\alpha {}_{}{}^{\mu }F_{\mu }^{}{}_{}{}^{\beta }F_{\beta }^{}{}_{}{}^{\alpha }+_\alpha J_1^\alpha ,`$ (17)
$`J_1^\alpha `$ $`=`$ $`F_{\nu \mu }^\mu F^{\nu \alpha }F_\nu {}_{}{}^{\alpha }_{\mu }^{}F^{\nu \mu }.`$ (18)
The last total derivative term in (17) is not important here as it can be cancelled by adding local counterterms , therefore we can drop it in our calculation.
For scalar bosons we have
$`F_{\mu \nu }=F_{\mu \nu }^AT^A,`$ $`E=0,Tr\left(F_{\mu \nu }F^{\mu \nu }\right)=C_\phi F_{\mu \nu }^AF^{A,\mu \nu },`$ (19)
$`Tr(F_\alpha {}_{}{}^{\mu }F_{\mu }^{}{}_{}{}^{\beta }F_{\beta }^{}{}_{}{}^{\alpha })`$ $`=`$ $`{\displaystyle \frac{1}{2}}C_\phi f^{ABC}F_\alpha ^{A,\mu }F_\mu ^{B,\beta }F_\beta ^{C,\alpha },`$ (20)
$`b_6^s`$ $`=`$ $`{\displaystyle \frac{C_\phi \left(^\alpha F_{\alpha \mu }^A\right)^2}{(4\pi )^360}}+{\displaystyle \frac{C_\phi f^{ABC}F_\alpha ^{A,\mu }F_\mu ^{B,\beta }F_\beta ^{C,\alpha }}{(4\pi )^3180}}.`$ (21)
The corresponding spinor contribution is
$`F_{\mu \nu }=F_{\mu \nu }^AT^A𝐈_\psi ,`$ $`E={\displaystyle \frac{1}{2}}F_{\mu \nu }^AT^A\gamma ^{\mu \nu },Tr\left(F_{\mu \nu }F^{\mu \nu }\right)=C_\psi (\mathrm{Tr}𝐈_\psi )F_{\mu \nu }^AF^{A,\mu \nu },`$ (22)
$`Tr(F_\alpha {}_{}{}^{\mu }F_{\mu }^{}{}_{}{}^{\beta }F_{\beta }^{}{}_{}{}^{\alpha })`$ $`=`$ $`{\displaystyle \frac{1}{2}}C_\psi (\mathrm{Tr}𝐈_\psi )f^{ABC}F_\alpha ^{A,\mu }F_\mu ^{B,\beta }F_\beta ^{C,\alpha },`$ (23)
$`b_6^f`$ $`=`$ $`{\displaystyle \frac{C_\psi (\mathrm{Tr}𝐈_\psi )\left(^\alpha F_{\alpha \mu }^A\right)^2}{(4\pi )^315}}+{\displaystyle \frac{C_\psi (\mathrm{Tr}𝐈_\psi )f^{ABC}F_\alpha ^{A,\mu }F_\mu ^{B,\beta }F_\beta ^{C,\alpha }}{(4\pi )^3180}}.`$ (24)
From (21) and (24) we can calculate the free-field result for the total trace anomaly of the $`d=6`$, $`(2,0)`$ tensor multiplet in the the presence of external vector fields and up to total derivatives terms we obtain
$`T_\mu {}_{}{}^{\mu }(x)`$ $`=`$ $`b_6^s(x)b_6^f(x)`$ (25)
$`=`$ $`{\displaystyle \frac{\left(^\alpha F_{\alpha \mu }^A\right)^2}{(4\pi )^315}}({\displaystyle \frac{C_\phi }{4}}+C_\psi (\mathrm{Tr}𝐈_\psi ))+{\displaystyle \frac{f^{ABC}F_\alpha ^{A,\mu }F_\mu ^{B,\beta }F_\beta ^{C,\alpha }}{(4\pi )^3180}}(C_\phi C_\psi (\mathrm{Tr}𝐈_\psi ))`$
$`=`$ $`{\displaystyle \frac{𝒞_{V,free}^{(6)}\pi ^3}{960}}\left(^\alpha F_{\alpha \mu }^A\right)^2+0f^{ABC}F_\alpha ^{A,\mu }F_\mu ^{B,\beta }F_\beta ^{C,\alpha }.`$
The get from the second line to the third in (25) we used the general expression for $`𝒞_{V,free}^{(6)}`$ and also the following selection rule obtained in
$`C_\phi =C_\psi (\mathrm{Tr}𝐈_\psi ),`$ (26)
for the free-field realization of the $`(2,0)`$ tensor multiplet. The value of $`\beta _V`$ read off from (25) coincides up to an overall $`4N^3`$ factor with (11), which is a consistency test for our calculation. We also obtain $`\alpha _V=0`$ as a direct result of the selection rule (26).
## 4 The strong-coupling result from AdS<sub>7</sub>/CFT<sub>6</sub> correspondence
To calculate the anomaly coefficients in the strong coupling limit we use the $`AdS`$ action (7) and follow the methods developed in . For that, we have to consider the on-shell dependence of (7) on the boundary value of the gauge field and then extract the logarithmic divergence. This is achieved by solving the equations of motion for the Yang-Mills field $`A_i(\widehat{x})A_i^AT^A=(A_0(x_0,x_\mu ),A_\mu (x_0,x_\mu ))`$ in $`AdS`$. These equations are significantly simplified by choosing to work on the gauge $`A_0=0`$, which is a natural gauge condition preserving gauge invariance on the boundary. In this gauge the equations of motions following from (7) are
$`\widehat{}_jF^{ji}`$ $``$ $`{\displaystyle \frac{1}{\sqrt{g}}}_j\left(\sqrt{g}F^{ji}\right)+[A_j,F^{ji}]=\mathrm{\hspace{0.17em}\hspace{0.17em}0},`$ (27)
$`\widehat{}_\mu F_{\mu x}`$ $`=`$ $`0,x_0^{d+1}_0\left(x_0^{d1+4}F_{x\mu }\right)+x_0^4\widehat{}_\nu F_{\nu \mu }=\mathrm{\hspace{0.17em}\hspace{0.17em}0}.`$ (28)
Then, following (see also ), we expand the vector fields around the conformal boundary in a power series as
$`A_\mu (x_0,x_\mu )`$ $`=`$ $`{\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}x_0^kA_\mu ^{(k)}(x_\mu ),F_{0\mu }=_0A_\mu ={\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}kx_0^{k1}A_\mu ^{(k)}(x_\mu ),`$ (29)
$`F_{\mu \nu }(x_0,x_\mu )`$ $`=`$ $`{\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}x_0^kF_{\mu \nu }^{(k)}(x_\mu )={\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}x_0^k\left(_\mu A_\nu ^{(k)}_\mu A_\nu ^{(k)}+{\displaystyle \underset{l=0}{\overset{k}{}}}[A_\mu ^{(l)},A_\nu ^{(kl)}]\right).`$ (30)
The boundary value of the vector field is $`A_\mu ^{(0)}(x)`$, consequently in order to be able to extract the logarithmic singularity of the action (7) we have in the following to introduce a suitable IR regulator in the $`x_0`$-integration. Here we just mention that a simple expansion such as (29) implemented with a suitable IR regularization of the $`x_0`$-integration can be easily shown to give the correct conformal anomaly for massless scalars in any dimension. Substituting (29) and (30) into (28) we obtain
$`{\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}(2+kd)kx_0^{k+2}A_\mu ^{(k)}`$ $`=`$ $`{\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}x_0^{k+3}_\nu F_{\nu \mu }^{(k1)}{\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}x_0^{k+4}{\displaystyle \underset{l=0}{\overset{k}{}}}[A_\nu ^{(l)},F_{\nu \mu }^{(kl)}],`$ (31)
$`{\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}kx_0^{k1}_\mu A_\mu ^{(k)}`$ $`=`$ $`{\displaystyle \underset{k=2}{\overset{\mathrm{}}{}}}x_0^{k1}{\displaystyle \underset{l=1}{\overset{k1}{}}}l[A_\mu ^{(kl)},A_\mu ^{(l)}],`$ (32)
where $`_\mu =_\mu +[A_\mu ^{(0)},\mathrm{}]`$. The equations (31), (32) can be recursively solved for all $`k`$. Here we consider only the relevant to us cases $`k=1,2,\mathrm{}`$ when we obtain the following set of solutions
$`A_\mu ^{(1)}`$ $`=`$ $`0,`$ (33)
$`A_\mu ^{(2)}`$ $`=`$ $`{\displaystyle \frac{1}{2(4d)}}_\nu F_{\nu \mu }^{(0)}\stackrel{\left(d=6\right)}{=}{\displaystyle \frac{1}{4}}_\nu F_{\nu \mu }(0),_\mu A_\mu ^{(2)}=0,`$ (34)
$`A_\mu ^{(3)}`$ $`=`$ $`0.`$ (35)
Plugging in (33)-(35) into (7) we can extract the logarithmic term as
$`{\displaystyle \frac{1}{4g_{SG_7}^2}}{\displaystyle _ϵ^{\mathrm{}}}{\displaystyle \frac{dx_0}{x_0}}{\displaystyle d^6x\left(8A_\mu ^{(2)}A_\mu ^{(2)}+4F_{\mu \nu }^{(0)}_\mu A_\nu ^{(2)}\right)}=\mathrm{ln}ϵ{\displaystyle \frac{1}{8g_{SG_7}^2}}{\displaystyle d^6x\left(^\mu F_{\mu \nu }^{(0)}\right)^2}.`$ (36)
In (36) we used the standard procedure to obtain the regularized generating functional for the boundary CFT<sub>6</sub> by considering the $`x_0`$-coordinate as an IR regulator in the bulk ($`x_0(\mathrm{},ϵ),ϵ0`$) which corresponds to an UV regulator in the boundary $`\mu =1/ϵ`$ dropping at the same time all bulk UV divergences. Then, from (8) we finally get by virtue of (10)
$`T_\mu ^\mu (x)_{AdS_7}=0F_\mu ^{A,\nu }F_\nu ^{B,\lambda }F_\lambda ^{C,\mu }f^{ABC}+{\displaystyle \frac{𝒞_V^{(6)}\pi ^3}{960}}\left(^\mu F_{\mu \nu }^A\right)^2.`$ (37)
## 5 Discussion
Our results (25) and (37) for the trace anomaly of the (2,0) multiplet in the presence of external gauge fields were obtained in the absence of external gravity. They are characterized both by the manifestation of the overall $`4N^3`$ as one goes from the weak-coupling (free-fields) to the strong-coupling regime and also by the vanishing of the one of the two possible structures (namely $`\alpha _V=0`$). The latter fact can be attributed, as seen from (25), to the selection rule (26) which in turn may be viewed as a manifestation of maximal supersymmetry both in the supergravity and also in the boundary CFT<sub>d</sub>.
Moreover, the results above are connected to the well-known results for the $`R`$-current anomaly . The structure of the supersymmetry algebra in $`d=6`$ is, however, quite involved and an explicit relation between the trace and the $`R`$-symmetry anomalies has not appeared in the literature as yet. Nevertheless, a separate discussion of the known results for the trace and $`R`$-current anomalies of the (2,0) multiplet might be useful.
In the absence of an external gravitational background the $`R`$-current anomaly is given by the following 8-form
$`I_8^{free}(F)`$ $`=`$ $`{\displaystyle \frac{1}{32^4}}\left[p_2(F)+{\displaystyle \frac{1}{4}}p_1(F)^2\right],`$ (38)
$`p_1(F)`$ $`=`$ $`{\displaystyle \frac{1}{2}}\mathrm{tr}\stackrel{~}{F}^2,p_2(F)={\displaystyle \frac{1}{4}}\left[\mathrm{tr}\stackrel{~}{F}^4{\displaystyle \frac{1}{2}}\mathrm{tr}\stackrel{~}{F}^2\mathrm{tr}\stackrel{~}{F}^2\right],\stackrel{~}{F}={\displaystyle \frac{\mathrm{i}}{2\pi }}F.`$ (39)
(38) gives the anomaly $`_6`$ in the six-dimensional theory via the descent equations $`d(\delta _6)=\delta I_7`$, $`I_8=\delta I_7`$. Notice that, at first sight, the descent formalism seems to give two possible linearly independent structures for the six-dimensional anomaly. This has to be compared with our trace anomaly results (25) and (37) which involve only one structure.
The $`R`$-symmetry anomaly of the strongly-coupled (2,0) multiplet has been also evaluated requiring the cancellation of the total anomaly of $`N`$ $`M5`$-branes when one takes into account the inflow anomaly . The result in the absence of external gravity is
$$I_8^{(2,0)}(F)=\frac{1}{32^4}\left[(2N^3N)p_2(F)+\frac{N}{4}p_1(F)^2\right].$$
(40)
Now, if AdS<sub>7</sub>/CFT<sub>6</sub> correspondence is valid, one should be able to recover the large-$`N`$ limit of the latter result by considering the maximally supersymmetric $`N=2`$ gauged SUGRA in $`d=7`$. Namely, one should find a result which lifted to 8-dimensions should read
$$I_8^{(AdS/CFT)}(F)=\frac{2N^3}{32^4}p_2(F).$$
(41)
One then observes that (38) and (41) seem to imply that the $`R`$-symmetry anomaly has different structure in the weak (free-fields) and the strong-coupling regimes. This in turn implies some kind of renormalization of the $`R`$-symmetry anomaly as one goes from the weak to the strong coupling regimes and it is reminiscent to the corresponding result for the trace anomaly of the (2,0) multiplet in an external gravitation background in which at least the Euler density term seems to be also renormalized . Moreover, by virtue of supersymmetry (38) and (41) seem to indicate that the corresponding result for the trace anomaly in the presence of external vector fields, but in the absence of external gravity, would also be renormalized as one goes from the weak to the strong-coupling regimes. Such a conclusion appears to be, at first sight, incompatible with our result (25) and (37) i.e. that the structure of the trace anomaly is the same in both the weak and the strong coupling regimes. It is conceivable that a better understanding of supersymmetry in $`d=6`$ might resolve this apparent puzzle.
## Acknowledgments
We wish to thank G. Arutyunov and R. Minasian for very enlightening discussions and also S. Frolov and especially A. Tseytlin for very useful correspondence. We would also like to thank the organizers of the Workshop “Duality, Strings and $`M`$-Theory” in E.S.I., Vienna, where this work was initiated, this for their worm hospitality and financial support. The work of R. M. was partially supported by Alexander von Humboldt Foundation and by the INTAS grant No:99-590. The work of A. C. P. was supported by Alexander von Humboldt Foundation.
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# Contents
## 1 Introduction
It has already been proven that there exists an algebra homomorphism between Yangian based on $`sl(N)`$ and finite $`𝒲(sl(Np),N.sl(p))`$-algebras. Such a connection plays a role in the study of physical models: for instance, in the case of the $`N`$-vectorial non-linear Schrödinger equation on the real line, the full symmetry is the Yangian $`Y(gl(N))`$, but the space of states with particle number less than $`p`$ is a representation of the $`𝒲(gl(Np),p.sl(N))`$ algebra .
The connection between Yangians and finite $`𝒲(sl(Np),N.sl(p))`$-algebras was proven in the Drinfeld presentation of the Yangian. Since the homomorphism is (obviously) not an isomorphism, it does not allow to carry the Yangian R-matrix ”down” to the finite $`𝒲`$-algebras. In this paper, we prove the correspondence in the ”RTT” presentation of the Yangian. The mentioned finite $`𝒲`$-algebras<sup>*</sup><sup>*</sup>*More precisely, it is the $`𝒲(gl(Np),N.sl(p))`$ algebras which are concerned, we will come back on this slight difference later on. appear to be ”truncation” of the Yangian, i.e. the resulting coset when modding out the Yangian ”high level” generators. These truncated Yangians were already introduced in under the name of Yangian of level $`p`$. Thanks to this presentation, we can deduce a R-matrix for the $`𝒲`$-algebras under consideration, as well as the complete classification of the finite-dimensional irreducible representations of these algebras. We also show that the Hopf structure of the Yangian cannot be carried by the homomorphism: although this is not a ”no-go theorem” for $`𝒲`$-algebras to be Hopf algebras, it severely constrains the possibilities to get this structure.
To prove our result, we need to combine three notions: Yangians, $`𝒲`$-algebras and cohomology. We have tried to be self-contained, and as such, we need to recall known results for these different fields: it is done in section 2 for Yangians, in section 3 for $`𝒲`$-algebras, and in the appendix B for cohomology. We collect our results in section 4 and then present applications in section 5. We conclude with a …conclusion, where possible generalizations and applications of our results are presented (section 6). Some calculations about $`gl(Np)`$ algebras are collected in the appendix A.
## 2 Yangians
Yangians can be seen as deformations of loop algebras (based on a simple Lie algebra) and associated to a rational solution to the Yang-Baxter equation. They have been extensively studied, and we refer to and references therein for more details. We will here focus on Yangians based on $`gl(N)`$, and recall the basic properties below.
### 2.1 The Yangian $`Y(gl(N))`$
There is essentially two presentations of $`Y(gl(N))`$: one based on generators and relations (Serre-Chevalley-type presentation), and the second (closer to integrable systems methods) using the $`R`$-matrix approach (see also and ref. therein). We use here the last one. The generators of the Yangian are gathered in a single matrix:
$$T(u)=\underset{n=0}{\overset{\mathrm{}}{}}\underset{i,j=1}{\overset{N}{}}u^nT_n^{ij}E_{ij}=\underset{n=0}{\overset{\mathrm{}}{}}u^nT_n=\underset{i,j=1}{\overset{N}{}}T^{ij}(u)E_{ij}\text{ with }T_0^{ij}=\delta ^{ij}$$
(2.1)
where $`u`$ is a spectral parameter and $`i,j`$ indices in the fundamental of $`gl(N)`$. $`E_{ij}`$ is the usual $`N`$x$`N`$ matrix with 1 at position $`(i,j)`$. The algebraic structure is encoded in the relation
$$R(uv)T_1(u)T_2(v)=T_2(v)T_1(u)R(uv)$$
(2.2)
with $`R(x)=11\frac{1}{x}P_{12}`$ and $`P_{12}`$ is the flip operator ($`P_{12}=_{i,j=1}^NE_{ij}E_{ji}`$ in representations). The commutation relations read in components
$$[T_m^{ij},T_n^{kl}]=\underset{r=0}{\overset{\text{min}(m,n)1}{}}(T_r^{kj}T_{m+nr1}^{il}T_{m+nr1}^{kj}T_r^{il})$$
(2.3)
Note that in (2.3), all the couples $`(r,s)`$, where $`s=m+n1r`$, satisfy $`s<`$min$`(m,n)`$ and $`r`$max$`(m,n)`$.
It is known that the Yangian $`Y(N)`$ is a deformation of a loop algebra based on $`gl(N)`$. The parameter $`\mathrm{}`$ can be recovered by multiplying the generators by an appropriate power of $`\mathrm{}`$:
$$T_n^{ij}\mathrm{}^{n1}T_n^{ij}$$
(2.4)
Then, the relations (2.3) can be rewritten as
$$[T_m^{ij},T_n^{kl}]=\delta ^{kj}T_{m+n1}^{il}\delta ^{il}T_{m+n1}^{kj}+o(\mathrm{})$$
(2.5)
which shows that $`Y(N)`$ is a deformation of a loop algebra (restricted to its positive modes). It can be proven that as soon as $`\mathrm{}0`$, all the Hopf algebras $`Y_{\mathrm{}}(N)`$ are isomorphic.
The Hopf structure is given by
$$\mathrm{\Delta }(T(u))=T(u)T(u);ϵ(T(u))=1;S(T(u))=T(u)$$
(2.6)
or in components:
$$\mathrm{\Delta }(T_m^{ij})=\underset{k=1}{\overset{N}{}}\underset{r=0}{\overset{m}{}}T_r^{ik}T_{mr}^{kj};ϵ(T_m^{ij})=\delta _{m,0}\delta _{i,j};S(T_m^{ij})=T_m^{ij}$$
(2.7)
For briefness, we will denote $`Y(N)Y(gl(N))`$.
### 2.2 Center of $`Y(N)`$ and associated Hopf subalgebras.
The center $`𝒟=𝒰(d_i,i)`$ of $`Y(N)`$ is generated by the quantum determinant:
$$\text{q-det}(T(u))\underset{\sigma \mathcal{\Sigma }}{}\text{sgn}(\sigma )T_{\sigma (1)1}(u)T_{\sigma (2)2}(u1)\mathrm{}T_{\sigma (N)N}(uN+1)=1+\underset{n=1}{\overset{\mathrm{}}{}}u^nd_n$$
(2.8)
The Hopf algebra $`Y(sl(N))`$ is the quotient of $`Y(N)`$ by the relation q-det$`T=1`$, i.e. $`Y(sl(N))Y(N)/𝒟`$.
We introduce
$$𝒟_r=𝒰(\{d_1,d_2,\mathrm{},d_r\})$$
It is not difficult to show that for any value of $`r`$, $`𝒟_r`$ is a Hopf ideal of $`Y(N)`$. It is obviously an algebra ideal (because it lies in the center of the Yangian), and from (2.7), one shows that
$$\mathrm{\Delta }(𝒟_r)𝒟_r𝒟_r\mathrm{\Delta }(𝒟_r)𝒟_rY(N)Y(N)𝒟_r$$
(2.9)
hence $`𝒟_r`$ is a coideal. Consequently, the coset $`Y(N)/𝒟_r`$ is also a Hopf algebra.
$$S_rY(N)=Y(N)/𝒟_r\text{ and }Y(N)S_rY(N)𝒟_r$$
(2.10)
This allows us to construct a series of Hopf subalgebras:
$`S_rY(N)=Y(N)/𝒟_r\text{ and }Y(N)S_rY(N)𝒟_rr`$
$`Y(N)S_0Y(N)S_1Y(N)\mathrm{}S_rY(N)\mathrm{}S_{\mathrm{}}Y(N)Y(sl(N))`$
where $`Y(sl(N))`$ is the only one which possesses a trivial center. The intermediate subalgebras will be of some use in the following.
### 2.3 Evaluation representations
The finite dimensional irreducible representations of $`Y(N)`$ have been classified , see also for more details. It uses the notion of evaluation representations:
###### Definition 2.1
Evaluation representations
An evaluation representation $`ev_\pi `$ is a morphism from the Yangian $`Y(gl(N))`$ to a highest weight irreducible representation $`\pi `$ of $`gl(N)`$. The morphism is given by
$$ev_\pi (T_{(1)}^{ij})=\pi (T_{(1)}^{ij})\text{ and }ev_\pi (T_{(n)}^{ij})=0,n>1$$
(2.11)
where we have identified the generators $`T_{(1)}^{ij}`$ with $`gl(N)`$ elements.
The evaluation representations form a very simple class of representations, since only one kind of Yangian generators is non-trivially represented. They are sufficient to get all finite-dimensional irreducible representations, through the tensor products of such representations:
###### Definition 2.2
Tensor product of evaluation representations
Let $`\{ev_{\pi _i}\}_{i=1,..,n}`$ be a set of evaluation representations. The tensor product of these $`n`$ representations $`ev_\stackrel{}{\pi }=ev_{\pi _1}..ev_{\pi _n}`$ is a morphism from the Yangian $`Y(gl(N))`$ to the tensor product of $`gl(N)`$ representations $`\stackrel{}{\pi }=_i\pi _i`$ given by
$$ev_\stackrel{}{\pi }(T_{(r)}^{ij})=\underset{r_1+r_2+..+r_n=r}{}\left(\underset{k=1}{\overset{n}{}}ev_{\pi _k}(T_{(r_k)}^{ij})\right)$$
(2.12)
It satisfies:
$$ev_\stackrel{}{\pi }(T_{(r)}^{ij})0\text{ if and only if }rn$$
(2.13)
Note that this definition follows from the Yangian coproduct (2.6). Tensor product of evaluation representations play an important role in the classification of finite dimensional irreducible representations of Yangians. This is reflected in the following theorems and corollary (proved in , see also for more details).
Theorem: Any finite dimensional irreducible representation of $`Y(N)`$ is highest weight and contains (up to multiplication by a scalar) a unique highest weight vector.
By highest weight vector, we mean a vector $`\eta `$ (in the representation) such that
$$\begin{array}{cc}t^{ij}(u)\eta =0\hfill & 1i<jN\hfill \\ t^{ii}(u)\eta =\lambda ^i(u)\eta \hfill & 1iN\hfill \end{array}$$
where $`\lambda ^i(u)=1+_{r>0}\lambda _{(r)}^iu^r`$, with $`\lambda _{(r)}^i`$, and $`t^{ij}(u)`$ represents $`T^{ij}(u)`$. As usual, $`\lambda (u)=(\lambda ^1(u),\mathrm{},\lambda ^N(u))`$ is called the weight of the representation.
Theorem: An irreducible highest weight representation of $`Y(N)`$ of weight $`\lambda (u)`$ is finite dimensional if and only if there exist $`(N1)`$ monic polynomials $`P_i(u)`$ such that
$$\frac{\lambda ^i(u)}{\lambda ^{i+1}(u)}=\frac{P_i(u+1)}{P_i(u)}$$
In that case, the representation is isomorphic to the subquotient of the tensor product of $`m=_im_i`$ evaluation representations, where $`m_i`$ is the degree of $`P_i(u)`$.
By monic polynomials we mean a polynomial of the form
$$P_i(u)=\underset{k=1}{\overset{m_i}{}}(u\gamma _k)\text{ with }\gamma _k$$
By subquotient, we mean the irreducible part of the highest weight submodule of the mentioned tensor product. More precisely, in the tensor product of evaluation representations (which are by definition highest weight representations), one considers the submodule generated by the tensor product of the highest weight vectors, and quotients it by all (sub)singular vectors which may appear.
Note that although generically the tensor product is irreducible (i.e. is equal to the mentioned submodule and has no singular vector), it is only for $`Y(2)`$ that it is always irreducible (see counter-example for $`Y(3)`$ in ).
A simpler characterization of the finite dimensional irreducible representations is given by the following corollary
Corollary The irreducible finite dimensional representations of $`Y(N)`$ are in one-to-one correspondence with the families $`\{P_1(u),\mathrm{},P_{N1}(u),\rho (u)\}`$ where $`P_i`$ are monic polynomials and $`\rho (u)=1+_{n>0}d_nu^n`$ encodes the values of the central elements.
### 2.4 Truncated Yangians
The notion of truncated Yangians has been already introduced in (although not named truncated, but Yangian of level $`p`$) as a tool in representation theory. They were also studied in . We now introduce the left ideal generated by $`𝒯_p=𝒰(\{T_n^{ij},n>p\})`$ :
$$_p=Y(N)𝒯_p$$
and the coset (truncation of the Yangian at order $`p`$)
$$Y(N)_p=Y(N)/_p$$
(2.14)
###### Property 2.3
The truncated Yangian $`Y(N)_p`$ is an algebra ($`N`$, $`p`$). $`\mathrm{\Delta }`$ is not a morphism of this algebra (for the structure induced by $`Y(N)`$).
Proof: We prove the Lie algebra structure of $`Y(N)_p`$ by showing that $`𝒯_p`$ is a bilateral ideal, i.e. that we have $`Y(N)_p_p`$. In fact, we will show a more stronger property, that is
$$[Y(N),𝒯_p]Y(N)𝒯_p\text{ and }[Y(N),𝒯_p]𝒯_pY(N)$$
(2.15)
We make the calculation for the first inclusion, the proof for the other inclusion being identical. Indeed, the relation (2.3) shows that $`[T_m^{ij},T_n^{kl}]`$ (for $`n>p`$) is the sum of two terms, the first being in $`Y(N)𝒯_p`$, the second belonging to $`𝒯_pY(N)`$. Focusing on the latter, one rewrites it as
$$\begin{array}{c}\underset{r=0}{\overset{\text{min}1}{}}T_{m+n1r}^{kj}T_r^{il}=\underset{r=0}{\overset{\text{min}1}{}}\left(T_r^{il}T_{m+n1r}^{kj}+\underset{s=0}{\overset{r1}{}}\left(T_s^{ij}T_{m+n2s}^{kl}T_{m+n2s}^{ij}T_s^{kl}\right)\right)\hfill \\ =\underset{r=0}{\overset{\text{min}1}{}}T_r^{il}T_{m+n1r}^{kj}+\underset{s=0}{\overset{\text{min}2}{}}(\text{min}s1)\left(T_s^{ij}T_{m+n2s}^{kl}T_{m+n2s}^{ij}T_s^{kl}\right)\hfill \end{array}$$
(2.16)
where min stands for min$`(m,n)`$. In (2.16), only the last term belongs to $`𝒯_pY(N)`$, with a summation which has one term less than the previous one: we can thus proceed recursively in a finite number of steps. The final result is an element of $`Y(N)𝒯_p`$.
As far as Hopf structure is concerned, the calculation
$$\mathrm{\Delta }(T_{p+1}^{ij})=T_{p+1}^{ij}1+1T_{p+1}^{ij}+\underset{n=1}{\overset{p}{}}T_n^{ik}T_{p+1n}^{kj}$$
shows that $`_p`$ is not a coideal, since we have
$$\mathrm{\Delta }(_p)Y(N)_p_pY(N)$$
Moreover, $`\mathrm{\Delta }`$ is not an algebra morphism anymore, since for instance
$$\mathrm{\Delta }\left([T_p^{ij},T_2^{kl}]\right)[\mathrm{\Delta }(T_p^{ij}),\mathrm{\Delta }(T_2^{kl})]=\underset{s+t=p}{}(T_{s+1}^{il}T_t^{kj}T_s^{il}T_{t+1}^{kj})0$$
(2.17)
Finally, we note that each $`Y(N)_p`$ is a deformation of a truncated loop algebra based on $`gl(N)`$. By truncated loop algebra, we mean the quotient of a usual $`gl(N)`$ loop algebra (of generators $`t_n^{ij}`$) by the relations $`t_n^{ij}=0`$ for $`n<0`$ and $`n>p`$. The construction is the same as for the complete Yangian.
### 2.5 Poisson Yangians
In the following we will deal with a Poisson version of the Yangian, where the commutator is replaced by Poisson bracket. It corresponds to the usual classical limit of quantum groups. One sets
$$T(u)=L(u);R_{12}(x)=\text{1I}+\mathrm{}r_{12}(x)+o(\mathrm{});[,]=\mathrm{}\{,\}+o(\mathrm{})$$
(2.18)
The relation (2.2) is then expanded as a series in $`\mathrm{}`$, the first non-trivial term being the $`\mathrm{}^2`$ coefficient. This new relation is the defining relation for the Poisson Yangian and reads:
$$\left\{L(u)\stackrel{}{,}L(v)\right\}=[r_{12}(uv),L(u)L(v)]\text{ with }r_{12}(x)=\frac{1}{x}P_{12}$$
(2.19)
where $`\{L(u)\stackrel{}{,}L(v)\}`$ is a matrix of component $`\{L^{ij},L^{kl}\}`$ in the basis $`E_{ij}E_{kl}`$. In components
$$\{T_m^{ij},T_n^{kl}\}=\underset{r=0}{\overset{\text{min}(m,n)1}{}}(T_r^{kj}T_{m+nr1}^{il}T_{m+nr1}^{ij}T_r^{ij})$$
(2.20)
Apart from the change from commutators to Poisson brackets (and the commutativity of the product), all the above algebraic properties still apply.
In particular, we can still define the truncated (Poisson) Yangian, with the same procedure as above.
## 3 $`𝒲`$-algebras
Such algebras can be constructed by symplectic reduction of finite dimensional Lie algebras in the same way the conformal (affine) $`𝒲`$-algebras arise as reduction of Kac-Moody (affine) Lie algebras , hence the name finite $`𝒲`$-algebras for the former . Some properties of such $`𝒲`$-algebras have been developed -. In particular, starting from a simple Lie algebra $`𝒢`$, a large class of $`𝒲`$-algebras can be seen as the commutant, in a localization of the enveloping algebra $`𝒰(𝒢)`$, of a $`𝒢`$-subalgebra . This feature has already been exploited in various physical contexts . A remarkable fact is that the involved $`𝒲`$-algebras are just of the type $`𝒲(sl(2n),n.sl(2))`$, a subclass of the $`𝒲[gl(Np),N.sl(p)]`$ algebras, in which we are interested here.
We note $`𝒲_p(N)𝒲[gl(Np),N.sl(p)]`$. This algebra is defined as the Hamiltonian reduction of the enveloping algebra of $`gl(Np)`$ (see below). In general, the $`𝒲`$-algebras are defined using semi-simple Lie algebras, but for $`gl(m)`$, we have the following property
$$𝒲[gl(m),]𝒲[sl(m)gl(1),]𝒲[sl(m),]gl(1)$$
which allows to extend the $`𝒲`$-algebra to $`gl(m)`$.
Note also that we are dealing with finite $`𝒲`$-algebra, i.e. the $`gl(m)`$ algebras we are speaking of are finite dimensional Lie algebras (not their affinization).
We use the notations introduced in the appendix A.
### 3.1 $`𝒲_p(N)`$ as an Hamiltonian reduction
Following the usual technic (see and for more details), we gather the generators of $`gl(Np)`$ in a $`(Np)\times (Np)`$ matrix:
$$𝕁=\underset{a,b=1}{\overset{N}{}}\underset{j=0}{\overset{p1}{}}\underset{m=j}{\overset{j}{}}J_{jm}^{ab}M_{ab}^{jm}$$
(3.1)
where $`M_{ab}^{jm}`$ are $`(Np)\times (Np)`$ matrices and $`J_{jm}^{ab}`$ are in the dual algebra of $`gl(Np)`$. They obey Poisson Brackets (PB) which mimic the commutation relations of $`gl(Np)`$:
$$\{J_{ab}^{j,m},J_{cd}^{\mathrm{},n}\}=\underset{r=|j\mathrm{}|}{\overset{j+\mathrm{}}{}}\underset{s=r}{\overset{r}{}}\left(\text{}\delta _{bc}<j,m;\mathrm{},n|r,s>J_{ad}^{r,s}\delta _{ad}<\mathrm{},n;j,m|r,s>J_{cb}^{r,s}\right)$$
(3.2)
On the dual algebra, we introduce first class constraints:
$$𝕁|_{const.}=ϵ_{}+\underset{a,b=1}{\overset{N}{}}\underset{j=0}{\overset{p1}{}}\underset{m=0}{\overset{j}{}}J_{jm}^{ab}M_{ab}^{jm}ϵ_{}+𝔹$$
(3.3)
Explicitly, these constraints are imposed on the negative grade generators $`J_{jm}^{ab}`$, $`m<0`$, $`j,a,b`$. They correspond to the vanishing of all these negative grade generators, but $`J_{1,1}^{00}`$ which is set to 1. We will denote them generically by $`\varphi _𝐱`$. Physically, these first class constraints generate gauge transformations, an infinitesimal form of which is:
$$\delta _\lambda J_{jm}^{ab}\underset{𝐱}{}\lambda _𝐱\{\varphi _𝐱,J_{jm}^{ab}\}$$
(3.4)
where the symbol $``$ means that one has to impose the constraints once the PB has been computed. The interesting quantities are the gauge invariant ones, and it can be shown that a way to construct a basis for them is to choose a gauge fixing for $`𝕁|_{const.}`$. In the present case, the gauge fixing is the highest weight gauge:
$$𝕁|_{g.f.}=ϵ_{}+\underset{a,b=1}{\overset{N}{}}\underset{j=0}{\overset{p1}{}}W_{jj}^{ab}M_{ab}^{jj}ϵ_{}+𝕎$$
(3.5)
where $`W_{jj}^{ab}`$ are the (unknown) generators of the gauge invariant polynomials.
In other words, there is a unique set of parameters $`\lambda _𝐱`$ such that the gauge transformations (3.4) leads $`𝕁|_{const.}`$ to $`𝕁|_{g.f.}`$. These parameters are polynomials in the original $`J_{jm}^{ab}`$, hence the generators $`W_{jj}^{ab}`$. Since they generate the gauge invariant polynomials, the $`W_{jj}^{ab}`$’s close (polynomially) under the PB: they generate the $`𝒲(gl(Np),N.sl(p))`$ algebra. The Lie algebra structure of this $`𝒲`$-algebra is given by the PB (3.2), together with the knowledge of the polynomials $`W_{jj}^{ab}`$. Unfortunately, the complete expression of these polynomials is difficult to obtain in the general case, so that different technics have been developed to compute the PB of the $`𝒲`$-algebra, without knowing the exact expression of the polynomials $`W_{jj}^{ab}`$.
There is essentially two different ways of defining the Poisson brackets of the $`𝒲_p(N)`$ algebras: through the Dirac brackets, or using the so-called soldering procedure. We will need them both, and describe them in the following.
### 3.2 Dirac brackets
It can be shown that the first class constraints together with the gauge fixing form a set of second class constraints, i.e. that if $`\mathrm{\Phi }=\{\varphi _\alpha \}_{\alpha I}`$ is the set of all constraints, we have
$$\mathrm{\Delta }_{\alpha \beta }=\{\varphi _\alpha ,\varphi _\beta \}\text{ is invertible: }\underset{\gamma I}{}\mathrm{\Delta }_{\alpha \gamma }\overline{\mathrm{\Delta }}^{\gamma \beta }=\delta _\alpha ^\beta \text{ where }\overline{\mathrm{\Delta }}^{\alpha \beta }(\mathrm{\Delta }^1)_{\alpha \beta }$$
(3.6)
Together with a set of second class constraints occurs the notion of Dirac brackets which are constructed in such a way that they are compatible with these constraints:
$$\{X,Y\}_{}\{X,Y\}\underset{\alpha ,\beta I}{}\{X,\varphi _\alpha \}\mathrm{\Delta }^{\alpha \beta }\{\varphi _\beta ,Y\}X,Y$$
(3.7)
where the symbol $``$ means that one has to apply the constraints on the right hand side once the Poisson Brackets have been computed. The compatibility of the Dirac brackets with the constraints reflects in the following property
$$\{X,\varphi _\alpha \}_{}0\alpha I,X$$
(3.8)
Then, the Poisson brackets of the $`𝒲`$-algebra are defined as the Dirac brackets of the unconstrained generators $`J_{jj}^{ab}`$:
$$\{W_j^{ab},W_{\mathrm{}}^{cd}\}\{J_{jj}^{ab},J_{\mathrm{}\mathrm{}}^{cd}\}_{}$$
(3.9)
In the case we are considering, the matrix $`\mathrm{\Delta }`$ take the form
$`\mathrm{\Delta }_{jm;k\mathrm{}}^{ab;cd}`$ $`=`$ $`\{J_{jm}^{ab},J_k\mathrm{}^{cd}\}a,b,c,d,j,k;m<j;\mathrm{}<k`$ (3.10)
$`=`$ $`(1)^m{\displaystyle \frac{j(j+1)m(m+1)}{2}}{\displaystyle \frac{\eta _j}{\eta _1}}\delta _{j,k}\delta _{m+\mathrm{}+1,0}\delta ^{bc}\delta ^{ad}+`$ (3.12)
$`+<j,m;k,\mathrm{}|t,t>\left(\delta ^{bc}J_{tt}^{ad}(1)^{j+m+k+\mathrm{}}\delta ^{ad}J_{tt}^{cb}\right)`$
$`=`$ $`(1)^m{\displaystyle \frac{j(j+1)m(m+1)}{2}}{\displaystyle \frac{\eta _j}{\eta _1}}\delta _{j,r}\delta _{m+s+1,0}\delta ^{be}\delta ^{af}\left(\mathrm{𝟏}\widehat{\mathrm{\Delta }}\right)_{k\mathrm{};rs}^{cd;ef}`$ (3.13)
$`\widehat{\mathrm{\Delta }}_{jm;k\mathrm{}}^{ab;cd}`$ $`=`$ $`{\displaystyle \frac{2\eta _1<j,m1;k,\mathrm{}|t,t>}{\eta _j(j(j+1)m(m+1))}}\left(\text{}(1)^m\delta ^{ac}J_{tt}^{bd}+(1)^{j+k+\mathrm{}}\delta ^{bd}J_{tt}^{ca}\right)`$ (3.14)
The form (3.13) shows that $`\mathrm{\Delta }`$ is invertible, for the matrix $`\widehat{\mathrm{\Delta }}`$ is nilpotent: due to the Clebsch-Gordan coefficient $`<j,m1;k,\mathrm{}|t,t>`$, we have $`(\widehat{\mathrm{\Delta }})^{2p1}=0`$. Hence, we deduce
$$\overline{\mathrm{\Delta }}_{ab;cd}^{jm;k\mathrm{}}=(1)^{m+1}\frac{j(j+1)m(m+1)}{2}\frac{\eta _1}{\eta _j}\underset{n=0}{\overset{2p1}{}}(\widehat{\mathrm{\Delta }}^n)_{j,m1;k\mathrm{}}^{ba;cd}$$
(3.15)
where we have set
$$(\widehat{\mathrm{\Delta }}^0)_{j,m;k\mathrm{}}^{ab;cd}=\delta _{j,k}\delta _{m,\mathrm{}}\delta ^{ac}\delta ^{bd}$$
(3.16)
Once $`(\mathrm{\Delta })^1`$ is known, one can compute the Dirac brackets. Unfortunately, in practice, (3.15) is difficult to achieve, and only partial results are obtained using the Dirac brackets.
### 3.3 Soldering procedure
The calculation of the Poisson brackets of the $`𝒲`$-algebra can be achieved through another way, called the soldering procedure, see also in the case of finite $`𝒲`$-algebras. It is not our aim to show the equivalence of this approach with the previous (Dirac) procedure. We give here jut a flavor of it in the context of $`𝒲_p(N)`$-algebras.
In the soldering procedure, the idea is to view the (adjoint) action of the $`𝒲_p(N)`$ algebra on itself as a ”residual” action of the whole $`gl(Np)`$ algebra on the currents, residual in the sense that is ”respects” the constraints that have been imposed. In other words, among all the transformations induced by the (enveloping algebra of) $`gl(Np)`$, we look for the ones that do not affect the form $`𝕁|_{g.f.}`$: these will be the transformations induced by the $`𝒲_p(N)`$-algebra.
In the present paper, thanks to the basis explicited in the appendix A, we will be able to synthetically present (and solve) this procedure in the case of $`𝒲_p(N)`$ algebras.
More precisely, the action of $`gl(Np)`$, with parameter $`\lambda =_{j,m;a,b}\lambda _{jm}^{ab}M_{ab}^{jm}`$, can be written
$$\delta _\lambda 𝕁=\{tr(\lambda 𝕁),𝕁\}=[\lambda ,𝕁]$$
(3.17)
where $`\{,\}`$ is the PB (on the $`J`$’s) and $`[,]`$ is the commutator (of $`Np`$x$`Np`$ matrices). Within all these transformations, we look for the ones which preserve the form of $`𝕁|_{g.f.}`$:
$$\delta _\lambda (𝕁|_{g.f.})=\delta _\lambda 𝕎=\underset{j;a,b}{}(\delta _\lambda J_{jj}^{ab})M_{ab}^{jj}$$
(3.18)
This constrains the parameters $`\lambda _{ab}^{jm}`$, and only $`N^2p`$ of them are left free: they correspond to the parameters of the $`𝒲`$-transformation.
Explicitly, the calculation $`[\lambda ,ϵ_{}+𝕎]=\delta _\lambda 𝕎`$ leads to
$`\lambda ^{j,m+1}`$ $`=`$ $`{\displaystyle \underset{k,r=0}{\overset{p1}{}}}{\displaystyle \underset{\mathrm{}=k}{\overset{k1}{}}}\left(\lambda ^k\mathrm{}W_r<k,\mathrm{};r,r|j,m>W_r\lambda ^k\mathrm{}<r,r;k,\mathrm{}|j,m>\right)`$
for $`jm<j`$ (3.19)
$`\delta _\lambda W_j`$ $`=`$ $`{\displaystyle \underset{k,r=0}{\overset{p1}{}}}{\displaystyle \underset{\mathrm{}=k}{\overset{k1}{}}}\left(\lambda ^k\mathrm{}W_r<k,\mathrm{};r,r|j,j>W_r\lambda ^k\mathrm{}<r,r;k,\mathrm{}|j,j>\right)`$ (3.20)
where $`\lambda _{j,m}=_{a,b}\lambda _{jm}^{ab}M_{ab}^{jm}`$, $`W_j=_{a,b}W_j^{ab}M_{ab}^{jj}`$ and the products are matricial products.
The system (3.19) is strictly triangular in $`\lambda _{jm}`$ with respect to the gradation $`gr(\lambda _{jm})=j+m`$. Indeed, the Clebsch-Gordan coefficients ensure that $`|jr|kj+r`$ and $`\mathrm{}+r=m`$, so that $`gr(\lambda _{k,\mathrm{}})=k+\mathrm{}j+m<j+m+1=gr(\lambda _{j,m+1})=`$ in (3.19). Thus, all the $`\lambda `$’s are expressible in terms of the $`\lambda _{j,j}`$ parameters.
#### 3.3.1 Calculation of $`\{W_0^{ab},W_j^{cd}\}`$
As a start up, we consider the variation of $`W_0`$. In that case, one has only to look at (3.20), which reads:
$`\delta _\lambda W_0`$ $`=`$ $`{\displaystyle \underset{k,r=0}{\overset{p1}{}}}\left(\lambda ^{k,k}W_r<k,k;r,r|0,0>W_r\lambda ^{k,k}<r,r;k,k|0,0>\right)`$ (3.21)
$`=`$ $`{\displaystyle \underset{k}{\overset{p1}{}}}(1)^k{\displaystyle \frac{\eta _k}{\eta _0}}[\lambda ^{k,k},W_r]`$ (3.22)
Thus, we get the equation:
$$\underset{j}{}\stackrel{~}{\lambda }_j\{W_j,W_0\}=\frac{1}{p}\underset{j}{}[\stackrel{~}{\lambda }_j,W_r]$$
(3.23)
where $`\stackrel{~}{\lambda }_j=(1)^j\eta _j\lambda _{j,j}`$. Hence, we are directly led to the PB:
$$\{W_0^{ab},W_j^{cd}\}=\frac{1}{p}(\delta ^{bc}W_j^{ad}\delta ^{ad}W_j^{cb})$$
(3.24)
#### 3.3.2 Calculation of $`\{W_1^{ab},W_j^{cd}\}`$
Now, focusing on the variation of $`W_1`$ and using the results (A.25-A.30), we are led to
$`\delta _\lambda W_1`$ $`=`$ $`{\displaystyle \underset{j}{}}c_j({\displaystyle \frac{1}{j(2j1)}}[\lambda _{j1,1j},W_j][\lambda _{j,1j},W_j]_++`$ (3.25)
$`{\displaystyle \frac{(j+1)(pj1)(p+j+1)}{2j+3}}[\lambda _{j+1,1j},W_j])`$
where $`c_j`$ has been defined in (A.25) and $`[,]`$ (resp. $`[,]_+`$) stands for the commutator (anti-commutator) of $`Np`$x$`Np`$ matrices. Then, solving the equation (3.19) for $`m=j,1j`$, and plugging the result into (3.25) gives
$`{\displaystyle \underset{j}{}}\stackrel{~}{\lambda }_j\{W_j,W_1\}`$ $`=`$ $`\frac{3}{p(p^21)}({\displaystyle \underset{j=1}{\overset{p1}{}}}{\displaystyle \frac{j(p^2j^2)}{2j+1}}[\stackrel{~}{\lambda }_{j1},W_j]+{\displaystyle \underset{j=1}{\overset{p1}{}}}{\displaystyle \underset{s=j}{\overset{p1}{}}}[[\stackrel{~}{\lambda }_s,W_{sj}],W_j]_+`$
$`+{\displaystyle \underset{j=0}{\overset{p1}{}}}{\displaystyle \underset{s=j+1}{\overset{p1}{}}}{\displaystyle \frac{sj1}{2j+1}}[[\stackrel{~}{\lambda }_{s1},W_{sj1}]_+,W_j]+`$
$`{\displaystyle \underset{j=0}{\overset{p1}{}}}{\displaystyle \underset{t=j+1}{\overset{p1}{}}}{\displaystyle \underset{s=t}{\overset{p1}{}}}{\displaystyle \frac{1}{t(2j+1)}}[[[\stackrel{~}{\lambda }_s,W_{st}],W_{t+1j}],W_j])`$
In component, we get the following PB:
$$\begin{array}{c}\{W_1^{ab},W_j^{cd}\}=\frac{3}{p(p^21)}[\text{}\frac{(j+1)(p^2(j+1)^2)}{2j+3}(\text{}\delta ^{cb}W_{j+1}^{ad}\delta ^{ad}W_{j+1}^{cb})+\hfill \\ \text{ }+j\left(\text{}\delta ^{cb}(W_0W_j)^{ad}\delta ^{ad}(W_jW_0)^{cb}+W_j^{cb}W_0^{ad}W_j^{ad}W_0^{cb}\right)+\hfill \\ \text{ }+\underset{s=1}{\overset{j}{}}(1+\frac{js}{2s+1})\left(\delta ^{cb}(W_sW_{js})^{ad}\delta ^{ad}(W_{js}W_s)^{cb}\text{}\right)+\hfill \\ \text{ }+\underset{s=1}{\overset{j}{}}(1\frac{js}{2s+1})\left(W_{js}^{ad}W_s^{cb}W_s^{ad}W_{js}^{cb}\text{}\right)+\hfill \\ \text{ }\underset{s=0}{\overset{j1}{}}\underset{t=s+1}{\overset{j}{}}\frac{1}{t(2s+1)}(\text{}\delta ^{cb}(W_sW_{ts1}W_{jt})^{ad}\delta ^{ad}(W_{jt}W_{ts1}W_s)^{cb}+\hfill \\ \text{ }+W_{jt}^{ad}(W_{ts1}W_s)^{cb}(W_sW_{ts1})^{ad}W_{jt}^{cb}+\hfill \\ \text{}\text{ }+W_{ts1}^{ad}(W_{jt}W_s)^{cb}(W_sW_{jt})^{ad}W_{ts1}^{cb}+\hfill \\ \text{ }+W_s^{ad}(W_{jt}W_{ts1})^{cb}(W_{ts1}W_{jt})^{ad}W_s^{cb}\text{})\text{}]\hfill \end{array}$$
(3.27)
## 4 Comparison between truncated Yangians and finite $`𝒲`$-algebras
We have seen that the truncated Yangians are a deformation of a truncated loop algebra based on $`gl(N)`$. We show below that $`𝒲_p(N)`$ is also a deformation of this algebra, and that these two deformations coincide. We use here the notions presented in appendix B.
We work at the classical (Poisson brackets) level.
### 4.1 $`𝒲_p(N)`$ as a deformation of a truncated loop algebra
To see that the $`𝒲_p(N)`$ is a deformation of a truncated loop algebra based on $`gl(N)`$, we modify the constraints to
$$𝕁=\frac{1}{\mathrm{}}ϵ_{}+\underset{a,b=1}{\overset{N}{}}\underset{j=0}{\overset{p1}{}}\underset{0mj}{}J_{jm}^{ab}M_{ab}^{jm}$$
(4.1)
These constraints are equivalent to the previous ones as soon as $`\mathrm{}0`$. With these new constraints, the matrix $`\mathrm{\Delta }`$ and its inverse read
$`(\mathrm{\Delta }_{\mathrm{}})_{jm;k\mathrm{}}^{ab;cd}`$ $`=`$ $`{\displaystyle \frac{1}{\mathrm{}}}(1)^m{\displaystyle \frac{j(j+1)m(m+1)}{2}}{\displaystyle \frac{\eta _j}{\eta _1}}\delta _{j,k}\delta _{m+\mathrm{}+1,0}\delta ^{bc}\delta ^{ad}\left(1\mathrm{}\widehat{\mathrm{\Delta }}_{k\mathrm{};rs}^{cd;ef}\right)`$
$`(\overline{\mathrm{\Delta }}_{\mathrm{}})_{ab;cd}^{jm;k\mathrm{}}`$ $`=`$ $`\mathrm{}(1)^{m+1}{\displaystyle \frac{j(j+1)m(m+1)}{2}}{\displaystyle \frac{\eta _1}{\eta _j}}{\displaystyle \underset{n=0}{\overset{2p1}{}}}\mathrm{}^n(\widehat{\mathrm{\Delta }}^n)_{j,m1;k\mathrm{}}^{ba;cd}`$ (4.2)
Then, computing the Dirac brackets associated to these new constraints, one finds
$`\{J_{jj}^{ab},J_{\mathrm{}\mathrm{}}^{cd}\}_{\mathrm{}}`$ $`=`$ $`\{J_{jj}^{ab},J_{\mathrm{}\mathrm{}}^{cd}\}{\displaystyle \underset{efgh;kmrs}{}}\{J_{jj}^{ab},J_{km}^{ef}\}(\overline{\mathrm{\Delta }}_{\mathrm{}})_{ef;gh}^{km;rs}\{J_{rs}^{gh},J_{\mathrm{}\mathrm{}}^{cd}\}`$ (4.3)
$`=`$ $`\delta ^{bc}J_{j+\mathrm{},j+\mathrm{}}^{ad}\delta ^{ad}J_{j+\mathrm{},j+\mathrm{}}^{cb}\mathrm{}P_{\mathrm{}}(J)`$ (4.4)
where $`P_{\mathrm{}}(J)`$ (polynomial in the $`J_{jj}^{ab}`$ which is computed using $`\overline{\mathrm{\Delta }}_{\mathrm{}}`$ as in section 3.2) has only positive (or null) powers of $`\mathrm{}`$. This clearly shows that the $`𝒲_p(N)`$ algebra is a deformation of the algebra generated by $`W_j^{ab}J_{jj}^{ab}`$ and with defining (undeformed) Poisson brackets:
$`\{W_j^{ab},W_{\mathrm{}}^{cd}\}_0`$ $`=`$ $`\delta ^{bc}W_{j+\mathrm{}}^{ad}\delta ^{ad}W_{j+\mathrm{}}^{cb}\text{ if }j+\mathrm{}<p`$ (4.5)
$`=`$ $`0\text{ if }j+\mathrm{}p`$ (4.6)
One recognizes in this algebra a (enveloping) loop algebra based on $`gl(N)`$ quotiented by the relations $`W_j^{ab}=0`$ if $`jp`$. In other words, this algebra is nothing but a truncated loop algebra, and the $`𝒲`$-algebra is a deformation of it.
### 4.2 Identification of $`𝒲_p(N)`$ and $`Y(N)_p`$
We have already seen that the truncated Yangians as well as the $`𝒲`$-algebras we consider are both a deformation of a truncated loop algebra:
$`\{W_j^{ab},W_{\mathrm{}}^{cd}\}_1`$ $`=`$ $`\{W_j^{ab},W_{\mathrm{}}^{cd}\}_0+{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}\mathrm{}^n\phi _n^W(W_j^{ab},W_{\mathrm{}}^{cd})0j,\mathrm{}p1`$ (4.7)
$`\{\overline{T}_m^{ij},\overline{T}_n^{kl}\}_2`$ $`=`$ $`\{\overline{T}_m^{ij},\overline{T}_n^{kl}\}_0+{\displaystyle \underset{r=1}{\overset{\mathrm{}}{}}}\mathrm{}^r\phi _r^T(\overline{T}_m^{ij},\overline{T}_n^{kl})0m,np1`$ (4.8)
where the cochains $`\phi _n^W`$ and $`\phi _r^T`$ obey (B.42-B.43). The undeformed PB $`\{,\}_0`$ are identical (via the identificationThe shift $`jj1`$ in the identification is due to a difference of convention between $`𝒲`$-algebras and Yangians: in the former case, the index $`j`$ denotes the underlying $`sl(2)`$ representation, while in the latter $`j`$ is the exponent of $`u`$ in the formal series (2.1). $`W_j^{ab}\overline{T}_{j1}^{ab}`$) and correspond to the truncated loop algebra. Thus, we have two deformed PB $`\{,\}_1`$ and $`\{,\}_2`$, and all we need is to show that the cochains $`\phi _n^W`$ and $`\phi _n^T`$ coincide $`n`$. To prove that it is indeed the case, we need the following properties:
###### Lemma 4.1
Let $`gl(N)_p`$ be the loop algebra based on $`gl(N)`$, truncated at order $`p`$, and $`u_j^{ab}`$ $`(j<p)`$ the corresponding generators. A 2-cocycle $`\phi `$ with values in $`𝒰(gl(N)_p)`$ is completely determined once one knows $`\phi (u_0^{ab},u_j^{cd})`$ and $`\phi (u_1^{ab},u_j^{cd})`$, $`a,b,c,d=1,\mathrm{},N`$ and $`j=0,\mathrm{},p1`$
Proof: We prove this lemma recursively. We write the cocycle condition for a triplet $`(u_1^A,u_j^B,u_k^C)`$, using indices $`A,B,C=1,\mathrm{},N^2`$ in the adjoint representation, and the commutation relations of $`gl(N)_p`$:
$`f_{}^{AB}{}_{D}{}^{}\phi (u_{1+j}^D,u_k^C)+f_{}^{BC}{}_{D}{}^{}\phi (u_{k+j}^D,u_1^A)+f_{}^{CA}{}_{D}{}^{}\phi (u_{1+k}^D,u_j^B)=`$
$`=\{u_1^A,\phi (u_j^B,u_k^C)\}+\{u_j^B,\phi (u_k^C,u_1^A)\}+\{u_k^C,\phi (u_1^A,u_j^B)\}`$ (4.9)
It can be rewritten as
$`\gamma _2\phi (u_{1+j}^A,u_k^B)`$ $`=`$ $`f_{A}^{}{}_{}{}^{CD}f_{DB}^{}{}_{}{}^{E}\phi (u_1^C,u_{k+j}^E)+f_{A}^{}{}_{}{}^{CD}f_{DB}^{}{}_{}{}^{E}\phi (u_j^C,u_{k+1}^E)+`$
$`+f_{A}^{}{}_{}{}^{CD}\left(\{u_k^B,\phi (u_j^C,u_1^D)\}+\{u_j^C,\phi (u_1^D,u_k^B)\}+\{u_1^D,\phi (u_k^B,u_j^C)\}\right)`$
where $`\gamma _20`$ is the value of the second Casimir operator in the adjoint representation.
For $`j=1`$, (4.2) allows to compute $`\phi (u_2^D,u_k^C)`$ $`C,D`$ and $`k1`$ once $`\phi (u_1^D,u_k^C)`$ $`C,D`$ and $`k`$ is known.
Suppose now that we know $`\phi (u_j^A,u_k^B)`$ for $`1j<\mathrm{}_0`$ and $`k`$. Then, (4.2) for $`j=\mathrm{}_01`$ allows to compute $`\phi (u_\mathrm{}_0^A,u_k^B)`$ $`k`$.
Thus, apart from the values $`\phi (u_0^A,u_k^B)`$ we are able to compute all the expressions $`\phi (u_j^A,u_k^B)`$. This ends the proof.
###### Property 4.2
There exist two sets of generators $`\{{}_{}{}^{\pm }\overline{W}_{j}^{ab}\}_{j=0,\mathrm{}}`$ in $`𝒲_p(N)`$ such that
$`\{{}_{}{}^{\pm }\overline{W}_{1}^{ab},{}_{}{}^{\pm }\overline{W}_{j}^{cd}\}=\delta ^{cb}{}_{}{}^{\pm }\overline{W}_{j+1}^{ad}\delta ^{ad}{}_{}{}^{\pm }\overline{W}_{j+1}^{cb}+\overline{W}_0^{cb}{}_{}{}^{\pm }\overline{W}_{j}^{ad}{}_{}{}^{\pm }\overline{W}_{j}^{cb}\overline{W}_0^{ad}`$
$`a,b,c,d=1,\mathrm{},N;j1`$ (4.11)
$`\{\overline{W}_0^{ab},{}_{}{}^{\pm }\overline{W}_{j}^{cd}\}=\delta ^{cb}{}_{}{}^{\pm }\overline{W}_{j}^{ad}\delta ^{ad}{}_{}{}^{\pm }\overline{W}_{j}^{cb}`$
The generators $`{}_{}{}^{\pm }\overline{W}_{j}^{ab}`$ are polynomial of degree $`(j+1)`$ in the original ones $`W_j^{ab}`$ and are recursively defined by
$`{}_{}{}^{\pm }\overline{W}_{j,\pm }^{ab}`$ $`=`$ $`{\displaystyle \underset{n=1}{\overset{j+1}{}}}{}_{}{}^{\pm }\overline{W}_{j,(n)}^{ab}={\displaystyle \underset{n=1}{\overset{j+1}{}}}{\displaystyle \underset{\stackrel{}{s}=j+1n}{}}{}_{}{}^{\pm }\alpha _{\stackrel{}{s}}^{n,j}(W_{s_1}\mathrm{}W_{s_n})^{ab}1<n\text{ and }1<j`$
$`{}_{}{}^{\pm }\overline{W}_{1}^{ab}`$ $`=`$ $`\pm {\displaystyle \frac{p(p^21)}{6}}W_1^{ab}+{\displaystyle \frac{p(p\pm 1)}{2}}(W_0W_0)^{ab}`$ (4.12)
$`\overline{W}_0^{ab}`$ $``$ $`{}_{}{}^{+}\overline{W}_{0}^{ab}={}_{}{}^{}\overline{W}_{0}^{ab}=pW_0^{ab}`$
for some numbers $`{}_{}{}^{\pm }\alpha _{\stackrel{}{s}}^{n,j}`$. The summation on $`\stackrel{}{s}`$ is understood as a summation on $`n`$ positive (or null) integers $`(s_1,\mathrm{},s_n)\stackrel{}{s}`$ such that $`|\stackrel{}{s}|=_{i=1}^ms_i=j+1n`$.
The subsets $`\{{}_{}{}^{\pm }\overline{W}_{j}^{ab}\}_{j=0,\mathrm{},p1}`$ form two bases of $`𝒲_p(N)`$, the other generators $`\{{}_{}{}^{\pm }\overline{W}_{j}^{ab}\}_{jp}`$ been polynomials in the basis elements.
Proof: We first remark that the form (4.12) clearly shows that the $`p`$ first generators are independent, and thus form a basis. The other ones must then be polynomials in any basis.
We prove the relation (4.11) by a recursion on $`j`$. It is easy to compute that the definitions are such that (4.11) is satisfied for $`j=1`$. For the recursion, we fix a basis $`{}_{}{}^{+}\overline{W}_{j}^{ab}`$ or $`{}_{}{}^{}\overline{W}_{j}^{ab}`$ (the proof is obviously independent of the choice), and write it $`\overline{W}_j^{ab}`$.
We suppose that we have found generators $`\overline{W}_j^{ab}`$ for $`jj_0`$ such that (4.11) is satisfied. This implies that we have:
$$N\overline{W}_{j_0+1}^{cd}=\{\overline{W}_1^{ca},\overline{W}_{j_0}^{ad}\}\overline{W}_0^{aa}\overline{W}_{j_0}^{cd}+\overline{W}_0^{cd}\overline{W}_{j_0}^{aa}+\delta _{cd}\overline{W}_{j_0+1}^{aa}$$
where we have used implicit summation on repeated $`gl(N)`$ indices. Then, we get
$`N\{\overline{W}_1^{ab},\overline{W}_{j_0+1}^{cd}\}`$ $`=`$ $`\{\overline{W}_1^{ab},\{\overline{W}_1^{ce},\overline{W}_{j_0}^{ed}\}\}\{\overline{W}_1^{ab},\overline{W}_0^{ee}\overline{W}_{j_0}^{cd}\overline{W}_0^{cd}\overline{W}_{j_0}^{ee}\}+\delta _{cd}\{\overline{W}_1^{ab},\overline{W}_{j_0+1}^{aa}\}`$
$`=`$ $`\{\{\overline{W}_1^{ab},\overline{W}_1^{ce}\},\overline{W}_{j_0}^{ed}\}+\{\overline{W}_1^{ce},\{\overline{W}_1^{ab},\overline{W}_{j_0}^{ed}\}\}+\delta _{cd}\{\overline{W}_1^{ab},\overline{W}_{j_0+1}^{aa}\}+`$
$`\overline{W}_0^{ee}\{\overline{W}_1^{ab},\overline{W}_{j_0}^{cd}\}+\overline{W}_0^{cd}\{\overline{W}_1^{ab},\overline{W}_{j_0}^{ee}\}+\overline{W}_{j_0}^{ee}\{\overline{W}_1^{ab},\overline{W}_0^{cd}\}`$
$`=`$ $`\{\overline{W}_1^{cb},\overline{W}_{j_0+1}^{ad}\}\{\overline{W}_2^{cb},\overline{W}_{j_0}^{ad}\}+\overline{W}_1^{ab}\overline{W}_{j_0}^{cd}\overline{W}_1^{cd}\overline{W}_{j_0}^{ab}`$
$`+N\left(\overline{W}_0^{cb}\overline{W}_{j_0+1}^{ad}\overline{W}_0^{ad}\overline{W}_{j_0+1}^{cb}\right)+\delta ^{cb}A^{ad}\delta ^{ad}B^{cb}+\delta ^{cd}C^{ab}`$
with the notation
$`A^{ad}`$ $`=`$ $`\{\overline{W}_2^{ae},\overline{W}_{j_0}^{ed}\}+\overline{W}_0^{ad}\overline{W}_{j_0+1}^{ee}\overline{W}_0^{ee}\overline{W}_{j_0+1}^{ad}+\overline{W}_1^{ad}\overline{W}_{j_0}^{ee}\overline{W}_1^{ee}\overline{W}_{j_0}^{ad}`$
$`B^{cb}`$ $`=`$ $`\{\overline{W}_1^{ce},\overline{W}_{j_0+1}^{eb}\}+\overline{W}_0^{cb}\overline{W}_{j_0+1}^{ee}\overline{W}_0^{ee}\overline{W}_{j_0+1}^{cb}`$
$`C^{ab}`$ $`=`$ $`\{\overline{W}_1^{ab},\overline{W}_{j_0}^{ee}\}+[\overline{W}_0,\overline{W}_{j_0+1}]^{ab}`$
It remains to compute $`\{\overline{W}_2^{cb},\overline{W}_{j_0}^{ad}\}`$. This is done using the same technics as above:
$$N\overline{W}_2^{cb}=\{\overline{W}_1^{ce},\overline{W}_1^{eb}\}\overline{W}_0^{ee}\overline{W}_1^{cb}+\overline{W}_0^{cb}\overline{W}_1^{ee}+\delta ^{cb}\overline{W}_2^{ee}$$
so that we have
$`N\{\overline{W}_2^{cb},\overline{W}_{j_0}^{ad}\}`$ $`=`$ $`\left(\text{}\{\overline{W}_1^{ab},\overline{W}_{j_0+1}^{cd}\}+\{\overline{W}_1^{cd},\overline{W}_{j_0+1}^{ab}\}\right)+N\left(\text{}\overline{W}_1^{ab}\overline{W}_{j_0}^{cd}\overline{W}_1^{cd}\overline{W}_{j_0}^{ab}\right)+`$
$`+N\left(\text{}\overline{W}_0^{ab}\overline{W}_{j_0+1}^{cd}\overline{W}_0^{cd}\overline{W}_{j_0+1}^{ab}\right)+\delta ^{ab}B^{cd}+`$
$`+\delta ^{cd}\left(\text{}\{\overline{W}_{j_0+1}^{ae},\overline{W}_1^{eb}\}+\overline{W}_0^{ee}\overline{W}_{j_0+1}^{ab}\overline{W}_0^{ab}\overline{W}_{j_0+1}^{ee}\right)+`$
$`+\delta ^{cb}\left(\text{}\{\overline{W}_2^{ee},\overline{W}_{j_0}^{ad}\}[\overline{W}_0,\overline{W}_{j_0+1}]^{ad}[\overline{W}_1,\overline{W}_{j_0}]^{ad}\right)`$
Then, a recurrent use of these two brackets leads to the result:
$`\{\overline{W}_1^{ab},\overline{W}_{j_0+1}^{cd}\}`$ $`=`$ $`\overline{W}_0^{cb}\overline{W}_{j_0+1}^{ad}\overline{W}_0^{ad}\overline{W}_{j_0+1}^{cb}+\delta ^{cb}\overline{W}_{j_0+2}^{ad}\delta ^{ad}\overline{W}_{j_0+2}^{cb}+`$
$`{\displaystyle \frac{\delta ^{ad}\delta ^{cb}}{N(N^21)}}\left(\{\overline{W}_1^{ee},\overline{W}_{j_0+1}^{ff}\}N\{\overline{W}_1^{ef},\overline{W}_{j_0+1}^{fe}\}\right)+`$
$`+{\displaystyle \frac{\delta ^{cd}}{N}}(\text{}\{\overline{W}_1^{ab},\overline{W}_{j_0+1}^{ee}\}+[\overline{W}_0,\overline{W}_{j_0+1}]^{ab}+`$
$`+{\displaystyle \frac{\delta ^{ab}}{N(N^21)}}(\text{}\{\overline{W}_1^{ee},\overline{W}_{j_0+1}^{ff}\}N\{\overline{W}_1^{ef},\overline{W}_{j_0+1}^{fe}\}))`$
for some polynomials $`\overline{W}_{j_0+2}^{ad}`$.
Finally, we remark that the forms (4.12) and the PB (3.27) clearly show that the PB $`\{\overline{W}_1^{ab},\overline{W}_j^{cd}\}`$ does not contain terms proportional to $`\delta ^{ab}`$ or $`\delta ^{cd}`$. Moreover, a direct calculation, using (3.27), shows that
$$\{\overline{W}_1^{aa},P_j^{bb}\}=0,P_j^{cd}(W)=\underset{n=1}{\overset{j+1}{}}\underset{\stackrel{}{s}=j+1n}{}\beta _\stackrel{}{s}^{n,j}(W_{s_1}\mathrm{}W_{s_n})^{cd}$$
This is enough to show that the two last lines in the PB (4.2) identically vanish.
Hence, we can deduce that the PB must be of the form
$$\{\overline{W}_1^{ab},\overline{W}_{j_0+1}^{cd}\}=\overline{W}_0^{cb}\overline{W}_{j_0+1}^{ad}\overline{W}_0^{ad}\overline{W}_{j_0+1}^{cb}+\delta ^{cb}\overline{W}_{j_0+2}^{ad}\delta ^{ad}\overline{W}_{j_0+2}^{cb}$$
(4.15)
which is exactly (4.11), so that the recursion on $`j`$ is proven.
We have computed the first and last terms ($`j0`$) that appear in the definition (4.12):
$`{}_{}{}^{\pm }\overline{W}_{j,(1)}^{ab}`$ $`=`$ $`(\pm 1)^j(j!)^2\left({\displaystyle \genfrac{}{}{0pt}{}{p+j}{2j+1}}\right)W_j^{ab}`$ (4.16)
$`{}_{}{}^{}\overline{W}_{j,(j+1)}^{ab}`$ $`=`$ $`\left({\displaystyle \genfrac{}{}{0pt}{}{p}{j+1}}\right)(\underset{j+1}{\underset{}{W_0\mathrm{}W_0}})^{ab}`$ (4.17)
$`{}_{}{}^{+}\overline{W}_{j,(j+1)}^{ab}`$ $`=`$ $`\left({\displaystyle \genfrac{}{}{0pt}{}{p+j}{j+1}}\right)(\underset{j+1}{\underset{}{W_0\mathrm{}W_0}})^{ab}`$ (4.18)
###### Corollary 4.3
The change of generators between $`\{{}_{}{}^{+}\overline{W}_{j}^{ab}\}`$ and $`\{{}_{}{}^{}\overline{W}_{j}^{ab}\}`$ is given by:
$${}_{}{}^{\pm }\overline{W}_{j}^{ab}=\underset{n=1}{\overset{j}{}}(1)^{j+n+1}\underset{\stackrel{}{s}=j+1n}{}({}_{}{}^{}\overline{W}_{s_1}^{}\mathrm{}{}_{}{}^{}\overline{W}_{s_n}^{})^{ab}$$
(4.19)
Proof: Using the expression (4.12) for $`j=1`$ and the PB (4.11), one computes that
$$\{{}_{}{}^{\pm }\overline{W}_{1}^{ab},{}_{}{}^{}\overline{W}_{j}^{cd}\}=\delta ^{bc}\left(\text{}(\overline{W}_0{}_{}{}^{}\overline{W}_{j}^{})^{ad}{}_{}{}^{}\overline{W}_{j+1}^{ad}\right)\delta ^{ad}\left(\text{}({}_{}{}^{}\overline{W}_{j}^{}\overline{W}_0)^{cb}{}_{}{}^{}\overline{W}_{j+1}^{cb}\right)$$
(4.20)
Then, a direct calculation shows that indeed the expression (4.19) satisfies (4.11).
###### Corollary 4.4
The basis $`\{{}_{}{}^{}\overline{W}_{j}^{ab}\}`$ is such that $`{}_{}{}^{}\overline{W}_{j}^{ab}=0`$ for $`jp`$. In the basis $`\{{}_{}{}^{+}\overline{W}_{j}^{ab}\}`$, all the $`{}_{}{}^{+}\overline{W}_{j}^{ab}`$ generators ($`jp`$) are not vanishing.
Proof: From the PB (4.11) it is clear that it is sufficient to show that $`\overline{W}_p^{ab}=0`$. Writing this PB for $`j=p`$ and using the form
$$\overline{W}_p^{ab}=\underset{n=2}{\overset{p+1}{}}\underset{\stackrel{}{s}=p+1n}{}\alpha _\stackrel{}{s}^{n,p}(\overline{W}_{s_1}\mathrm{}\overline{W}_{s_n})^{ab}$$
(4.21)
one gets only two possibilities for the $`\alpha `$’s:
$$\alpha _\stackrel{}{s}^{n,p}=(1)^nA\text{ with }A=0\text{ or }1$$
(4.22)
If $`A=0`$, then $`\overline{W}_p^{ab}=0`$ while if $`A=1`$, the change of basis given in the corollary 4.19 shows that in the other basis we have $`\overline{W}_p^{ab}=0`$. Hence, we have to determine which basis corresponds to $`\overline{W}_p^{ab}=0`$.
Looking at the expressions (4.18) and (4.17), one concludes that $`{}_{}{}^{}\overline{W}_{p}^{ab}=0`$ while $`{}_{}{}^{+}\overline{W}_{j}^{ab}0,j`$.
In the following, we choose for $`𝒲_p(N)`$ the $`\{{}_{}{}^{}\overline{W}_{j}^{ab}\}`$ basis and omit the superscript $``$ for the generators. Now, from above, it is easy to show:
###### Theorem 4.5
The $`𝒲_p(N)`$ algebra is the truncated Yangian $`Y(N)_p`$.
Proof: Let us first remark that the two algebras have identical (in fact undeformed) PB on the couples $`(\overline{W}_0^{ab},\overline{W}_j^{cd})`$, which proves that the cochains $`\phi _n^W`$ and $`\phi _n^T`$ coincide (in fact vanish) on these points.
Moreover, the property 4.2 shows that the cochains $`\phi _n^W`$ and $`\phi _n^T`$ coincide on the couples $`(\overline{W}_1^{ab},\overline{W}_j^{cd})`$. Since $`\phi _1`$ is a cocycle, this is enough (using lemma 4.1) to prove that $`\phi _1^T`$ and $`\phi _1^W`$ are identical.
Now, suppose that we have proven that $`\phi _n^W`$ and $`\phi _n^T`$ are identical for $`n<n_0`$. Then, eq. (B.43) fixes $`\phi _{n_0}^W`$ and $`\phi _{n_0}^T`$, up to a cocycle:
$`\phi _{n_0}^W`$ $`=`$ $`\phi _{n_0}+\xi _{n_0}^W`$
$`\phi _{n_0}^T`$ $`=`$ $`\phi _{n_0}+\xi _{n_0}^T`$
where $`\phi _{n_0}`$ is a function of the cochains $`\phi _n^W=\phi _n^T`$, $`n<n_0`$. But property 4.2 shows that the two cocycles $`\xi _{n_0}^W`$ and $`\xi _{n_0}^T`$ coincide on the couples $`(\overline{W}_1^{ab},\overline{W}_j^{cd})`$, which proves that they are identical (due to lemma 4.1).
Thus, $`\phi _{n_0}^W`$ and $`\phi _{n_0}^T`$ are identical, and we have proven recursively the property.
#### 4.2.1 Quantization
We have shown that truncated Yangians and $`𝒲`$-algebras coincide at the classical level. It remains to show that it is still true at the quantum level. Fortunately, an algebra morphism between Yangians and $`𝒲`$-algebras has already been given in , at classical and quantum levels. This relation was not sufficient to establish the identification between $`𝒲`$-algebras and truncated Yangians, since all the horizontal arrows involved in the diagram
$$\begin{array}{ccc}Y(N)& & 𝒲_p(N)\\ & & \mathrm{?}\\ Y(N)& & Y_p(N)\end{array}$$
(4.23)
are not isomorphisms. Hence the calculations done in this paper.
However, once the relation (between $`Y_p(N)`$ and $`𝒲_p(N)`$) has been established at the classical level, we can use the result of to promote it at the quantum level. More precisely, now that we can identify the $`𝒲_p(N)`$ algebra with $`Y_p(N)`$ at the classical level, we can use the results of at the quantum level: it has been established that any quantization of $`𝒲_p(N)`$ still obey to the Drinfeld relation, and hence the homomorphism still exists at the quantum level.
Thus, theorem 4.5 is valid both at classical and quantum level, and the figure 4.23 is correct (without question mark).
Let us remark that in the proof we have establish, we have constructed $`𝒲`$-algebras as deformations of a truncated loop algebra and identified them with the truncated Yangians,i.e. truncations of deformed loop algebras. Denoting by $`(gl(N))`$ the loop algebra defined on $`gl(N)`$, and by $`(gl(N))_p`$ its truncation, the above sentence can be pictured as the following commutative diagram:
$$\begin{array}{ccccc}& & Y(N)& & \\ & _{\mathrm{}}& & ^p& \\ (gl(N))& & & & Y_p(N)𝒲_p(N)\\ & ^p& & _{\mathrm{}}& \\ & & (gl(N))_p& & \end{array}$$
(4.24)
where $`_{\mathrm{}}`$ stands for a deformation, and $`^p`$ for a truncation (at level $`p`$).
## 5 Applications
### 5.1 $`R`$-matrix for $`𝒲`$-algebras
The above construction allows us to associate the $`𝒲`$-algebras to the $`R`$-matrix of the Yangian, the difference between these two algebras lying in the modes development of $`T(u)`$: in both cases, the development is done in powers of $`u^1`$, but for the Yangian it is an infinite series, while the development is truncated to a polynomial for the $`𝒲`$-algebra. Explicitly, the presentation of the $`𝒲_p(N)`$-algebra take the form:
$$R(uv)T_1(u)T_2(v)=T_2(v)T_1(u)R(uv)\text{ with }\{\begin{array}{c}T(u)=1+\underset{n=1}{\overset{p}{}}\underset{a,b=1}{\overset{N^2}{}}u^nE_{ab}T_n^{ab}\hfill \\ R(x)=11\frac{1}{x}P_{12}\hfill \end{array}$$
Let us remark that this procedure is similar to the ”factorization procedure” which leads from the elliptic algebra $`𝒜_{q,p}(N)`$ to the Sklyanin algebra $`S_{q,p}(N)`$ (see also for more examples about factorizations). In all cases, one chooses for $`T(u)`$ a special dependence in $`u`$ to get a finite algebra: this special dependence is nothing but a coset by some of the modes of $`T(u)`$. In all the examples, the Hopf structure of the starting algebra does not survive to this quotient.
Note that the $`R`$-matrix presentation of the $`𝒲`$-algebras provides an exhaustive set of commutation relations among the $`𝒲_p(N)`$ generators for generic $`N`$ and $`p`$, while, up to now, a complete set of commutation relation was known only for a small number of $`𝒲`$-algebras.
Let us also remark that the $`R`$-matrix presentation allows to define the $`𝒲`$-algebras without any reference to the underlying $`gl(Np)`$ algebra, and thus is a more ”abstract” definition.
### 5.2 Irreducible representations of $`𝒲_p(N)`$-algebras
Once again, the $`R`$-matrix presentation provides a very natural framework for the classification of $`𝒲`$-representationsWe thank P. Sorba for drawing our attention to this point.. It is based on the notion of evaluation representations, as it appears in the Yangian context (see section 2.3). In fact, this classification was done in , in the context of (truncated) Yangians. We have the following theorem:
###### Theorem 5.1
Finite dimensional irreducible representations of $`𝒲_p(N)`$
Any finite dimensional irreducible representation of the $`𝒲[gl(Np),N.gl(p)]`$ algebra is isomorphic to an evaluation representation or to the subquotient of tensor product of at most $`p`$ evaluation representations.
Proof:
By evaluation representations for $`𝒲_p(N)`$ algebra, we mean the definitions 2.1 and 2.2 with the change $`T_r^{ab}W_{r1}^{ab}`$ (i.e. the evaluation representations of the truncated Yangian). The property (2.13) clearly shows that the (subquotient of) tensor product of $`n`$ evaluation representations is a representation of the truncated Yangian as soon as $`np`$. It also shows that if it is irreducible for the Yangian, then it is also irreducible for the truncated Yangian and that they are finite dimensional.
Now conversely, an irreducible representation $`\pi `$ of the $`𝒲_p(N)`$ algebra can be lifted to a representation of the whole Yangian by setting $`\pi (T_{(r)}^{ij})=0`$ for $`r>n`$. It is then obviously irreducible for the Yangian, and thus is isomorphic to the tensor product of evaluation representations.
We remark that the theorem 5.1 allows to construct any (finite dimensional) representation of $`𝒲_p(N)`$ in term of $`p`$ representations of $`gl(N)`$ (including trivial representations). This is exactly what one obtains from the so-called ”Miura transformation” that appears in the context of $`𝒲`$-algebras. Indeed, this transformation allows to construct a representation of the $`𝒲(𝒢,)`$-algebra using representations of $`𝒢_0`$, the zero-grade subalgebra of $`𝒢`$. In the case of $`𝒲_p(N)`$, we get $`𝒢_0=N.gl(p)`$, and hence need $`N`$ representations of $`gl(p)`$, as it is stated in theorem 5.1.
Finally, as for Yangians, we have the following characterization (proved using above theorem and the characterization for Yangians):
###### Corollary 5.2
The irreducible finite-dimensional representations of $`𝒲_p(N)`$ are in one-to-one correspondence with the families $`\{P_1(u),\mathrm{},P_{N1}(u),\rho (u)\}`$ where $`P_i`$ are monic polynomials of degree $`m_i`$ such that $`_im_ip`$, and $`\rho (u)=1+_{n>0}d_nu^n`$.
### 5.3 Generalization to $`S_rY(N)_p`$ truncated Yangians.
As well as we have defined truncated Yangians based on $`Y(N)`$, the same construction can be done for each of the $`S_rY(N)`$ Hopf algebras, to construct $`S_rY(N)_p`$ algebras ($`rp`$): these truncated Yangians will correspond to the quotient of $`𝒲_p(N)`$ by $`𝒟_r`$, which is a part of the $`𝒲_p(N)`$-center (see below).
In particular, the $`𝒲(sl(Np),N.sl(p))`$ algebra usually encountered in the literature is nothing but the truncation $`S_1Y(N)_p`$. For this algebra, one sees that there exist two algebra homomorphisms: $`Y(gl(N))𝒲(sl(Np),N.sl(p))`$ and $`Y(sl(N))𝒲(sl(Np),N.sl(p))`$. The second one corresponds to the case given in .
More generally, we have the following property:
###### Property 5.3
There is an algebra homomorphism from $`S_rY(N)_{p+q}`$ to $`S_{r+s}Y(N)_p`$, for any values of $`p,q,r,s=0,1,\mathrm{},\mathrm{}`$.
Proof: It is a trivial composition of algebra and Hopf algebra homomorphisms, as it is visualized in figure 1.
#### 5.3.1 Finite dimensional irreducible representations of $`S_rY(N)_p`$ algebras
Starting from the theorem 5.1 and using cosets by central elements, it is easy to get
###### Corollary 5.4
Any finite-dimensional irreducible representation of the $`S_rY(N)_p`$ algebra is obtained from the subquotient of tensor product of at most $`p`$ evaluation representations, quotiented by $`r`$ constraints on the generators of $`𝒟_r`$.
The finite-dimensional irreducible representations of the $`S_rY(N)_p`$ algebra are in one-to-one correspondence with the families $`\{P_1(u),\mathrm{},P_{N1}(u),\rho (u)\}`$ where $`P_i`$ are monic polynomials of degree $`m_i`$ such that $`_im_ip`$, and $`\rho (u)=1+_{n=0}^rd_nu^n`$.
In particular, in the case of the $`𝒲(sl(Np),N.sl(p))`$ algebra, we obtain the result given in for $`N=2`$.
### 5.4 Center of $`𝒲_p(N)`$ algebras.
From the definition of $`𝒲(gl(Np),N.sl(p))`$ algebras, one already knows that their center contains the Casimir operators of $`gl(Np)`$, since, being central, these operators are obviously gauge invariant. Hence the dimension of the center is at least $`Np`$. However, it was not proved (to our knowledge) that its dimension is exactly $`Np`$. Fortunately, the center of the truncated Yangians $`Y_p(N)`$ has been determined in :
Property: A basis of the $`Y_p(N)`$ center is given by all the coefficients of the principal part of the following generating function
$$H(x)=\underset{wS_N}{}\underset{i=1}{\overset{N}{}}\underset{r_i=0}{\overset{p1}{}}(1)^{sg(w)}T_{r_1}^{w(1)1}T_{r_2}^{w(2)2}\mathrm{}T_{r_N}^{w(N)N}\underset{j=1}{\overset{N}{}}\left(\frac{(xj)^{p1r_j}}{_{k=1}^p(xju_k)}\right)$$
(5.25)
where $`S_N`$ is the symmetric group and $`T_r^{ab}`$ are the Yangian generators.
Looking at the poles of $`H(x)`$, it is easy to see that there are exactly $`Np`$ poles (including multiplicities). A basis for this center (using quantum determinant) was also given in . Hence, using this property and the above remark, we can deduce
###### Corollary 5.5
The center of $`𝒲_p(N)`$ is $`Np`$-dimensional and is given by $`𝒟_{Np}/_p`$. A basis of this center is canonically associated to the Casimir operators of $`gl(Np)`$.
Let us remark that the $`p`$ first Casimir operators can be chosen as elements of the $`𝒲_p(N)`$ basis, while the next $`p(N1)`$ ones are polynomials in the basis generators. Note that a different way to get these central generators has been given in . It uses a determinant formula for $`gl(Np)`$ expressed for $`𝕁_{gf}`$, namely:
$$\text{det}(𝕁_{gf}\lambda 𝕀)=(1)^{Np}\lambda ^{Np}+\underset{n=0}{\overset{Np1}{}}C_{Npn}\lambda ^n$$
(5.26)
More generally, the same reasoning leads to the following center for $`S_rY(N)_p`$:
$$Z(S_rY(N)_p)=𝒰(d_{r+1},\mathrm{},d_{pN})/𝒯_p$$
(5.27)
It is generated by the last $`(Npr)`$ independent Casimirs of $`gl(Np)`$.
## 6 Conclusion
We have shown that the finite $`𝒲(gl(Np),N.sl(p))`$ algebras are nothing but truncated Yangians $`Y(gl(N))_p`$, i.e. coset of the Yangian $`Y(gl(N))`$ by the relations $`T_{(n)}^{ab}=0`$ for $`np`$. The resulting coset is an algebra, but the Yangian Hopf structure does not survive to the quotient. This property enlightens the algebra homomorphism between Yangians and finite $`𝒲`$-algebras, and which was given in . Using this property, we have been able to present these $`𝒲`$-algebras as exchange algebras, with the help of the Yangian R-matrix. This more abstract presentation is not linked to an Hamiltonian reduction, as were usually defined the $`𝒲`$-algebras. It could be of some help in the seek of a geometrical interpretation of $`𝒲`$-algebras. As a consequence, we have also given a complete classification of the finite dimensional irreducible representations for these $`𝒲`$-algebras. This classification completes the one given in for $`𝒲(sl(2n),2.sl(n))`$ algebras. Physically, one can hope to construct lattice models associated to $`𝒲`$-algebras, starting from models with Yangian symmetry.
Now that the relation between Yangians and $`𝒲`$-algebras is well-understood, one can hope to construct R-matrices for general $`𝒲`$-algebras: work is in progress in this direction. Conversely, one can think of generalizing the notion of Yangian as certain limits of $`𝒲(𝒢,)`$ algebras in which a (quasi) Hopf structure can be recovered. This would provide a wide class of new types of quantum groups.
Let us also remark that two other approaches for Yangians and $`𝒲`$-algebras could be related. On the one hand, one can construct Yangians as the projective limit of the centralizer of $`gl(n)`$ in $`𝒰(gl(m+n))`$ (see also ), and on the other hand, some finite $`𝒲`$-algebras (of type $`𝒲(gl(2n),n.sl(2))`$) have been realized as commutants of a $`gl(2n)`$ parabolic subalgebra in a certain localization of $`𝒰(gl(2n))`$ . It seems to us quite natural to look for a global description of these two point of view.
Of course, the case of conformal $`𝒲`$-algebras (i.e. extensions of the Virasoro algebra) has to be considered. It could be related to a multi-parametric generalization of Yangians. Would such a generalization be possible, one could think of an ”RTT presentation” of Virasoro algebra: this would allow to relate ”usual” $`𝒲`$-algebras with the deformed $`𝒲`$-algebras presentation, a link which is not clear up to now, since two different deformed algebras can be constructed . Note finally that the construction of some conformal $`𝒲`$-algebras (such as the Virasoro and the $`𝒲_3`$ algebras) as commutant in a localization of an affine Kac-Moody algebra (see above paragraph) as been already achieve : this could be a way to generalize the notion of Yangians, using the centralizer construction.
Acknowledgments
We warmly thank Daniel Arnaudon, Michel Bauer and Paul Sorba for fruitful and clarifying discussions and Alexander Molev for his enlightening comments about representations of Yangians.
## Appendix A General settings on $`gl(Np)`$
We have gathered here the notations and properties we need about $`gl(Np)`$ algebra.
We consider the $`gl(Np)`$ algebra in its fundamental representation ($`Np`$x$`Np`$ matrices), and take a basis adapted to the decomposition with respect to the $`sl(2)`$ algebra principal in $`N.sl(p)\underset{N}{\underset{}{sl(p)\mathrm{}sl(p)}}`$. This decomposition makes naturally appear the ”factorization” $`gl(Np)=gl(N)gl(p)`$, valid in the fundamental representation, and the $`sl(2)`$ principal in $`sl(p)`$.
### A.1 The principal embedding of $`sl(2)`$ in $`sl(p)`$
We will denote by $`M_{j,m}`$ (with $`jmj`$ and $`1jp`$) the $`p`$x$`p`$ matrices resulting from the decomposition in $`sl(2)`$ multiplets:
$`[e_+,M_{j,m}]`$ $`=`$ $`{\displaystyle \frac{j(j+1)m(m+1)}{2}}M_{j,m+1}`$ (A.1)
$`[e_{},M_{j,m}]`$ $`=`$ $`M_{j,m1}`$ (A.2)
$`[e_0,M_{j,m}]`$ $`=`$ $`mM_{j,m}`$ (A.3)
$`[e_0,e_\pm ]`$ $`=`$ $`\pm e_\pm \text{ and }[e_+,e_{}]=e_0`$ (A.4)
where $`e_{\pm ,0}`$ are the generators of the $`sl(2)`$ algebra principal in $`sl(p)`$. The normalizations in (A.1-A.2), although not symmetric, are adapted to the $`𝒲`$-algebra framework. When working with $`gl(p)`$ instead of $`sl(p)`$, we will add the $`j=0`$ generator, proportional to the identity matrix.
The decomposition of $`M_{j,m}`$ in terms of the $`p`$x$`p`$ matrices $`E_{ab}`$ reads
$$\begin{array}{ccc}M_{j,m}=\underset{k=1}{\overset{pm}{}}a_{j,m}^kE_{k,k+m}\hfill & \text{ with }\hfill & a_{j,m}^k=\underset{i=0}{\overset{jm}{}}(1)^{i+j+m}\left(\genfrac{}{}{0pt}{}{jm}{i}\right)a_{j,j}^{ki}\hfill \\ \text{ for }0mj\hfill & & \end{array}$$
(A.5)
$$\begin{array}{ccc}M_{j,m}=\underset{k=1}{\overset{p+m}{}}a_{j,m}^kE_{km,k}\hfill & \text{ with }\hfill & a_{j,m}^k=\underset{i=0}{\overset{jm}{}}(1)^{i+j+m}\left(\genfrac{}{}{0pt}{}{jm}{i}\right)a_{j,j}^{kim}\hfill \\ \text{ for }jm0\hfill & & \end{array}$$
(A.6)
$$a_{j,j}^k=\frac{(k+j1)!(pk)!}{(k1)!(pkj)!}$$
(A.7)
The generators $`e_{\pm ,0}`$ of the $`sl(2)`$ algebra are proportional to the $`M_{1,m}`$ generators:
$`e_+`$ $`=`$ $`{\displaystyle \underset{k=1}{\overset{p1}{}}}{\displaystyle \frac{k(pk)}{2}}E_{k,k+1}={\displaystyle \frac{1}{2}}M_{1,1}`$
$`e_0`$ $`=`$ $`{\displaystyle \underset{k=1}{\overset{p}{}}}({\displaystyle \frac{p+1}{2}}k)E_{k,k}={\displaystyle \frac{1}{2}}M_{1,0}`$ (A.8)
$`e_{}`$ $`=`$ $`{\displaystyle \underset{k=1}{\overset{p1}{}}}E_{k+1,k}={\displaystyle \frac{1}{2}}M_{1,1}`$
Let us remark that we have the following generating function for the coefficients $`{}_{}{}^{(p)}a_{j,m}^{k}`$ (where $`(p)`$ refers to the $`gl(p)`$ algebra under consideration):
$`{}_{}{}^{(p)}a_{j,m}^{k}={\displaystyle \frac{j!}{p!(k1)!(jm)!}}\left[{\displaystyle \frac{d^p}{du^p}}{\displaystyle \frac{d^j}{dz^j}}{\displaystyle \frac{d^{jm}}{dy^{jm}}}{\displaystyle \frac{d^{k1}}{dx^{k1}}}a(x,y,z;u)\right]_{\begin{array}{c}x=y=0\hfill \\ z=u=0\hfill \end{array}}`$ (A.11)
$`\text{with }a(x,y,z,u)={\displaystyle \frac{u}{\left(\text{}1+y(1x)\right)\left(\text{}1u[1+z+x(1u)]\right)}}`$ (A.12)
The scalar product is given by
$`\eta _{j,m;\mathrm{},n}`$ $`=`$ $`(M_{j,m},M_{\mathrm{},n})=tr(M_{j,m}M_{\mathrm{},n})=(1)^m\eta _j\delta _{j,\mathrm{}}\delta _{m+n,0}`$ (A.14)
$`\text{with }\eta _j=(2j)!(j!)^2\left({\displaystyle \genfrac{}{}{0pt}{}{p+j}{2j+1}}\right)`$
where the dot stands for the matrix product and $`tr`$ is the trace of matrices (in the $`p`$-dimensional representation).
In the following, we will need the Clebsch-Gordan like coefficients given by:
$$M_{j,m}M_{\mathrm{},n}=\underset{r=|j\mathrm{}|}{\overset{j+\mathrm{}}{}}\underset{s=r}{\overset{r}{}}<j,m;\mathrm{},n|r,s>M_{r,s}$$
(A.15)
As for usual Clebsch-Gordan coefficients, one can prove (using commutators by $`e_{\pm ,0}`$) that $`r`$ must be in $`[|j\mathrm{}|,j+\mathrm{}]`$ and that $`s`$ must be equal to $`mn`$. However, since we are in the fundamental of $`sl(p)`$, the coefficients will be truncated in such a way that only the values $`rp`$ are kept in the decomposition (A.15). We will still call them Clebsch-Gordan coefficients.
Using the scalar product, one can compute these coefficients to be
$$<j,m;\mathrm{},n|r,s>=\frac{(1)^s}{\eta _r}tr\left(M_{j,m}M_{\mathrm{},n}M_{r,s}\right)$$
(A.16)
### A.2 Few results about the Clebsch-Gordan like coefficients
Using the cyclicity of the trace, one shows
$$<j,m;\mathrm{},n|r,s>=(1)^{s+m}\frac{\eta _j}{\eta _r}<\mathrm{},n;r,s|j,m>=(1)^{s+n}\frac{\eta _{\mathrm{}}}{\eta _r}<r,s;j,m|\mathrm{},n>$$
(A.17)
We will also use the property
$$<j,m;\mathrm{},n|r,s>=\frac{(jm)!(\mathrm{}n)!(r+s)!}{(j+m)!(\mathrm{}+n)!(rs)!}<\mathrm{},n;j,m|r,s>$$
(A.18)
where the coefficients are due to the non-symmetric basis we have chosen.
With these two properties, one can compute:
$`<r,r;k,k|j,j>`$ $`=`$ $`(1)^k{\displaystyle \frac{\eta _{j+k}}{\eta _j}}\delta _{r,j+k}`$ (A.19)
$`<k,k;r,r|j,j>`$ $`=`$ $`(1)^k{\displaystyle \frac{\eta _{j+k}}{\eta _j}}\delta _{r,j+k}`$ (A.20)
$`<r,1r;k,k|j,1j>`$ $`=`$ $`(1)^k{\displaystyle \frac{j}{j+k}}{\displaystyle \frac{\eta _{j+k}}{\eta _j}}\delta _{r,j+k}`$ (A.21)
$`<k,k;r,1r|j,1j>`$ $`=`$ $`(1)^k{\displaystyle \frac{j}{j+k}}{\displaystyle \frac{\eta _{j+k}}{\eta _j}}\delta _{r,j+k}`$ (A.22)
$`<r,r;k,k|j,1j>`$ $`=`$ $`(1)^{k+1}kj{\displaystyle \frac{\eta _{j+k1}}{\eta _j}}\delta _{r+1,j+k}`$ (A.23)
$`<k,k;r,r|j,1j>`$ $`=`$ $`(1)^kkj{\displaystyle \frac{\eta _{j+k1}}{\eta _j}}\delta _{r+1,j+k}`$ (A.24)
We will also need the following coefficients:
$`<k,k;k,1k|1,1>`$ $`=`$ $`(1)^{k+1}{\displaystyle \frac{\eta _k}{\eta _1}}c_k`$ (A.25)
$`<k,1k;k,k|1,1>`$ $`=`$ $`c_k`$ (A.26)
$`<k,k;k1,1k|1,1>`$ $`=`$ $`{\displaystyle \frac{1}{k(2k1)}}c_k`$ (A.27)
$`<k1,1k;k,k|1,1>`$ $`=`$ $`{\displaystyle \frac{1}{k(2k1)}}c_k`$ (A.28)
$`<k+1,1k;k,k|1,1>`$ $`=`$ $`{\displaystyle \frac{(k+1)(p^2(k+1)^2)}{2k+3}}c_k`$ (A.29)
$`<k,k;k+1,1k|1,1>`$ $`=`$ $`{\displaystyle \frac{(k+1)(p^2(k+1)^2)}{2k+3}}c_k`$ (A.30)
### A.3 Basis for $`gl(Np)`$
We can use the above basis of $`gl(p)`$ to construct a basis for $`gl(Np)`$. Using the $`N`$x$`N`$ matrices $`E_{ab}`$, the generators $`\mathrm{{\rm Y}}_{ab}^{jm}`$ of $`gl(Np)`$ in the fundamental will be represented by
$$\pi _F(\mathrm{{\rm Y}}_{ab}^{jm})=M_{ab}^{jm}=E_{ab}M^{jm}$$
(A.31)
The generators of the $`sl(2)`$ algebra principal in $`N.sl(p)`$ are then
$$ϵ_{\pm ,0}=1_Ne_{\pm ,0}$$
(A.32)
where $`e_{\pm ,0}`$ are the $`p`$x$`p`$ matrices defined above. We have the following commutation relations
$`[ϵ_+,M_{ab}^{j,m}]`$ $`=`$ $`{\displaystyle \frac{1}{2}}(j(j+1)m(m+1))M_{ab}^{j,m+1}`$ (A.33)
$`[ϵ_{},M_{ab}^{j,m}]`$ $`=`$ $`M_{ab}^{j,m1}`$ (A.34)
$`[ϵ_0,M_{ab}^{j,m}]`$ $`=`$ $`mM_{ab}^{j,m}`$ (A.35)
$`[ϵ_0,ϵ_\pm ]`$ $`=`$ $`\pm ϵ_\pm \text{ and }[ϵ_+,ϵ_{}]=ϵ_0`$ (A.36)
together with
$$[M_{ab}^{00},M_{cd}^{00}]=\delta _{bc}M_{ad}^{00}\delta _{ad}M_{cb}^{00}$$
(A.37)
This last commutator reveals the $`gl(N)`$ algebra which commutes with the $`sl(2)`$ subalgebra under consideration.
More generally, the product law (in the fundamental representation) reads
$$M_{ab}^{j,m}M_{cd}^{\mathrm{},n}=\delta _{bc}\underset{r=|j\mathrm{}|}{\overset{j+\mathrm{}}{}}\underset{s=r}{\overset{r}{}}<j,m;\mathrm{},n|r,s>M_{ad}^{r,s}$$
(A.38)
which leads to the following commutation relations (valid in the abstract algebra):
$$[\mathrm{{\rm Y}}_{ab}^{j,m},\mathrm{{\rm Y}}_{cd}^{\mathrm{},n}]=\underset{r=|j\mathrm{}|}{\overset{j+\mathrm{}}{}}\underset{s=r}{\overset{r}{}}\left(\text{}\delta _{bc}<j,m;\mathrm{},n|r,s>\mathrm{{\rm Y}}_{ad}^{r,s}\delta _{ad}<\mathrm{},n;j,m|r,s>\mathrm{{\rm Y}}_{cb}^{r,s}\right)$$
(A.39)
The scalar product is
$$\eta _{ab,cd}^{j,m;\mathrm{},n}=(\mathrm{{\rm Y}}_{ab}^{j,m},\mathrm{{\rm Y}}_{cd}^{p,n})=tr(M_{ab}^{j,m}M_{cd}^{p,n})=\delta _{a,d}\delta _{b,c}\eta ^{j,m;\mathrm{},n}$$
(A.40)
## Appendix B Deformations and cohomology
We include here some definitions (in the context of Chevalley cohomology) to be self-content. For more details about deformations and their relation to cohomology, we refer to and ref. therein.
### B.1 Few words about Chevalley cohomology
We begin with an algebra $`𝒜`$, and first introduce the space $`C_n(𝒜,𝒜)`$ of $`n`$-cochains with values in $`𝒜`$, i.e. skew-symmetric linear maps from $`^n𝒜`$ to $`𝒜`$. The Chevalley derivation $`\delta `$ maps $`n`$-cochains to $`(n+1)`$-cochains as:
$`(\delta \chi _n)(u_0,\mathrm{},u_n)`$ $`=`$ $`{\displaystyle \underset{i=0}{\overset{n}{}}}(1)^i\{\text{}u_i,\chi _n(u_0,u_1,\mathrm{},\widehat{u_i},\mathrm{},u_n)\}+`$
$`+{\displaystyle \underset{0i<jn}{}}(1)^{i+j}\chi _n(\{u_i,u_j\},u_0,u_1,\mathrm{},\widehat{u_i},\mathrm{},\widehat{u_j},\mathrm{},u_n)`$
where, as usual, $`\widehat{u_i}`$ means that $`u_i`$ has to be discarded in the list (or product, or sum, or whatever) we consider.
It can be shown that $`\delta `$ squares to zero:
$$(\delta (\delta \chi _n))(u_1,u_0,u_1,\mathrm{},u_n)=0u_1,u_0,u_1,\mathrm{},u_n;\chi _n;n$$
(B.41)
Thus, we introduce the cohomology associated to $`\delta `$, i.e. we focus on Ker$`\delta `$. Elements of Ker$`\delta `$ are called cocycles, and we will see that they play a direct role in the deformation of Lie algebras. The space of $`n`$-cocycles (with values in $`𝒜`$) is denoted $`Z_n(𝒜,𝒜)`$: Ker$`\delta =_nZ_n(𝒜,𝒜)`$. Since $`\delta ^2=0`$, we have Im$`\delta `$Ker$`\delta `$: each $`n`$-cochain provides a $`(n+1)`$-cocycle. The elements $`\delta \chi _n`$ correspond to ”trivial” cocycles: they are called coboundaries, and the corresponding space denoted $`B_n(𝒜,𝒜)`$. The cohomology describes the non-trivial cocycles, i.e. it is the space $`H_n(𝒜,𝒜)=Z_n(𝒜,𝒜)/B_n(𝒜,𝒜)`$, $`H(𝒜,𝒜)=_nH_n(𝒜,𝒜)=`$Ker$`\delta /`$Im$`\delta `$. Due to its definition, the Chevalley cohomology is naturally associated to Lie algebras. When the cochains take values in $``$ instead of $`𝒜`$, the space $`H_2(𝒜,)`$ classifies the non-trivial central extensions of $`𝒜`$: see for instance where central extensions of generalized loop algebras are classified and computed. In the case we are considering, $`C(𝒜,𝒜)`$ is related to deformations of $`𝒜`$.
### B.2 Deformations
We start again with an algebra, with generators $`u_\alpha `$ ($`\alpha \mathrm{\Gamma }`$).
$$\{u_\alpha ,u_\beta \}=f_{}^{\alpha \beta }{}_{\gamma }{}^{}u_\gamma $$
Actually, we will consider its enveloping algebra $`𝒜`$, and introduce a deformation of it
$$\{u_\alpha ,u_\beta \}_{\mathrm{}}=f_{}^{\alpha \beta }{}_{\gamma }{}^{}u_\gamma +\underset{n=1}{\overset{\mathrm{}}{}}\mathrm{}^n\phi _n(u_\alpha ,u_\beta )$$
where the antisymmetric bilinear forms $`\phi _n`$ take values in $`𝒜`$: they are all elements of $`C_2(𝒜,𝒜)`$.
Asking the bracket $`\{.,.\}_{\mathrm{}}`$ to obey the Jacobi identity leads to the following equations:
$`\delta \phi _1`$ $`=`$ $`0`$ (B.42)
$`\delta \phi _n`$ $`=`$ $`{\displaystyle \underset{j+k=n}{}}\left(\text{}\phi _j(\phi _k(u,v),w)+\phi _j(\phi _k(v,w),u)+\phi _j(\phi _k(w,u),v)\right)\text{ for }n>1`$ (B.43)
where the operation $`\delta `$ defined in section B.1 has naturally appeared.
These equations indicate that $`\phi _1`$ is a cocycle, while $`\phi _n`$ is determined by the $`\phi _p`$’s ($`p<n`$) up to a cocycle. Note that $`\delta \phi _n`$ is a coboundary, so that the $`\phi _p`$, $`p<n`$, must be such that the r.h.s. of (B.43) is also a coboundary (it can be proven that this r.h.s. is indeed a cocycle, i.e. is annihilated by $`\delta `$). If the third cohomological space is not trivial, the r.h.s. of (B.43) may be a cocycle while being not a coboundary: this leads to the usual assumption that the third cohomological space classify the obstructions to deformations. In other words, it could appear that, in the attempt to construct a deformation, the chosen $`\phi _p,p<n`$ are such that the l.h.s. of (B.43) is a non-trivial cocycle, so that one cannot solve this equation at level $`n`$. In that case, the deformation would be ill-defined.
Fortunately, in the case we will consider below, we already know that we have well-defined deformations, and we have not to deal with a possible obstruction.
Note also that if $`\phi _n`$ is a coboundary
$$\phi _n(u_\alpha ,u_\beta )=\delta \chi _n(u_\alpha ,u_\beta )=\{u_\alpha ,\chi _n(u_\beta )\}\{u_\beta ,\chi _n(u_\alpha )\}\chi _n(\{u_\alpha ,u_\beta \})$$
we can perform a change of basis
$$\stackrel{~}{u}_\alpha =u_\alpha \mathrm{}^n\chi _n(u_\alpha )$$
such that in this new basis, the term in $`\mathrm{}^n`$ has disappeared:
$$\{\stackrel{~}{u}_\alpha ,\stackrel{~}{u}_\beta \}_{\mathrm{}}=f_{}^{\alpha \beta }{}_{\gamma }{}^{}\stackrel{~}{u}_\gamma +\underset{m=1}{\overset{n1}{}}\mathrm{}^m\phi _m(\stackrel{~}{u}_\alpha ,\stackrel{~}{u}_\beta )+\underset{m=n+1}{\overset{\mathrm{}}{}}\mathrm{}^m\stackrel{~}{\phi }_m(\stackrel{~}{u}_\alpha ,\stackrel{~}{u}_\beta )$$
where $`\stackrel{~}{\phi }_m,m>n`$ are new cochains resulting from the change of variables. In that sense, a coboundary leads to a trivial deformation. However, one has to be careful that to ”trivialize” the full deformation, the change of basis has to be done recursively and the coboundarity at level $`n`$ has to be checked once the change of basis at level $`n1`$ has been done (since the cochains are modified at higher order).
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# HST/STIS spectroscopy of the optical outflow from DG Tau: structure and kinematics on sub-arcsecond scales1
## 1 Introduction
One of the most interesting questions in young stellar object (YSO) research is how their jets, e.g. Camenzind (1997), Ray (1996) or Eislöffel and Mundt (1997), are collimated and accelerated. To address this problem one must obtain not only high spatial resolution but in addition kinematic information as close as possible to their source. Ground-based long-slit spectroscopic studies of optically visible jet sources have shown that the structure and kinematics of the outflow region on scales $`\stackrel{<}{}`$1<sup>′′</sup> is rather complex (e.g. Solf and Böhm (1993); Hirth, Mundt, and Solf (1997)). In a significant fraction of classical T Tauri stars (CTTSs) the forbidden emission lines (FELs) show two (or more) blueshifted velocity components which have very different properties (note that in many CTTSs the corresponding redshifted part of the flow is occulted by a circumstellar disk, at least close to the source). In the case of double-peaked FEL profiles, the emission consists of a so-called high-velocity component (HVC), having typical radial velocities of -60 to -200 km s<sup>-1</sup> and a low-velocity component (LVC), having typical radial velocities of -5 to -20 km s<sup>-1</sup> with respect to the systemic velocity of the star (see Hartigan, Edwards, and Ghandour (1995) and Solf (1997) for comprehensive discussions). While the HVC can be spatially very extended (at least in the outflow direction) and is often identified with a jet, the LVC is much more compact ($`\stackrel{<}{}`$1<sup>′′</sup>). Other differences between the two components have also been noted, for example the LVC appears to be of higher density but lower excitation than the HVC (Hirth, Mundt, and Solf, 1997) and the LVC shows a clear inverse correlation of velocity with increasing critical density. The latter effect has been interpreted as evidence of acceleration in the LVC with increasing distance from the star.
A number of theories have been put forward to explain the origin of the various FEL components. For example Kwan and Tademaru (1988, 1995) have suggested that the LVC and HVC are separate flows. According to their model the LVC is a poorly collimated wind coming from the outer periphery of the YSO disk while the HVC is a separate jet launched from closer to the star. Others have sought to explain the observations in terms of a one component, non-isotropic, wind model in which the appearence of separate FEL components is due to projection effects (Hartmann and Raymond, 1989; Safier, 1993; Ouyed and Pudritz, 1994). In any event, Calvet (1997) has shown that the luminosity of the two components is tightly correlated, which seems to imply that they are not independent flows, and that the apparent dominance of one component over the other may be a density effect.
In order to examine in more detail the nature of the compact FEL region, we have observed the CTTS DG Tau with STIS on-board the HST. Multiple overlapping slit positions parallel to the outflow from this star were chosen so as to build up a 3-D spatial intensity-velocity “cube”. Our target was picked not only because it is one of the closest CTTSs, but more significantly because its FEL region has a broad range of velocities, probably also a result of the relatively small angle between the jet axis and the line of sight ($`38^{}`$, see Eislöffel and Mundt (1998)). This is important given the moderate wavelength resolution of STIS. Historically, DG Tau was amongst the first T Tauri stars from which a jet-like outflow (HH 158) was discovered (Mundt and Fried, 1983) and it has been imaged by the HST prior to the installation of the telescope’s correcting optics (Kepner et al., 1993). On large scales ($``$ 10<sup>′′</sup>) the jet seems to terminate in a bow shock (Eislöffel and Mundt, 1998) while high resolution spectro-imaging ground-based studies (Lavalley et al., 1997) have shown that the flow close to the star contains at least two resolved knots (at 2.7<sup>′′</sup> and 4<sup>′′</sup>, epoch 1994.8). The outermost of these has a morphology and velocity gradient which is also consistent with it being a bow shock.
In §2 we describe our observational technique and give details of our STIS data. Our primary results are described in §3 and discussed in §4. A more detailed analysis will be presented in Bacciotti et al. (2000).
## 2 Observations
STIS spectra of DG Tau were taken on January 14 1999 using the G750M grating and a central wavelength setting of 6581 Å. The spectral range of 562 Å included several strong forbidden lines (\[OI\] $`\lambda \lambda `$6300,6363, \[NII\] $`\lambda \lambda `$6548,6583, \[SII\] $`\lambda \lambda `$6716,6731) as well as H$`\alpha `$. The STIS/CCD detector, a 1024$`\times `$1024 pixel array, had a spectral resolution of 0.56 Å pixel<sup>-1</sup> and a nominal sampling of 0.05<sup>′′</sup> pixel<sup>-1</sup> in the dispersion and spatial directions respectively. The actual spatial resolution, however, was limited by the PSF of the HST in the red to a FWHM of approximately 0.1<sup>′′</sup>. The slit aperture was 52$`\times `$0.1 arcsec<sup>2</sup>. Seven different long-slit spectra were taken, keeping the slit parallel to the outflow axis (P.A. 226), but with steps of 0.07<sup>′′</sup> in the transverse direction, i.e. with offsets southeast and northwest of the jet axis. In this way we built-up a 3-D flux density/radial velocity data cube of the optical outflow from DG Tau, with a total spatial width of about 0.5<sup>′′</sup> in the direction perpendicular to the jet axis. The spectra are labelled S1,S2,…S7 going from the southeast to the north-west; the central slit position (S4) coincided with the star (RA 4<sup>h</sup> 27<sup>m</sup> 04<sup>s</sup>.71, Dec +26 6 16.8<sup>′′</sup>.) Exposure times were split in order to facilitate removal of cosmic rays. Total exposure times were approximately 2740 s for S1 – S5, 1930 s for S6 and 2050 s for S7. Note that the somewhat shorter exposure times for the last two spectra were due to minor problems with the HST.
The pipeline spectra were found to be contaminated by a large number of “hot” pixels and so the raw data were first fully re-calibrated using the CALSTIS suite of programs and new reference files made available by the STScI. Subsequent reduction and data analysis was carried out using standard IRAF routines. In order to minimise contrast problems close to the source, the stellar continuum, and that of a faint reflection nebula extending up to 2<sup>′′</sup> from the source, were carefully subtracted from all images. This operation proved to be particularly critical for the central spectra S3, S4, and S5, due to the appearance of artificial undulations in the stellar continuum, caused by sampling effects (see the STIS Handbook v3.0).
From the acquired spectra we created synthetic 2-D images for the various lines in four broad distinct radial velocity intervals (Figs. 1–2). These images were constructed by adjoining the seven row-averaged columns of pixels obtained from S1–S7 for the corresponding velocity bin. Averaging was done over five columns, corresponding to approximately 2.83Å, in the dispersion direction. In order to study the lowest velocity emission in finer detail, we also (Fig. 3) carried out a similar procedure in \[OI\]$`\lambda `$6300 and \[SII\]$`\lambda `$6731 but this time using just one column in the dispersion direction. All velocities quoted in this Letter are systemic using v$`{}_{,hel}{}^{}`$ +16.5 km s<sup>-1</sup>, as derived from the Li $`\lambda `$6707 photospheric absorption line.
## 3 Results
Our composite images of the outflow close to DG Tau in the four broad velocity bins, and in each of four lines (H$`\alpha `$, \[NII\]$`\lambda `$6583, \[SII\]$`\lambda `$6731, and \[OI\]$`\lambda `$6363), are shown in Figs. 1 and 2. Note that we did not include \[OI\]$`\lambda `$6300 in these figures as the blueward portion of this line was missing from our spectra due to a restricted choice for the STIS central wavelength. The velocity range is approximately +50 km s<sup>-1</sup> to -450 km s<sup>-1</sup> and the four velocity bins (low, medium, high, and very high), with widths of about 125 km s<sup>-1</sup>, represent increasingly blueshifted velocities. On all figures the position of the star (i.e. the peak in the continuum subtracted light) is marked. Similar images, for the \[OI\]$`\lambda `$6300 and \[SII\]$`\lambda `$6731 lines, but only for the lowest velocities (from about +60 km s<sup>-1</sup> to -70 km s<sup>-1</sup>) and with bin widths of only about 25 km s<sup>-1</sup>, are shown in Fig. 3.
First of all, we note that the appearence of the outflow changes, remarkably in some cases, from line to line and from velocity bin to velocity bin (see in particular Figs. 1 and 2). The jet from DG Tau is, as one might expect, most evident at the highest velocities and is visible as a narrow feature up to about 0.7<sup>′′</sup> from the star. We will discuss it in more depth shortly. Aside from the jet a limb-brightened bubble-like structure can also be observed that extends between 0.4<sup>′′</sup> and 1.5<sup>′′</sup> from DG Tau. The “bubble”, which is most evident at intermediate velocities, in combination with the jet, gives an overall pear-shaped morphology to the emission.
Comparison of Figs. 1, 2 and 3 lead to a number of important results:
(i) The lowest velocity emission close to the star ($`\stackrel{<}{}`$0.3<sup>′′</sup>) is spatially broad, i.e. it is well resolved in the transverse direction to the outflow, but it does not extend far along the flow.
(ii) Within $`\stackrel{<}{}`$0.5<sup>′′</sup> from the star, there appears to be a more or less gradual increase in the degree of collimation with velocity: at high and very high velocities, the flow is primarily confined to the central axis.
(iii) Beyond 0.5<sup>′′</sup> from the star there are at least two well-defined structures (which we have labelled A1 and A2 in the medium velocity H$`\alpha `$ image following the knot nomenclature of Eislöffel and Mundt (1998)). The outermost one, A1, is bow-shaped and observed in all lines but the innermost one, A2, is more easily seen in H$`\alpha `$. There is very high velocity gas immediately behind A1, which is seen in all lines but which is not symmetrically distributed around the jet axis. The opposite effect is, incidentally, seen near A2 where the higher velocity gas appears in the downstream direction. The emission from A2 is also distributed asymmetrically with respect to the outflow axis.
(iv) Comparison of the \[OI\], \[SII\], and \[NII\] data (in Figs. 1 and 2) show that the central jet can be traced much closer to the source in \[OI\] and \[NII\] than in \[SII\]. This result is probably due to a combination of quenching and excitation effects (see Bacciotti et al. (2000) for details). Note that the \[OI\] line has a 100 (10) times higher critical density than the \[SII\] (\[NII\]) line, and that the \[NII\] line usually traces gas of higher excitation. The emission peak (centroid) in the \[SII\] and \[OI\] images moves outwards with increasing velocity. This effect is visible at all velocities, but it is most obvious at the highest velocities where virtually no \[SII\] emission and only faint \[OI\] emission is seen in the region $`\stackrel{<}{}`$1<sup>′′</sup>. Since the high velocity jet can be traced right back to the source in H$`\alpha `$, the lack of high-velocity and the presence of low-velocity FEL emission at $`<`$ 1<sup>′′</sup> can probably only be explained by quenching effects in both the longitudinal and lateral directions. In other words our data suggest that the jet density increases not only with proximity to the star (in the longitudinal direction) but also transversely toward the outflow axis. Without the latter effect one would not be able to explain why there is so much spatially extended and easily quenchable low-velocity FEL emission seen so close to the source.
(v) Quenching effects likewise manifest themselves at low velocities (see Fig. 3) in that the centroid of the low velocity \[OI\]$`\lambda `$6300 emission is closer to the star than the centroid of the \[SII\]$`\lambda `$6731 emission. This was expected on the basis of groundbased measurements (Hirth, Mundt, and Solf, 1997).
(vi) In the region close to the star ($`\stackrel{<}{}`$0.7<sup>′′</sup>) \[NII\] emission comes primarily from the high velocity jet implying this is the region with the highest excitation.
(vii) The medium velocity bow-shaped structure A1 is less evident in \[SII\] and \[OI\] than in H$`\alpha `$, suggesting that its excitation level may be high.
(viii) There is a strong low velocity peak in H$`\alpha `$ emission coinciding with the star. This is almost certainly scattered emission from a region close to DG Tau, including possibly a magnetospheric contribution (Edwards, 1997) and therefore does not constitute part of the outflow per se.
## 4 Discussion
The STIS observations presented here clearly show that the HVC emission in DG Tau comes from the most highly focused part of the outflow i.e. the jet. The jet at high velocities can be traced back to at least 0.1<sup>′′</sup> (15 AU) from the star where quenching effects become important even in the case of the \[OI\]$`\lambda \lambda `$6300,6363 lines, which have the highest critical density ($`10^6`$ cm<sup>-3</sup>) among the studied FELs. Obviously it is likely that the jet is collimated on even smaller scales. The shifts of the emission centroids in the different lines and in the different velocity bins can be explained quite naturally if there is an increase in jet density not only with proximity to the star (in the longitudinal direction) but also to the central outflow axis. That is to say the high velocity “core” of the jet is denser than its periphery. This is confirmed by an inspection of the \[SII\]$`\lambda \lambda `$6716,6731 doublet ratio (see Fig. 2 and Bacciotti et al. 2000). Note also that structures such as A1 and A2 are reminiscent of the “bubbles” recently seen by Krist et al. (1999) in the case of XZ Tauri. These are almost certainly internal working surfaces caused by temporal variations in the outflow from DG Tau. Certainly there is plenty of evidence for strong jet velocity variations in DG Tau (and in many other CTTS stars) on timescales of years. For example, the data of Mundt et al. (1987) and Solf and Böhm (1993) show an increase in the radial velocity of the jet at D $``$ 0.5<sup>′′</sup> by a factor of 2 within 8 years. Also the proper motion data of Eislöffel and Mundt (1998) indicate large velocity variations. Finally, several of the knots in the DG Tau jet show bow shock-like structures (see also the HST Archive data presented in Bacciotti et al. (2000)) and these provide indirect evidence for strong velocity variations.
Turning now to the LVC, its nature still remains somewhat enigmatic. It was already clear from groundbased observations that the LVC and HVC differ in many properties such as density and excitation. A new difference is reported here i.e. the rather large spatial width W of the LVC perpendicular to the jet at distances D of about 0.1<sup>′′</sup>–0.3<sup>′′</sup> from the source. A comparison of W (at FWHM) between the high and low-velocity emission in \[SII\] and \[OI\] shows that the average W(LVC) $``$ 0.18<sup>′′</sup> while the HVC is hardly resolved at a distance D of 0.2<sup>′′</sup>. We note that this comparison of spatial widths can only be done using our \[SII\] and \[OI\] “images” as the H$`\alpha `$ data is heavily contaminated by stellar H$`\alpha `$ emission and in \[NII\] the LVC is very weak. Another interesting result of our study is the smaller velocity of the LVC at the edges of the flow (for full details see Bacciotti et al. (2000)). For example at D = 0.2<sup>′′</sup> the LVC peaks at $``$ -95 km s<sup>-1</sup> in \[OI\] for the central slit position (S4) while in the two outermost slit positions (S1,S7) it peaks at about -18 km s<sup>-1</sup>. The corresponding values for the \[SII\] line are -60 km s<sup>-1</sup> and -40 km s<sup>-1</sup>, respectively. Since these lines are optically thin, such observations clearly point to a rise in the average LVC velocity as the central outflow axis is approached.
To what degree the observations of DG Tau presented here, particularly of the LVC, are representative of other CTTSs is an open question. DG Tau is one of the most active CTTSs known and we caution that the LVC of DG Tau is unusual in that it has the highest absolute velocity of all the CTTSs listed by Hartigan, Edwards, and Ghandour (1995), as well as one of the highest accretion rates. That said it shares the typical properties of other LVCs and, in particular, the ratio of its luminosity (L) to that of the DG Tau HVC is in perfect agreement with the L(HVC) v. L(LVC) relationship noted by Calvet (1997). Thus DG Tau may simply be displaying the higher activity tail of the distribution of outflow properties amongst CTTSs. In conclusion our STIS observations show, for the first time, a quasi-continuous variation in the outflow velocity close to a YSO in the transverse direction to the flow. Detailed studies are required, however, to test whether these observations can constrain models for the generation of the LVC and HVC.
We thank the anonymous referee for his/her very helpful comments. FB was supported by an European Space Agency contract at the Dublin Institute for Advanced Studies during the course of this work.
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# Asymptotic behaviour and the moduli space of doubly-periodic instantons
### Introduction and statement of the results
The aim of this paper is to understand the analytical properties of certain finite energy solutions of the Yang-Mills anti-self-dual equations over $`T\times `$. These so-called extensible doubly-periodic instantons have been studied by the second author in , where they were shown to be equivalent to certain singular solutions of Hitchin’s equations over an elliptic curve via a construction known as the Nahm transform. The present paper grew from questions raised in the works mentioned above.
More precisely, consider an $`SU_2`$ bundle $`ET\times `$. The instanton connections $`A`$ considered in satisfied the following hypothesis:
1. *quadratic curvature decay*: $`|F_A|=O(r^2)`$ with respect to the Euclidean metric on $`T\times `$;
2. *extensibility*: there is a holomorphic rank two vector bundle $`T\times ^1`$ with trivial determinant such that $`|_{T\times (^1\{\mathrm{}\})}(E,\overline{}_A)`$, where $`\overline{}_A`$ is the holomorphic structure on $`E`$ induced by $`A`$;
where $`w`$ is a coordinate in the complex line, and by the notation $`O(|w|^\gamma )`$ we mean the set of functions on $``$ such that: $`lim_{|w|\mathrm{}}|f(w)|/|w|^\gamma <\mathrm{}`$.
One of the goals of this paper is to prove that the technical hypothesis of extensibility is actually a consequence of the anti-self-duality equation, and more generally to understand completely the behaviour at infinity of all instantons with quadratic curvature decay.
Model solutions. Special solutions of the anti-self-duality equations may be obtained by restricting to torus invariant connections. Such instantons come from solutions $`(B,\psi )`$ of Hitchin’s equations on $``$
$$\{\begin{array}{c}F_B+[\psi ,\psi ^{}]=0\hfill \\ \overline{}_B\psi =0\hfill \end{array}$$
in the following way. Recall that $`B`$ is a $`SU_2`$-connection on $``$, and $`\psi `$ is a (1,0)-form with values in $`𝔰𝔩_2`$. Let $`\psi =\frac{1}{2}(\psi _0+i\psi _1)dw`$, and consider the connection (where $`x`$ and $`y`$ are coordinates on $`T`$):
$$A_0=B+\psi _0dx+\psi _1dy$$
which is a torus invariant instanton. Assuming that $`|F_{A_0}|=O(r^2)`$, the asymptotic behavior of solutions $`(B,\psi )`$ is given by one of the following models:
$`B=d`$ $`\psi =\left(\begin{array}{cc}\lambda & 0\\ 0& \lambda \end{array}\right)dw`$ (3)
$`B=d+\left(\begin{array}{cc}\alpha & 0\\ 0& \alpha \end{array}\right)d\theta `$ $`\psi =\left(\begin{array}{cc}\mu & 0\\ 0& \mu \end{array}\right){\displaystyle \frac{dw}{w}}`$ (8)
$`B=d+\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right){\displaystyle \frac{d\theta }{\mathrm{ln}r^2}}`$ $`\psi =\left(\begin{array}{cc}0& 1\\ 0& 0\end{array}\right){\displaystyle \frac{dw}{w\mathrm{ln}r^2}}`$ (13)
where $`\lambda ,\mu `$ and $`\frac{1}{2}\alpha <\frac{1}{2}`$. The solutions of examples (3) & (8) can be superimposed, and such superpositions are called the semisimple solutions. On the other hand, solutions of example (13) cannot be superimposed with the others; these are called the nilpotent solutions, and can only exist when $`\lambda =\mu =\alpha =0`$. The torus invariant instanton is then given by, in the semisimple case:
$$A_0=d+i\left(\begin{array}{cc}a_0& 0\\ 0& a_0\end{array}\right)$$
with
$$a_0=\lambda _1dx+\lambda _2dy+(\mu _1\mathrm{cos}\theta \mu _2\mathrm{sin}\theta )\frac{dx}{r}+(\mu _1\mathrm{sin}\theta +\mu _2\mathrm{cos}\theta )\frac{dy}{r}+\alpha d\theta ;$$
while in the nilpotent case, we have:
$$A_0=d+i\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right)\frac{d\theta }{\mathrm{ln}r^2}+\frac{1}{r\mathrm{ln}r^2}\left(\begin{array}{cc}0& e^{i\theta }(dxidy)\\ e^{i\theta }(dx+idy)& 0\end{array}\right)$$
and note that the curvature is $`O(r^2|\mathrm{ln}r^2|^2)`$.
Remark that the connection $`A_0`$ has a flat limit over the torus at infinity,
$$d+i\left(\begin{array}{cc}\lambda _1dx+\lambda _2dy& 0\\ 0& \lambda _1dx\lambda _2dy\end{array}\right),$$
and one can prove that such flat limit for a connection $`A`$ exists as soon as $`|F_A|=O(r^{1ϵ})`$; the flat limit underlies a holomorphic vector bundle $`L_{\xi _0}L_{\xi _0}`$, where the elements of the dual torus $`\pm \xi _0\widehat{T}`$ are called the *asymptotic states* of the connection.
We show that the three standard examples above completely describe the behavior at infinity of doubly-periodic instantons with quadratic curvature decay:
###### Theorem 0.1.
Let $`A`$ be a doubly-periodic instanton with curvature $`O(r^2)`$. Then there is a gauge near infinity such that
$$A=A_0+a,$$
where $`A_0`$ is one of the previous models, and, for some $`\delta >0`$, in the semisimple case:
$$|a|=O\left(\frac{1}{r^{1+\delta }}\right),|_{A_0}a|=O\left(\frac{1}{r^{2+\delta }}\right);$$
in the nilpotent case:
$$|a|=O\left(\frac{1}{r(\mathrm{ln}r)^{1+\delta }}\right),|_{A_0}a|=O\left(\frac{1}{r^2(\mathrm{ln}r)^{2+\delta }}\right).$$
In the case where the limit at infinity of $`A`$ is non trivial, one can prove the theorem under the weaker assumption that the curvature is $`O(r^{1ϵ})`$; this condition is very close to the finite energy condition, and it is natural to suppose that the theorem actually describes the behaviour of all finite energy instantons. The instantons we will use (for example, those coming from the inverse Nahm transform) have quadratic curvature decay, so that this hypothesis is sufficient for our applications.
The theorem, to be proved in section 4, provides a complete characterization of the instanton parameters which are invariant under $`L^2`$ deformations. The parameter $`\lambda `$ is equivalent to the asymptotic states $`\pm \xi _0`$. The two remaining parameters are new: $`\alpha `$ is called the limiting holonomy of the instanton $`A`$, while $`\mu `$ is called the residue. The motivation for the latter nomenclature will be made clear latter on. Notice that, in contrast with the instanton number (see below) and the asymptotic states, the limiting holonomy and the residues are defined only for anti-self-dual connections.
Instantons and holomorphic bundles. We are now ready to state our second main result, which in particular solves the extensibility problem. Recall that the instanton number $`k`$ of the doubly-periodic instanton $`A`$ is defined by the formula:
$$k=\frac{1}{8\pi ^2}_{T\times }|F_A|^2$$
as usual.
###### Theorem 0.2.
There is a 1-1 correspondence between the following objects:
* $`SU_2`$-doubly-periodic instanton connections with quadratic curvature decay and fixed asymptotic parameters $`(k,\pm \xi _0,\alpha )`$;
* $`\alpha `$-stable, rank two holomorphic vector bundles $`T\times ^1`$ with trivial determinant such that $`c_2()=k`$ and $`|_{T\times \{\mathrm{}\}}=L_{\xi _0}L_{\xi _0}`$.
The stability condition of the statement is a variant of the stability condition for *parabolic bundles*; the degree is calculated with respect to a non ample class (the fundamental class of the torus). The precise definition will be given in section 5, where this result is proved.
In a broader context, theorem 0.2 can be seen as the analog of Donaldson’s correspondence between instantons on $`^4`$ and framed holomorphic bundles over $`^2`$ . In this last case, no stability condition is needed in order to produce an instanton, while in the case of a compact surface, stability (with respect to an ample class) is necessary. Thus, in some sense, our stability criterion goes midway between these two situations.
Moduli space. We then pass to the analytical construction of the moduli space of doubly-periodic instantons. We prove:
###### Theorem 0.3.
The moduli space of doubly-periodic instantons with fixed instanton number $`k`$ and asymptotic parameters $`(\pm \xi _0,\alpha ,\mu )`$ is a smooth hyperkähler manifold of real dimension $`8k4`$.
Of course, this theorem is interesting only if the moduli space is not empty. Fortunately, as mentioned in , existence of doubly-periodic instantons for generic values of the parameters $`(k,\pm \xi _0,\alpha ,\mu )`$ is guaranteed via the Nahm transform (see below) of meromorphic Higgs bundles over $`\widehat{T}`$, whose existence follows from Simpson among others; theorem 0.1 puts these instantons in our moduli spaces. Another equivalent, probably more direct, way for guaranteeing existence is of course theorem 0.2. See also section 5 for some cases where the moduli space is empty, and section 6 for a description of the $`k=1`$ moduli space.
Nahm transform. Finally we revisit the Nahm transform of doubly-periodic instantons defined in with two main objectives in mind. Before explaining what these objectives are, let us say a few words about the Nahm transform.
Here we restrict to the semisimple case, since Nahm transform was defined only in this case. Recall from (see also part III) the Nahm Transform is a 1-1 correspondence between irreducible, doubly-periodic instantons and certain meromorphic Higgs pairs $`(B,\mathrm{\Phi })`$ on a bundle $`V`$ over the dual torus $`\widehat{T}`$. The rank of $`V`$ is given by the instanton number. The Higgs field $`\mathrm{\Phi }`$ has simple poles at the two points corresponding to the asymptotic states $`\pm \xi _0`$. Moreover, $`\mathrm{\Phi }`$ has semisimple residues of rank one if $`\xi _0\xi _0`$, and rank two otherwise. We denote by $`\mathrm{Res}\mathrm{\Phi }(\pm \xi _0)`$ the residue of the Higgs field at the singular point $`\pm \xi _0`$.
Thus, it is natural to ask how are the new asymptotic parameters defined by theorem 0.1 interpreted in terms of the Nahm transformed Higgs pair. This question in answered in section 7, and the precise statement is given in theorem 7.1. As expected from the general principle Nahm transform is a non-linear Fourier transform, the asymptotic behavior is converted into singularity behavior.
It is well known that the moduli space of Higgs pairs on a Riemann surface is hyperkähler ; for the moduli space of Higgs pairs with fixed singularities at the punctures, this follows from . The second goal can now be summarized in our last result:
###### Theorem 0.4.
The Nahm transform of doubly-periodic instantons is a hyperkähler isometry.
Note that similar results have been proved for the other well-known examples of Nahm transform: the ADHM construction, see ; the duality between monopoles and solutions of Nahm equations, see ; and the Fourier-Mukai transform of instantons over 4-tori, see . Indeed, it is reasonable to expect that such result holds for any Nahm transform.
Outline. The paper is divided in three parts. The first part is technical: we study the asymptotic behavior of connections on $`E`$ with quadratic curvature decay, but which are not necessarily anti-self-dual; the technical goal is the construction of a partial Coulomb gauge (theorem 0.5). In the second part, we obtain theorems 0.1, 0.2 and 0.3. Finally, the third part deals with the Nahm transform aspects of the paper.
Acknowledgements. The second author would like to thank the École Polytechnique for its support, and Antony Maciocia for useful conversations.
## Part I Connections with quadratic curvature decay
In this part, we study the behaviour at infinity of (not necessarily anti-self-dual) connections with quadratic curvature decay on a $`SU_2`$-bundle $`E`$ on $`T\times ^2`$. Such connections will have a limit flat connection $`\mathrm{\Gamma }`$ on the torus at infinity $`T_{\mathrm{}}`$, which decomposes $`E|_T_{\mathrm{}}`$ as a sum of two flat line bundles $`L_{\xi _0}L_{\xi _0}`$; when $`L_{\xi _0}^2=0`$, we can reduce to the case where $`L_{\xi _0}=0`$ by globally tensoring $`E`$ with $`L_{\xi _0}`$; therefore we will always suppose that
$$\text{either }L_{\xi _0}^20,\text{ or }L_{\xi _0}=0.$$
(14)
Over any torus $`T`$, we consider the $`L^2`$-orthogonal decomposition
$$End(E)=(\mathrm{ker}_\mathrm{\Gamma })(\mathrm{ker}_\mathrm{\Gamma })^{}$$
(15)
and we decompose accordingly any section $`u`$ of $`End(E)`$ as
$$u=u_\mathrm{\Gamma }+u_{}.$$
(16)
If we write explicitly $`\mathrm{\Gamma }=d+\gamma `$, with
$$\gamma =\left(\begin{array}{cc}\lambda _1& 0\\ 0& \lambda _1\end{array}\right)dx+\left(\begin{array}{cc}\lambda _2& 0\\ 0& \lambda _2\end{array}\right)dy,$$
then, in view of (14), $`\mathrm{ker}_\mathrm{\Gamma }`$ is described as the $`T`$-invariant sections of $`\mathrm{ker}\gamma `$; if $`\gamma `$ is nontrivial, these are reduced to $`T`$-invariant diagonal matrices.
The technical goal of this part is a partial Coulomb gauge on the $`a_{}`$ part of a connection $`A=\mathrm{\Gamma }+a`$ with curvature $`O(r^2)`$. More precisely, let $`V_R`$ denote the complement of a disc of radius $`R`$ centered at the origin.
###### Theorem 0.5.
Given a constant $`\eta >0`$, there exists $`R`$ sufficiently large such that if $`A`$ is a doubly-periodic connection satisfying $`sup_{rR}\left(r^2|F_A|\right)\eta `$, then there is a gauge $`g`$ on $`T\times V_R`$ such that $`g(A)=\mathrm{\Gamma }+a_\mathrm{\Gamma }+a_{}`$, with:
(i) $`d_{\mathrm{\Gamma }+a_\mathrm{\Gamma }}^{}a_{}=0;`$
(ii) $`_r\mathrm{}a_{}(r=R)=0;`$
(iii) $`r^2F_{\mathrm{\Gamma }+a_\mathrm{\Gamma }}_{C^0}+r^{2ϵ}a_{}_{C^0}Cr^2F_A_{C^0}.`$
Note that gauge transformations $`g=g_\mathrm{\Gamma }`$ preserve the Coulomb gauge constructed in this theorem. This kind of partial gauge fixing reminds of Råde’s fibered Hodge gauge .
###### Remark 0.6.
Actually, if $`\mathrm{\Gamma }`$ is nontrivial, the proof gives a Coulomb gauge under a weaker bound on the curvature, namely $`|F|=O(r^{(1+ϵ)})`$; this condition is very close to the finite energy condition, since $`r^\delta `$ is in $`L^2`$ when $`\delta >1`$.
### 1 Limit flat connection
Our first task is to establish the existence of a flat limit connection $`\mathrm{\Gamma }`$ for every connection $`A`$ with quadratic curvature decay:
###### Proposition 1.1.
Suppose that the connection $`A`$ on $`T\times ^2`$ satisfies
$$|F_A|\frac{c_1}{r^2}.$$
Then $`A`$ has a flat limit $`\mathrm{\Gamma }`$ on $`T`$ at infinity, and there exists a sequence of connections $`A_j`$, such that
1. $`|F_{A_j}|c_2/r^2`$ ;
2. $`A_j`$ is gauge equivalent to $`A`$ on $`\{rj\}`$ ;
3. $`A_j=\mathrm{\Gamma }`$ on $`\{r2j\}`$.
###### Remark 1.2.
This proposition remains true if the curvature is $`O(r^{(1+ϵ)})`$.
###### Proof.
We begin by proving the existence of the flat limit $`\mathrm{\Gamma }`$. Take a radial gauge
$$A=d+a_\theta d\theta +a_xdx+a_ydy$$
for $`A`$; from the bound on the curvature, we deduce
$$|_ra_x|+|_ra_y|=O(r^2),|_ra_\theta |=O(r^1);$$
(17)
from this we deduce that $`a_x`$ and $`a_y`$ have limits $`a_x^{\mathrm{}}(\theta ,x,y)`$ and $`a_y^{\mathrm{}}(\theta ,x,y)`$ when $`r`$ goes to infinity; moreover, the bound on the curvature implies that for each $`\theta `$, the connection $`d+a_x^{\mathrm{}}(\theta )dx+a_y^{\mathrm{}}(\theta )`$ is flat on $`T`$. It remains to see that it is independent of $`\theta `$: for this we pick a base point in $`T=S^1\times S^1`$ and prove that the monodromies along the two circles remain conjugate when $`\theta `$ varies; this is a consequence of the bound on the curvature and the following lemma (see for example \[1, lemma 1\]):
###### Lemma 1.3.
Suppose we have a connection $`A`$ on $`[0,1]\times S^1`$, and $`m(t)`$ is the monodromy of $`A`$ along the circle $`\{t\}\times S^1`$; note $`h(t)`$ the parallel transport from the point $`(0,0)`$ to the point $`(t,0)`$; then
$$\left|_t\left(h(t)^1m(t)h(t)\right)\right|_{\{t\}\times S^1}|F_A|.$$
Therefore we have constructed a flat limit $`\mathrm{\Gamma }`$ on $`T`$ for the connection $`A`$. Now pass to the approximation statement.
###### Claim.
On $`\{r\}\times S^1\times T`$, there exists a gauge so that $`A=\mathrm{\Gamma }+a`$, $`|a|c/r`$.
This statement (a $`C^0`$ gauge only), can be proven by elementary means and is left to the reader. Now, we extend radially this gauge on $`\{j\}\times S^1\times T`$ to $`[j,2j]\times S^1\times T`$, and the bounds (17) imply that $`A=\mathrm{\Gamma }+a`$ with still $`|a|c/r`$ on $`[j,2j]`$; then we choose a cutoff function $`\chi =\chi (r)`$ so that
$$\chi (rj)=1,\chi (r2j)=0,|_r\chi |2/j,$$
and define a connection $`A_j`$ by
$`rj,`$ $`A_j=A,`$
$`rj,`$ $`A_j=\mathrm{\Gamma }+\chi a;`$
on $`rj`$, the curvature of $`A_j`$ is
$$F_{A_j}=\chi F_A+d\chi a+(\chi ^2\chi )aa$$
and this remains bounded by $`c/r^2`$ on $`[j,2j]`$, which means that $`|F_{A_j}|`$ is uniformly bounded by $`c/r^2`$. ∎
###### Remark 1.4.
Actually, it is not difficult to go a bit further and to prove that there is a global gauge in which $`A=\mathrm{\Gamma }+a`$ and $`|a|=O(\mathrm{ln}r/r)`$; this gives a result used without proof in . Of course, the result will also be a consequence of theorem 0.5.
In the case of a torus invariant connection, we need a stronger statement.
###### Proposition 1.5.
Under the hypotheses of proposition 1.1, if $`A=d+a`$ with $`a=a_\mathrm{\Gamma }`$ (in particular $`A`$ is torus invariant), then there is a gauge such that
$$A=\mathrm{\Gamma }+a+b,$$
where $`d+a`$ is a connection on $`^2`$ and $`b=b_xdx+b_ydy`$ a 1-form along $`T`$, satisfying
$$|b|\frac{c_3}{r},|_{\mathrm{\Gamma }+a}b|\frac{c_2}{r^2},$$
and
$$a=i\left(\begin{array}{cc}\alpha (r)& 0\\ 0& \alpha (r)\end{array}\right)d\theta +b$$
with
$$|_r\alpha |+|b|c_3/r,\underset{j}{sup}j^{22/p}b_{L^p(jr2j)}c_3.$$
The meaning is that we want a gauge with not only a $`C^0`$ bound, but also a $`C^1`$ bound; actually this is not possible (because elliptic regularity does not hold in $`C^k`$ spaces) and this explains why we use $`L^p`$ derivatives instead. So the proposition must be considered as a regularization of the connection. The standard way to obtain this is to use Hodge gauges in order for the curvature to become an elliptic equation: locally Uhlenbeck’s theorem provides the required statement, but the glueing is not easy, especially on a non simply connected manifold. We present here a proof based on the following lemma, which is a consequence of the Hodge gauge constructed in \[1, theorem 1\]:
###### Lemma 1.6.
Any connection $`A`$ on $`[0,1]\times S^1`$, with $`F_A_{L^p}`$ sufficiently small, is gauge equivalent to a connection $`d+i\alpha d\theta +a`$ with $`a_{L^{1,p}}cF_A_{L^p}`$, where $`\alpha `$ is a diagonal matrix, with coefficients in $`[0,1[`$, such that $`\mathrm{exp}(2\pi i\alpha )`$ is the monodromy of $`A`$ along the circle $`\{0\}\times S^1`$. ∎
###### Proof of proposition 1.5.
If $`A`$ is torus invariant, then a torus invariant gauge transformation $`g`$ acts on $`b`$ only by $`gbg^1`$, and the bounds on the curvature immediately imply the required bounds on $`b`$. Therefore, we are reduced to look at a connection $`d+a`$ on $`^2`$.
Now note that the region $`r1`$ is conformally equivalent to the half-cylinder $`_+\times S^1`$ (with coordinate $`t=\mathrm{ln}r`$); in the rest of the proof we will use only the flat metric on the cylinder. The bound on the curvature becomes $`|F_A|c_2`$; eventually pulling back $`A`$ using the transformation $`t\lambda t`$ with $`\lambda `$ sufficiently small, we may suppose that $`c_2`$ is very small. This means that we are now able to use lemma 1.6, for some $`p`$ very big, to produce on each $`[j1,j+1]\times S^1`$ a gauge $`g_j`$ so that
$$g_j(A)=d_{\alpha _j}+a_j,a_j_{L^p}+(+i\alpha _jd\theta )a_j_{L^p}cc_2.$$
We perform recursively diagonal gauge transformations with coefficients of type $`\mathrm{exp}(ik\theta )`$ ($`k`$ integer) so that we have
$$|\alpha _{j+1}\alpha _j|<c_2;$$
this is possible because of lemma 1.3, and the operation does not affect the bound on $`a_j`$ (but we have only $`|\alpha _j|c_2j`$). We want to glue together these local gauges: the transition $`h_j=g_{j+1}g_j^1`$ satisfies
$$dh_j+[\alpha _j,h_j]=h_ja_j(\alpha _{j+1}\alpha _j+a_{j+1})h_j;$$
the RHS is controled by $`cc_2`$, and this implies that $`h_j`$ is very close to some $`\stackrel{~}{h}_j(\theta )`$ in the kernel of $`d+\alpha _j`$; replacing $`g_{j+1}`$ by $`\stackrel{~}{h}_jg_{j+1}`$, we now may suppose that the transition $`g_{j+1}g_j^1`$ is close to the identity (in $`L^{2,p}`$ norm), and a standard argument now enables us to glue together all these gauges: for a similar argument, see \[1, pages 447–8\]. If we choose diagonal matrices $`\alpha (t)`$ so that
$$\alpha (j)=\alpha _j,|_t\alpha |2c_2,$$
we finally get a gauge $`d+i\alpha (t)d\theta +b`$, with
$$b_{L^p([j1,j+1])}+(+i\alpha (t)d\theta )b_{L^p([j1,j+1])}cc_2.$$
Sobolev embedding implies that $`b_{C^0}`$ is controled as well; translating back these bounds in the metric of $`^2`$, we get the proposition. ∎
###### Remark 1.7.
The proof of proposition 1.5 becomes certainly easier if $`A`$ is abelian (which is the case if the limit $`\mathrm{\Gamma }`$ is regular), since in this case, it is easy to produce a global Hodge gauge.
###### Remark 1.8.
In general, we are unable to prove proposition 1.5 if the curvature is only $`O(r^{(1+ϵ)})`$: this is because, in order to use get a controled gauge on $`_+\times S^1`$, we need the curvature to be bounded; if $`\mathrm{\Gamma }`$ is nontrivial, the problem becomes abelian, and then it is easy to construct a global Hodge gauge on $`_+\times S^1`$, from which the proposition follows easily (and one gets a bound in $`O(r^ϵ)`$ on $`a`$).
### 2 The linear problem
In this section we study the linear analysis on the $`(\mathrm{ker}_\mathrm{\Gamma })_{}`$ part for the Laplacian operator $`d_\mathrm{\Gamma }^{}d_\mathrm{\Gamma }`$ acting on 0-forms and the deformation operator $`d_\mathrm{\Gamma }^++d_\mathrm{\Gamma }^{}`$ acting on 1-forms, with fixed boundary conditions.
For this analysis, we will use the Sobolev spaces $`L^{p,k}`$ of functions with $`k`$ derivatives in $`L^p`$; the weighted Sobolev spaces $`L_\delta ^{p,k}`$ of functions $`f`$ such that $`(1+r^2)^{\delta /2}fL^{p,k}`$.
The basis of the analysis is the following simple lemma, which is an immediate consequence of the decomposition (15).
###### Lemma 2.1.
There is a constant $`c`$, depending on $`p`$, such that on each torus T, for any section $`u`$ of $`End(E)`$, we have:
$$_T|_\mathrm{\Gamma }u_{}|^pc_T|u_{}|^p.$$
(18)
Analysis on 0-forms
###### Lemma 2.2.
The Neumann problem on sections of $`End(E)_{}`$ on $`rR`$,
$$\{\begin{array}{c}\mathrm{\Delta }_\mathrm{\Gamma }u=v\hfill \\ _ru(r=R)=0\hfill \end{array}$$
(19)
is an isomorphism $`L_\delta ^{2,2}L_\delta ^2`$.
Proof. The solution $`u`$ of the Neumann problem is obtained by minimizing the functional
$$\frac{1}{2}|_\mathrm{\Gamma }u|^2u,v$$
in the space $`L^{1,2}`$; the minimization is possible because of the estimate (18); local elliptic regularity gives that the $`L^{1,2}`$-solution actually lives in $`L^{2,2}`$, and this gives the statement when there is no weight.
In the case we have a weight $`\delta `$, the following estimate holds:
$`{\displaystyle \mathrm{\Delta }_\mathrm{\Gamma }u,ur^{2\delta }}`$ $`=`$ $`{\displaystyle |_\mathrm{\Gamma }u|^2r^{2\delta }}+2{\displaystyle \frac{\delta }{r}}__ru,ur^{2\delta }`$
$``$ $`{\displaystyle (1\frac{\delta }{r})|_\mathrm{\Gamma }u|^2r^{2\delta }}{\displaystyle \frac{\delta }{r}}|u|^2r^{2\delta }`$
and using (18) we get, if $`R`$ is large enough,
$`\mathrm{\Delta }_\mathrm{\Gamma }u_{L_\delta ^2}u_{L_\delta ^2}`$ $``$ $`{\displaystyle \mathrm{\Delta }_\mathrm{\Gamma }u,ur^{2\delta }}`$
$``$ $`Cu_{L_\delta ^2}^2`$
and therefore
$$Cu_{L_\delta ^2}\mathrm{\Delta }_\mathrm{\Gamma }u_{L_\delta ^2}$$
which proves that the isomorphism persists between weighted $`L^2`$-spaces, at least if $`R`$ is large enough.
This would be enough for our applications, but one can prove easily that the statement remains true for any $`R`$: because $`\mathrm{\Delta }_\mathrm{\Gamma }`$ is an isomorphism for $`R`$ big enough, it remains a Fredholm operator for any $`R`$ (just glue the inverse near infinity with a parametrix on the compact part); the index is locally constant and therefore does not depend on the weight $`\delta `$; this means that it is equal to the $`L^2`$-index, that is $`0`$; now, because the $`L^2`$-kernel is zero, the $`L_\delta ^2`$-kernel is zero if $`\delta >0`$; for general $`\delta `$, the kernel is the $`L_\delta ^2`$-kernel, while the cokernel is the $`L_\delta ^2`$-kernel: as at least one of them is trivial and the index is $`0`$, both are trivial. ∎
We now want to deduce the same result in $`L^p`$ spaces. We need an estimate on the solution of problem (19) when $`v`$ is $`L^p`$. After a conformal change in the Euclidean metric $`g_E`$, we can pass to the cusp metric ($`r=e^t`$):
$$g_C=dt^2+d\theta ^2+e^{2t}(dx^2+dy^2)=\frac{1}{r^2}g_E$$
The operator $`\mathrm{\Delta }_\mathrm{\Gamma }`$ now has singular coefficients, but is basically of the type studied in , where Hölder and $`L^p`$ estimates are deduced from the $`L^2`$-estimates. Here, the same techniques lead to the desired result:
###### Lemma 2.3.
The Neumann problem (19) for 0-forms on $`rR`$ is an isomorphism $`L_\delta ^{2,p}L_\delta ^p`$ for all weights $`\delta `$.
Proof. For the convenience of the reader, we give here a sketch of proof for the statement, inspired from \[3, section 6\], but written with respect to the Euclidean metric. The proof below works for $`p>2`$ (the case we will use), but the statement remains true for general $`p`$.
The first step is to give an elliptic estimate
$$u_{L_\delta ^{2,p}[r,2r]}c\left(\mathrm{\Delta }_\mathrm{\Gamma }u_{L_\delta ^p([\frac{1}{2}r,3r])}+u_{L_{\delta 1+2/p}^2([\frac{1}{2}r,3r])}\right).$$
(20)
The weight
$$\delta _2=\delta 1+2/p$$
(21)
chosen for the $`L^2`$ space corresponds to functions with the same order of decreasing in $`r^{\delta 2/p}`$ as in the weighted $`L^p`$ space, but actually the proof below will give more. In order to prove this, we remark that
$$_\mathrm{\Gamma }=e^{i(ax+by)}e^{i(ax+by)}$$
so that if we consider $`x`$ and $`y`$ as coordinates on $`^2`$, the equation $`\mathrm{\Delta }_\mathrm{\Gamma }u=v`$ becomes equivalent to
$$\mathrm{\Delta }u^{}=e^{i(ax+by)}v,u^{}=e^{i(ax+by)}u.$$
In the domain $`[1,2]\times S^1\times [1,1]^2^2\times ^2`$, we have an elliptic estimate
$$u^{}_{L^{2,p}}c\left(u^{}_{L^2}+\mathrm{\Delta }u^{}_{L^p}\right)$$
which implies on the homothetic domain $`[R,2R]\times S^1\times [R,R]^2^2\times ^2`$
$$R^{24/p}^2u^{}_{L^p}c\left(R^2u^{}_{L^2}+R^{24/p}\mathrm{\Delta }u^{}_{L^p}\right)$$
and therefore on $`[R,2R]\times S^1\times T`$
$$R^{22/p}_\mathrm{\Gamma }^2u_{L^p}c\left(R^1u_{L^2}+R^{22/p}\mathrm{\Delta }_\mathrm{\Gamma }u_{L^p}\right)$$
which we can rewrite, still on $`[R,2R]\times S^1\times T`$,
$$_\mathrm{\Gamma }^2u_{L_\delta ^p}c\left(u_{L_{\delta 3+2/p}^2}+\mathrm{\Delta }_\mathrm{\Gamma }u_{L_\delta ^p}\right)$$
now the estimate (18) implies
$$_\mathrm{\Gamma }^ku_{L^p}cu_{L^p};$$
this, with local elliptic regularity, gives the estimate (20).
The second step now consists in going from the $`L^2`$-estimates with weights to the $`L^p`$-estimate. Basically, one can do the following: let $`P`$ be the inverse obtained by the $`L^2`$-resolution; decompose
$$v=v_i$$
(22)
where $`v_i`$ has support in $`\mathrm{exp}(i/2)<r<\mathrm{exp}(3i)`$; by the $`L^2`$-resolution for the weight $`\delta _2`$ defined by (21), one has
$$Pv_i_{L_{\delta _2}^2}cv_i_{L_{\delta _2}^2}cv_i_{L_\delta ^p};$$
on the other hand, we decompose similarly $`u_i=Pv_i`$ as
$$u_i=\underset{j}{}u_{ij},$$
and we note that the $`L^2`$ resolution gives the estimate
$`u_{ij}_{L_{\delta _2}^2}`$ $``$ $`ce^{ϵi}u_{ij}_{L_{\delta _2+ϵ}^2}`$
$``$ $`ce^{ϵi}v_i_{L_{\delta _2+ϵ}^2}`$
$``$ $`ce^{ϵ(ij)}v_i_{L_{\delta _2}^2};`$
if we choose $`ϵ`$ to be $`\pm ϵ`$ according to the sign of $`ij`$, we get the estimate
$`u_{ij}_{L_{\delta _2}^2}`$ $``$ $`ce^{ϵ|ij|}v_i_{L_{\delta _2}^2}`$
$``$ $`ce^{ϵ|ij|}v_i_{L_\delta ^p}`$
now, note $`\kappa _{ij}=1`$ if $`|ij|1`$ and $`0`$ otherwise; using (20), we deduce
$`u_{ij}_{L_\delta ^p}`$ $``$ $`c\left(\kappa _{ij}v_i_{L_\delta ^p}+e^{ϵ|ij|}v_i_{L_{\delta _2}^2}\right)`$
$``$ $`ce^{ϵ|ij|}v_i_{L_\delta ^p}`$
from which we deduce immediately
$$u_{L_\delta ^p}cv_{L_\delta ^p},$$
which proves, with the help of local elliptic regularity, that the operator is an isomorphism $`L_\delta ^{2,p}L_\delta ^p`$. ∎
###### Remark 2.4.
Actually, the proof gives a bit more, namely the norm of the inverse operator is bounded by a constant which is independent of $`R`$ ($`R`$ big enough); this is because we have explicit constants for the $`L^2`$ inverse, and the constants in the above proof do not depend on $`R`$.
###### Remark 2.5.
The same proof works in Hölder spaces, and gives an isomorphism between Hölder weighted spaces. In $`C^k`$ spaces, we have no more elliptic regularity; nevertheless, if $`v`$ is in $`C_\delta ^0`$, one can still deduce from the above proof the estimate
$$r^{\delta ϵ}u_{C^0}cr^\delta v_{C^0};$$
(23)
this estimate is not a consequence of the $`L^p`$ estimate, because the Sobolev embedding (which can be proven like the elliptic estimate (20) by a homothety argument),
$$u_{C_\delta ^0}c\left(u_{L_{\delta 2/p}^p}+u_{L_{\delta +12/p}^p}\right),$$
(24)
implies $`L_{\delta +12/p}^{1,p}C_\delta ^0`$, so that there is a loss of weight, since $`L_{\delta +12/p}^{1,p}`$ corresponds to functions $`O(r^{\delta 1})`$ when $`C_\delta ^0`$ corresponds to functions $`O(r^\delta )`$. Note also that in the case where $`v`$ lies in the component where $`\gamma `$ acts non trivially, the maximum principle provides directly the estimate (23) without the $`ϵ`$.
Analysis on 1-forms
In the next few lemmas, we assume that $`a`$ is a 1-form with values in $`End(E)`$ such that $`_r\mathrm{}a=0`$ on $`r=R`$. Again we suppose that $`a`$ is reduced to its component $`a_{}`$. All Sobolev norms are taken over the set $`T\times V_R=\{rR\}`$.
###### Lemma 2.6.
One has the identity
$$d_\mathrm{\Gamma }^{}a_{L^2}^2+d_\mathrm{\Gamma }a_{L^2}^2=_\mathrm{\Gamma }a_{L^2}^2_{r=R}\left|\frac{1}{r}\frac{}{\theta }\mathrm{}a\right|^2𝑑x𝑑y𝑑\theta $$
(25)
with respect to the Euclidean metric.
Proof. The equality follows from the Weitzenböck formula in the Euclidean metric:
$$d_\mathrm{\Gamma }^{}d_\mathrm{\Gamma }+d_\mathrm{\Gamma }d_\mathrm{\Gamma }^{}=_\mathrm{\Gamma }^{}_\mathrm{\Gamma }$$
Just integrate by parts and check the boundary terms. ∎
###### Lemma 2.7.
For any real function $`f`$ and any $`R>0`$, one has:
$$f(R)^2\frac{2}{R}_R^{R+1}(|_rf|^2+|f|^2)r𝑑r$$
(26)
The proof is left to the reader.
###### Lemma 2.8.
If $`R`$ is sufficiently large, then for some constant $`c`$:
$`d_\mathrm{\Gamma }^{}a_{L_\delta ^p}+d_\mathrm{\Gamma }a_{L_\delta ^p}`$ $``$ $`c_\mathrm{\Gamma }a_{L_\delta ^p}`$
$`d_\mathrm{\Gamma }^{}a_{C_\delta ^0}+d_\mathrm{\Gamma }a_{C_\delta ^0}`$ $``$ $`ca_{C_{\delta ϵ}^0}`$
with respect to the Euclidean metric.
###### Remark 2.9.
Remind that on the component $`a=a_{}`$ we look at, $`_\mathrm{\Gamma }a`$ controls $`a`$ by (18).
###### Proof.
From lemma 2.7 and lemma 2.1, we have:
$`{\displaystyle _{r=R}}|a|^2𝑑x𝑑y`$ $``$ $`{\displaystyle \frac{2}{R}}{\displaystyle _{[R,R+1]}}(|__ra|^2+|a|^2)r𝑑r𝑑x𝑑y`$
$``$ $`{\displaystyle \frac{C_1}{R}}{\displaystyle _{[R,R+1]}}|_\mathrm{\Gamma }a|^2𝑑x𝑑yr𝑑r`$
for some constant $`C_1`$; in particular
$$_{r=R}\left|\frac{1}{r}\frac{}{\theta }\mathrm{}a\right|^2𝑑x𝑑y𝑑\theta \frac{C_1}{R}_{[R,R+1]}|_\mathrm{\Gamma }a|^2𝑑x𝑑yr𝑑r𝑑\theta $$
and we deduce from lemma 2.6, for $`R`$ big enough,
$$d_\mathrm{\Gamma }^{}a_{L^2}^2+d_\mathrm{\Gamma }a_{L^2}^2\frac{1}{2}_\mathrm{\Gamma }a_{L^2}^2$$
which proves the $`L^2`$-estimate of the lemma.
The $`L^2`$-estimate with weights is proven in the same way. In the integration by parts, new terms appear because of the weight $`r^{2\delta }`$. However, as in the proof of lemma 2.2, these terms have all a coefficient $`O(r^1)`$ and therefore are a small perturbation if $`R`$ is large enough (note that we can take the same $`R`$ if the weight remains bounded).
Finally, one may deduce the $`L^p`$ and $`C^0`$ estimates from the $`L^2`$ estimates as in lemma 2.3 and remark 2.5, since the operator $`d_\mathrm{\Gamma }^{}+d_\mathrm{\Gamma }`$ has injective symbol, and the boundary condition $`_r\mathrm{}a=0`$ is an elliptic boundary condition. The proof is a slightly more complicated, because one has to compose the decomposition (22) with a $`L^2`$-projection on the image of the operator. ∎
###### Lemma 2.10.
The operator $`2d_\mathrm{\Gamma }^{}d_\mathrm{\Gamma }^++d_\mathrm{\Gamma }d_\mathrm{\Gamma }^{}`$ on 1-forms lying in $`\mathrm{\Omega }^1End(E)_{}`$, with Dirichlet condition on $`r=R`$, is an isomorphism in weighted Sobolev or Hölder spaces for all weights $`\delta `$.
###### Proof.
Again the Weitzenböck formula
$$2d_\mathrm{\Gamma }^{}d_\mathrm{\Gamma }^++d_\mathrm{\Gamma }d_\mathrm{\Gamma }^{}=_\mathrm{\Gamma }^{}_\mathrm{\Gamma }$$
gives the $`L^2`$-estimate (for forms vanishing on the boundary)
$`((2d_\mathrm{\Gamma }^{}d_\mathrm{\Gamma }^++d_\mathrm{\Gamma }d_\mathrm{\Gamma }^{})u,u)_{L^2}`$ $`=`$ $`_\mathrm{\Gamma }u_{L^2}^2`$
$``$ $`cu_{L^2}^2`$
from which the $`L^2`$-statement (without weight) follows immediately. One can then deduce weighted statements as in the proofs of lemmas 2.2 and 2.3. ∎
### 3 Existence of a Coulomb gauge
After the technical work of the previous section, we are finally in a position to establish theorem 0.5, the key analytical result of this paper. The first step is the nonlinear version of the Hölder estimate in lemma 2.8; the exponent $`p`$ is fixed, near infinity.
###### Lemma 3.1.
Given $`\eta _1`$ sufficiently small, if a connection $`A=\mathrm{\Gamma }+a_\mathrm{\Gamma }+a_{}`$ on $`rR`$ satisfies:
1. $`d_{\mathrm{\Gamma }+a_\mathrm{\Gamma }}^{}a_{}=0`$ ,
2. $`_r\mathrm{}a(r=R)=0`$ ,
3. $`r^ϵa_{C^0}\eta _1`$ ,
then:
$$r^2F_{\mathrm{\Gamma }+a_\mathrm{\Gamma }}_{C^0}+r^{2ϵ}a_{}_{C^0}+(_\mathrm{\Gamma }+a_\mathrm{\Gamma })a_{}_{L_{22/pϵ}^p}cr^2F_A_{C^0}.$$
(27)
###### Proof.
First, note that:
$$F_A=F_{\mathrm{\Gamma }+a_\mathrm{\Gamma }}+d_{\mathrm{\Gamma }+a_\mathrm{\Gamma }}a_{}+\frac{1}{2}[a_{},a_{}].$$
(28)
Therefore, using the decomposition in (15), we have:
$`\left(F_A\right)_\mathrm{\Gamma }`$ $`=`$ $`F_{\mathrm{\Gamma }+a_\mathrm{\Gamma }}+{\displaystyle \frac{1}{2}}\left([a_{},a_{}]\right)_\mathrm{\Gamma },`$ (29)
$`\left(F_A\right)_{}`$ $`=`$ $`d_{\mathrm{\Gamma }+a_\mathrm{\Gamma }}a_{}+{\displaystyle \frac{1}{2}}\left([a_{},a_{}]\right)_{},`$ (30)
from which the the estimates below follow:
$`r^2\left(F_A\right)_\mathrm{\Gamma }_{C^0}`$ $``$ $`r^2F_{\mathrm{\Gamma }+a_\mathrm{\Gamma }}_{C^0}ra_{}_{C^0}^2,`$ (31)
$`r^2\left(F_A\right)_{}_{C^0}`$ $``$ $`r^2d_\mathrm{\Gamma }a_{}_{C^0}r^2[a_\mathrm{\Gamma },a_{}]_{C^0}ra_{}_{C^0}^2.`$ (32)
Using $`C_2^0L_{22/pϵ}^p`$ and the estimate in lemma 2.8, we get:
$`r^2F_A_{C^0}`$ $``$ $`c\left(r^2F_{\mathrm{\Gamma }+a_\mathrm{\Gamma }}_{C^0}+r^{2ϵ}a_{}_{C^0}+(_\mathrm{\Gamma }+a_\mathrm{\Gamma })a_{}_{L_{22/pϵ}^p}\right)`$
$`c^{}\left(ra_{}_{C^0}^2+r^2[a_\mathrm{\Gamma },a_{}]_{C^0}\right);`$
from the third hypothesis, we have
$$ra_{}_{C^0}^2+r^2[a_\mathrm{\Gamma },a_{}]_{C^0}\eta _1r^{2ϵ}a_{}_{C^0};$$
if $`\eta _1`$ is small enough, these two inequalities give the required estimate. ∎
###### Lemma 3.2.
Given $`\eta `$, there exists $`R`$ such that if $`A`$ is a connection over $`T\times V_R`$ such that $`A\mathrm{\Gamma }`$ is compactly supported and $`|F_A|\eta r^2`$, then there is a gauge $`g`$ such that $`g(A)=_\mathrm{\Gamma }+a_\mathrm{\Gamma }+a_{}`$, with:
(i) $`d_{\mathrm{\Gamma }+a_\mathrm{\Gamma }}^{}a_{}=0,`$
(ii) $`_r\mathrm{}a(r=R)=0,`$
(iii) $`r^2F_{\mathrm{\Gamma }+a_\mathrm{\Gamma }}_{C^0}+r^{2ϵ}a_{}_{C^0}+(_\mathrm{\Gamma }+a_\mathrm{\Gamma })a_{}_{L_{22/pϵ}^p}cr^2F_A_{C^0}.`$
Proof. We now have all the necessary ingredients for a proof by continuity. Consider the homothety $`\varphi _t(r)=e^tr`$ and the connections $`A_t=\varphi _t^{}A`$. We have $`A_0=A`$ and, for $`t`$ big enough, say $`tT`$, $`A_t=d_\mathrm{\Gamma }`$ because of the assumption on compact support. Moreover, it is clear from the form of the metric that
$$|F_{A_t}|=|\varphi _t^{}F_A|\varphi _t^{}|F_A|\frac{ce^{2t}}{r^2}$$
so that the whole path of connections $`(A_t)`$ satisfies the hypothesis of the lemma. Moreover, after gauge transformation, we can also assume that $`A_t=\mathrm{\Gamma }+a_t`$ with $`_r\mathrm{}a_t(r=R)=0`$ for all $`t`$.
We prove that the subset $`S[0,T]`$ containing all the values of $`t`$ for which the theorem holds for $`A_t`$ is both closed and open. Since $`S`$ is nonempty (it contains $`t=T`$), $`S`$ must be the whole interval and the result holds for $`t=0`$.
The closedness is trivial, since the estimate on the connection provides all the needed bounds.
For openness, first remark that proposition 1.5 provides a gauge in which
$`r^2F_{d_\mathrm{\Gamma }+a_\mathrm{\Gamma }}_{C^0}`$ $``$ $`c{\displaystyle \frac{r}{\mathrm{ln}r}}a_\mathrm{\Gamma }_{C^0}`$
$``$ $`c{\displaystyle \frac{R^{1ϵ}}{\mathrm{ln}R}}r^ϵa_\mathrm{\Gamma }_{C^0};`$
on the other hand, from (iii),
$`r^2F_A_{C^0}`$ $``$ $`cr^{2ϵ}a_{}_{C^0}`$
$``$ $`cR^{22ϵ}r^ϵa_{}_{C^0};`$
we deduce
$$r^ϵa_{C^0}c^1R^{(12ϵ)}\eta ;$$
(33)
taking $`R`$ big enough so that the RHS is smaller than $`\eta _1`$ of lemma 3.1, we see that (i) and (ii) imply (iii).
It remains to solve problem (i)-(ii) near a solution. Fix some $`t`$ and suppose that $`g_t(A_t)=\mathrm{\Gamma }+b`$ with $`\mathrm{\Gamma }+b`$ satisfying (i), (ii) and (iii). If we have a connection $`\mathrm{\Gamma }+b+\varpi `$ with $`_r\mathrm{}\varpi (r=R)=0`$, we want to find a gauge $`g`$ such that:
$$\{\begin{array}{c}g(\mathrm{\Gamma }+b+\varpi )=\mathrm{\Gamma }+c_\mathrm{\Gamma }+c_{}\hfill \\ d_{\mathrm{\Gamma }+c_\mathrm{\Gamma }}^{}c_{}=0\hfill \end{array}$$
Looking at solutions of the form $`g=e^u_{}`$, the equation to be solved is:
$$L(u_{},\varpi )=d_{\mathrm{\Gamma }+c_\mathrm{\Gamma }}^{}\left(e^u(\mathrm{\Gamma }+c_\mathrm{\Gamma }+c_{})e^ud_{\mathrm{\Gamma }+c_\mathrm{\Gamma }}(e^u)e^u\right)=0;$$
we would like to solve this equation with $`u_{}`$ in $`C^2`$, but $`C^k`$ spaces are not suitable for elliptic analysis; instead, we use weighted $`L^p`$ spaces with $`p`$ very big; since we have the freedom to apply a $`\mathrm{\Gamma }`$-invariant gauge transformation, using proposition 1.5, we can choose a gauge in which the derivatives of $`b_\mathrm{\Gamma }`$ are also controled, and therefore the operator $`L`$ is well defined; its linearization along the first variable is given by the operator:
$$ud_\mathrm{\Gamma }^{}d_\mathrm{\Gamma }u+\text{perturbation};$$
if $`R`$ is big enough, the perturbation is sufficiently small and we get an isomorphism by lemma 2.2.
This completes the proof. ∎
#### Completing the proof of theorem 0.5.
Our final task is to remove from lemma 3.2 the assumption that $`A\mathrm{\Gamma }`$ is compactly supported.
Using proposition 1.1, we approximate the connection $`A`$ by a sequence $`A_i`$ such that $`\mathrm{\Gamma }A_i`$ is compactly supported, and $`r^2F_{A_i}_{C^0}`$ remains bounded.
We can apply lemma 3.2 to each connection $`A_i`$, thus obtaining a gauge $`g_i`$ such that $`g_i(A_i)=d_\mathrm{\Gamma }+a_i`$, and $`a_i`$ satisfies (i)–(iii) of lemma 3.2. Using proposition 1.5 for the $`(a_i)_\mathrm{\Gamma }`$ part, the $`(a_i)`$ converge (weakly) to a limit $`a`$ still satisfying (i)–(iii), such that $`d_\mathrm{\Gamma }+a`$ is gauge equivalent to $`A`$. ∎
## Part II Instantons, holomorphic bundles, and the moduli space
So far, $`A`$ has simply been a connection on $`ET\times `$ with quadratic curvature decay. From now on, we shall assume that $`A`$ is also an instanton.
### 4 Asymptotic behavior: proof of theorem 0.1
Let us now assume that $`A`$ is a doubly-periodic instanton connection. Using theorem 0.5, if $`R`$ is big enough, we can put it in a Coulomb gauge on $`rR`$, so that $`A=\mathrm{\Gamma }+a_\mathrm{\Gamma }+a_{}`$, with $`a_\mathrm{\Gamma }`$ and $`a_{}`$ satisfying the Coulomb gauge equation,
$$d_{\mathrm{\Gamma }+a_\mathrm{\Gamma }}^{}a_{}=0,$$
and the anti-self-duality equation,
$$d_\mathrm{\Gamma }^+a+\frac{1}{2}[a,a]^+=0.$$
These can be rewritten as follows:
$`d_\mathrm{\Gamma }^{}a_{}`$ $`=`$ $`a_\mathrm{\Gamma }^{}a_{}`$ (34)
$`d_\mathrm{\Gamma }^+a_{}`$ $`=`$ $`[a_\mathrm{\Gamma },a_{}]^+{\displaystyle \frac{1}{2}}[a_{},a_{}]_{}^+`$ (35)
$`d^+a_\mathrm{\Gamma }+{\displaystyle \frac{1}{2}}[a_\mathrm{\Gamma },a_\mathrm{\Gamma }]^+`$ $`=`$ $`[a_{},a_{}]_\mathrm{\Gamma }^+`$ (36)
Now let $`\chi =\chi (r)`$ be a smooth cut-off function supported on $`T\times V_R`$; we have, using equations (35) and (36):
$$\left(d_\mathrm{\Gamma }^++d_\mathrm{\Gamma }^{}\right)(\chi a_{})=\chi (a_\mathrm{\Gamma }a_{}+a_{}a_{})+d\chi a_{}$$
(37)
where $``$ denotes some bilinear operations.
From theorem 0.5 and proposition 1.5, we already know that $`|a_{}|=O(r^{2+ϵ})`$ and that we can choose a gauge such that $`|a_\mathrm{\Gamma }|=O(\mathrm{ln}r/r)`$. We now apply lemma 2.10 to the equation (37): a priori the lemma applies to the laplacian $`(d_\mathrm{\Gamma }^+)^{}d_\mathrm{\Gamma }^++d_\mathrm{\Gamma }d_\mathrm{\Gamma }^{}`$ but the estimates also imply estimates for the first order elliptic operator $`d_\mathrm{\Gamma }^++d_\mathrm{\Gamma }^{}`$ (alternatively one may take one derivative of equation (37) and use the bounds on the derivatives of $`a_{}`$ and $`a_\mathrm{\Gamma }`$); the RHS of equation (37) is $`O(r^{3+ϵ})`$, therefore $`|a_{}|=O(r^{3+ϵ_2})`$, where $`ϵ_2>ϵ`$; by the same argument, we have that $`|a_{}|=O(r^{4+ϵ_3})`$, etc. Therefore, $`|a_{}|=O(r^\delta )`$ for any $`\delta >0`$.
Now come back to equation (34): it now means that $`d+a_\mathrm{\Gamma }`$ satisfies the instanton equation up to a term which goes very quickly to $`0`$ at infinity; as $`a_\mathrm{\Gamma }`$ is translation invariant, this means, by dimensional reduction, that $`d+a_\mathrm{\Gamma }`$ is a solution of Hitchin’s equations for Higgs bundles on $`^2`$ near infinity, up to a term decaying quicker than any $`O(r^\delta )`$. The behavior of the solutions of Hitchin’s equations near a singularity has been studied by Simpson , Biquard . The arguments in these papers are not affected by a very quickly decaying perturbation. Moreover, the bounds in proposition 1.5 implies that the Higgs field is $`O(1/r)`$ at infinity, so that the Higgs bundle is “tame” in Simpson’s terminology. Finally, we deduce from these articles that $`d+a_\mathrm{\Gamma }`$ is close to one of the examples described in the introduction, in the sense of theorem 0.1. ∎
### 5 Holomorphic extension
The theorem 0.1 proves that any instanton $`A`$ with quadratic curvature decay can be put in a gauge near infinity so that
$$A=A_0+a,$$
where $`A_0`$ is one of the model torus invariant instantons induced by model Higgs bundles, and $`a`$ is a small perturbation.
Local aspects
Let us now restrict to the semisimple case. Therefore, we have
$$A_0=d+i\left(\begin{array}{cc}a_0& 0\\ 0& a_0\end{array}\right)$$
with
$$a_0=\lambda _1dx+\lambda _2dy+(\mu _1\mathrm{cos}\theta \mu _2\mathrm{sin}\theta )\frac{dx}{r}+(\mu _1\mathrm{sin}\theta +\mu _2\mathrm{cos}\theta )\frac{dy}{r}+\alpha d\theta ;$$
observe that the (0,1)-part of this form is
$$a_0^{0,1}=\lambda d\overline{z}+\mu \frac{d\overline{z}}{w}\frac{\alpha }{2}\frac{d\overline{w}}{\overline{w}},\lambda =\frac{\lambda _1+i\lambda _2}{2},\mu =\frac{\mu _1+i\mu _2}{2},$$
so there is a singularity in the direction of transverse disks to the torus at infinity. We first reduce to a normal form on transverse disks.
###### Lemma 5.1.
Near the torus at infinity, there exists a continuous complex gauge transformation $`g`$, such that
1. $`g|_T_{\mathrm{}}=1`$ ;
2. $`|_{A_0}gg^1|=O(r^{(1+\delta )})`$ (and $`g_{}`$ is $`O(r^\delta )`$ for any $`\delta `$);
3. $`g(\overline{}_A)=A_0+bd\overline{z}`$, with $`b=O(r^{(1+\delta )})`$.
###### Proof.
We give a concise proof, since this is parallel to \[3, section 9\]. Remark that
$$\overline{}_\alpha =\overline{}\frac{\alpha }{2}\frac{d\overline{w}}{\overline{w}}=r^\alpha \overline{}r^\alpha ;$$
(38)
now the problem to be solved is
$$\frac{\overline{}_{A_0}}{\overline{w}}gga=0,$$
that is, using $`g=1+u`$,
$$\left(\frac{}{\overline{w}}\frac{1}{2}\left(\begin{array}{cc}\alpha & 0\\ 0& \alpha \end{array}\right)\right)uua=a;$$
this is a $`\overline{}`$-problem on small disks near infinity; for the model problem (38) the Cauchy formula gives us an explicit solution; in general, with the small perturbation $`a`$, the solution is produced by a fixed point theorem, and we even have an estimate
$$supr^\delta |u|csupr^{1+\delta }|a|;$$
one can then deduce the regularity statement on $`u`$. ∎
Note $`b_{jk}`$ the coefficients of the matrix $`b`$ above. Let $`(e_1,e_2)`$ be the orthonormal basis for the trivialisation of the bundle near infinity. From the lemma and equation (38), we deduce that the sections
$$(\sigma _1=r^\alpha g(e_1),\sigma _2=r^\alpha g(e_2))$$
(39)
are holomorphic on transverse disks, and, moreover, in the basis $`(\sigma _1,\sigma _2)`$, we now have
$$\overline{}_A=\overline{}+\left(\begin{array}{cc}\lambda & 0\\ 0& \lambda \end{array}\right)d\overline{z}+\left(\begin{array}{cc}\mu & 0\\ 0& \mu \end{array}\right)\frac{d\overline{z}}{w}+\left(\begin{array}{cc}b_{11}& r^{2\alpha }b_{12}\\ r^{2\alpha }b_{21}& b_{22}\end{array}\right)d\overline{z},$$
(40)
with all coefficients of the last matrix holomorphic in $`w`$. From this, we see immediately that in the basis $`(\sigma _1,\sigma _2)`$, the operator (40) defines a holomorphic extension $``$ over $`T\times ^1`$.
Since
$$|\sigma _1|r^\alpha ,|\sigma _2|r^\alpha ,$$
we see that, from an intrinsic point of view, if $`\alpha <1/2`$, the local holomorphic sections of $``$ are characterized as the local holomorphic sections $`\sigma `$ outside $`T_{\mathrm{}}`$ satisfying the growth condition
$$|\sigma |=O(r^\alpha ).$$
(41)
When $`0<\alpha <1/2`$, this global extension has a subbundle $``$ over the torus at infinity, given by the values of the local holomorphic sections $`\sigma `$ satisfying the growth condition
$$|\sigma |=O(r^\alpha ).$$
(42)
Therefore, the growth of the holomorphic sections at infinity determine a “parabolic structure”
$$0,$$
with weights $`\alpha <\alpha `$ (the sign is changed because the local coordinate near infinity is $`w^1`$).
Actually one can say more : over $`T_{\mathrm{}}`$, the $`\overline{}`$-operator (40) is
$$\overline{}+\left(\begin{array}{cc}\lambda & 0\\ 0& \lambda \end{array}\right)d\overline{z},$$
which means that
$$|_T_{\mathrm{}}=L_{\xi _0}L_{\xi _0}.$$
Of course, if $`\alpha `$ is nontrivial, then $`=L_{\xi _0}`$ is canonically determined by the growth condition (42).
Actually, the decomposition $`L_{\xi _0}L_{\xi _0}`$ can almost always be made canonical: this is clear if $`\xi _00`$, and in this case, since the off-diagonal components of the connection decay quicker than any $`O(r^\delta )`$, we deduce from equation (40) that, still in the basis $`(\sigma _1,\sigma _2)`$,
$$\overline{}_A=\overline{}+\left(\begin{array}{cc}\lambda & 0\\ 0& \lambda \end{array}\right)d\overline{z}+\left(\begin{array}{cc}\mu & 0\\ 0& \mu \end{array}\right)\frac{d\overline{z}}{w}+O(r^2);$$
(43)
this gives the asymptotic behavior of $`|_{T_w}`$ when $`w`$ goes to infinity.
Moreover, when $`\xi _0=0`$, we still get something from (40): since the coefficients are holomorphic in $`w`$, we note $`b_{12}^{}`$ the coefficient of $`r^{2\alpha }b_{12}`$ on $`w^1`$ (in the case $`\alpha =0`$, we simply have $`b_{12}^{}=0`$), so
$$\overline{}_A=\overline{}+\left(\begin{array}{cc}\mu & b_{12}^{}\\ 0& \mu \end{array}\right)\frac{d\overline{z}}{w}+O(r^2);$$
(44)
if $`\mu 0`$, the matrix appearing above can always been diagonalized with eigenvalues $`\pm \mu `$, which means that up to changing $`\sigma _2`$ by some multiple of $`\sigma _1`$, we are reduced to (43) so that a supplementary subspace of $``$ is still well defined (and when $`\alpha =0`$, the decomposition $``$ still makes sense, as the eigenspaces of this matrix).
Note also that, as a consequence of (39), since $`g`$ is continuous, the unitary extension (given by the basis $`(e_1,e_2)`$ of the Coulomb gauge) and the holomorphic extension are topologically isomorphic.
Therefore, we have proven the following proposition.
###### Proposition 5.2.
In the semisimple case, for $`\alpha <1/2`$, if $`A`$ is a doubly-periodic instanton connection satisfying $`|F_A|=O(r^2)`$, then $`A^{0,1}`$ has a unique holomorphic extension $``$ over $`T\times ^1`$, whose holomorphic sections satisfy the growth condition (41). Moreover, one has $`c_2()=k`$ and a decomposition (if $`\lambda `$ or $`\mu `$ is nonzero) $`|_T_{\mathrm{}}=L_{\xi _0}L_{\xi _0}`$.∎
###### Remark 5.3.
Note that when $`\alpha =1/2`$, we cannot get a $`Sl_2`$-extension this way: indeed we could equally well choose the sections $`(w\sigma _1,\sigma _2/w)`$, giving a different extension. One way to construct a canonical extension is to use (41) with $`\alpha =1/2`$, which furnishes a $`Gl_2`$-extension where all nonzero sections have norm $`O(r^{1/2})`$. Also, a $`Sl_2`$-extension can be constructed if $`\xi _0\xi _0`$, by *deciding* that sections with nonzero values in $`L_{\pm \xi _0}`$ have norm like $`r^{1/2}`$.
In the sequel we will ignore this case, but all the statements can be easily adapted to it.
###### Remark 5.4.
In the nilpotent case (then $`\lambda `$, $`\mu `$, and $`\alpha `$ are trivial), the result is the same, but (as in the case of Higgs bundles) the growth of the holomorphic sections at infinity is now logarithmic:
$$|\sigma |=O\left((\mathrm{ln}r)^{\frac{1}{2}}\right),$$
(45)
and there is a line subbundle $``$ defined by the growth condition
$$|\sigma |=O\left((\mathrm{ln}r)^{\frac{1}{2}}\right).$$
(46)
The subbundle $``$ has no canonical supplementary subspace. The tools in \[3, section 9\] handle this situation as well.
Also observe that the $`\overline{}`$-operator for the model instanton (13) is (in an orthonormal basis $`(e_1,e_2)`$)
$$\overline{}+\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right)\frac{d\overline{w}}{2\overline{w}\mathrm{ln}r^2}+\frac{1}{r\mathrm{ln}r^2}\left(\begin{array}{cc}0& e^{i\theta }d\overline{z}\\ 0& 0\end{array}\right)$$
which gives, in the basis $`(e_1/(\mathrm{ln}r^2)^{\frac{1}{2}},e_2(\mathrm{ln}r^2)^{\frac{1}{2}})`$,
$$\overline{}+\left(\begin{array}{cc}0& \frac{d\overline{z}}{w}\\ 0& 0\end{array}\right);$$
in particular, $`|_{T_w}`$ is the nontrivial extension of $`\underset{¯}{}`$ by $`\underset{¯}{}`$; it is easy to see that this remains true for instantons, asymptotic to this nilpotent model.
Non-existence results
The proposition 5.2 gives obstructions for the existence of instantons. Here are some examples.
###### Lemma 5.5.
There are no instantons with $`\xi _0=\xi _0`$ and $`k=1`$.
###### Proof.
For a contradiction, let $`A`$ be an instanton with $`\xi _0=\xi _0`$ and $`k=1,2`$, and consider the extended holomorphic bundle $``$ given by theorem 0.2. The restriction of $``$ to the elliptic fibres $`T_p`$ must be semistable for all $`p^1`$ (see ). Moreover, $`|_{T_p}`$ cannot be generically the nontrivial extension of $`\underset{¯}{}`$ by itself, since this would give a non-constant map from $`^1`$ to $``$ (which parametrises the extensions of $`\underset{¯}{}`$ by itself).
Therefore, as shown in , index theory tells us that for each $`\xi \widehat{T}`$:
$$\mathrm{\Sigma }_{w^1}h^0(T_w,L_\xi |_{T_w})=k$$
(47)
But if $`|_T_{\mathrm{}}=L_{\xi _0}L_{\xi _0}`$, then $`h^0(T_{\mathrm{}},L_{\xi _0}|_T_{\mathrm{}})=2`$, thus contradicting the assumption that $`k=1`$. ∎
###### Lemma 5.6.
There are no instantons with $`\xi _0\xi _0`$ and $`\mu =0`$.
###### Proof.
The lemma is a consequence of the Nahm transform of doubly-periodic instantons defined in , more exactly of its holomorphic aspects; we anticipate a bit here, but see the introduction to Part III for a summary of the construction.
Again for a contradiction, let $`A`$ be an instanton with $`\mu =0`$ and asymptotic state $`\pm \xi _0`$ not of order two. The corresponding Nahm transformed Higgs field $`\mathrm{\Phi }`$ has simple poles at $`\pm \xi _0`$; its residues have rank one. However, as we shall see in the proof of theorem 7.1, the non-zero eigenvalues of the residues of $`\mathrm{\Phi }`$ are exactly $`\pm \mu `$, and more generally, the eigenvalues of $`\mathrm{\Phi }`$ at $`\xi \widehat{T}`$ are the $`w`$ such that $`H^0(T_w,L_\xi )0`$; hence, the vanishing of $`\mu `$ implies that the eigenvalues of $`\mathrm{\Phi }`$ remain bounded when $`\xi `$ goes to $`\xi _0`$.
Now if $`\xi _0\xi _0`$ then $``$ remains isomorphic to some $`L_\xi L_\xi `$ on each torus near infinity. It is then clear (again, see the proof of theorem 7.1) that the eigenvalues of $`\mathrm{\Phi }`$ must go to infinity and we get a contradiction. ∎
Global aspects, stability
More subtil obstructions come from stability properties. We investigate this for the extension $``$ of an instanton $`A`$ with quadratic curvature decay. Notice that by theorem 0.1, in the semisimple case, the curvature is only $`O(r^2)`$, but
$$|\iota _{\{\}\times }F_A|+|\iota _{T\times \{\}}F_A|=O\left(r^{(2+ϵ)}\right);$$
(48)
in the nilpotent case, we have
$$|F_A|=O\left(r^2(\mathrm{ln}r)^2\right);$$
(49)
the point here is that these two controlling factors are in $`L^1`$, whence $`F_A`$ itself is not $`L^1`$: this will enable us to define a degree.
The degree of a saturated subsheaf $`L`$ of $``$ with respect to the Euclidean Kähler form $`\omega `$ is \[17, lemma 3.2\]
$$2\pi \mathrm{deg}L=i\text{tr}(\pi F_A)\omega |\overline{}\pi |^2$$
(50)
where $`\pi `$ is orthogonal projection on $`L`$; from (48) and (49), this can be $`\mathrm{}`$ or a real number; in the last case, $`\overline{}\pi `$ is in $`L^2`$: this condition must be analyzed more precisely.
Again, we now restrict to the semisimple case (see remark 5.10 for the nilpotent case), so that $`|_T_{\mathrm{}}=L_{\xi _0}L_{\xi _0}`$, with weights $`\alpha `$ and $`\alpha `$, and behavior (43). In this case, we have near infinity
$$|_{T_w}=L_{\xi (w)}L_{\xi (w)}.$$
(51)
###### Lemma 5.7.
Suppose $`\alpha 0`$, then the degree of a subsheaf $`L`$ of $``$ is finite if and only if
1. $`L|_T_{\mathrm{}}`$ is flat; in particular, if $`L_{\xi _0}L_{\xi _0}`$, this means that $`LL_{\pm \xi _0}`$ ;
2. if $`L|_T_{\mathrm{}}=L_{\xi _0}`$, then $`L|_{T_w}L_{\xi (w)}`$ up to first order near infinity.
Now suppose $`\alpha =0`$, then the degree of a subsheaf $`L`$ of $``$ is finite if and only if $`L|_{T_w}L_{\pm \xi (w)}`$ up to first order near infinity.
###### Remark 5.8.
The first order condition can be seen as a reminiscence of the approximating Higgs bundle at infinity; indeed the Higgs field has eigenspaces $`L_{\pm \xi (w)}`$ and for Higgs bundle stability, one looks only at subsheaves stable under the action of the Higgs field.
###### Proof.
We analyze the situation locally near infinity; in the decomposition (51), the metric is approximately
$$\left(\begin{array}{cc}r^{2\alpha }& 0\\ 0& r^{2\alpha }\end{array}\right),$$
and we will simplify the problem by using this metric to make the calculations (the correction term can be easily bounded); at a point on $`T_{\mathrm{}}`$ where $`L`$ is a subbundle, we suppose for example that $`L`$ is not contained in $`L_{\xi _0}`$; choose a local flat section $`\sigma `$ for $`L_{\xi _0}`$, and note $`\sigma ^t`$ the dual flat section of $`L_{\xi _0}`$; extend $`\sigma `$ near $`T_{\mathrm{}}`$, keeping it parallel on $`T_w`$ (this is possible with our approximation for the metric); locally, $`L`$ is generated by $`s=\sigma +f\sigma ^t`$, where $`f`$ is holomorphic, and an orthogonal section is given by $`t=r^{2\alpha }\sigma \overline{f}\sigma ^tr^{2\alpha }`$, and
$$\overline{}_Tt=(\overline{}_T\overline{f})\sigma ^tr^{2\alpha },$$
from which we deduce
$$\pi (\overline{}_Tt)=\frac{\overline{}_T\overline{f}}{r^{2\alpha }+|f|^2r^{2\alpha }}s,$$
and finally, since our choice of $`t`$ satisfies $`|s|=|t|`$, and $`f`$ is holomorphic,
$$|\overline{}_T\pi |=\frac{|d_Tf|}{r^{2\alpha }+|f|^2r^{2\alpha }};$$
in order for $`\overline{}_T\pi `$ to be in $`L^2`$, it is necessary that $`d_Tf=0`$ on $`T_{\mathrm{}}`$, and therefore $`L`$ is constant.
Now restrict to the case of nontrivial decomposition $`L_{\xi _0}L_{\xi _0}`$ (the other cases are similar); therefore we may suppose that $`f=0`$ on $`T_{\mathrm{}}`$; if the first order term of $`d_Tf`$ does not vanish, then
$$|\overline{}_T\pi |r^{1+2\alpha }$$
this still is not in $`L^2`$ if $`\alpha 0`$ (but it is in $`L^2`$ if $`\alpha <0`$, which corresponds to the case $`L|_T_{\mathrm{}}L_{\xi _0}`$); this means that we need $`d_Tf`$ to vanish up to first order.
Concerning $`\overline{}_{}\pi `$, it is easy to verify that the $`L^2`$-condition is always satisfied. ∎
Recall that
$$F_L=\pi F_A\pi +\overline{}\pi \pi .$$
(52)
When the degree is finite, that is when $`\overline{}\pi `$ is $`L^2`$, the restriction of $`\omega `$ to $``$ does not contribute: indeed, $`dwd\overline{w}=\overline{}|w|^2`$, and this leads to
$$_{rR}F_Ldwd\overline{w}=_{r=R}w𝑑wF_L$$
but using (48) and (52), we see that this goes to zero as $`R`$ goes to infinity. Then we can rewrite the degree (denoting $`\overline{}_{}`$ the $`\overline{}`$ operator in the $``$ direction)
$$2\pi \mathrm{deg}L=i\pi F_A\omega _T|\overline{}_{}\pi |^2,$$
(53)
and this in turn is easily interpreted \[2, (4.1)\] as a “parabolic degree”:
$$\mathrm{deg}L=\{\begin{array}{c}c_1(L)[t]+\alpha [\omega _T],[t]\text{ if }L_T_{\mathrm{}}L_{\xi _0},\hfill \\ c_1(L)[t]\alpha [\omega _T],[t]\text{ if }L_T_{\mathrm{}}L_{\xi _0},\hfill \end{array}$$
(54)
where $`[t]`$ is the fundamental class of $`T`$ and $`\omega _T`$ the given Kähler form on $`T`$; of course this is not a degree in the usual sense on $`T\times ^1`$, since we use the non ample class $`[t]`$.
Define $`\alpha `$-*stability* of $``$ as the fact that any subsheaf satisfying the condition of lemma 5.7 has negative degree (we shall forget the $`\alpha `$ when there is no ambiguity); standard arguments give us
###### Proposition 5.9.
If $`A`$ is an instanton with quadratic curvature decay, then the holomorphic extension $``$ is $`\alpha `$-stable. ∎
###### Remark 5.10.
In the nilpotent case, the proposition remains true; here $`\alpha =0`$, and, following the proof of lemma 5.7, the degree is finite for all subsheaves with flat restriction to $`T_{\mathrm{}}`$.
###### Remark 5.11.
It is important to note that the stability condition just defined is not an empty one. Indeed, $`\alpha `$-unstable bundles $`T\times ^1`$ can be obtained as extensions in the following way:
$$0p_1^{}L_{\xi _0}p_2^{}𝒪_^1(b)p_1^{}L_{\xi _0}p_2^{}𝒪_^1(b)_k0$$
where $`b>0`$ and $`_k`$ is the ideal sheaf of $`k>0`$ points in $`T\times ^1`$, and we assume that none of these points are in $`T_{\mathrm{}}`$. Every sheaf $``$ so obtained is locally-free, since the sheaf on the LHS is locally-free and the one on the RHS is torsion-free. Clearly, $``$ has trivial determinant, instanton number $`k`$ and asymptotic states $`\pm \xi _0`$.
To finish the proof of theorem 0.2, it remains to prove the following proposition.
###### Proposition 5.12.
Every $`\alpha `$-stable, holomorphic $`S\mathrm{}_2`$-bundle $``$ over $`T\times ^1`$ restricting to $`L_{\xi _0}L_{\xi _0}`$ on $`T_{\mathrm{}}`$ can be obtained as the holomorphic extension of an instanton on $`T\times `$ with asymptotic states $`\pm \xi _0`$, and whose monodromy around the torus at infinity has eigenvalues $`\mathrm{exp}(\pm 2\pi i\alpha )`$.
###### Proof.
We will give two different ideas to prove the proposition, but we will not give the proofs, because they follow essentially well known arguments.
The first idea is direct construction: construct a Hermitian-Einstein metric on $`|_{T\times }`$ (so that the Chern connection is anti-self-dual); for this, one has first to build a metric $`h_0`$ on $``$ which gives asymptotically at infinity an instanton: this is possible because $`\alpha `$ and the behavior of $``$ near infinity (see (43)) give all the parameters at infinity of the instanton; then one wants to deform $`h_0`$ to a solution $`h`$ of the Hermitian-Einstein equation, mutually bounded with $`h_0`$; Simpson’s method cannot be used, because $`T\times `$ has infinite volume, but one can apply the method in , using precise analysis at infinity, which will be explained in the next section for the study of the moduli space.
The second idea, giving a different proof, consists in using the Nahm transform of instantons. Recall that our instantons are in correspondence with Higgs bundles with singularities on the dual torus $`\widehat{T}`$, with a harmonic metric. Actually, the correspondence has a purely holomorphic interpretation, and this is an occurrence of the so-called Fourier-Mukai transform. Stability is ususally preserved by such a correspondence, so that an $`\alpha `$-stable bundle on $`T\times ^1`$ would transform into a stable parabolic Higgs bundle on $`\widehat{T}`$; then one can apply Simpson’s theorem to construct a harmonic metric, whose inverse Nahm transform provides an instanton with quadratic curvature decay, and by theorem 0.1 this instanton has exactly the desired behavior at infinity. ∎
### 6 Moduli spaces
We now proceed to the differential geometric construction of the moduli space. The $`L^2`$ metric will then provide a hyperkähler structure on it.
We will restrict to the semisimple case; this choice simplifies the construction, because theorem 0.1 says that it is enough to look at functional spaces with weights which are powers of $`r`$; the analysis in the nilpotent case is possible, as in , but requires functional spaces with logarithmic weights.
Recall the model connection on the bundle $`E`$, trivialized near infinity:
$`A_0=d`$ $`+`$ $`i\left(\begin{array}{cc}\lambda _1& 0\\ 0& \lambda _1\end{array}\right)dx+i\left(\begin{array}{cc}\lambda _2& 0\\ 0& \lambda _2\end{array}\right)dy`$
$`+`$ $`i\left(\begin{array}{cc}\mu _1\mathrm{cos}\theta \mu _2\mathrm{sin}\theta & 0\\ 0& \mu _1\mathrm{cos}\theta +\mu _2\mathrm{sin}\theta \end{array}\right){\displaystyle \frac{dx}{r}}`$
$`+`$ $`i\left(\begin{array}{cc}\mu _1\mathrm{sin}\theta +\mu _2\mathrm{cos}\theta & 0\\ 0& \mu _1\mathrm{sin}\theta \mu _2\mathrm{cos}\theta \end{array}\right){\displaystyle \frac{dy}{r}}+`$
$`+`$ $`i\left(\begin{array}{cc}\alpha & 0\\ 0& \alpha \end{array}\right)d\theta .`$
Note that in order to get $`L^2`$ deformations, we cannot move the parameters $`\lambda `$, $`\mu `$ and $`\alpha `$; in view of theorem 0.1, it is natural to consider connections $`A_0+a`$, such that
$$|a|=O(r^{(1+\delta )}),|_{A_0}a|=O(r^{(2+\delta )});$$
actually, this $`C^1`$ space is not good for analysis, and we have the choice to substitute either a Hölder space $`C^{1,\eta }`$ or a Sobolev space $`L^{1,p}`$; we make the last choice, for $`p`$ big enough, and this leads to the technical definitions
$`\mathrm{\Omega }_\delta ^1`$ $`=`$ $`\{a\mathrm{\Omega }^1(𝔰𝔲(E)),aL_{12/p+\delta }^p,_{A_0}aL_{22/p+\delta }^p\}`$
$`𝒜`$ $`=`$ $`A_0+\mathrm{\Omega }_\delta ^1`$
$`𝒢`$ $`=`$ $`\{gSU(E),_{A_0}gg^1\mathrm{\Omega }_\delta ^1\}`$
$``$ $`=`$ $`\{F\mathrm{\Omega }_+^2(𝔰𝔲(E)),FL_{22/p+\delta }^p\}.`$
The Lie algebra of $`𝒢`$ is
$$T_1𝒢=\{u𝔰𝔲(E),_{A_0}u\mathrm{\Omega }_\delta ^1\}.$$
Note that for $`a\mathrm{\Omega }_\delta ^1`$, lemma 2.1 implies that actually $`a_{}L_{22/p+\delta }^p`$, so that this Sobolev space is the same as the one considered is part I. Also, the Sobolev embedding (24) implies $`\mathrm{\Omega }_\delta ^1C_\delta ^0`$, and an important property is that the embedding $`\mathrm{\Omega }_\delta ^1C_\delta ^{}^0`$ is compact if $`\delta ^{}<\delta `$; gauge transformations $`g𝒢`$ can be continuously extended over $`T_{\mathrm{}}`$, so that
$$g|_T_{\mathrm{}}=\left(\begin{array}{cc}u& 0\\ 0& u^1\end{array}\right),$$
where $`uS^1`$ is fixed. Also, $`𝒢`$ acts smoothly on $`𝒜`$ and the curvature is a smooth map from $`𝒜`$ to $``$.
Remark that there is no reducible connection in $`𝒜`$, since a reduction would decompose the bundle $`E`$ as $`LL^1`$, with $`L`$ topologically trivial on the torus at infinity; but then we would get $`c_2(E)=0`$.
Now we need the following proposition; the proof is given at the end of the section.
###### Proposition 6.1.
For $`k>0`$ and $`A𝒜`$, we have:
1. the laplacian $`\mathrm{\Delta }_A:T_1𝒢L_{22/p+\delta }^p`$ is an isomorphism; therefore there is a slice at $`A`$ to the action of $`𝒢`$ on $`𝒜`$, given by $`\{A+a,d_A^{}a=0\}`$;
2. if $`A`$ is an instanton, then the map $`d_A^+d_A^{}:\mathrm{\Omega }_\delta ^1L_{22/p+\delta }^p`$ is Fredholm surjective; the kernel coincide with the $`L^2`$-kernel.
Note that in the first statement of the proposition, it was crucial to allow gauge transformations to take non trivial values on $`T_{\mathrm{}}`$, otherwise one cannot obtain the slice $`d_A^{}a=0`$.
Define the moduli space $``$ as the space of instantons $`A𝒜`$ modulo the gauge group $`𝒢`$. As is well-known, $`F_A^+`$ is a hyperkähler moment map for the action of $`𝒢`$ on $`𝒜`$ with respect to the three complex structures on $`T\times `$:
$`I_1(z_1,z_2,w_1,w_2)`$ $`=`$ $`(z_2,z_1,w_2,w_1)`$
$`I_2(z_1,z_2,w_1,w_2)`$ $`=`$ $`(w_1,w_2,z_1,z_2)`$ (59)
$`I_3(z_1,z_2,w_1,w_2)`$ $`=`$ $`(w_2,w_1,z_2,z_1)`$
where $`z=z_1+iz_2`$ and $`w=w_1+iw_2`$. With the help of the previous proposition, standard theory now gives us:
###### Proposition 6.2.
The moduli space $``$ is a smooth hyperkähler manifold; the tangent space at $`[A]`$ is isomorphic to the $`L^2`$-kernel of $`d_A^+d_A^{}`$ acting on $`\mathrm{\Omega }^1(𝔰𝔲(E))`$. It has dimension $`8k4`$.
###### Proof of proposition 6.1.
First, we have to understand the behavior of the laplacian $`\mathrm{\Delta }_A`$ acting on sections of $`End(E)`$. We want to prove that it is Fredholm. This property is not changed by a perturbation in $`\mathrm{\Omega }_\delta ^1`$ (this adds to $`\mathrm{\Delta }_A`$ a compact operator), and we can therefore restrict to the case when $`A=A_0`$ on $`rR`$. On this domain $`rR`$, the laplacian preserves the decomposition $`u_\mathrm{\Gamma }u_{}`$.
The case of $`u_{}`$ is easier: since we have seen that $`_{A_0}^2u_{L_{22/p+\delta }^p}`$ controls $`u_{L_{22/p+\delta }^p}`$, it follows that $`A_0\mathrm{\Gamma }`$, which is $`O(r^1)`$, is small if $`R`$ is big enough; therefore lemma 2.3 proves that $`_A`$ is an isomorphism on $`rR`$ for the Neumann boundary condition (the same is true for Dirichlet boundary condition).
The case of $`u_\mathrm{\Gamma }`$ is more complicated, but can be reduced to standard theory: recall that $`u_\mathrm{\Gamma }`$ is torus invariant, so that the operator now reduces to an operator on $`^2`$; the action of $`\mathrm{\Delta }_\mathrm{\Gamma }`$ on off-diagonal coefficients (which exist only when $`\mathrm{\Gamma }`$ is trivial) is by
$$\frac{1}{r^2}\left((r_r)^2+(_\theta ^2\pm 2i\alpha )^2+|\mu |^2\right),$$
and the action on diagonal coefficients is the standard laplacian on $`^2`$ (that we obtain by making $`\alpha =\mu =0`$ in the previous formula); now $`r^2\mathrm{\Delta }_A`$ becomes the translation invariant laplacian
$$_t^2(_\theta ^2\pm 2i\alpha )^2+|\mu |^2$$
on the conformal cylinder $`_+\times S^1`$, so that standard theory now applies: such operator (say, with Dirichlet boundary condition on $`r=R`$) is Fredholm for all weights, except a discrete set of critical weights $`\delta `$ (they are characterized by the existence at infinity of solutions of type $`\mathrm{exp}(\delta t)t^k`$); moreover, as the operator is self-adjoint, its index is 0 at the weight 0 if it is noncritical, or $`1`$ for small positive weights if 0 is critical; in our situation, $`uT_1𝒢`$ corresponds to the decay $`uL_{\delta 2/p}^p`$, and this becomes exactly the weight $`\delta `$ on the cylinder; there are two cases: if $`\alpha `$ or $`\mu `$ is non zero (off-diagonal coefficients), then the weight 0 is not critical, and the operator remains Fredholm for nearby $`\delta `$, with index 0: actually is is an isomorphism, because it easy to verify that is has no kernel; if $`\alpha `$ and $`\mu `$ are zero, then the laplacian has index $`1`$ for small weights $`\delta >0`$, so that it becomes an isomorphism if we add the possibility to consider solutions $`u`$ of $`\mathrm{\Delta }u=v`$ with $`u`$ having some nonzero limit at infinity (and this is exactly our definition of $`𝒢`$). All these results can also be checked by direct calculation, after decomposing $`u`$ into Fourier series along each circle.
Finally, we deduce from these considerations that the laplacian $`\mathrm{\Delta }_{A_0}`$ is an isomorphism $`T_1𝒢L_{22/p+\delta }^p`$ for the Dirichlet boundary conditions on $`rR`$, and gluing this isomorphism with a parametrix on the compact part, it follows that $`\mathrm{\Delta }_A`$ is Fredholm on $`T\times ^2`$.
In order to calculate the index, if $`\mathrm{\Gamma }`$ is nontrivial, we have seen that the index is not changed if we modify $`A`$ so that $`A=\mathrm{\Gamma }`$ near infinity; the index of a self-adjoint operator on a compact manifold is zero; by an excision principle, this has the consequence that the index comes only from the contribution at infinity; therefore, it is equal to the index of the operator $`\mathrm{\Delta }_\mathrm{\Gamma }`$ acting on the trivial bundle $`𝔰𝔲(^2)`$; now this operator is completely explicit: on the $`u_{}`$ component, it is an isomorphism, and on the $`u_\mathrm{\Gamma }`$ component (that is, diagonal, torus invariant, matrices), it is simply the standard laplacian in $`^2`$, and its index between the spaces that we have defined is again 0, with 1-dimensional kernel and cokernel equal to constant diagonal matrices.
If $`\mathrm{\Gamma }`$ is trivial, we cannot reduce to the operator of flat space, but we can reduce to $`\mathrm{\Delta }_{A_0^{}}`$, with $`A_0^{}`$ the diagonal connection
$$A_0^{}=\chi (r)A_0+(1\chi (r))d,$$
(60)
where $`\chi (r)`$ is a cutoff function which equals 1 for $`r>R`$ and 0 for $`r<R1`$; then, as above, it is not difficult to prove that $`\mathrm{\Delta }_{A_0}^{}`$ is an isomorphism on non-diagonal components (and the operator on the diagonal components is the same as above).
Finally, the operator $`\mathrm{\Delta }_A`$ has no kernel in $`T_1𝒢`$, since an element in the kernel would decompose $`A`$, which is impossible. This finishes the proof of the first part of the proposition.
If $`A𝒜`$ is an instanton, observe that the operator $`d_A^+d_A^{}`$ acting on self-dual 2-forms, by the Weitzenböck formula, equals the laplacian $`_A^{}_A`$; this means that the above results remain true for $`d_A^+d_A^{}`$, and we deduce that the operator
$$d_A^+d_A^{}:\mathrm{\Omega }_\delta ^1L_{22/p+\delta }^p$$
is surjective; its kernel equals the kernel of the laplacian $`2(d_A^+)^{}d_A^++d_Ad_A^{}`$; again one can prove (in particular using lemma 2.10) that this operator is Fredholm (for the weight $`\delta `$); remark that the $`L^2`$ condition corresponds to a critical weight (on diagonal components, where the operator is asymptotically the standard laplacian of $`^2`$), when $`\mathrm{\Omega }_\delta ^1`$ corresponds to a slightly greater weight; nevertheless, it remains true that the $`L^2`$-kernel equals the kernel for slightly greater weights (the possible new solutions in the kernel at the critical weight are never $`L^2`$). ∎
###### Proof of proposition 6.2.
It remains only to calculate the dimension, which, by proposition 6.1, is the index of the operator $`d_A^+d_A^{}`$. If the limit flat connection $`\mathrm{\Gamma }`$ is non trivial, this is simple to calculate by comparison to the same operator for $`\mathrm{\Gamma }`$: actually, by the excision principle,
$$ind(d_A^+d_A^{})=ind(d_\mathrm{\Gamma }^+d_\mathrm{\Gamma }^{})+8k;$$
now for the flat connection $`\mathrm{\Gamma }`$, the operator $`d_\mathrm{\Gamma }^+d_\mathrm{\Gamma }^{}`$ has no kernel (by the Weitzenböck formula), but its cokernel equals the cokernel of the operator $`d_\mathrm{\Gamma }^{}d_\mathrm{\Gamma }+d_\mathrm{\Gamma }^+d_\mathrm{\Gamma }^{}`$ acting on $`\mathrm{\Omega }^0(𝔰𝔲(E))\mathrm{\Omega }_+^2(𝔰𝔲(E))=^4𝔰𝔲(E)`$; we have seen above that the cokernel of this operator on $`𝔰𝔲(E)`$ is the $`L^2`$-orthogonal of constant, diagonal matrices. This proves the formula for the index.
If $`\mathrm{\Gamma }`$ is trivial, the same result holds, but one must compare with the operator $`d_{A_0^{}}^+d_{A_0^{}}^{}`$ defined in (60). ∎
#### Fibration structure.
It was shown in that the moduli space of rank two holomorphic vector bundles over $`T\times ^1`$ with trivial determinant and instanton number $`k`$ contains an open set $`_k^{}`$ (corresponding to the so-called *regular bundles*) which has the structure of a fibration:
$$𝕋\mathrm{}_k^{}\mathrm{\Sigma }_k$$
The fibres are complex tori of complex dimension $`2k1`$, and the base can be interpreted as the set of rational maps $`^1^1`$ of degree $`k`$, so that $`\mathrm{dim}\mathrm{\Sigma }_k=2k+1`$.
Fixing the splitting of $``$ at $`T_{\mathrm{}}`$, i.e. fixing the asymptotic state of the corresponding instanton connection $`A`$, amounts to fixing the value of these rational maps at $`\mathrm{}^1`$. Moreover, as we will see in the next section, fixing the residue of $`A`$ amounts to fixing the first derivative at $`\mathrm{}^1`$.
Therefore, according to theorem 0.2, we conclude that $`_{(k,\pm \xi _0,\mu )}`$, the moduli space of $`SU_2`$ doubly-periodic instantons with fixed instanton number $`k`$, asymptotic states $`\pm \xi _0`$ and residue $`\mu `$ with the complex structure induced from the complex structure $`I_1`$ on $`T\times ^2`$, is a fibration over $`\mathrm{\Sigma }_{(k,\pm \xi _0,\mu )}`$, the space of rational maps $`f:^1^1`$ with fixed $`f(w=\mathrm{})`$ and $`f^{}(w=\mathrm{})`$, with fibres given complex tori of dimension $`2k1`$.
Moreover, it is possible to show that the such fibres are lagrangian with respect to complex symplectic structure on $`_{(k,\pm \xi _0,\mu )}`$ induced from the complex symplectic structure $`\omega _{I_2}+I_1\omega _{I_3}`$ on $`T\times `$ (see for the proof of a similar result for elliptic K3 and abelian surfaces).
#### An example: $`𝐤=\mathrm{𝟏}`$.
We shall now give an explicit model for the moduli space of doubly-periodic instantons with $`k=1`$; clearly, we also assume that $`\xi _0\xi _0`$ and $`\mu 0`$.
Our approach is based on the observations made above, that is, we shall study the set of rational maps $`f:^1^1`$ of degree 1; in a neighbourhhod of $`\mathrm{}^1`$, such maps can be written as follows:
$$f(w)=\frac{w+b}{cw+d},\mathrm{where}w=0\mathrm{corresponds}\mathrm{to}\mathrm{}^1.$$
As we discussed above, we must still fix $`f(0)`$ and $`f^{}(0)`$. This means that $`b/d`$ and $`(dcb)/d^20`$ are fixed. Thus, $`\mathrm{\Sigma }_{(1,\pm \xi _0,\mu )}=`$, so that $`_{(k,\pm \xi _0,\mu )}`$ is an elliptic fibration over $``$.
Actually, one can say more: there is an action of $`T\times `$ on the moduli space (by translations), so the moduli space is exactly $`T\times `$, and the metric is flat.
## Part III Nahm transform
We now shift our attention to the Nahm transform of doubly-periodic instanton connections . Note that this transform was defined in only for instantons such that the restriction of the underlying holomorphic bundle to a generic torus is $`L_\xi L_\xi `$ (this is what we called the semisimple case). In this part, we shall restrict to this case.
Throughout this part, we assume familiarity with , but let us quickly recall how Nahm transform is defined. Given an instanton $`A`$ on a $`SU_2`$-bundle $`E`$ on $`T\times ^2`$, one may twist $`A`$ by a flat connection on $`T`$; these twists $`A_\xi `$ are parameterized by $`\xi \widehat{T}`$. Now there is a coupled Dirac operator
$$D_{A_\xi }:\mathrm{\Gamma }(S^+E)\mathrm{\Gamma }(S^{}E)$$
and one can show that the bundle of $`L^2`$-cokernels of $`D_{A_\xi }`$ is a rank $`k`$ vector bundle $`V`$ over $`\widehat{T}\{\pm \xi _0\}`$; there is a natural connection $`B`$ on $`V`$ obtained by projection, and one can define an endomorphism $`\mathrm{\Phi }`$ of $`V`$ by taking an element $`\beta \mathrm{ker}D_{A_\xi }^{}`$ to the projection of $`w\beta `$ on this kernel; the pair $`(B,\mathrm{\Phi })`$ satisfies Hitchin’s equations on $`\widehat{T}\{\pm \xi _0\}`$.
From the holomorphic point of view, the picture is very clear: the spinor bundle $`S`$ is identified $`\mathrm{\Lambda }^{0,}`$, so that the $`L^2`$-kernel of $`D_{A_\xi }^{}`$ is exactly the $`L^2`$-kernel of $`\overline{}_{A_\xi }\overline{}_{A_\xi }^{}`$ on $`\mathrm{\Omega }^{0,1}E`$. It can be proven that this $`L^2`$-kernel coincides with $`H^1(T\times ^1,L_\xi )`$, where $``$ is the holomorphic extension of $`A`$ on $`T\times ^1`$; this provides a holomorphic extension of $`V`$ on the whole $`\widehat{T}`$; this extension has degree $`2`$, as can be checked by Riemann-Roch theorem for families. Moreover, there is a natural interpretation for the Higgs field: one has the identification
$$V_\xi =H^1(T\times ^1,L_\xi )=\underset{w}{}H^0(T_w,EL_\xi )$$
(61)
where of course there is only a finite number of points $`w`$ (actually $`k`$, counted with multiplicity) such that $`H^0(T_w,EL_\xi )0`$. Now the Higgs field $`\mathrm{\Phi }`$ is multiplication by $`w`$ on $`H^0(T_w,EL_\xi )`$. From this description, one can see that the Higgs field has a simple pole at $`\pm \xi _0`$ with semisimple residue, and the residue has only one nonzero eigenvalue if $`\xi _0\xi _0`$, two otherwise.
We first study how the new asymptotic parameters of doubly-periodic instantons introduced in Part II behave under Nahm Transform. This will prepare the way for the proof of theorem 0.4, our last result.
### 7 Asymptotic parameters
Following the general philosophy that the Nahm Transform is a sort of nonlinear Fourier Transform, it is reasonable to expect the asymptotic behavior of the instanton to be translated into further singularity data for the Higgs field.
Recall that $`|_T_{\mathrm{}}=L_{\xi _0}L_{\xi _0}`$. From (61) we deduce a holomorphic splitting of $`V`$ on a small neighborhood of $`\pm \xi _0`$:
$$V_\xi =B_\xi R_\xi $$
(62)
where $`B_\xi `$ corresponds to the points in $``$ that remain bounded as $`\xi \xi _0`$ and $`R_\xi `$ corresponds to the points that go off to infinity. Clearly, $`B_\xi `$ approaches the kernel of $`\mathrm{Res}_{\pm \xi _0}\mathrm{\Phi }`$ as $`\xi \xi _0`$, while $`R_\xi `$ approaches the eigenspace of the nontrivial eigenvalues of the residue.
The behaviour of the Higgs bundle with harmonic metric near the singularities $`\pm \xi _0`$ is completely determined by the following theorem.
###### Theorem 7.1.
Let $`A`$ be a doubly-periodic instanton with limiting holonomy $`\alpha `$ and residue $`\mu `$; let $`(B,\mathrm{\Phi })`$ be its Nahm transformed Higgs pair. The unique nonzero eigenvalue of $`\mathrm{Res}_{\pm \xi _0}\mathrm{\Phi }`$ is given by $`\pm \mu `$. In the decomposition (62), the harmonic metric on $`V`$ remains bounded on $`B`$, but behaves like $`|\xi \pm \xi _0|^{1\pm \alpha }`$ on $`R`$.
###### Remark 7.2.
The sum of the degree of $`V`$, that is $`2`$, and of the weights $`1\pm \alpha `$, equals 0, as must be for a solution of Hitchin’s equations. The monodromy of the connection $`B`$ near the punctures is semisimple, with only one nontrivial eigenvalue $`\mathrm{exp}(2\pi i\alpha )`$ on $`R`$ (or two if $`\xi _0=\xi _0`$).
We first prove the statement concerning the residues. The argument to establish the statement concerning the limiting holonomy is much more technical, and will involve a series of lemmas.
#### Residues.
Let $`\rho =r^1`$ and let $`w^{}=w^1=\rho e^{i\theta }`$ be a coordinate near $`\mathrm{}^1`$. Clearly, the holomorphic structure on the restriction $`E|_{T_w^{}}`$ is given by the $`(0,1)`$-part of the $`A|_{T_w^{}}`$. Rewriting equation (43) in terms of $`w^{}`$, we obtain:
$$\overline{}_A|_{T_w^{}}=\overline{}+\left(\begin{array}{cc}\lambda & 0\\ 0& \lambda \end{array}\right)d\overline{z}+\left(\begin{array}{cc}\mu & 0\\ 0& \mu \end{array}\right)w^{}d\overline{z}+O(\rho ^2)$$
so that:
$$\frac{d}{dw^{}}\left(\overline{}_A|_{T_w^{}}\right)|_{w^{}=0}=\left(\begin{array}{cc}\mu & 0\\ 0& \mu \end{array}\right)$$
In other words, the residue $`\mu `$ can be regarded as the infinitesimal variation of the holomorphic bundle $`|_{T_w}`$ at $`w=\mathrm{}`$.
Since for every $`w^{}`$ sufficiently close to $`\mathrm{}^1`$ we can assume that $`E|_{T_w^{}}=L_{\xi (w^{})}L_{\xi (w^{})}`$, the above expression implies that:
$$\frac{d}{dw^{}}\xi (w^{})|_{w^{}=0}=\mu .$$
The eigenvalue of $`\mathrm{\Phi }`$ going to infinity is $`w(\xi )=1/w^{}(\xi )`$ by (61); the statement follows.∎
#### Limiting holonomy.
Let us now look at the coupled Dirac laplacian $`\mathrm{\Delta }_{A_\xi }`$ acting on sections of $`S^+E`$; since $`A`$ is an instanton, we have that $`D_{A_\xi }^{}D_{A_\xi }=_{A_\xi }^{}_{A_\xi }`$, i. e. the Dirac laplacian coincides with the trace laplacian. This laplacian is inversible in $`L^2`$ for $`\xi \pm \xi _0`$ (see ; this is also a consequence of the lemmas below), and we note its inverse by $`G_{A_\xi }`$. Such inverse is useful to produce harmonic representative of elements of $`H^1(T\times ^1,L_\xi )`$. Indeed, if we have a compactly supported (0,1)-form $`\beta `$ with values in $`E`$ such that $`\overline{}_{A_\xi }\beta =0`$, then the $`L^2`$-harmonic representative of the class $`[\beta ]`$ is given by
$$\beta \overline{}_{A_\xi }G_{A\mathrm{\_}xi}\overline{}_{A_\xi }^{}\beta .$$
We now want to understand the inverse $`G_{A_\xi }`$ when $`\xi `$ approaches the asymptotic states $`\pm \xi _0`$. For simplicity, assume that $`\xi _0=0`$ in the next three lemmas; the general case can be obtained by substituting $`\xi \xi _0`$ for $`\xi `$ in the expressions below.
We know from theorem 0.1, where $`\lambda =\xi `$:
$$A_\xi =(A_0)_\xi +a\mathrm{with}|a|=O(r^{1ϵ})$$
and
$`(A_0)_\xi `$ $`=`$ $`d+i\left(\begin{array}{cc}\alpha & 0\\ 0& \alpha \end{array}\right)d\theta +i\left(\begin{array}{cc}\lambda _1dx+\lambda _2dy& 0\\ 0& \lambda _1dx\lambda _2dy\end{array}\right)+`$
$`{\displaystyle \frac{i}{r}}\left(\begin{array}{cc}\mu _1dx+\mu _2dy& 0\\ 0& \mu _1dx\mu _2dy\end{array}\right).`$
We assume also that at either $`\mu _1`$ or $`\mu _2`$ is nonzero; however, the proofs below will also work if $`\mu _1=\mu _2=0`$, but $`\alpha 0`$.
###### Lemma 7.3.
Let $`\sigma `$ is a section of $`ET\times `$; if $`\lambda `$ is sufficiently small and $`|w|`$ is large enough, then:
$$_{T_w}|_{(A_0)_\xi }\sigma |^2\left|\lambda +\frac{\mu }{w}\right|^2_{T_w}|\sigma |^2.$$
###### Proof.
Consider the Fourier expansion $`\sigma =\mathrm{\Sigma }\sigma _{nm}e^{i(nx+my)}`$. Then on the torus $`T_w`$, we have:
$`{\displaystyle _{T_w}}|\sigma |^2`$ $`=`$ $`{\displaystyle _{T_w}}\left|\left(_x+i\lambda _1+i{\displaystyle \frac{\mu _1}{|w|}}\right)\sigma \right|^2+\left|\left(_y+i\lambda _2+i{\displaystyle \frac{\mu _2}{|w|}}\right)\sigma \right|^2`$
$`=`$ $`{\displaystyle }|(n+im+\lambda +{\displaystyle \frac{\mu }{|w|}}|^2|\sigma _{nm}|^2.`$
However, under the hypothesis above,
$$\left|n+im+\lambda +\frac{\mu }{|w|}\right|\left|\lambda +\frac{\mu }{w}\right|$$
for all $`n,m`$, which proves the lemma. ∎
###### Lemma 7.4.
Under the hypothesis of lemma 7.3, we have:
$$_{rR}|_{(A_0)_\xi }\sigma |^2c|\mu |^2_{rR}\frac{|\sigma |^2}{r^2}$$
(65)
###### Proof.
By the previous lemma, the estimate holds away from the region where $`|\lambda +\mu /w|`$ is small, that is:
$$\frac{1}{2}\frac{|\lambda |}{|\mu |}|w|2\frac{|\lambda |}{|\mu |}$$
Actually, we claim that if the estimate of the lemma is satisfied outside this region, then it must be satisfied everywhere. Indeed, one has the inequality for any function $`f:^2`$, and a constant $`c`$ independent of $`\rho `$,
$$_{\rho r2\rho }\frac{f^2}{r^2}c\left(_{2\rho r4\rho }\frac{f^2}{r^2}+_{\rho r4\rho }|_rf|^2\right)$$
(66)
and the lemma follows by applying (66) to $`f=|\sigma |`$ and $`\rho =|\lambda /\mu |`$. The proof of (66) is left to the reader. ∎
Note that an estimate similar to (65) remains valid if $`\mu =0`$, but $`\alpha 0`$. In fact, the proof is even simpler, since one has the estimate:
$$_{r=R}|_{(A_0)_\xi }\sigma |^2\frac{|\alpha |^2}{r^2}_{rR}|\sigma |^2$$
from which one immediately obtains:
$$_{rR}|_{(A_0)_\xi }\sigma |^2|\alpha |^2_{rR}\frac{|\sigma |^2}{r^2}.$$
(67)
###### Lemma 7.5.
The solution of the Poisson equation $`\mathrm{\Delta }_{A_\xi }u=v`$ satisfies:
$$r^1u_{L^2}+_{A_\xi }u_{L^2}crv_{L^2}$$
$$\mathrm{and}|\xi |^2u_{L^2}+|\xi |_{A_\xi }u_{L^2}crv_{L^2}$$
###### Proof.
First, note that:
$$|_{A_\xi }\sigma |^2c\left(|\xi |^2|\sigma |^2+\frac{|\sigma |^2}{r^2}\right).$$
(68)
Near infinity, this a consequence of lemma 7.4 and of the fact that $`A=A_0+O(r^{1ϵ})`$. Globally, the estimate follows from the Poincaré-type inequality:
$$_{rR}|\sigma |^2c\left(_{rR}|\sigma |^2+_{R/2rR}|\sigma |^2\right)$$
To prove the lemma itself, we have that:
$`_{A_\xi }u_{L^2}^2`$ $`=`$ $`{\displaystyle \mathrm{\Delta }_{A_\xi }u,u}={\displaystyle v,u}`$
$``$ $`rv_{L^2}r^1u_{L^2}crv_{L^2}_{A_\xi }u_{L^2}`$
by (68). Thus, we conclude that $`_{A_\xi }u_{L^2}crv_{L^2}`$, and again by (68) we have $`r^1u_{L^2}crv_{L^2}`$. The second estimate is obtained in a similar way. ∎
We are now finally ready to complete the proof of theorem 7.1. Let us first analyze the behavior of the harmonic metric on the local sub-bundle $`BV`$ with fibers given by $`B_\xi `$. Let $`\beta `$ be a section of $`B`$. Then, for each $`\xi \xi _0`$, we know from (61) that $`\beta (\xi )`$ can be represented as a section of $`\mathrm{\Lambda }^{0,1}EL_\xi `$ supported on $`rR`$ for some $`R`$ sufficiently large. Furthermore, its harmonic representative in $`H^1(T\times ,EL_\xi )`$ is given by $`\beta (\xi )\overline{}_{A_\xi }G_{A_\xi }\overline{}_{A_\xi }^{}\beta (\xi )`$. By lemma 7.5, we have:
$$\overline{}_{A_\xi }G_{A_\xi }\overline{}_{A_\xi }^{}\beta (\xi )_{L^2}cr\overline{}_A^{}\beta (\xi )_{L^2}cR\overline{}_A^{}\beta (\xi )_{L^2}$$
which remains bounded even as $`\xi \xi _0`$. This means that the limit
$$\beta (\xi _0)=\underset{\xi \xi _0}{lim}\beta (\xi )$$
has a square-integrable harmonic representative, so that the harmonic metric restricted to the sub-bundle $`B`$ extends across $`\pm \xi _0`$.
Now let $`RV`$ be a local sub-bundle with fibers given by $`R_\xi `$; remind that near infinity, we have $`|_{T_w}=L_{\xi (w)}L_{\xi (w)}`$; take a section $`\beta (\xi )`$ of $`R_\xi `$ coming by (61) from sections of $`|_{T_{w(\xi )}}L_\xi `$ converging to a section of $`|_T_{\mathrm{}}L_{\xi _0}=L_{2\xi _0}`$. Here we have to be more specific: say that a section $`\sigma H^0(T_{w(\xi )},L_\xi )`$ corresponds to the class in $`H^1(T\times ^1,L_\xi )`$ represented by the (0,1)-current
$$\sigma (z)\delta _{w(\xi )}(w)d\overline{w},$$
(69)
where $`\delta _{w(\xi )}`$ is the Dirac function at the point $`w(\xi )`$. From this description, we see that, for each $`\xi \xi _0`$, the representative $`\beta (\xi )`$ can be chosen with compact support near $`r=|w(\xi )|`$, and bounded in $`L^{1,2}`$. Now lemma 7.5 gives, as above,
$$\overline{}_{A_\xi }G_{A_\xi }\overline{}_{A_\xi }^{}\beta (\xi )_{L^2}cr\overline{}_A^{}\beta (\xi )_{L^2}\frac{c}{|\xi \xi _0|}\overline{}_A^{}\beta (\xi )_{L^2}.$$
This means that the norm of the harmonic representative of $`\beta (\xi )`$ is bounded by $`|\xi \xi _0|^1`$.
This result must be interpreted, since (69) actually does not extend to $`w=\mathrm{}`$, so that our $`[\beta (\xi )]`$ is not a section of $`R`$ which extends over the puncture $`\xi _0`$. There are two changes to make; first, note that a (0,1)-form smooth on $`^1`$ near infinity is $`d\overline{w}/\overline{w}^2`$, so we see that we must consider $`\beta (\xi )/\overline{w}(\xi )^2`$ instead of $`\beta (\xi )`$. The second change to be made is that we want $`\beta (\xi )`$ holomorphic in $`\xi `$. This involves a constraint on the choice of $`\sigma `$: from the growth of the holomorphic sections of $``$ at infinity studied in section 5, it follows that $`|\sigma ||w(\xi )|^\alpha `$, and we can finally conclude that the norm of a holomorphic section of $`R`$ is bounded by $`|\xi \xi _0|^{1\alpha }`$.
From these results, it follows that the harmonic metric of the Higgs bundle $`V`$ extends on $`B`$, and is bounded by $`|\xi \pm \xi _0|^{1\pm \alpha }`$ on $`R`$. This gives a bound $`1\pm \alpha `$ for the weights of the parabolic structure of $`V`$. However, the “parabolic degree” of the bundle must be zero, and $`V`$ has degree $`2`$, so that the weights must be exactly equal to $`1\pm \alpha `$. ∎
#### Reformulating the Nahm transform theorem
Together with , theorem 7.1 allows us to state a complete version of the Nahm transform theorem, including the new asymptotic parameters defined in Part II:
###### Theorem 7.6.
The Nahm transform is a correspondence between the following objects:
* $`SU(2)`$ doubly-periodic instantons with instanton number $`k>0`$ and asymptotic parameters $`(\pm \xi _0,\alpha ,\mu )`$;
* rank $`k`$ logarithmic Higgs bundles with harmonic metric over $`\widehat{T}`$ with singularity behavior as described in theorem 7.1.
### 8 The hyperkähler property
Our final task is to prove that the Nahm transform of doubly-periodic instantons define a hyperkähler isometry between $``$, the moduli space of doubly-periodic instanton constructed in section 6, and $`\widehat{}`$, the moduli space of meromorphic Higgs pairs satisfying the conditions of theorem 7.6. To do that, we shall follow the following strategy. First, we compute the derivative of the map:
$`N:`$ $``$ $`\widehat{}`$
$`A`$ $``$ $`(B,\mathrm{\Phi })`$
defined by the Nahm transform, verifying that it is indeed well-defined. We then show that $`D_{[A]}N`$ preserves the three complex structures in each space. The last step is to show that $`D_{[A]}N`$ preserve the metrics in each space.
#### Computing the derivative.
Recall the definition of the tangent space $`T_{[A]}`$ at the gauge equivalence class of an instanton $`A`$ can be characterized as follows:
$$T_{[A]}=\{aL^2(\mathrm{\Omega }^1𝔰𝔲(E))\mathrm{s}.\mathrm{t}.\begin{array}{cc}\hfill (i)& d_A^{}a=0\hfill \\ \hfill (ii)& d_A^+a=0\hfill \end{array}\}$$
(70)
The 1-form $`a`$ is regarded as a infinitesimal variation of the instanton connection $`A`$, inducing a 1-parameter family of connections $`A_t=A+ta`$, which are anti-self-dual up to first order.
Now let $`\{\mathrm{\Psi }(\xi )^j\}_{j=1}^k`$ be an orthonormal base for coupled adjoint Dirac operator $`\mathrm{ker}D_{A_\xi }^{}`$. In order to compute the derivative $`D_{[A]}N`$, we must understand the infinitesimal change on harmonic spinors induced by the infinitesimal change on the instanton. We are looking for negative spinors $`\phi (\xi )^j`$ such that the 1-parameter family $`\mathrm{\Psi }_t(\xi )^j=\mathrm{\Psi }(\xi )^j+t\phi (\xi )^j`$ satisfies $`D_{(A_\xi )_t}^{}\mathrm{\Psi }_t(\xi )^j=0`$ up to first order. In other words,
$$\frac{d}{dt}D_{(A_\xi )_t}^{}\mathrm{\Psi }_t(\xi )^j|_{t=0}=D_{A_\xi }^{}\phi (\xi )^j+a\mathrm{\Psi }(\xi )^j=0$$
where $``$ means Clifford multiplication. Therefore, the infinitesimal variations on harmonic spinors are given by:
$$\phi (\xi )^j=D_{A_\xi }G_{A_\xi }(a\mathrm{\Psi }(\xi )^j)$$
(71)
Recall from that the Nahm transformed Higgs pair is defined as follows:
$$\begin{array}{ccc}B(\xi )^{ij}=\mathrm{\Psi }(\xi )^i,\widehat{d}\mathrm{\Psi }(\xi )^j& \mathrm{and}& \mathrm{\Phi }(\xi )^{ij}=\mathrm{\Psi }(\xi )^i,w\mathrm{\Psi }(\xi )^jd\xi \end{array}$$
(72)
where $`\widehat{d}`$ means differentiation with respect to $`\xi `$, the coordinate on the dual torus $`\widehat{T}`$, and the inner products are taken in $`L^2(ES^{})`$. Thus, the infinitesimal change in the Nahm transformed Higgs pair $`(B,\mathrm{\Phi })`$ is given by:
$`b(\xi )^{ij}`$ $`=`$ $`{\displaystyle \frac{d}{dt}}\mathrm{\Psi }_t(\xi )^j,\widehat{d}\mathrm{\Psi }_t(\xi )^j|_{t=0}=`$ (73)
$`=`$ $`G_{A_\xi }\mathrm{\Psi }(\xi )^i,\mathrm{\Omega }a\mathrm{\Psi }(\xi )^j\mathrm{\Omega }a\mathrm{\Psi }(\xi )^i,G_{A_\xi }\mathrm{\Psi }(\xi )^j`$
and
$`\varphi (\xi )^{ij}`$ $`=`$ $`{\displaystyle \frac{d}{dt}}\mathrm{\Psi }_t(\xi )^j,w\mathrm{\Psi }_t(\xi )^j|_{t=0}=`$ (74)
$`=`$ $`G_{A_\xi }\mathrm{\Psi }(\xi )^i,dwa\mathrm{\Psi }(\xi )^jd\xi `$
where $`\mathrm{\Omega }=i\left(d\xi _1dz_1+d\xi _2dz_2\right)`$ is the curvature of the Poincaré bundle over $`T\times \widehat{T}`$.
The tangent space $`T_{[(B,\mathrm{\Phi })]}\widehat{}`$ at the gauge equivalence class of a Higgs pair $`(B,\mathrm{\Phi })`$, can described as follows (see for instance ):
$$T_{[(B,\mathrm{\Phi })]}\widehat{}=\{\begin{array}{c}bL^2(\mathrm{\Omega }^1𝔲(V))\hfill \\ \varphi L^2(\mathrm{\Omega }^{1,0}𝔤𝔩(V))\hfill \end{array}\mathrm{s}.\mathrm{t}.\begin{array}{cc}\hfill (i)& d_Bb+[\mathrm{\Phi },\varphi ^{}]+[\varphi ,\mathrm{\Phi }^{}]=0\hfill \\ \hfill (ii)& \overline{}_B\varphi +[b^{0,1},\mathrm{\Phi }]=0\hfill \\ \hfill (iii)& d_B^{}b+\mathrm{Re}[\mathrm{\Phi }^{},\varphi ]=0\hfill \end{array}\}$$
(75)
Again, $`(b,\varphi )`$ define a 1-parameter family of pairs $`(B_t=B+tb,\mathrm{\Phi }_t=\mathrm{\Phi }+t\varphi )`$ which satisfy Hitchin’s equations up to first order.
Therefore, it is clear from (73) and (74) that the pair $`(b,\varphi )`$ satisfies the linearized Hitchin’s equations ($`(i)`$ and $`(ii)`$ in (75)).
We must only verify that $`(b,\varphi )`$ are transversal to infinitesimal changes in $`(B,\mathrm{\Phi })`$ arising from infinitesimal gauge transformations, i.e. must check equation $`(iii)`$ in (75). To do that, denote by $`\stackrel{~}{B}`$ and $`\stackrel{~}{b}`$ the $`(^2)^{}`$-invariant 1-forms on $`\widehat{T}\times (^2)^{}`$ obtained from $`(B,\mathrm{\Phi })`$ and $`(b,\varphi )`$, respectively. Clearly, $`\stackrel{~}{B}`$ is anti-self-dual and
$$d_B^{}b+\mathrm{Re}[\mathrm{\Phi }^{},\varphi ]=0d_{\stackrel{~}{B}}^{}\stackrel{~}{b}=0$$
The following result completes our first step towards the proof of theorem 0.4
###### Lemma 8.1.
If $`d_A^{}a=0`$, then $`d_{\stackrel{~}{B}}^{}\stackrel{~}{b}=0`$.
###### Proof.
See proposition 3.1 in . ∎
###### Remark 8.2.
Using the ideas above, one can easily compute the derivative of the inverse Nahm transform, thus showing that $`N:\widehat{}`$ is a diffeomorphism. Noting that, since$``$ is smooth, the diffeomorphism type of the moduli space of instantons does not depend on the choice of asymptotic parameters $`(\alpha ,\lambda ,\mu )`$, one concludes that the diffeomorphism type of the moduli of Higgs bundles is independent not only of the singularity data (residues and parabolic structure), as it was observed by Nakajima in , but also of the position of the singularities.
#### Commuting with the complex structures.
Consider coordinates $`(\xi _1,\xi _2,\omega _1,\omega _2)`$ on $`(^4)^{}`$, which are dual to $`(z_1,z_2,w_1,w_2)`$. Each of the complex structures (59) in $`^4`$ naturally induces a similar complex structures $`\widehat{I}_j`$ on $`(^4)^{}`$. Thus, we have maps:
$$\mathrm{\Lambda }^1^4\stackrel{I_j}{}\mathrm{\Lambda }^1^4\mathrm{and}\mathrm{\Lambda }^1(^4)^{}\stackrel{\widehat{I}_j}{}\mathrm{\Lambda }^1(^4)^{}$$
The complex structures on $`\widehat{}`$ can be then defined as follows. As above, let $`\stackrel{~}{b}\mathrm{\Lambda }^1(^4)^{}`$ be the $`(^2\times ^2)^{}`$-invariant 1-form obtained from $`(b,\varphi )`$. Then $`\widehat{I}_j(\stackrel{~}{b})`$ is also a $`(^2\times ^2)^{}`$-invariant 1-form on $`(^4)^{}`$, which can then be interpreted as an element of (75). It is easy to see that these coincide with the complex structures originally defined by Hitchin in . Therefore, we have to show that the following diagram:
(76)
commutes. The horizontal maps are defined as follows:
$$D_{[A]}N(a)=\stackrel{~}{b}=G_{A_\xi }\mathrm{\Psi }(\xi )^i,\stackrel{~}{\mathrm{\Omega }}a\mathrm{\Psi }(\xi )^j\stackrel{~}{\mathrm{\Omega }}a\mathrm{\Psi }(\xi )^i,G_{A_\xi }\mathrm{\Psi }(\xi )^j$$
(77)
with $`\stackrel{~}{\mathrm{\Omega }}=i\left(d\xi _1dz_1+d\xi _2dz_2+d\omega _1dw_1+d\omega _2dw_2\right)`$.
Each $`I_j`$ induces an isomorphism $`l_j:^4^2`$ satisfying the following commutative diagram:
(78)
where the map on the left hand side is multiplication by $`i=\sqrt{1}`$. Of course, a similar diagram holds for $`\widehat{l_j}:(^4)^{}(^2)^{}`$.
The key point is to note that each map:
$$D_{[A]}N_{}=\widehat{l_j}D_{[A]}Nl_j^1:\mathrm{\Lambda }^{(1,0)}^2\mathrm{\Lambda }^{(1,0)}(^2)^{}$$
$$D_{[A]}N_{}(\alpha )=G_{A_\xi }\mathrm{\Psi }(\xi )^i,\stackrel{~}{\mathrm{\Omega }}_{}\overline{\alpha }\mathrm{\Psi }(\xi )^j\stackrel{~}{\mathrm{\Omega }}_{}\alpha \mathrm{\Psi }(\xi )^i,G_{A_\xi }\mathrm{\Psi }(\xi )^j$$
is $``$-linear, where $`\stackrel{~}{\mathrm{\Omega }}_{}=\widehat{l_j}\times l_j(\stackrel{~}{\mathrm{\Omega }})`$. Therefore, we conclude:
$`\widehat{I}_j(D_{[A]}N(a))`$ $`=`$ $`\widehat{l_j}^1(i)\widehat{l_j}D_{[A]}N(a)=D_{[A]}Nl_j^1(i)l_j(a)=`$
$`=`$ $`D_{[A]}N(I_j(a))`$
as desired.
#### The Nahm transform is an isometry.
Again, the fact that the Nahm transform is an isometry is actually a property of the underlying four-dimensional transform. The calculations of Braam and van Baal are quite precise and also apply to the present situation.
Recall that the metric on the instanton moduli space is given by the $`L^2`$ norm of the tangent vectors, that is:
$$g(a_1,a_2)=_{T\times }\mathrm{Tr}(a_1a_2)$$
while the metric on the Higgs moduli space is given by
$$\widehat{g}((b_1,\varphi _1),(b_2,\varphi _2))=_{\widehat{T}}\mathrm{Tr}(b_1^{}b_2+\varphi _1\varphi _2^{})$$
or, equivalently, in terms of the 4-dimensional 1-forms $`\stackrel{~}{b_1}`$ and $`\stackrel{~}{b_2}`$:
$$\widehat{g}(\stackrel{~}{b_1},\stackrel{~}{b_2})=_^2^{}\mathrm{Tr}(\stackrel{~}{b_1}\stackrel{~}{b_2})$$
where integration is now done only with respect to the two coordinates on $`(^4)^{}`$ on which $`\stackrel{~}{b_1}`$ and $`\stackrel{~}{b_1}`$ depend.
Let $`(b,\varphi )=D_{[A]}N(a)`$; it is enough to show that:
$$\widehat{g}(D_{[A]}N(a),(b,\varphi ))=g(a,D_{[A]}N^1(b,\varphi ))$$
This can be done exactly as proposition 3.2 of .
Alternatively, we can reduce the isometry property to a purely algebraic statement as follows.
Fix the complex structure $`I_1`$ on $`T^2\times ^2`$. The instanton moduli space $``$ is then identified with the moduli space of $`\alpha `$-stable holomorphic vector bundles $`T\times ^1`$ as a Kähler manifold. Moreover, its tangent space becomes identified with $`H^1(T\times ^1,\mathrm{End})`$. One can define a complex symplectic structure on $``$ via the bilinear pairing:
$$H^1(T\times ^1,\mathrm{End})\times H^1(T\times ^1,\mathrm{End})\stackrel{\omega }{}H^2(T\times ^1,\mathrm{End})=$$
On the other hand, the moduli space of Higgs pairs $`\widehat{}`$ becomes identified, as a Kähler manifold, with the moduli space of stable parabolic Higgs bundles. The tangent is then given by the hypercohomology $`^1`$ of the following complex of sheaves:
$$ParEnd(𝒱)\stackrel{[,\varphi ]}{}\mathrm{\Lambda }^1ParEnd(𝒱)$$
where $`ParEnd(𝒱)`$ is the sheaf of parabolic endomorphism of the holomorphic Higgs bundle $`𝒱`$, see for a detailed explanation. A complex symplectic structure on $`\widehat{}`$ can be defined via the bilinear pairing
$$^1\times ^1\stackrel{\widehat{\omega }}{}^2=$$
In order to show that the Nahm transform is an isometry, it is enough to prove that the holomorphic version of the Nahm transform (see ) preserves the bilinear pairings above. This is an algebraic statement, which one can hope to prove using spectral sequences. Indeed, as we mentioned before, the holomorphic version of the Nahm transform of doubly-periodic instantons is an example of a Fourier-Mukai transform, which usually preserves this type of pairings.
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# NEW DETREMINATION METHOD OF PRIMORDIAL Li ABUNDANCE
## 1. Introduction
Recent spectral and photometric observations of Type Ia supernovae at high redshifts (Riess et al. 1998; Perlmutter et al. 1999) have raised a possibility that the cosmic expansion is accelerated. For a flat cosmology these data have $`\chi ^2`$-minimum around $`\mathrm{\Omega }_00.3`$ and $`\mathrm{\Omega }_\mathrm{\Lambda }0.7`$, allowing Hubble time $``$15Gyr which is not inconsistent with the age of the Milky Way constrained from the observations of the oldest globular clusters.
Cosmological model for primordial nucleosynthesis provides independent method to determine $`\mathrm{\Omega }_0`$. The Big-Bang nucleosynthesis model (Copi et al. 1995) predicts $`0.04\mathrm{\Omega }_bh_{50}^20.08`$. Combining this value with X-ray observations of rich clusters that indicate 0.3$`h_{50}^{3/2}\mathrm{\Omega }_b`$/$`\mathrm{\Omega }_0`$ (Bahcall et al. 1995; White et al. 1993), total $`\mathrm{\Omega }_0`$ turns out to be $`\mathrm{\Omega }_0h_{50}^{1/2}`$ 0.1$``$0.3, which is consistent with flat cosmology.
However, in the determination of $`\mathrm{\Omega }_b`$, a difficulty has been imposed by recent detections of a low deuterium abundance, 2.9$`\times `$10<sup>-5</sup> $``$ D/H $``$4.0$`\times `$10<sup>-5</sup>, in Lyman-$`\alpha `$ clouds along the line of sight to high red-shift quasars (Burles & Tytler 1998ab). Primordial abundance of <sup>7</sup>Li is constrained from the observed ”Spite plateau”, 0.91$`\times `$10<sup>-10</sup> $``$ <sup>7</sup>Li/H$``$1.91$`\times `$10<sup>-10</sup> (Ryan et al. 2000a), and the <sup>4</sup>He abundance by mass, 0.226$`Y_p`$0.247 (Olive et al. 1999), from the observations in the HII regions. In order to satisfy these abundance constraints by single $`\mathrm{\Omega }_b`$ value, one has to assume an appreciable depletion in the observed abundance of <sup>7</sup>Li, which is still controversial both theoretically and observationally. We are now forced to critically study the uncertainty. An independent method to determine the primordial <sup>7</sup>Li is also desirable.
## 2. Primordial Nucleosynthesis
### 2.1. <sup>4</sup>He vs. D and Neutrino Degeneracy
Shown in Fig. 1 is the comparison between the observed abundance constraints on <sup>4</sup>He, D/H, and <sup>7</sup>Li/H and the calculated curves in the homogeneous Big-Bang model as a function of $`\eta `$, where $`\eta =n_B/n_\gamma `$ and $`\mathrm{\Omega }_bh_{50}^2`$ = $`\eta `$ 1.464$`\times `$10<sup>8</sup>. Solid curves display the theoretical prediction of primordial abundances in the standard particle model for neutrino, which preserves the lepton symmetry $`L_\nu `$=0.
There is now a good collection of abundance information on the <sup>4</sup>He mass fraction, $`Y_p`$, in over 50 extragalactic HII regions, from which the upper limit on primordial abundance, $`Y_p0.240`$, and a systematic error, $`\mathrm{\Delta }Y_{sys}`$=0.005, were extracted. Unfortunately, for this upper limit one cannot find $`\mathrm{\Omega }_b`$ to satisfy both abundance constraints on <sup>4</sup>He and D/H. (See the solid curves in Fig. 1.)
It has been recognized that $`\mathrm{\Delta }Y_{sys}`$ may even be larger (Izatov et al. 1994; Thuan 2000), making the upper limit as large as $`Y_p`$0.247. If this upper limit is adopted, the Universe model with $`\eta `$5$`\times `$10<sup>-10</sup> is marginally consistent with both abundance constraints. However, since even smaller value, $`Y_p`$=0.235$`\pm `$0.003, in low-metallicity extragalactic HII regions has been reported by Peimbert & Peimbert (2000), this potential conflict is to be studied more carefully.
One possible solution is to introduce a lepton asymmetry. Theoretically, it is natural to assume that both baryon and lepton symmetries are simultaneously broken, $`B`$0 and $`L_\nu `$0, due to the CP violation in baryogenesis. $`L_\nu `$0 is fulfilled by neutrino degeneracy with non-zero $`\xi _{\nu _e}`$, where $`\xi _{\nu _e}=\mu _{\nu _e}/kT_\nu `$ and $`\mu _{\nu _e}`$ is the chemical potential of electron neutrino. Since neutrinos had energy density comparable to the densities due to photons and charged leptons in the early Universe, even a small degeneracy 0$`<\xi _{\nu _e}`$1 leads to an appreciable decrease in the neutron-to-proton number ratio, slightly faster acceleration of the Universal expansion, and a small increase of the weak-decoupling temperature. As a net result, <sup>4</sup>He abundance decreases with increasing $`\xi _{\nu _e}`$, as shown in Fig. 1, while keeping D/H and <sup>7</sup>Li/H almost unchanged in logarithmic scale (Kajino & Orito 1998). Since the abundance constraint on primordial <sup>4</sup>He is more accurate than the other light elements, this helps determine the most likely $`\xi _{\nu _e}`$. $`\xi _{\nu _e}0.05`$ can best fit the <sup>4</sup>He abundance as well as low deuterium abundance D/H$``$10<sup>-5</sup>, leaving inevitable requirement that the observed abundance level of Spite plateau, <sup>7</sup>Li/H$``$10<sup>-10</sup>, should be the result of depleted primordial abundance.
### 2.2. <sup>7</sup>Li vs. D
There are several input parameters in the primordial nucleosynthesis calculation. As the number of light neutrino families $`N_\nu `$ = 3 and the neutron lifetime $`\tau _n`$ = 886.7 $`\pm `$ 1.9 s are known, the remaining major uncertainty arises from input nuclear reaction data. We did not take account of the effects of sterile neutrino which is a hypothetical particle for interpreting flavor mixing.
Laboratory cross section measurements ever done provide rather precise thermonuclear reaction rates for the production of D, T, <sup>3</sup>He, and <sup>4</sup>He. It however was claimed in literature (Smith et al. 1993) that the <sup>7</sup>Li abundance is strongly subject to large error bars associated with the measured cross sections for <sup>4</sup>He(<sup>3</sup>H,$`\gamma `$)<sup>7</sup>Li at $`\eta <`$ 2$`\times 10^{10}`$ and <sup>4</sup>He(<sup>3</sup>He,$`\gamma `$)<sup>7</sup>Be at 3$`\times 10^{10}<\eta `$. There are in fact several inconsistent data with one another, leading to large uncertainty in the primordial <sup>7</sup>Li, as displayed by long-dash-dotted curves in Fig. 1.
We studied these two reactions very carefully and concluded that the proper 2$`\sigma `$ error bars could be 1/4$``$1/3 of the previous ones (Kajino et al. 2000). This improvement owes mostly to, first, the new precise measurement (Brune et al. 1994) of the cross sections for <sup>4</sup>He(<sup>3</sup>H,$`\gamma `$)<sup>7</sup>Li and, second, the systematic theoretical studies of both reaction dynamics and quantum nuclear structures of <sup>7</sup>Li and <sup>7</sup>Be, whose validity is critically tested by electromagnetic form factors measured by high-energy electron scattering experiments.
When our recommended error estimate is applied to the determination of $`\mathrm{\Omega }_b`$ in Fig. 1, we lose $`\mathrm{\Omega }_b`$ value to explain both D/H and <sup>7</sup>Li/H simultaneously. If we allow for larger primordial <sup>7</sup>Li abundance in Population II halo stars because
of possible lithium depletion for diffusion or rotation-induced mixing of matter (Deliyannis et al. 1998 ; Pinsonneault et al. 1992) or some systematic uncertainty in the model atmospheres (Kurucz 1995), we can recover the concordance. Taking depletion factor $``$ 2.5, $`\mathrm{\Omega }_bh_{50}^2`$ 0.075 best fits all abundance constraints in the homogeneous Big-Bang model. Note that larger $`\mathrm{\Omega }_bh_{50}^2`$ 0.2 is allowed in the inhomogeoeus Big-Bang model (Kajino & Orito 1998).
## 3. <sup>7</sup>Li/<sup>6</sup>Li Ratio in the Interstellar Medium (ISM)
The lithium in the ISM is almost free from the complicated stellar processes. A diffuse cloud along the line of sight to $`\zeta `$Oph was observed to show the lithium abundance depleted by 1.58dex from the meteoritic solar-system value 12.3. This is due to dust grain formation (Savage & Sembach 1996). The isotopic ratio is free from such condensation effects and represents the real ratio of chemical compositions in the gas phase. The D/H (Wannier 1980) and <sup>3</sup>He/H (Rood et al. 1995) abundance ratios in the ISMs have been observed over wide Galactocentric distance range 0$``$R$``$12kpc and used to constrain the primordial abundance of D/H (Dearborn et al. 1996), but the distribution of <sup>7</sup>Li/<sup>6</sup>Li was poorly known.
### 3.1. Observation
Observations of isotopic abundance ratio, <sup>7</sup>Li/<sup>6</sup>Li, have been performed by several groups (Ferlet & Dennefeld 1983, Lemoine et al. 1993, 1995, Meyer et al. 1993) only for the ISMs in our solar neighborhood. The observed ratio is less than 12.3 and larger than 2.1 which is a predicted GCR abundance ratio.
Using the Coude spectrograph of the 74-inch telescope at Okayama Astrophysical Observatory, Japan, we have succeeded for the first time in the determination of <sup>7</sup>Li/<sup>6</sup>Li in the diffuse cloud along the line of sight to $`\chi ^2`$Ori, which is a member of OB association Gem-OB1, being located at R = 10kpc (Kawanomoto et al. 2000). The telescope performance was R=43,000 (with slit width of 100 $`\mu `$m), exp=50hours, and S/N=2,800.
We found a decreasing gradient of <sup>7</sup>Li/<sup>6</sup>Li, as shown in Figs. 2 & 3. It is interpreted as a result of gradual extinction of the stellar production of <sup>7</sup>Li.
### 3.2. A New Method to Determine Primordial <sup>7</sup>Li
In order to study the sensitivity to the primordial lithium abundance, we have calculated Galactic chemical evolution (GCE) of lithium (Kawanomoto 2000). We adopted a hybrid model (Ryan et al. 2000b) of the inhomogeneous GCE model (Suzuki et al. 1999), which was constructed for the early evolution of metal-deficient stars, being smoothly connected with a simple one-zone GCE model for later evolution. Five different sources of lithium production are included in this model: Primordial nucleosynthesis, GCR interactions with ISM, $`\nu `$-induced nucleosynthesis in Type II SNe, AGB star nucleosynthesis, and nova nucleosynthesis. We took an approximation that each ring having different Galactocentric distance evolves independently so that the observed present day star-formation-rate and the gas fraction are reproduced very well.
The calculated time variation of <sup>7</sup>Li/<sup>6</sup>Li is shown as a function of R in Fig. 2. Remarkable decrease of <sup>7</sup>Li/<sup>6</sup>Li in the inner region is caused by faster gas consumption for the star formation. It is discussed in literature that the meteoretic chemical compositions are peculier and different from those of ISM because they were possibly polluted by nearby AGB star. One might speculate another possibility that the solar-system might ever have moved outward over hundreds of turns of the Galactic disc, keeping high <sup>7</sup>Li/<sup>6</sup>Li = 12.3 as it was in the original position when the solar system was isolated from viscous gas component at $`t_G26`$Gyr.
Figure 3 displays sensitivity of the <sup>7</sup>Li/<sup>6</sup>Li ratio at the present time $`t_G`$=12Gyr to the primordial abundance of <sup>7</sup>Li. It is very sensitive to <sup>7</sup>Li<sub>p</sub>. Except for old data point at $`\rho `$-Oph (Ferlet & Dennefeld 1983), which has the largest error bar among all data for the solar neighborhood, the observed ratios look more consistent with <sup>7</sup>Li<sub>p</sub> = (1.4$``$3.5)$`\times `$10<sup>-10</sup> than <sup>7</sup>Li<sub>p</sub> = 1.5$`\times `$10<sup>-9</sup>. More data with smaller error bars are highly desirable in order to convince the gradient of the <sup>7</sup>Li/<sup>6</sup>Li ratio and to determine the primordial abundance of <sup>7</sup>Li in this method.
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Kajino, T., Orito, M., Sakai, K., & Deliyannis, P.C. 2000, in preparation.
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# An Extragalactic H i Cloud with No Optical Counterpart?
## 1 INTRODUCTION
Until recently, there has been no extensive survey of the extragalactic sky in neutral hydrogen (H i). The mass distribution in the Universe has therefore been traced by galaxy surveys in other wavebands, particularly in the optical, but also in the infrared using the IRAS satellite. Because of the limited nature of existing H i surveys, it is unclear whether there is a significant local population of H i-rich objects not seen in optical and infrared surveys (away from the Galactic Plane). In particular, the existence of H i clouds without optical counterparts, or ’H i protogalaxies’ might significantly add to the local H i mass density of the Universe. Recent observations by Zwaan et al. (1997) and Schneider, Spitzak & Rosenberg (1998) arrive at somewhat differing conclusions regarding the low H i-mass galaxy population, leading to some uncertainty in the contribution of these objects to the HI mass in the local Universe.
Limits to the presence of a population of H i-rich clouds and ‘protogalaxies’ have been set by a number of ‘blind’ or semi-‘blind’ H i surveys conducted in a range of environments: clusters (Barnes et al., 1997; Dickey, 1997), groups (Kraan-Korteweg et al., 1999; Banks et al., 1999; Fisher & Tully, 1981a), voids (Krumm & Brosch, 1984), and in the general field (Henning, 1995; Sorar, 1994; Spitzak & Schneider, 1998; Shostak, 1977). In addition, there have been numerous observations in off-source calibration regions for galaxies which have been optically-selected. For example, Fisher & Tully (1981b) serendipitously detected a couple of dozen galaxies in the calibration scans from their survey, all with optical counterparts.
A number of H i clouds without optical counterparts have been found. However, many appear to be closely associated with massive galaxies (Williams & Van Gorkom, 1988; Giovanelli et al., 1995) and perhaps the product of tidal interactions (Hibbard & van Gorkom, 1996). A number of potential H i protogalaxies have later been classified as High Velocity Clouds (Mathewson et al., 1975), and argument continues over the extragalactic nature of a specific class of Compact High Velocity Clouds (Braun & Burton, 1999; Zwaan et al., 2000). One of the larger blind H i surveys to date (Spitzak & Schneider, 1998) detected 75 galaxies in 55 deg<sup>2</sup> of sky. One detection was found not to have any obvious optical counterpart, though a bright star just near the peak of the H i emission may obscure a low surface brightness optical companion.
The closest and best examples of H i clouds without prominent counterparts are the Virgo Cloud HI 1225+01 (Giovanelli & Haynes, 1989), and the Leo ring (Schneider et al., 1983). Both of these H i clouds were later found to be associated with optical galaxies. HI 1225+01 was found serendipitously during routine calibration observations. It comprises two regions of H i emission: the largest has a Magellanic-type dwarf irregular galaxy at the position of the highest H i column density, while the smaller, possibly infalling cloud shows no optical emission to a faint limiting magnitude (Chengalur et al., 1995). The larger component has a total dynamical mass of $`1.0\times 10^{10}d_{20}_{}`$, and a gas mass of $`3\times 10^9d_{20}^2_{}`$ where $`d_{20}`$ is the distance in units of $`20\text{Mpc}`$. The heliocentric velocity of the cloud is 1275 km s<sup>-1</sup> and the extent of the H i emission is about $`200d_{20}\text{kpc}`$. The Leo ring is an intergalactic cloud found during H i observations of the Leo group of galaxies. It comprises a large ring of H i emission surrounding NGC 3384 and M105, with an H i mass of $`10^9_{}`$. The velocity structure suggests that it is probably not a tidal tail from the nearby galaxies, but a primordial gas cloud which has not started forming stars (Schneider, 1989). The cloud is gravitationally bound to the optical galaxies in the group. Both of these objects have been used in H$`\alpha `$ studies to determine limits on the local ionising background (Donahue et al., 1995), finding limits which suggest that quasar light, not galactic light, dominates the local ionizing background at low redshift.
The H i Parkes<sup>1</sup><sup>1</sup>1The Parkes telescope is part of the Australia Telescope which is funded by the Commonwealth of Australia for operation as a National Facility managed by CSIRO. All Sky Survey (HIPASS) has been underway since 1997 and, when complete, will be the largest blind H i survey to date, surveying a volume at least two orders of magnitude larger than any previous survey. In early HIPASS observations of 600 deg<sup>2</sup> in the Centaurus region, Banks et al. (1999) found 10 new members of the nearby Centaurus A group. However, all of these have optical companions, although 5 are faint and were previously uncatalogued. Other early observations include the south celestial cap (Putman et al., 1998; Kilborn, Webster & Staveley-Smith, 1999). During routine inspection of this data, a resolved H i detection, HIPASS J1712-64, was selected for follow-up observations as it showed no catalogued galaxy in the vicinity and no optical counterpart on the Digitised Sky Survey. It also had a low systemic velocity (though well-separated from Galactic gas and High Velocity Clouds) and a significant angular size.
The observations of HIPASS J1712-64 are discussed in § 2; the H i structure, dynamics and optical limits are discussed in § 3. Alternative High Velocity Cloud and Local Group hypotheses are discussed in § 4, the nature of HIPASS J1712-64 is discussed in § 5, and the results are summarised in § 6. Throughout we assume a Hubble constant of 75 km s<sup>-1</sup>Mpc<sup>-1</sup>.
## 2 OBSERVATIONS AND DATA REDUCTION
### 2.1 HIPASS
HIPASS is a blind H i survey of the sky $`\delta <2\mathrm{°}`$. The survey uses the Parkes 64 m telescope which is equipped with a 13-beam receiver (Staveley-Smith et al., 1996). The observations began in February 1997, and will be finished by the start of 2000. A northern extension is being contemplated from Parkes, and parts of a complementary northern survey are currently being conducted at Jodrell Bank Observatory’s Lovell Telescope. The velocity range covered is $`1200`$ to $`12700\text{km\hspace{0.17em} s}\text{-1}`$. Observations are taken by scanning the multibeam receiver in declination strips of length $``$ 8°. Each declination scan is separated by 35′ in R.A., and each area of sky is scanned five times, resulting in a final scan separation of 7′. The final integration time is approximately 460 s beam<sup>-1</sup>. The mean telescope FWHM beam is 14$`\stackrel{}{\mathrm{.}}`$3, although the gridding process increases this to $`15\stackrel{}{\mathrm{.}}5`$ (Barnes et al., 2000).
Bandpass calibration, spectral smoothing and Doppler correction of the data are applied in real time at the telescope (Barnes et al., 1998). Once all data are collected, the spectra are gridded into data cubes which have an rms noise level of $`13\text{mJy}`$ beam<sup>-1</sup>. The channel spacing in the final cubes is $`13.2\text{km\hspace{0.17em} s}\text{-1}`$, the FWHM resolution is $`18.0\text{km\hspace{0.17em} s}\text{-1}`$ and the pixel size is $`4\mathrm{}\times 4\mathrm{}`$.
### 2.2 H i Observations
HIPASS J1712-64 was discovered during a routine visual inspection of data cubes. The detection is very significant: it has a peak flux of $`150\text{mJy}`$ beam<sup>-1</sup>, corresponding to a signal-to-noise ratio of about 12. It is also significantly resolved in the beam of the Parkes telescope, extending over at least two beamwidths. On 1999 June 2, a further 2.5 hour observation was made at the Parkes telescope, using the multibeam receiver, to confirm the detection. The observation was of a $`3\mathrm{°}\times 4\mathrm{°}`$ field centered on the original detection. This observation also resulted in a clear positive detection. The rms noise in the follow-up observation is $`19\text{mJy}`$ beam<sup>-1</sup>. The combined, and spatially integrated, Parkes spectrum is shown in Fig. An Extragalactic H i Cloud with No Optical Counterpart?. The total flux density peaks at 250 mJy at a heliocentric velocity of $``$ 455 km s<sup>-1</sup>.
Higher resolution observations in the H i line were made at the Australia Telescope Compact Array (ATCA) on 1999 June 30. Approximately 12 hours of data was obtained using the 375 m array, and the pointing center for the observations was R.A. $`17^\mathrm{h}12^\mathrm{m}13^\mathrm{s}`$, Decl. $`64^\mathrm{d}39^\mathrm{m}12^\mathrm{s}`$ (J2000). The secondary calibrator PKS 1814-637 was observed every 30 minutes for 4 minutes, to calibrate the amplitude and phase of the visibility data. The primary flux scale was set using the primary calibrator PKS 1934-638. The bandwidth of 8 MHz was divided into 512 channels, resulting in a channel spacing of 3.3 km s<sup>-1</sup> and a FWHM resolution, before any smoothing, of 4.0 km s<sup>-1</sup>.
The data were reduced using the miriad reduction package. The data were edited, calibrated in amplitude and phase and bandpass-corrected. Several bright continuum sources were then removed by fitting spectral baselines to the line-free channels in the visibility domain. Natural weighting was used to form the data cube, which was cleaned until the absolute maximum residual in each plane fell below 10 mJy. The final FWHM beam is $`2\stackrel{}{\mathrm{.}}0\times 1\stackrel{}{\mathrm{.}}9`$, and the rms noise is 3.7 mJy beam<sup>-1</sup>. No correction for the primary beam pattern was applied.
J1712-64 was again detected in the ATCA data. The spatially integrated spectrum (smoothed to the same velocity resolution) is overlaid on the Parkes/HIPASS data in Fig. An Extragalactic H i Cloud with No Optical Counterpart?. All except $`20`$ % of the Parkes flux density (mainly at higher velocities) is recovered. This implies that some diffuse or low-level emission is being missed by the ATCA. The ATCA channel images (Hanning-smoothed to 6.6 km s<sup>-1</sup> resolution) are shown in Fig. An Extragalactic H i Cloud with No Optical Counterpart?.
A combined cube of the Parkes and ATCA data set was made by spectrally smoothing the ATCA data to match the HIPASS data, then spatially resampling the HIPASS data to match the ATCA data. The data were then combined using the miriad task immerge. It was assumed that both the ATCA and HIPASS flux density scales were accurate to a few per cent. The final column density image, with contours overlaid is shown in Fig. An Extragalactic H i Cloud with No Optical Counterpart?. The ATCA velocity field, derived from the gipsy task moments, is shown in Fig. An Extragalactic H i Cloud with No Optical Counterpart? overlaid on the column density image.
### 2.3 Optical Observations
A search for an optical counterpart using blue and red film copies of the UKST/ESO southern sky survey revealed no optical counterpart. Subsequently, two deep original plates were obtained from the UK Schmidt Telescope (UKST) plate archive and subjected to a photographic enhancement technique (Malin, 1978) which is capable of revealing extended features over 5 mag fainter than the night sky (Malin & Hadley, 1997). At the site of the UKST (Siding Spring) this is typically $`\mu (B)=22.5`$ mag arcsec<sup>-2</sup>.
Only two suitable plates were available, and an image was made by combining photographically enhanced derivatives from both of them: J1659 (IIIa-J/GG 395, 395–530 nm) and OR15092 (IIIa-F/RG 590, 590-690 nm). Both were plates of excellent quality but in both cases the object centre was close to the vignetting region towards the edge of the UKST focal plane.
The IIIa-J passband is the most sensitive to extended stellar light (assuming a solar spectrum) and would reveal extended features at least as faint as 27.5 mag arcsec<sup>-2</sup>. The derivative from the red-sensitive plate achieves a detection limit of about $``$ 26.8 mag arcsec<sup>-2</sup> because of the brighter night sky in the red. Nothing unusual is seen on deep derivatives from either plate, or on the combined image. However, it should be mentioned that large, featureless low surface brightness objects are particularly difficult to detect in crowded fields. For an optical companion of 1′ in diameter and lying just beneath our optical limit, these limits suggest an absolute magnitude $`M_B>8.8`$ mag, and an H i mass-to-light ratio $`M_{HI}/L_B>24`$ (Table 2) .
The source was also imaged in the optical $`I`$-band at the ANU 40-inch telescope at Siding Springs Observatory on 1999 April 15. It was also imaged in the B and R-bands by the CTIO 40-inch telescope on 1999 May 4. The data were reduced using the iraf<sup>2</sup><sup>2</sup>2iraf is distributed by the National Optical Astronomy Observatories, which is operated by the Association of Universities for Research in Astronomy, Inc., under contract to the National Science Foundation. package and calibration for the $`B`$ and $`R`$-band images used photometric observations of the Landolt Standard field 104 (Landolt, 1983). The optical point source magnitude limits are $`B19.8`$ and $`R19.2`$ from the CCD observations. None of the images show any evidence of extended optical emission near the location of the H i detection. Fig. An Extragalactic H i Cloud with No Optical Counterpart? shows the H i contours of HIPASS J1712-64 overlayed on the combined $`B`$ and $`R`$-band CCD images, showing a dense starfield but no extended objects.
## 3 PHYSICAL PARAMETERS
### 3.1 H i Structure
HIPASS J1712-64 is well-resolved in each of the velocity channels with significant emission (see Fig. An Extragalactic H i Cloud with No Optical Counterpart?). The object appears to have two major components. The prominent one (hereafter the NE component) is centered at R.A. $`17^\mathrm{h}12^\mathrm{m}35^\mathrm{s}`$, Decl. $`64^{}38\mathrm{}12\mathrm{}`$ (J2000) and has a velocity spread from 438 to 464 km s<sup>-1</sup> in the ATCA data. The systemic velocity is $`451\text{km\hspace{0.17em} s}\text{-1}`$ and the velocity width is $`W_{50}=11\text{km\hspace{0.17em} s}\text{-1}`$ (Table 1). The velocities appear to increase from southeast to northwest across this component, but the difference is only a few km s<sup>-1</sup> at the most (Fig. An Extragalactic H i Cloud with No Optical Counterpart?). The other component (hereafter the SW component) is centered at R.A. $`17^\mathrm{h}11^\mathrm{m}35^\mathrm{s}`$, Decl. $`64^{}45\mathrm{}11\mathrm{}`$ (J2000) and has a velocity spread from 451 to 477 km s<sup>-1</sup>. The systemic velocity is $`464\text{km\hspace{0.17em} s}\text{-1}`$ and the velocity width is $`W_{50}=20\text{km\hspace{0.17em} s}\text{-1}`$ (Table 1). Any gradient across this component is again only a few km s<sup>-1</sup> (Fig. An Extragalactic H i Cloud with No Optical Counterpart?).
A bridge of emission appears to join the NE and SW components in a manner reminiscent of the Virgo Cloud HI 1225+01 (Chengalur et al., 1995). Velocities appear to increase smoothly from the NE to the SW components. The overall H i structure appears to be that of two separate, but interacting components.
### 3.2 Optical Limits
The source is located close to the Galactic plane at $`(l,b)=(326\stackrel{}{\mathrm{.}}6,14\stackrel{}{\mathrm{.}}6`$), and therefore extinction may be important. We have estimated this using three methods. The local H i column density from the HIPASS data (reprocessed to avoid any spatial filtering) is $`484\text{K\hspace{0.17em}km\hspace{0.17em} s}\text{-1}`$, which corresponds to $`N_{\mathrm{HI}}=8.8\times 10^{20}\text{atoms\hspace{0.17em} cm}\text{-2}`$. Following the conversion formula of Burstein & Heiles (1978), this implies $`E(BV)=0.12`$ mag. The extinction determined from the COBE/DIRBE and IRAS/ISSA maps based on dust temperature is essentially identical: $`E(BV)=0.11`$ mag (Schlegel et al., 1998). Finally, the Burstein & Heiles (1982) extinction maps suggest a similar value $`E(BV)=0.09`$ mag. Taking the mean value, $`E(BV)=0.11`$, this implies a blue-light absorption of $`A_B=4.0E(BV)=0.44`$ mag. Thus although the stellar density in the optical images is quite high (see Fig. An Extragalactic H i Cloud with No Optical Counterpart?), there is no evidence for large optical extinction in the direction of the H i detection.
An optical counterpart could be overlooked if the J1712-64 is very close to the Milky Way and resolved into stars. One method of detecting a resolved galaxy in a crowded field of stars is to search a colour-colour plot of the stars in a field, to look for a separate stellar population. This method was used to detect the Sagittarius dwarf galaxy (Ibata, Gilmore & Irwin, 1995). We applied this method to the $`B`$ and $`R`$-band CCD images. These images cover an area of about 1 deg<sup>2</sup>, which is substantially larger than the H i detection. Fig An Extragalactic H i Cloud with No Optical Counterpart? shows a histogram of stellar magnitudes for a 20′ field centred on the H i detection, and other nearby fields of the same size. These fields are statistically equivalent, thus there is no evidence for a resolved galaxy in the direction of the H i detection.
### 3.3 Distance and Mass
The systemic velocity of the H i detection is $`v_{\mathrm{sys}}=451`$ km s<sup>-1</sup> (NE component). The velocity with respect to the Local Group can be determined using $`v_{\mathrm{LG}}=v_{\mathrm{sys}}79\mathrm{cos}l\mathrm{cos}b+296\mathrm{sin}l\mathrm{cos}b36\mathrm{sin}b`$ (Yahil, Tammann & Sandage, 1977). This gives $`v_{\mathrm{LG}}=239`$ km s<sup>-1</sup> which, if we assume HIPASS J1712-64 to be extragalactic and $`H_{}=75`$ km s<sup>-1</sup> Mpc<sup>-1</sup>, corresponds to a distance of 3.2 Mpc (see Table 2). This would place HIPASS J1712-64 at some distance beyond the Local Group. This distance is subject to considerable uncertainty. Alternate Local Group velocities range from $`v_{\mathrm{LG}}=291\text{km\hspace{0.17em} s}\text{-1}`$ (using $`300\mathrm{sin}\mathrm{}\mathrm{cos}b`$) to $`v_{\mathrm{LG}}=172`$ km s<sup>-1</sup> (potent, Bertschinger (1990)), giving distances of 3.9 Mpc and 2.3 Mpc, respectively. Without the aid of optical distance measures, subsequent calculations involving the distance to J1712-64 will therefore be based on the value of $`3.2(\pm 1)`$ Mpc. At this distance, the projected separation of the NW and SW components is $`9\text{kpc}`$, the overall H i size is $`15\text{kpc}`$ and the H i mass is $`1.7\times 10^7_{}`$. For a substantially smaller distance, say 100 kpc (see § 4), the separation, size and mass are $`300`$ pc, $`500`$ pc and $`1.7\times 10^4_{}`$, respectively.
The velocities within J1712-64 are very low. The velocity separation of the NE and SW components is only 12 km s<sup>-1</sup>, and the velocity dispersions within each component are small. For the NE component, the mean velocity dispersion from the ATCA data is $`\sigma 4\text{km\hspace{0.17em} s}\text{-1}`$ (see Table 1), possibly slightly more for the SW component, although the S/N ratio is too low for a reliable estimate.
If the velocity difference between the NE and SW components is indicative of the rotation of a bound, binary system at 3.2 Mpc, the minimum system mass is $`1.5\times 10^8_{}`$ (Faber & Gallagher, 1979), an order of magnitude above the H i mass. The ‘virial’ distance, at which the total mass is 1.3 times the H i mass (to account for He), is implausibly distant ($`>20`$ Mpc) implying that the system contains substantial dark matter ($`9M_{\mathrm{HI}}`$, which would be a reasonable value for disk galaxies), or else is not bound. The orbital period of the two components corresponding to the minimum total mass is $`4(d/3.2\mathrm{Mpc})`$ Gyr.
The NE component alone can be modelled as a pressure-supported sphere of gas of total enclosed mass
$$M_T=\frac{\alpha \sigma ^2r}{G}$$
(1)
where $`\sigma `$ is the line-of-sight velocity dispersion ($`4\text{km\hspace{0.17em} s}\text{-1}`$), $`r`$ is the radius ($`4`$ kpc) and $`\alpha 13`$ depending on the mass distribution. For an isothermal sphere, $`\alpha 2`$, which gives $`M_T(\mathrm{NE})=3.0\times 10^7_{}`$, which is only $`1.8`$ times higher than the H i mass for this component. Again, for a distance of 100 kpc, the mass drops to $`M_T(\mathrm{NE})=10^6_{}`$, but this is now $`60`$ times greater than the H i mass if in hydrostatic equilibrium.
### 3.4 Stability
The critical column density above which local axisymmetric instabilities can occur in a uniformly rotating gaseous disk is given by the Toomre stability criterion:
$$\mathrm{\Sigma }_c=\frac{2v_s\mathrm{\Omega }}{\pi G},$$
(2)
where $`v_s`$ is the sound speed and $`\mathrm{\Omega }`$ is angular frequency (Binney and Tremaine, 1987). Kennicutt (1989) has shown that similar relations for differentially rotating stellar disks approximately predict the column density above which star formation occurs in spiral galaxies. If we use $`\mathrm{\Omega }=v/r`$ with $`v\sigma v_s`$ and $`r=4`$ kpc for the NE component of J1712-64 (assuming it to be extragalactic), then we obtain $`\mathrm{\Sigma }_c8\times 10^{19}\text{atoms\hspace{0.17em} cm}\text{-2}`$. This is a factor of two above the peak measured column density of just $`3.5\times 10^{19}\text{atoms\hspace{0.17em} cm}\text{-2}`$ (Fig An Extragalactic H i Cloud with No Optical Counterpart?), consistent with the absence of visible star-formation.
## 4 ALTERNATIVES
### 4.1 A Local Group Galaxy?
Nearby, faint galaxies which are resolved into stars can be especially hard to spot even far from the crowding in the Galactic Plane (e.g. Sextans – Irwin et al. (1990)). In § 3.2 we have shown that there is no evidence in color-magnitude diagrams for this. In addition, the Local Group velocity, $`v_{\mathrm{LG}}=239`$ km s<sup>-1</sup> is well beyond the $`\pm 60`$ km s<sup>-1</sup> which appears to encompass most known Local Group members (Grebel, 1997; van den Bergh, 1999). This is demonstrated in Fig. An Extragalactic H i Cloud with No Optical Counterpart?, where J1712-64 has the highest heliocentric and Local Group velocity of all the probable and possible Local Group members plotted. It has also has a higher velocity than the class of ‘Local Group outliers’ which have distances of $`12`$ Mpc and are regarded by Grebel (1997) as only potential members.
In addition, the irregular galaxy IC 4662 lies 3$`\stackrel{}{\mathrm{.}}`$7 on the sky from HIPASS J1712-64, has a velocity of $`v_{\mathrm{sys}}=302\pm 1\text{km\hspace{0.17em} s}\text{-1}`$, almost 150 km s<sup>-1</sup> lower, yet is not classified as belonging to the Local Group (van den Bergh, 1999). Therefore, from its velocity, there is no evidence for J1712-64 being a galaxy lying within the Local Group.
### 4.2 A High Velocity Cloud?
The most likely alternative explanation for the classification of HIPASS J1712-64 is that it is an unusual high-velocity cloud (HVC) lying at the edge of our own Galaxy. HVCs are H i objects which do not fit into a simple model of Galactic rotation and do not have optical counterparts (Wakker & van Woerden, 1997; Braun & Burton, 1999). They are widespread across the sky, with a covering factor estimated to be 0.2 – 0.4 for $`N_{\mathrm{HI}}>7\times 10^{17}\text{atoms\hspace{0.17em} cm}\text{-2}`$ (Wakker & van Woerden, 1997). There are two good reason to exclude the possibility of J1712-64 being a normal HVC. The first is its velocity, which is 40 km s<sup>-1</sup> higher than any currently catalogued HVC (e.g. those around the Large Magellanic Cloud – Morras et al. (2000)). The second is J1712-64’s HI structure. Most HVCs have a core-halo morphology with only 10–50% of the total flux able to be recovered by interferometry (Wakker & van Woerden, 1997). As shown in Fig. An Extragalactic H i Cloud with No Optical Counterpart?, however, the ATCA observations were able to recover 80% of the total flux. It may, however, be an extreme version of the compact HVCs studied in detail by Braun & Burton (2000). These objects exhibit narrow linewidths and also show apparent velocity gradients.
J1712-64 lies in a 2400 deg<sup>2</sup> region around the south celestial cap (Putman et al., 1998) where many of the HVCs appear to have originated from the interaction of the Magellanic Clouds with the Milky Way. It is possible that a compact HVC similar to J1712-64 could have originated from gravitational scattering off the three-body Magellanic Clouds/Galaxy system. However, HIPASS J1712-64 does not lie close to the Clouds (it is 46° from the LMC), and the Leading Arm, which includes HVCs 334 and 352 of Wakker & van Woerden (1991) (see Figure 3 of Putman et al. (1998)), is $`20\mathrm{°}`$ away. A velocity separation of $`300\text{km\hspace{0.17em} s}\text{-1}`$ and a spatial separation from the LMC of $`60`$ kpc (or $`80`$ kpc from the Galaxy) would be consistent with an ejection event coinciding with the last LMC/SMC encounter which probably occurred $`2\times 10^8`$ yr ago (Gardiner & Noguchi, 1996). An examination of the spatial and velocity distribution of future objects, similar to HIPASS J1712-64, will ultimately demonstrate their relationship, if any, with the Magellanic System.
## 5 DISCUSSION
Although there have been many previous searches for extragalactic H i in optically blank fields (Fisher & Tully, 1981a; Briggs, 1990; Barnes et al., 1997; Schneider, Spitzak & Rosenberg, 1998), there has been little success, implying that massive H i ‘protogalaxies’ are rare at the present epoch. However, HIPASS J1712-64 has a very low H i mass, $`1.7\times 10^7`$ $`_{}`$, implying that similar objects may exist in abundance and may have escaped previous detection through lack of sensitivity. Some limits exist from the Cen A HIPASS results of Banks et al. (1999), who find 10 new group members, all with corresponding optical counterparts, in the 600 deg<sup>2</sup> they surveyed around the Cen A group. However, their sensitivity limit is $`10^7`$ $`_{}`$ so it is possible that slightly deeper observations (Disney et al., 1999) may reveal more such objects. The HIPASS observations of the south celestial cap, Decl. $`<62\mathrm{°}`$ (Kilborn, Webster & Staveley-Smith, 1999), comprise 2400 deg<sup>2</sup> of sky, but contains only a few other candidate ‘protogalaxies’, which are awaiting follow-up CCD data. However, HIPASS J1712-64 is the most significant detection in terms of flux density and angular size. Nevertheless, for there to exist large numbers of similar objects would require them to have H i masses $`<10^7`$ $`_{}`$.
Interestingly, a system with a similar velocity and H i morphology, ZOA J1616-55, was found in a search for galaxies in the Zone of Avoidance (ZOA) near $`\mathrm{}=325\mathrm{°}`$ (Staveley-Smith et al., 1998). This system has a similar distance and H i column density to HIPASS J1712-64 but a higher H i mass ($`9\times 10^7`$ $`_{}`$) and diameter (86 kpc). Although lying in the ZOA, the estimated blue absorption is only $`A_B=2.7`$ mag, so the lack of an optical or IR counterpart led Staveley-Smith et al. (1998) to conclude it was a galaxy pair of low optical surface brightness. J1712-64 is only 11° (projected separation 600 kpc) from J1616-55 and 19° (projected separation 1.1 Mpc) from the massive Circinus galaxy. Possibly, all three are part of a loose southern extension to the Cen A group. In the HVC hypothesis, the alternative is that Circinus is unrelated and both J1712-64 and J1616-55 have a Magellanic origin.
From a study of the Ly$`\alpha `$ forest lines and a VLA H i study, Shull, Penton & Stocke (1999) suggest that, at $`z=0`$, there may exist clouds, or sheets, of ionized gas $`1`$ Mpc in extent which contain $`30`$ times more baryons than the neutral gas confined to known galaxies. Objects such as J1712-64 may be the most easily visible manifestations of these large primordial clouds. So, although the neutral gas content in J1712-64 is hardly enough to make a star cluster, let alone a galaxy, there may be significantly more baryons in the region available to form a galaxy if they were able to cool. The possibility that the dynamical mass of the J1712-64 system is $`10`$ times greater than the H i mass suggests there may be dark matter of some sort associated with the system.
Isolated H i clouds can be used to set stringent limits on the local metagalactic background radiation. Deep narrow-band CCD images of the H$`\alpha `$ line in HI 1225+01 and the Leo ring place limits of $`J_0<38\times 10^{23}\text{ergs s}\text{-1}\text{ cm}\text{-2}\text{ sr}\text{-1}\text{ Hz}\text{-1}`$ for quasar-like light (Donahue et al., 1995). H$`\alpha `$ observations of J1712-64 may also help provide strong constraints on this background. Such observations would also be useful in providing a lower limit to the distance to J1712-64 in light of the observation that many nearby HVCs have faint but detectable H$`\alpha `$ emission probably resulting from the Galactic ionization field (Bland-Hawthorn & Maloney, 1999).
The mean velocity dispersion of the NE component of J1712-64 appears to be low, $`\sigma =4`$ km s<sup>-1</sup> (Table 1). This gives an upper limit to the kinetic temperature of $`T<2000`$ K. The existence of cool gas would be interesting as this requires that there are densities high enough to allow cooling and may therefore indicate that star-formation is about to commence. In order to search for cool gas, we examined the ATCA H i spectrum of PKS 1708-648 which is 19$`\stackrel{}{\mathrm{.}}`$4 from the center of the J1712-64. However, we were unable to detect any unambiguous evidence for absorption. This is probably not significant as the continuum source is well beyond the last measured contour of HIPASS J1712-64.
## 6 SUMMARY
An isolated H i cloud, HIPASS J1712-64, has been found during the course of the HIPASS survey. It appears to be a nearby extragalactic cloud lying beyond the outer fringes of the Local Group at a distance of $`3.2`$ Mpc, though an alternative explanation is that of an HVC with unique properties. Optical limits imply a surface brightness, corrected for obscuration of $`\mu (B)>27`$ mag arcsec<sup>-2</sup> and an H i mass-to-light ratio $`M_{HI}/L_B>24_{}/L_{}`$ for a putative optical counterpart of diameter 1′. The system appears to have two components, separated by 9 kpc, with a bridge of gas between them. The mean velocity dispersion of the NE component is low, $`4`$ km s<sup>-1</sup> implying that HIPASS J1712-64 contains cool gas ($`T<2000`$ K). However, it appears to be dynamically stable against star-formation.
Although the H i mass is low, $`1.7\times 10^7`$ $`_{}`$ (for a distance of 3.2 Mpc), the velocity differences within the system suggest that there may be an extended dark matter halo if the system is bound. Is HIPASS J1712-64 important from a cosmological point of view? This depends on whether similar objects exist in great numbers, whether they represent the cool, high-density baryonic peaks of low-amplitude overdensities, whether such objects are the building blocks of galaxy and star cluster formation, whether they are simply the left-overs of these processes, and whether they are indeed extragalactic. Future HIPASS observations will cast some light on the frequency of their occurrence, and their distribution on the sky.
We are grateful to Michael Brown for assistance with analysis of the CCD data, using SExtractor, to Emma Ryan for assistance with the ATCA observing, to the staff at the Parkes telescope, and to members of the ZOA observing team for assistance with observations. P. Knezek acknowledges partial support by a grant from NASA administered by the American Astronomical Society.
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# A Dynamical Principle For The Salpeter Equation
## 1. Introduction: Historical Background
The problem of probability interpretation of the well-known Bethe-Salpeter Equation (BSE) has for long been a much discussed issue in the Physics literature, and has led to various attempts \[2-5\] towards its 3D reduction, as well as for intrinsic 3D formulations of various kinds . On the other hand, the issue should be viewed against the theoretical perspective of the BSE as only an $`approximate`$ description (in the so-called ladder approximation) of the equations of motion that follow strictly from the standard QED Lagrangian, viz., the Schwinger-Dyson Equations (SDE) which are an infinite chain of equations connecting successively higher order vertex functions. Thus a conceptual problem like lack of probability interpretation for a $`truncated`$ two-body system need hardly cause undue surprise. Nevertheless, the approximate nature of the BSE has never been a detracting issue from its practical usefulness which has lent considerable credence to the attempts at its 3D reduction designed among other things to restore the probability interpretation to the BS amplitude. Of these attempts, the Salpeter Equation is perhaps the oldest, and has claimed considerable attention in the contemporary literature, from the atomic two-body problem to the QCD context of heavy quarkonia. Another possible motivation for the 3D reduction of the BSE is the observed $`O(3)`$-like spectra of the respective energy levels.
To meet this dual requirement of both the probability interpretation as well as the observed spectra, the ”Instantaneous Approximation” was historically the earliest ansatz to be invoked for a 3D reduction of the full 4D BSE, leading to the Salpeter equation . However, the ansatz (which also suffered from lack of Lorentz Covariance) seemed to beg a solid underlying principle. On the other hand, the extensive use of the Salpeter equation in the contemporary literature demands a firmer theoretical foundation.
### 1.1 The Markov-Yukawa Transversality Principle
In the quest for a theoretical basis for the Instantaneous Approximation, success seems to have come from a rather unexpected quarter, viz., the half-century old Markov-Yukawa Transversality Principle (MYTP) , as was to be discovered by the Dubna Group . Now the MYTP presupposes a dependence of the field variable on both $`x`$ and $`p`$. While this is unacceptable for an elementary particle description (and is probably the reason why the Yukawa non-local field thery \[7b\] was found unattractive), it seems to be ideally suited to a $`composite`$ particle description, wherein the momentum dependence comes from the direction of the total 4-momentum $`P_\mu `$. For a bilocal field $`(z,X)`$, the Transversality condition was shown by Lukierski et al to be equivalent to a ‘gauge principle’ which expresses the redundance of the $`longitudinal`$ component of the relative momentum for the physical interaction between the two constituents. As will be shown in Sect.2, this condition amounts to a covariant 3D support to the input 4-quark Lagrangian, whence follows the 4D BSE with a 3D kernel support governed by Covariant Instantaneity. This Principle was first invoked by the Dubna group to show that the 3D Salpeter equation follows as an exact consequence of the covariant 3D support to the Bethe-Salpeter kernel, with the preferred direction as $`P_\mu `$. This gave a firm theoretical basis to the 3D Salpeter equation.
The other side of the coin, apparently missed by the Dubna Group , concerns the question of whether the information on the 4D content of the original BSE is retrievable, after the 3D reduction. As was to be found soon afterwards , the inbuilt structure of MYTP ensures that the original 4D BSE is exactly recovered by retracing the steps ! This two-way interconnection between the 3D and 4D BSE forms was initially proved for an idealized spinless fermion problem, but, as will be shown in this paper, the logic goes through equally well for spinor fermions, thus facilitating an exact reconstruction of the original 4D BS amplitude in terms of the 3D ingredients of the Salpeter equation. This is, surprisingly enough, a $`new`$ result, considering the fact that this aspect of the Salpeter equation has never seen the light of the day despite its half century old existence. A generalization of the 3D-4D interconnection of Bethe-Salpeter amplitudes under MYTP to the 3-body problem has been given recently .
The paper is organized as follows. In Sect.2, starting with the logic of the MYTP which mandates a covariant 3D support to the kernel of a BSE, we first recapulate the main steps that lead to an exact 3D-4D interconnection between the corresponding BS amplitudes for a ‘spinless’ two-particle system. In Sect.3 we outline a corresponding derivation for the Salpeter equation by recalling the main steps leading from the 4D to the 3D form , and reversing these steps. Sect.4 concludes with some comments on the significance of this results vis-a-vis the Markov-Yukawa Principle, especially the applicability of the ‘Salpeter Vertex fn’ to transition amplitudes as 4D loop integrals.
## 2. MYTP As A Gauge Principle
The logic of MYTP may be traced to Yukawa’s non-local field theory \[7b\], characterized by the field dependence on both coordinate and momentum. As this violates local micro-causality, this concept as a basic theory of elementary particles did not find much favour within the physics community. However this (limited) perspective had to change with the advent of QCD which pushed the status of hadrons from the elementary to a composite level, and gave rise to the concept of bilocal fields . Within such a bilical scenario, the total 4-momentum $`P_\mu `$ of the composite hadron provides a naturally preferred direction which forms the basis for a covariant 3D support to the interaction kernel .
An important feature of bilocal dynamics is the redundance of the relative ‘time’ variable $`x_0`$, ($`x=x_1x_2`$), whose covariant definition is just the longitudinal component of $`x_\mu `$ in the direction of $`P_\mu `$, viz., $`x.PP_\mu /P^2`$. This ‘redundance’ is expressed by the statement that a translation of the relative coordinate $`x_\mu x_\mu ^{}+\xi P_\mu `$ on the bilocal field $`(x,P)`$:
$$(x_\mu ,P_\mu )_\xi (x_\mu ,P_\mu )=(x_\mu +\xi P_\mu ,P_\mu )$$
, which is a sort of ‘gauge transformation’ of the bilocal field , should leave this quantity $`invariant`$. This invariance is just the content of the Markov-Yukawa subsidiary condition which, under an interchange of the relative coordinates and the momenta reads as \[9, 8b\]
$$P_\mu \frac{}{x_\mu }(x_\mu ,P_\mu )=0$$
(1)
where the direction $`P_\mu `$ guarantees an irreducible representation of the Poincare’ group for the bilocal field $``$ . An equation of type has been used in other approaches to bilocal field dynamics (see ref for other references), but this ‘gauge’ interpretation of the subsidiary condition provides a more transparent view of the same condition which we have abbreviated as MYTP above.
Eq.(1) amounts to an effective 3D support to the interaction between the constituents of the bilocal field, which may be alternatively postulated directly for the pairwise BSE kernel $`K`$ by demanding that it be a function of only $`\widehat{q}_\mu =qq.PP_\mu /P^2`$, which implies that $`\widehat{q}.P0`$. In this approach, the propagators are left untouched in their full 4D forms. This is somewhat complementary to the 3D BSE reduction methods \[4-6\] (propagators manipulated but kernel left untouched), so that the resulting equations look rather unfamiliar vis-a-vis 3D BSE’s \[4-6\], but it has the advantage of allowing a $`simultaneous`$ use of both 3D and 4D BSE forms via their interlinkage. Indeed what distinguishes the Covariant Instantaneity Ansatz from the more familiar 3D BSE reductions \[4-6\] is its capacity for a 2-way linkage: an exact 3D BSE reduction, and an equally exact reconstruction of the original 4D BSE form without extra charge . In contrast the more familiar methods of BSE reduction \[4-6\] give at most a one-way connection, viz., a $`4D3D`$ reduction, but not vice versa. This is a plus point for MYTP, and may well have a wider significance than the mere BSE context above, as an effective dynamics for strong interactions.
### 2.1 3D-4D Interconnection: Spinless Particles
To demonstrate the basic 3D-4D interconnection under MYTP , consider a system of two dentical spinless particles, with the BSE
$$i(2\pi )^4\mathrm{\Phi }(q,P)=(\mathrm{\Delta }_1\mathrm{\Delta }_2)^1d^3\widehat{q}^{}M𝑑\sigma ^{}K(\widehat{q},\widehat{q}^{})\mathrm{\Phi }(q^{},P);[\mathrm{\Delta }_{1,2}=m_q+p_{1,2}^2]$$
(2)
where the 3D support to the kernel $`K`$ is implied in its ‘hatted’ structure:
$$\widehat{q}_\mu =q_\mu \sigma P_\mu ;\sigma =q.P/P^2;\widehat{q}.P0.$$
(3)
The relative and total 4-momenta are related by
$$p_1+p_2=P=p_1^{}+p_2^{};2q=p_1p_2;2q^{}=p_i^{}p_2^{}.$$
The 3D wave function $`\varphi (\widehat{q})`$ is defined by
$$\varphi (\widehat{q})=M𝑑\sigma \mathrm{\Phi }(q,P)$$
(4)
When (4) is substituted on the RHS of (2) one gets
$$i(2\pi )^4\mathrm{\Phi }(q,P)=(\mathrm{\Delta }_1\mathrm{\Delta }_2)^1d^3\widehat{q}^{}K(\widehat{q},\widehat{q}^{})\varphi (\widehat{q}^{})$$
(5)
Now integrate both sides of this equation wrt $`\sigma `$ to get an explicit 3D equation
$$(2\pi )^3D(\widehat{q})\varphi (\widehat{q})=d^3\widehat{q}^{}K(\widehat{q},\widehat{q}^{})\varphi (\widehat{q}^{})$$
(6)
where the 3D denominator function is given by
$$2i\pi D^1(\widehat{q})=M𝑑\sigma (\mathrm{\Delta }_1\mathrm{\Delta }_2)^1$$
(7)
A comparison of (5) with (6) via (7) gives the 3D-4D interconnection
$$21\pi \mathrm{\Delta }_1\mathrm{\Delta }_2\mathrm{\Phi }(q,P)=D(\widehat{q})\varphi (\widehat{q})$$
(8)
which directly identifies the RHS as the hadron-quark Vertex Function
$$\mathrm{\Gamma }=D\times \varphi /2i\pi .$$
(9)
## 3. Salpeter Eqn: 3D-4D Interlinkage
Let us now look at the Salpeter equation for the relativistic hydrogen atom problem, which in the notation of the original paper reads as
$$i\pi ^2F(q_\mu )\psi (q)=\alpha d^4k𝐤^2\psi (q+k)$$
(10)
A comparison of this equation with eq.(6) shows a precise correspondence, except for certain technicalities arisin from its fermionic content. Indeed it stems from an equation of the form (2), where the 3D kernel support is due to the (non-covariant) instantaneous (adiabatic) assumption , manifesting from its dependence on the 3-vector $`𝐤`$, while the quantity $`F(q_\mu )`$ plays just the role of the product of the two 4D propagators $`\mathrm{\Delta }_1`$ and $`\mathrm{\Delta }_2`$ in (2):
$$F(q)=(\mu _1EH_1(𝐪)+ϵ)(\mu _2EH_1(𝐪)ϵ)$$
(11)
with the time-like components identified as the $`ϵ`$ terms ! Next, define the 3D wave function $`\varphi (𝐪)`$ by
$$\varphi (𝐪)=𝑑ϵ\psi (𝐪,ϵ)$$
(12)
which is the counterpart of (4), and use this result to integrate both sides of (10) wrt $`ϵ`$, after dividing by $`F(q)`$, so as to get the 3D Salpeter equation
$$[EH_1(𝐪)H_2(𝐪)]\varphi _{\pm \pm }=\pm \mathrm{\Lambda }_\pm ^{(1)}\mathrm{\Lambda }_\pm ^{(2)}(2i\pi \mathrm{\Gamma }(𝐪)=(4i\alpha )d^3k𝐤^_2\varphi (𝐪+𝐤)$$
(13)
where the $`\pm `$ components are associated with the energy projection operators $`\mathrm{\Lambda }`$ which however do not involve the time-like $`ϵ`$.
The new aspect, on the other hand, is the 3D-4D interconnection which is obtained by substituting the second part of eq.(13) on the RHS of (10), after making use of (12):
$$F(q)\psi (q)=\mathrm{\Gamma }(𝐪)$$
(14)
where $`\mathrm{\Gamma }(𝐪)`$ is the 3D BS vertex function. It is the precise fermionic counterpart of the scalar eq.(9), since the $`F(q)`$ function is the product of the two 4D propagators. The form (14) is not formally covariant, but this is a mere technicality which can be remedied by standard methods; see e.g., ref .
The more interesting thing about this demonstration is the exciting prospect of using the reconstructed 4D ‘Salpeter vertex function’ (14) as a basic ingredient for the calculation of various types of transition amplitudes as 4D loop integrals by standard Feynman techniques without having to face the usual problems of probabilty interpretation and/or spectroscopy, both of which are now subsumed in the 3D equation (13). This gives a sort of ‘two-tier’ description, the 3D form (13) just right for spectroscopy, energy levels, etc, while the 4D form (14) provides the proper vehicle for 4D loop integrals. It is only the first (3D) part of the Salpeter Equation that has so far been evidenced in the contemporary literature, but the second (4D) aspect is entirely $`new`$.
## 4. Retrospect And Summary
In retrospect, we have attempted to project an aspect of the well-known Salpeter equation , which had remained obscured from view for decades, viz., a theoretical basis for its underlying Instantaneous Approximation, offered by the Markov-Yukawa Transversality Principle : An in-built MYTP in the 4D BSE structure leads to an exact 3D reduction which, as first shown by the Dubna group , is a covariant generalization of the Salpeter Equation . The (new) complementary aspect of MYTP is its in-built capacity to $`reconstruct`$ with equal ease , the 4D vertex function, (9) or (14), in terms of 3D ingredients, which allows access to transition amplitues of diverse types as 4D loop integrals. This offers a two-tier description for the Salpeter Equation, analogously to the quark-level hadronic BSE problem that has been under study for several years, with its 3D form providing access to spectroscopy , and the 4D form offering applications to processes like e.m. form factors , within a single framework. This dual feature distinguishes MYTP from most other 3D approaches to strong interaction dynamics \[4-6\] which give at most a one-way connection (4D to 3D). This remarkable property of 3D-4D interlinkage enjoyed by the Salpeter equation , by virtue of its compliance with MYTP, should hopefully offer new incentives for its (second stage) applications to 4D loop integrals in a covariant manner.
The 3D-4D interlinkage offered by MYTP is also generalizable to a 3-body BSE with pairwise kernels under covariant 3D support . A second type of generalization of MYTP is to the covariant null-plane which facilitates trouble-free evaluation of form factors with triangle loops.
To summarise, the instantaneous approximation which characterizes the Salpeter equation, comes as a mere consequence of the Markov-Yukawa Transversality Principle which by its very definition gives a precise 3D support to the BSE kernel. Secondly, MYTP allows reconstruction of the 4D Salpeter amplitude in terms of 3D ingredients, a property which had remained obscured from view so far. Thus the Salpeter equation is amenable to a two-tier formalism, the 3D form for spectroscopy, and the reconstructed 4D vertex function for 4D loop integrals. Finally, the non-covariance of the Salpeter equation is a mere technicality which is easily remedied by standard techniques .
A preliminary version of this work was reported at the XVI Few-Body Conf at Taiwan.
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# A class of perfect-fluid cosmologies with polarised Gowdy symmetry and a Kasner-like singularity
## 1 Introduction
In this paper we use the Fuchsian algorithm to construct a family of perfect-fluid cosmologies with polarised Gowdy symmetry and with Kasner like asymptotics at early times. This family depends on the maximum number of free functions for spacetimes within the symmetry class. The technique used is to perturb exact Bianchi I solutions in one space-direction and to solve a Fuchsian system of equations for the perturbation. The results obtained are a generalisation of those in in the sense that the symmetry requirement has been relaxed a little, but here we require the free data for the field equations to be analytic rather than merely $`C^{\mathrm{}}`$.
The main result is stated as Theorem 6.1 at the end of the paper.
## 2 Exact solutions
The general Bianchi I solution of the Einstein-perfect fluid equations, as given in is
$$ds^2=B^{2(\gamma 1)}d\tau ^2+\tau ^{2p_1}B^{2q_1}dx^2+\tau ^{2p_2}B^{2q_2}dy^2+\tau ^{2p_3}B^{2q_3}dz^2$$
(1)
where
$$B^{2\gamma }=\alpha +m^2\tau ^{2\gamma }\alpha 0,m>0$$
(2)
$$p_1+p_2+p_3=1,p_1^2+p_2^2+p_3^2=1,q_i=\frac{2}{3}p_i$$
(3)
(The case $`\alpha =0`$ is an FRW model and is excluded in what follows.)
The density of the fluid is given by
$$\mu =\frac{4m^2}{3\tau ^\gamma B^\gamma }$$
(4)
When we come to write down Einstein’s equations for a metric with polarised Gowdy symmetry we will use conformal coordinates, i.e coordinates for which the metric takes the form
$$ds^2=e^{2A(t,x)}(dt^2+dx^2)+R(t,x)e^{W(t,x)}dy^2+R(t,x)e^{W(t,x)}dz^2$$
(5)
The model solution (1) may be written in conformal coordinates by making the transformation
$$\tau t=_0^\tau B^{p_1+\gamma \frac{5}{3}}(s)s^{p_1}𝑑s$$
(6)
Then the metric (1) takes the form (5) with
$$e^A=\tau ^{p_1}B^{\frac{2}{3}p_1},e^W=\tau ^{p_2p_3}B^{p_3p_2},R=\tau ^{1p_1}B^{(p_1+\frac{1}{3})}$$
(7)
and $`\tau `$ given implicitly as function of $`t`$ by the relation (6).
## 3 The Einstein-perfect fluid equations
We will assume that spacetime is filled with a polytropic perfect fluid, so that the stress tensor takes the form
$$T^{\alpha \beta }=(\rho +P)u^\alpha u^\beta +Pg^{\alpha \beta }$$
(8)
with $`u^\alpha u_\alpha =1`$ and $`P=(\gamma 1)\rho ,1<\gamma <2`$.
With this stress tensor the Einstein evolution equations for the metric (5) are
$$R_{tt}R_{xx}=Re^{2A}\rho (2\gamma )$$
(9)
$$W_{tt}W_{xx}+R^1(R_tW_tR_xW_x)=0$$
(10)
$$A_{tt}A_{xx}\frac{1}{4}R^2(R_t^2R_x^2)+\frac{1}{4}(W_t^2W_x^2)=\frac{1}{2}\gamma e^{2A}\rho $$
(11)
while the constraints read
$$R^1R_{xx}\frac{1}{4}R^2(R_t^2+R_x^2)R^1R_tA_tR^1R_xA_x+\frac{1}{4}(W_t^2+W_x^2)$$
$$=e^{2A}\rho (1+\gamma (v^1)^2)$$
(12)
$$R^1R_{tx}+R^1(R_tA_x+A_tR_x)+\frac{1}{2}R^2R_tR_x\frac{1}{2}W_tW_x=e^{2A}\gamma \rho v^0v^1$$
(13)
where $`v^i=e^Au^i`$ and hence $`(v^0)^2(v^1)^2=1`$.
The Euler equations $`^\alpha T_{\alpha \beta }=0`$ take the following explicit form
$$\gamma v^0(v^0\rho _t+v^1\rho _x)+(1\gamma )\rho _t$$
$$+\gamma \rho \{v^0(2v_t^0+A_xv^1)+v^1(v_x^1v^1A_x+A_xv^0+A_tv^1)\}$$
$$+\gamma \rho v^0\{v^0A_t+v_x^1+v^1A_x+R^1(v^0R_t+v^1R_x)\}=0$$
(14)
and
$$\gamma v^1(v^0\rho _t+v^1\rho _x)+(1\gamma )\rho _x+\gamma \rho \{v^0(v_t^1+v^0A_x)\}$$
$$+\gamma \rho \{v^1(v_t^0+2v_x^1+v^0(2A_t+R^1R_t)+v^1(A_x+R^1R_x)\}=0$$
(15)
## 4 Inhomogeneous perturbations
We now seek solutions $`(\rho ,v^1,A,R,W)`$ of (9)-(15) which depend on the maximum number of free analytic functions, namely four, and which approach the model forms (1)-(4) for each $`x`$ as $`t0`$.
To be specific we make the following ansatz
$$A=p_1\mathrm{log}\tau +\left(\frac{2}{3}p_1\right)\mathrm{log}B+t\stackrel{~}{A}$$
(16)
$$R=\tau ^{1p_1}B^{(p_1+\frac{1}{3})}(1+\stackrel{~}{R})$$
(17)
$$W=c(x)+(p_2p_3)\mathrm{log}\tau +(p_3p_2)\mathrm{log}B+t\stackrel{~}{W}$$
(18)
$$\mathrm{log}\rho =\mathrm{log}\mu +t^{(\gamma 1p_1)/(1p_1)+ϵ}\varphi $$
(19)
$$v^1=\tau ^{\gamma 1p_1}(G(x)+t^ϵ\stackrel{~}{\psi })$$
(20)
where
$$B^{2\gamma }=\alpha (x)+m^2(x)\tau ^{2\gamma }$$
(21)
$$G(x)=\frac{3}{4m^2\gamma }\alpha ^{(p_1+\frac{1}{3})/(2\gamma )}\times $$
(22)
$$\left\{\frac{1}{2}(p_2p_3)c_x+\frac{\alpha _x}{\alpha }\frac{(3\gamma (1p_1)+2(3p_11))}{3(2\gamma )}+(p_1)_x\left(\frac{2(1p_1)}{(2\gamma )}\mathrm{log}\alpha +1\right)\right\}$$
(23)
and $`\tau (t,x)`$ is determined by the relation (6).
$`\alpha (x)`$ and $`m(x)`$ are chosen as strictly positive analytic functions, $`c(x)`$ is an analytic function and the Kasner exponent $`p_1(x)`$ is subject to the restriction $`p_1(x)<\gamma 1k`$ for some arbitrarily small constant $`k`$. The analytic function $`ϵ(x)`$ is chosen to satisfy
$$0<ϵ(x)<\text{min}(\frac{2\gamma }{1p_1},\frac{\gamma 1p_1}{1p_1})$$
(24)
We are looking for analytic solutions of the Einstein-perfect fluid equations for which $`\stackrel{~}{A},\stackrel{~}{R},\stackrel{~}{W},\varphi `$ and $`\stackrel{~}{\psi }`$ tend to zero as $`t`$ tends to zero. The choice of $`G(x)`$ ensures that the momentum constraint is satisfied at $`t=0`$.
It is convenient to write the field equations for $`\stackrel{~}{A},\stackrel{~}{R},\stackrel{~}{W},\varphi `$ and $`\stackrel{~}{\psi }`$ in first order form. To do this we introduce the following new variables:
$$U=\stackrel{~}{R}_t,Q=(t\stackrel{~}{A})_t,X=\stackrel{~}{R}_x,Y=t\stackrel{~}{A}_x,V=(t\stackrel{~}{W})_t,Z=t\stackrel{~}{W}_x,S=t^1\stackrel{~}{R}$$
In terms of these variables the evolution equations (9)-(11) take the form
$$tX_t=tU_x$$
(25)
$$tY_t=tQ_x$$
(26)
$$t\stackrel{~}{R}_t=tU$$
(27)
$$t\stackrel{~}{A}_t+\stackrel{~}{A}Q=0$$
(28)
$$tZ_t=tV_x$$
(29)
$$t\stackrel{~}{W}_t+\stackrel{~}{W}V=0$$
(30)
$$tS_t+SU=0$$
(31)
$$tU_t+2U=tX_x+2U\left(1\left((1p_1)B^{2\gamma }+m^2\left(p_1+\frac{1}{3}\right)\tau ^{2\gamma }\right)B^{(p_1+\frac{1}{3})}t\tau ^{p_11}\right)$$
$$+t\tau ^{p_11}B^{(p_1+\frac{1}{3})}\{2X(\tau ^{1p_1}B^{p_1+\frac{1}{3}})_x+(1+\stackrel{~}{R})(\tau ^{1p_1}B^{p_1+\frac{1}{3}})_{xx}\}$$
$$+\frac{4m^2}{3}(2\gamma )(1+tS)t\tau ^{2p_1\gamma }B^{\frac{4}{3}2p_1\gamma }(\text{exp}(2t\stackrel{~}{A}+t^{ϵ+(\gamma 1p_1)(1p_1)^1}\varphi )1)$$
(32)
$$tQ_t\frac{1}{2}U+\frac{1}{2}\frac{(p_2p_3)}{(1p_1)}V=tY_x+t\{(p_1\mathrm{log}\tau )_x+(\mathrm{log}B^{\frac{2}{3}p_1})_{xx}\}$$
$$+\frac{1}{4}(1+\stackrel{~}{R})^2tU^2\frac{1}{2}U\left(1t\tau ^{p_11}B^{(p_1+\frac{1}{3})}(1+tS)^1\left((1p_1)B^{2\gamma }+m^2\left(p_1+\frac{1}{3}\right)\tau ^{2\gamma }\right)\right)$$
$$+\frac{1}{4}t\tau ^{2(p_11)}B^{2(p_1+\frac{1}{3})}\left(\tau ^{1p_1}B^{p_1+\frac{1}{3}}X+(1+\stackrel{~}{R})(\tau ^{1p_1}B^{p_1+\frac{1}{3}})_x\right)^2$$
$$\frac{1}{4}tV^2+\frac{1}{2}V(p_2p_3)\left((1p_1)^1t\tau ^{p_11}B^{\frac{5}{3}+p_1\gamma }+m^2t\tau ^{1+p_1\gamma }B^{\frac{1}{3}p_1}\right)$$
$$+\frac{1}{4}t\left(c_x+((p_2p_3)\mathrm{log}\tau )_x+((p_3p_2)\mathrm{log}B)_x+Z\right)^2$$
$$\frac{2m^2}{3}\gamma B^{\frac{4}{3}2p_1\gamma }t\tau ^{2p_1\gamma }\left(\text{exp}(2t\stackrel{~}{A}+t^{ϵ+(\gamma 1p_1)(1p_1)^1}\varphi )1\right)$$
(33)
$$tV_t+V+\frac{(p_2p_3)}{(1p_1)}U=tZ_x+t\{c_{xx}+((p_2p_3)\mathrm{log}\tau )_{xx}+((p_3p_2)\mathrm{log}B)_{xx}\}$$
$$+\frac{t\tau ^{p_11}B^{p_1\frac{1}{3}}}{(1+\stackrel{~}{R})}(c_x+((p_2p_3)\mathrm{log}\tau )_x+((p_3p_2)\mathrm{log}B)_x+Z)(\tau ^{1p_1}B^{p_1+\frac{1}{3}}X+(1+\stackrel{~}{R})(\tau ^{1p_1}B^{p_1+\frac{1}{3}})_x)$$
$$+U\left(\frac{(p_2p_3)}{(1p_1)}(1+tS)^1((p_2p_3)t\tau ^{p_11}B^{\frac{5}{3}p_1\gamma }+m^2(p_3p_2)t\tau ^{1+p_1\gamma }B^{\frac{1}{3}p_1}+tV)\right)$$
$$+V(1B^{p_1\frac{1}{3}}t\tau ^{p_11}((1p_1)B^{2\gamma }+m^2\tau ^{2\gamma }(p_1+1/3)))$$
(34)
The Euler equations , after some rearrangement, may be written
$$\left(\frac{1+2\tau ^{2\beta }\psi ^2}{v^0}\frac{2\gamma v^0\tau ^{2\beta }\psi ^2}{1+\gamma \tau ^{2\beta }\psi ^2}\right)(t\stackrel{~}{\psi }_t+ϵ\stackrel{~}{\psi })=$$
$$t^{1ϵ}\psi \left(\frac{\gamma v^0(v^0+\tau ^\beta )}{1+\gamma \tau ^{2\beta }\psi ^2}2\right)(\tau ^\beta \psi _x+\psi (\tau ^\beta )_x)$$
$$+t^{1ϵ}\tau ^\beta \left(\frac{\gamma 1}{\gamma }+\frac{\gamma (v^0)^2\tau ^{2\beta }\psi ^2}{1+\gamma \tau ^{2\beta }\psi ^2}\right)((m^2)_x\gamma (\tau _x+B_x)+t^{ϵ+\beta (1p_1)^1}(\varphi _x+\varphi (\beta (1p_1)^1)_x\mathrm{log}t))$$
$$+t^{1ϵ}v^0\psi \left(\frac{\gamma ((v^0)^2+\tau ^{2\beta }\psi ^2)}{1+\gamma \tau ^{2\beta }\psi ^2}2\right)((2/3p_1)B^{\frac{1}{3}p_1}m^2\tau ^\beta +Q)$$
$$t^{1ϵ}\tau ^\beta \left(1+2\tau ^{2\beta }\psi ^2\frac{\gamma \tau ^\beta \psi v^0}{1+\gamma \tau ^{2\beta }\psi ^2}(3v^0\tau ^\beta \psi \tau ^{2\beta }\psi ^2)\right)((p_1\mathrm{log}\tau )_x+(\mathrm{log}B^{\frac{2}{3}p_1})_x+Y)$$
$$+t^{1ϵ}\psi \left(\frac{\gamma (v^0)^3}{1+\gamma \tau ^{2\beta }\psi ^2}v^0\right)((1+\stackrel{~}{R})^1U+m^2(p_1+1/3)B^{p_1\frac{1}{3}}\tau ^\beta )$$
$$+t^{1ϵ}\tau ^{\gamma 2}\psi ^2\left(\frac{\gamma (v^0)^2}{1+\gamma \tau ^{2\beta }\psi ^2}1\right)B^{p_1\frac{1}{3}}(1+\stackrel{~}{R})^1(\tau ^{1p_1}B^{p_1+\frac{1}{3}}X+(1+\stackrel{~}{R})(\tau ^{1p_1}B^{p_1+\frac{1}{3}})_x)$$
$$+t^{1ϵ}\tau ^{p_11}\psi B^{\frac{5}{3}p_1\gamma }\{\beta (1+\frac{2\gamma v^0\tau ^{2\beta }\psi ^2}{1+\gamma \tau ^{2\beta }\psi ^2}\frac{1+2\tau ^{2\beta }\psi ^2}{v^0})$$
$$+p_1(\gamma (\frac{v^0((v^0)^2+\tau ^{2\beta }\psi ^2)}{1+\gamma \tau ^{2\beta }\psi ^2}1)+2(1v^0))+(1p_1)(\gamma (\frac{(v^0)^2}{1+\gamma \tau ^{2\beta }\psi ^2}1)+1v^0)\}$$
(35)
and
$$\left(1+\gamma \tau ^{2\beta }\psi ^2\frac{2\gamma (v^0)^2\tau ^{2\beta }\psi ^2}{1+2\tau ^{2\beta }\psi ^2}\right)t\varphi _t+\left(ϵ+\frac{\beta }{1p_1}\right)\varphi =$$
$$=\gamma \{(ϵ+\beta (1p_1)^1)\varphi (\frac{1}{\gamma }+\frac{2(v^0)^2\tau ^{2\beta }\psi ^2}{1+2\tau ^{2\beta }\psi ^2}\frac{1+\gamma \tau ^{2\beta }\psi ^2}{\gamma })$$
$$+\gamma t^{1+\frac{\beta }{p_11}ϵ}\tau ^{p_11}\tau ^{2\beta }\psi ^2\left(1\frac{2(v^0)^2}{1+2\tau ^{2\beta }\psi ^2}\right)(B^{\frac{5}{3}p_1\gamma }+m^2\tau ^{2\gamma }B^{\frac{1}{3}p_1})$$
$$(t\varphi _x+t\varphi \mathrm{log}t(ϵ+\beta (1p_1)^1)_x+t^{1+\frac{\beta }{p_11}ϵ}((m^2)_x\gamma (\tau _x+B_x)))v^0\tau ^\beta \psi $$
$$\times \left(1+\frac{2}{1+2\tau ^{2\beta }\psi ^2}\left(\frac{\gamma 1}{\gamma }\tau ^{2\beta }\psi ^2\right)\right)$$
$$+t^{1+\frac{\beta }{p_11}ϵ}\left(\frac{4\tau ^{2\beta }\psi ^2v^0}{1+2\tau ^{2\beta }\psi ^2}v^0\tau ^\beta \psi \right)((\tau ^\beta G)_x+\tau ^\beta t^ϵ\stackrel{~}{\psi }_x+(\tau ^\beta t^ϵ)_x\stackrel{~}{\psi })$$
$$+t^{1+\frac{\beta }{p_11}ϵ}\left(\frac{4(v^0)^2\tau ^{2\beta }\psi ^2}{1+2\tau ^{2\beta }\psi ^2}12\tau ^{2\beta }\psi ^2\right)Q$$
$$+t^{1+\frac{\beta }{p_11}ϵ}\tau ^{2\beta }\psi ^2\left(\frac{4(v^0)^2}{1+2\tau ^{2\beta }\psi ^2}2\right)(p_1\tau ^{p_11}B^{\frac{5}{3}p_1\gamma }+((2/3)p_1)B^{\frac{1}{3}p_1}m^2\tau ^\beta )$$
$$+t^{1+\frac{\beta }{p_11}ϵ}\left(\tau ^\beta \psi (\tau ^\beta \psi v^0)((p_1\mathrm{log}\tau )_x+(\mathrm{log}B^{\frac{2}{3}p_1})_x+Y)+\frac{(v^0\tau ^\beta \psi (v^0)^2)}{1+\stackrel{~}{R}}U\right)$$
$$+t^{1+\frac{\beta }{p_11}ϵ}\tau ^{2\beta }\psi ^2\left(\frac{2(v^0)^2}{1+2\tau ^{2\beta }\psi ^2}1\right)\tau ^{p_11}B^{p_1\frac{1}{3}}((1p_1)B^{2\gamma }+m^2(p_1+1/3)\tau ^{2\gamma })$$
$$+t^{1+\frac{\beta }{p_11}ϵ}\tau ^{\gamma 2}\psi (\frac{2v^0\tau ^{2\beta }\psi ^2}{1+2\tau ^{2\beta }\psi ^2}v^0)\frac{B^{p_1\frac{1}{3}}}{1+\stackrel{~}{R}}(\tau ^{1p_1}B^{p_1+\frac{1}{3}}X+(1+\stackrel{~}{R})(\tau ^{1p_1}B^{p_1+\frac{1}{3}})_x)\}$$
(36)
## 5 Existence and uniqueness of solutions
A careful inspection of the field equations (25)-(36) shows that they may be written in the form
$$t_tu+N(x)u=t^\delta H(t,x,u,u_x)$$
(37)
where $`u`$ stands for $`(U,\stackrel{~}{A},Q,X,Y,\stackrel{~}{R},V,Z,S,\varphi ,\stackrel{~}{\psi })`$$`\delta `$ is a strictly positive constant and $`H`$ is continuous in $`t`$ and analytic in its other arguments. The matrix $`N(x)`$ has positive eigenvalues.
It follows by that (37) has a unique analytic solution $`u`$ with $`u(0)=0`$
## 6 The constraints
Define constraint quantities $`C_1,C_0`$ by
$$C_0=R^1R_{xx}\frac{1}{4}R^2(R_t^2+R_x^2)R^1R_tA_tR^1R_xA_x+\frac{1}{4}(W_t^2+W_x^2)$$
$$+e^{2A}\rho (1+\gamma (v^1)^2)$$
(38)
$$C_1=R^1R_{tx}+R^1(R_tA_x+A_tR_x)+\frac{1}{2}R^2R_tR_x\frac{1}{2}W_tW_x+e^{2A}\gamma \rho v^0v^1$$
(39)
If the evolution equations (25)-(36) are satisfied then a calculation shows that the following hold
$$_tC_0=_xC_1\frac{R_t}{R}C_0\frac{R_x}{R}C_1$$
(40)
$$_tC_1=_xC_0\frac{R_t}{R}C_1\frac{R_t}{R}C_0$$
(41)
One also calculates that the quantities $`\stackrel{~}{C}_0=tC_0,\stackrel{~}{C}_1=tC_1`$ tend to zero as $`t`$ tends to zero. These quantities satisfy the following
$$t_t\stackrel{~}{C}_0+\left(\frac{tR_t}{R}1\right)\stackrel{~}{C}_0=t_x\stackrel{~}{C}_1\frac{tR_x}{R}\stackrel{~}{C}_1$$
(42)
$$t_t\stackrel{~}{C}_1+\left(\frac{tR_t}{R}1\right)\stackrel{~}{C}_1=t_x\stackrel{~}{C}_0\frac{tR_x}{R}\stackrel{~}{C}_0$$
(43)
Now $`R_x/R`$ is $`O(1)`$ and $`tR_t/R=1+O(t^\delta )`$ for some $`\delta >0`$. It thus follows that $`\stackrel{~}{C}_0`$ and $`\stackrel{~}{C}_1`$ are identically zero and thus the constraints are satisfied.
Summarising the results of sections (4)-(6) we have proved the following
Theorem 6.1 Given two strictly positive analytic functions $`\alpha (x),m(x)`$, an analytic function $`c(x)`$ and an analytic function $`p_1(x)`$ satisfying
$$\frac{1}{3}p_1(x)<\gamma 1k$$
(44)
for some small $`k>0`$ , there exists a unique solution $`(g_{\mu \nu },\rho ,u^\alpha )`$ of the Einstein equations coupled to a $`\gamma `$-law perfect fluid on $`^3\times (0,T)`$ satisfying
$$ds^2=e^{2A}(t,x)(dt^2+dx^2)+R(t,x)e^{W(t,x)}(dy)^2+R(t,x)e^{W(t,x)}(dz)^2$$
$$A=p_1\mathrm{log}t+\mathrm{log}B^{\frac{2}{3}p_1}+O(t)$$
$$R=\tau ^{1p_1}B^{p_1+\frac{1}{3}}(1+O(t))$$
$$W=c(x)+(p_2p_3)\mathrm{log}t+(p_3p_2)\mathrm{log}B+O(t)$$
$$\mathrm{log}\rho =\mathrm{log}\frac{4m^2}{3\tau ^\gamma B^\gamma }+O(t^{(\gamma 1p_1)(1p_1)^1})$$
$$v^1=O(t^{(\gamma 1p_1)(1p_1)^1})$$
where $`\tau `$ is implicitly given by
$$t=_0^\tau B^{p_1+\gamma \frac{5}{3}}(s)s^{p_1}𝑑s.$$
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# Quantum Time-Frequency Transforms
## 1 Introduction
The Fourier transform is an operator that expresses a time-dependent signal as a sum (or integral) of periodic signals. In other words the Fourier transform changes a function of time $`s(t)`$ into a function of frequency $`S(\omega )`$. If a signal is a function of time it said to be in the “time domain” and if it is a function of frequency it is said to be in the “frequency domain”. For signals whose spectrum is changing in time, i.e. nonstationary signals, sometimes the best description is a mixture of the time and frequency components. Signal representations which mix the time and frequency domains are called, naturally enough, “time-frequency representations” and are often used to describe time-varying signals for which the pure frequency or Fourier representation is inadequate ,. A familiar example of a time-frequency representation is a musical score, which describes when (time) certain notes (frequency) are to be played.
Formally speaking for our present purposes, a quantum signal is simply a quantum state $`|\psi `$ where the Hilbert space is the group algebra $`[G]`$ of a finite abelian group $`G`$. The Quantum Fourier Transform (QFT) is central to the important quantum algorithms for factoring and discrete logarithm. Mathematically speaking, the Quantum Fourier Transform is a linear operator on the Hilbert Space $`[G]`$ which is a change of basis from the basis of group elements $`\{|g_1,\mathrm{}.,|g_{|G|}\}`$ to the basis of characters of $`G`$, $`\{|\chi _1,|\chi _2,\mathrm{},|\chi _{|G|}\}`$.
We present efficient algorithms for quantum versions of the Zak and Weyl-Heisenberg transforms. Both these time-frequency tranforms can be seen as generalizations of Fourier transforms and the quantum algorithms make heavy use of the Quantum Fourier Transform. We follow the theory and notation of and recommend this book as background to this material.
## 2 Zak Transforms
### 2.1 Background
Let $`A`$ be a finite, abelian group, $`A^{}`$ the group of characters of $`A`$ (note: in this paper does not mean conjugation), $`BA`$ a subgroup of $`A`$, $`B_{}=\{a^{}A^{}:a^{}(b)=1,bB\}`$ the dual to $`B`$, and $`fC[A]`$, the group algebra of $`A`$. Define
$$Z(B)fC[A\times A^{}]$$
by the formula
$$Z(B)f(a,a^{})=\underset{bB}{}f(a+b)\overline{a^{}(b)}.$$
$`F=Z(B)f`$ is called the Zak transform of $`f`$ over $`B`$. A simple calculation shows that $`F(a+b,a^{}+b_{})=a^{}(b)F(a,a^{})`$ where $`bB`$ and $`b_{}B_{}`$. Therefore $`F`$ is determined by its values on a set of coset representatives of $`B\times B_{}`$ in $`A\times A^{}`$ and thus conceptually we may think of $`F`$ as a function on $`T`$ where $`T`$ is a set of coset representatives. Since
$$\frac{|A\times A^{}|}{|B\times B_{}|}=\frac{|A|^2}{|B||B_{}|}=|A|$$
we have the same number of degrees of freedom with which we started. Notice that if $`B`$ contains only the identity, i.e. is the trivial subgroup, then $`Z(B)f(a,a^{})=f(a)`$ and is basically the identity map. Also notice that if $`B=A`$ then $`Z(A)f(0,a^{})=a^{}|f`$ and therefore $`Z(A)f`$ is basically the Fourier transform of $`f`$. So the Zak transform mediates between the time domain and frequency domain depending on the subgroup $`B`$.
Consider the function $`f(a_0)=\delta (xa_0)`$ which is $`1`$ on $`a_0`$ and $`0`$ otherwise. Applying the above formula for the Zak transform yields $`F(a,a^{})=a^{}(aa_0)`$ for $`aa_0+B,a^{}A^{}`$ and $`0`$ otherwise. But since $`F`$ is determined by its values on a set of coset representatives of $`B\times B_{}`$ in $`A\times A^{}`$ let us introduce such a set of representatives $`T=T_1\times T_2=\{(x_i,a_j^{})\}`$ where $`T_1=\{x_i\}`$ is a set of coset representatives of $`B`$ in $`A`$ and $`T_2=\{a_j^{}\}`$ is a set of coset representatives of $`B_{}`$ in $`A^{}`$. Bearing in mind the above transformation of a delta function, we now offer our definition of the Quantum Zak Transform (QZT) (with respect to $`T`$) by
$$|a\frac{1}{\sqrt{|B|}}\underset{a_j^{}T_2}{}a_j^{}(x_aa)|x_a|a_j^{}.$$
where $`x_aT_1`$ is the coset representative of $`a`$. Now notice that $`x_aaB.`$ Therefore $`a_j^{}`$ is restricted to $`B`$ and therefore can be considered to be a character of $`B`$, i.e. an element of $`B^{}`$, and this restriction is independent of the choice of coset representative, i.e. it is natural or canonical. Therefore an equivalent formulation of the QZT is given by
$$|a\frac{1}{\sqrt{|B|}}\underset{b^{}B^{}}{}b^{}(x_aa)|x_a|b^{}.$$
The only difference in these two formulations is in the interpretation of the observed content of the second register.
### 2.2 The Quantum Algorithm
We now show that the QZT is efficiently implementable. Define $`P(B)`$ to be the transform
$$P(B)|a=|x_a|x_aa$$
which decomposes $`a`$ into its coset representative and the corresponding element of $`B`$. $`P`$ is clearly unitary and efficiently implementable. After applying $`P(B)`$ we apply the Quantum Fourier Transform (over the group $`B`$, denoted $`F_B`$) to the second register. This results in the state
$$\frac{1}{\sqrt{|B|}}\underset{b^{}B^{}}{}b^{}(x_aa)|x_a|b^{}.$$
Therefore the QZT is simply $`Z(B)=(IF_B)P(B).`$
## 3 Weyl-Heisenberg Transforms
### 3.1 Background
Define $`g_{(x,x^{})}(a)=g(ax)x^{}(a)`$ to be the time-frequency translate of $`g`$ by $`(x,x^{})`$ where $`gC[A]`$. We will use time-frequency translates to form orthonormal bases so we also require $`|g|=1`$. Let $`\mathrm{\Delta }=B\times B_{}`$ and $`(g,\mathrm{\Delta })=\{g_{(x,x^{})}:(b,b_{})\mathrm{\Delta }\}.`$ We call $`(g,\mathrm{\Delta })`$ a W-H system over $`\mathrm{\Delta }`$ with window $`g`$. A basic result (, Theorem 12.1 corrected version) is that $`(g,\mathrm{\Delta })`$ is an orthonormal basis of $`C[A]`$ if and only if for all $`(a,a^{})A\times A^{}`$ we have $`|G(a,a^{})|=\sqrt{\frac{|B|}{|A|}}`$ where the Zak tranform is taken over $`B`$. Because von Neumann measurements must be unitary we will restrict our attention to window functions $`g`$ which satisfy this constraint. Utilizing POVMs one could consider implementing nonorthonormal W-H systems but we will not address this in this note. This orthogonality constraint together with the earlier observation that $`G`$ is determined by its values on a set of coset representatives of $`B\times B_{}`$ in $`A\times A^{}`$ implies that orthonormal W-H systems are in bijective correspondence with the set of all $`|A|`$-tuples of complex numbers with modulus $`\sqrt{\frac{|B|}{|A|}}`$. In this note we will restrict the W-H systems under consideration by assuming that for each $`(a,a^{})A\times A^{}`$ the phase of $`G(a,a^{})`$ is a rational fraction of $`2\pi `$ which we can compute in polynomial time. Whether or not this last assumption is excessively restrictive would depend on the intended application. Notice that if $`g`$ is the constant function $`g=\frac{1}{\sqrt{|A|}}`$ and $`\mathrm{\Delta }=\{0\}\times A^{}`$ then $`(g,\mathrm{\Delta })`$ is the (normalized) Fourier basis, $`G(a,a^{})=\frac{1}{\sqrt{|A|}}`$ and this restriction holds trivially.
We define the Quantum Weyl-Heisenberg Transform (QWHT) by
$$|\psi \underset{(b,b_{})\mathrm{\Delta }}{}\psi |g_{(b,b_{})}|b,b_{}.$$
In other words, the QWHT expresses $`|\psi `$ in the orthonormal basis of time-frequency translates of the window function.
### 3.2 The Quantum Algorithm
Let
$$f=\underset{(b,b_{})\mathrm{\Delta }}{}\alpha (b,b_{})g_{(b,b_{})}$$
i.e. the $`\alpha `$’s are the coeffients of the WH-expansion of $`f`$. Define
$$P(a,a^{})=\underset{(b,b_{})\mathrm{\Delta }}{}\alpha (b,b_{})b_{}(a)\overline{a^{}(b)}.$$
Notice that $`P`$ is $`\mathrm{\Delta }`$-periodic and that the $`\alpha `$’s are, by definition, the Fourier coefficents (over $`A\times A^{}`$) of $`P`$. A fundamental result (, Theorem 7.5) states that $`F=GP`$. This result suggests an algorithm for computing the WH-coefficients of $`f`$, namely compute the Fourier coefficients of $`P=\frac{F}{G}.`$
Define $`\mathrm{\Phi }(g)`$ to be the unitary transformation which acts on the Hilbert space $`C[T]`$ (recall $`T`$ is the set of coset representatives of $`B\times B_{}`$ in $`A\times A^{}`$) by
$$|x_i|a_j^{}\frac{1}{G(x_i,a_j^{})}|x_i|a_j^{}.$$
Since the phase of $`G(x_i,a_j^{})`$ is, by assumption, a rational fraction of $`2\pi `$ computable in polynomial time we may efficiently implement $`\mathrm{\Phi }(g)`$ by the phase kickback technique described in . Finally in order to complete our description of the algorithm, we must assume that we are given an explicit isomorphism between $`A`$ and $`A^{}`$. These groups are isomorphic, though not canonically so. Therefore in any computational situation we provide an explicit isomorphism by choosing an explicit computational representation of the groups $`A`$ and $`A^{}`$. This isomorphism induces explicit isomorphisms between $`B`$ and $`B^{}`$ and between the factor group $`A/B`$ and $`B_{}`$. We will see shortly how we will employ these three interrelated isomorphisms. We will highlight this interrelation, and abuse notation, by using the symbol $`\varphi `$ to refer to all three of these isomorphisms, allowing for context to make the usage clear. As in the case of the Zak transformtion, these isomorphisms are simply reinterpretations of the contents of the registers.
Our QWHT is the sequence $`F_{B_{}\times B}\mathrm{\Phi }(g)Z(B).`$ Let us see how this unitary transformation acts on $`|a`$. We have
$$Z(B)|a=\frac{1}{\sqrt{|B|}}\underset{a_j^{}T_2}{}a_j^{}(x_aa)|x_a|a_j^{}$$
and then after applying $`\mathrm{\Phi }(g)`$ we obtain:
$$\frac{1}{\sqrt{|B|}}\underset{a_j^{}T_2}{}\frac{a_j^{}(x_aa)}{G(x_a,a_j^{})}|x_a|a_j^{}$$
which by the fundamental result discussed above equals:
$$\frac{1}{\sqrt{|B|}}\underset{b^{}}{}P(x_a,b^{})|x_a|b^{}$$
where we are now considering the contents of the second register to be an element of $`B^{}`$. We now utilize our explicit isomorphisms to reinterpret the contents of the first register as an element of $`B_{}`$ and the contents of the second register as an element of $`B`$:
$$\frac{1}{\sqrt{|B|}}\underset{b}{}P(b_{},b)|b_{}|b=\frac{1}{\sqrt{|B|}}\underset{\varphi (a_j^{})}{}P(\varphi (x_a),\varphi (b^{}))|\varphi (x_a)|\varphi (b^{}).$$
By applying the final transformation in the sequence $`F_{B_{}\times B}`$ we obtain our desired expansion:
$$\underset{(b,b_{})\mathrm{\Delta }}{}a|g_{(b,b_{})}|b|b_{}.$$
## Acknowledgements
We thank Myoung An, Richard Cleve, Peter Hoyer, Michele Mosca, and Richard Tolimieri for helpful conversations.
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# Adiabatic Elimination in Compound Quantum Systems with Feedback
## I Introduction
The quantum theory of continuous Markovian feedback is now well understood . Continuous feedback arises in a situation where a system continuously interacts with its environment, and the environment is deliberately engineered such that the influence of the system on the environment acts back on the system at a later time. This can be described as a Markovian process when (a) the natural coupling of the system to the environment is approximately Markovian, and (b) the effective time-delay in the feedback process is negligible compared to any relevant time scale of the system. If the Markovian approximation is appropriate, this leads to the great simplification that the system evolution may be described by a master equation of the Lindblad form .
It is possible to divide quantum feedback into two categories, which we may call coherent and incoherent, following Lloyd (but without being limited by his definitions). In the latter case of incoherent feedback, it is not necessary to use a quantum description of the entire feedback loop. Rather, at some point, it is permissible to change from a quantum to a classical description by invoking a measurement step. In a quantum optical context, this corresponds to electro-optical feedback where a photocurrent derived from detecting the light radiated by the system is used to control electro-optical devices which change the behavior of the system. In the former case of coherent feedback, a quantum description of the entire feedback loop is necessary. In a quantum optical context this corresponds to all-optical feedback in which the light radiated by the system is reflected so that in interacts with the system again, perhaps via some other system.
Continuous quantum feedback may be non-Markovian for a number of reasons. The coupling to the environment may be non-Markovian. The time delay in the feedback loop may be non-negligible. The feedback may act via a second system, the ancilla. In this paper we are concerned with the last possibility. This is of interest because it arises very naturally in quantum optics in both all-optical and electro-optical contexts. In principle, this sort of feedback can be described as a Markovian process in the larger state space of the system plus ancilla. In practice, this procedure is often not useful, because of the critical word larger in the previous sentence. If the required basis size of the system and ancilla are $`N`$ and $`M`$ respectively, then the Liouvillian for the compound system has of order $`N^4M^4`$ elements. Clearly for $`M`$ large, this is much larger than a Liouvillian for the system alone.
Consequently it would be an advantage to obtain a master equation for the system alone, without the ancilla. This is possible if the ancilla can be adiabatically eliminated. That is, if the ancilla has a decay rate much faster than any relevant system rate, so that it is always in a steady state determined by the system state. It is the purpose of this paper to determine numerically the conditions under which this is possible, and to derive the resultant master equations under those conditions, for a variety of general feedback systems.
Previous work in this area has left the situation somewhat confused. Wiseman and Milburn considered all-optical feedback via an ancilla system, and adiabatically eliminated the ancilla. This was shown to be equivalent to electro-optical feedback for quadrature feedback. However, for intensity feedback it was the same only to second order in the feedback strength. Moreover, the master equation derived (to second order) was not of the Lindblad form.
Slosser and Milburn considered electro-optic feedback of the photocurrent from the idler mode of a non-degenerate parametric oscillator onto the pump mode. Here the signal and idler mode formed the system and the pump mode was the ancilla. The procedure they adopted for deriving a master equation for the system was as follows. They expanded the feedback master equation for the compound system to first order in the feedback strength, adiabatically eliminated the pump mode, but the final result presented for the system master equation contained first and second order terms. As in Ref. , this second-order master equation was not of the Lindblad form. Furthermore, the steady-state field averages were calculated using an unstated all-order master equation (which was of the Lindblad form). There are other problems with this paper , but they are not relevant to the present work.
In this work we show how adiabatic elimination can be done rigorously in compound quantum feedback systems such as those of Refs. . As well as being of interest in the field of quantum feedback, the methods we use for adiabatic elimination are of more general interest. While adiabatic elimination of an ancilla mode which is linearly coupled to the system is well understood, adiabatic elimination with a nonlinear (e.g. proportional to the intensity) coupling is not. In particular, the methods we use here put the results obtained by Doherty and co-workers on the motion of an atom coupled to a damped optical cavity mode on a more rigorous footing.
This paper is organized as follows. In Sec. II we consider simple direct-detection feedback, and the four types of analogous feedback in compound systems: electro-optic feedback via a two-level atom, electro-optic feedback via an optical mode, all-optical feedback via a two-level atom, and all-optical feedback via a mode. We show that in all four cases it is possible to eliminate the ancilla under suitable conditions, giving a master equation for the system alone. In Sec. III we compare the stationary state of these master equations with the solution of the full dynamics of the compound systems. For this test we choose the free dynamics of the system to be that of a below-threshold parametric oscillator, the quantity being fed back to be the intensity, and the quantity being controlled by the feedback to be the detuning. We also compare the results of all five feedback mechanisms with that caused by an analogous “reversible feedback” generated by a $`\chi ^{(3)}`$ nonlinearity. In Sec. IV we conclude with a discussion of our results, and present a generalization of the all-optical case to multiple optical modes.
## II Adiabatic Elimination
### A Simple Feedback
In order to discern how the dynamics of a system are affected by a feedback loop that includes an ancilla, it is useful to know the master equation for simple feedback. By simple feedback it is meant that the measurement results, based on continuous observation of a source system, are immediately used to alter the evolution of the source without the involvement of any other quantum system. To use an example from quantum optics, a photodetector may register photon arrivals from a cavity at discrete times and, at these times, some specified change to the system may be made (see Fig. 1). Types of changes include altering the optical path length or damping rate of the cavity. In the remainder of this paper we will often use quantum optics terminology, but it should be remembered that the theory is not restricted to optical physics.
The most general form of the simple feedback master equation has been derived by Wiseman . Consider a system with Hamiltonian $`H`$ and some dissipation at rate $`\gamma `$ and with lowering operator $`c`$. With $`\mathrm{}`$ set equal to unity, the master equation is
$$\dot{\rho }(t)=i[H,\rho ]+\gamma 𝒟[c]\rho ,$$
(1)
where the Lindblad superoperator is
$$𝒟[c]=𝒥[c]𝒜[c],$$
(2)
where for arbitrary operators $`A`$ and $`B`$,
$$𝒥[A]B=ABA^{};𝒜[A]B=\frac{1}{2}\{A^{}A,B\}.$$
(3)
It is the dissipation which allows for continuous observation, the result of which is a current $`I(t)`$. In this paper we are concerned with what is known as direct detection where
$$I(t)=dN(t)/dt,$$
(4)
where $`dN(t)`$ is the point process (the increment in the number of photons counted) defined by
$`[dN(t)]^2`$ $`=`$ $`dN(t)`$ (5)
$`\mathrm{E}[dN(t)]`$ $`=`$ $`\gamma dt\mathrm{Tr}[c^{}c\rho _\mathrm{c}(t)]`$ (6)
Here E denotes a classically probabilistic expectation value, while the c subscript denotes that the state $`\rho _c`$ is conditioned on the previous measurement results. We have assumed that the detection is perfectly efficient; the generalization to inefficient detectors is trivial .
Simple feedback arises from adding a Hamiltonian to the system evolution of the form
$$H_{\mathrm{fb}}(t)=I(t)Z$$
(7)
where $`Z`$ is an Hermitian system operator. Taking into account the singularity of $`I(t)`$, and the fact that the feedback must act after the measurement, it is possible to derive a master equation for the system with feedback, averaging over all realizations of the stochastic measurement record $`I(t)`$. The result is
$$\dot{\rho }=i[H,\rho ]+\gamma 𝒟[e^{iZ}c]\rho .$$
(8)
To compare this master equation with those obtained later it is useful to expand the exponentials to third order
$`\dot{\rho }`$ $``$ $`𝒞[H]\rho +\gamma 𝒟[c]\rho `$ (10)
$`+\gamma \left\{𝒞[Z]+{\displaystyle \frac{1}{2}}\left(𝒞[Z]\right)^2+{\displaystyle \frac{1}{6}}\left(𝒞[Z]\right)^3\right\}𝒥[c]\rho ,`$
where $`𝒞[A]B=i[A,B]`$ for arbitrary operators $`A`$ and $`B`$.
The derivation outlined above for the feedback master equation treats the photocurrent $`I(t)`$ as a classical stochastic process, which causes the conditioned system state $`\rho _\mathrm{c}`$ to undergo stochastic evolution (known as a quantum trajectory ). There is an alternative derivation which treats the photocurrent $`I(t)`$ as an operator. This derivation works in the Heisenberg picture, where the system evolution is described by stochastic operator differential equations known as quantum Langevin equations . This method is useful for adiabatic elimination, so we will briefly review its features.
Quantum Langevin equations (QLE) are constructed without using the concept of measurement. The dissipative evolution of Eq. (1) can be derived in a quantum optical context from a linear coupling (in a rotating frame and with the rotating wave approximation)
$$V=i\sqrt{\gamma }[v^{}(t)cc^{}v(t)]$$
(11)
between the system and a bath of harmonic oscillators. Here $`v(t)`$ is the bath annihilation operator at the point at which it interacts with the system. Just before this point, the bath is an input vacuum, with field operator $`v_{\mathrm{in}}(t)`$ satisfying
$$[v_{\mathrm{in}}(t),v_{\mathrm{in}}^{}(t^{})]=\delta (tt^{}),$$
(12)
and all normally-ordered moments vanishing. Just after this point, the bath is an output (non-vacuum) with field operator
$$v_{\mathrm{out}}(t)=v_{\mathrm{in}}(t)+\sqrt{\gamma }c(t).$$
(13)
The photocurrent operator $`I(t)`$ is simply the intensity of the output field
$$I(t)=v_{\mathrm{out}}^{}(t)v_{\mathrm{out}}(t).$$
(14)
Adding together the evolution due to $`H`$, $`V`$, and $`H_{\mathrm{fb}}`$, and again noting that the feedback must act after the interaction, one can derive the following quantum Langevin equation for an arbitrary system operator $`s`$
$`ds`$ $`=`$ $`[v_{\mathrm{in}}^{}+\sqrt{\gamma }c^{}](e^{iZ}se^{iZ}s)[v_{\mathrm{in}}+\sqrt{\gamma }c]dt`$ (17)
$`+\gamma (c^{}sc{\displaystyle \frac{1}{2}}sc^{}c{\displaystyle \frac{1}{2}}c^{}cs)dt`$
$`\sqrt{\gamma }[dV_{\mathrm{in}}^{}cc^{}dV_{\mathrm{in}},s]+i[H,s]dt,`$
where $`dV_{\mathrm{in}}=v_{\mathrm{in}}dt`$. All operators have time argument $`t`$. When the expectation value of this equation is taken an equation is obtained that can be converted to the master equation (8) for simple feedback. If $`Z`$ is set to zero then the Langevin equation describes damping alone.
### B Electro-optic Feedback via an Atom
The simplest possible ancilla system is a two-level atom (TLA). In this section we consider incoherent (electro-optic) feedback via this ancilla. The output from the system is monitored by direct detection, the results of which are used to affect the evolution of the two level atom which is coupled to the system, as shown in Fig. 2. The system and ancilla are assumed to have approximately the same resonant frequency. If the atom is to be adiabatically eliminated, it must be heavily damped, in which case it will mostly be in the ground state. Then the most natural form of feedback involves flipping the state of the TLA whenever the photodetector monitoring the system makes a detection. This can be achieved with a feedback Hamiltonian of the form
$$H_{\mathrm{fb}}=\frac{\pi }{2}\sigma _xI(t).$$
(18)
Here $`\sigma _x`$ is the usual Pauli spin matrix for describing an atomic state . It could be realized experimentally by very briefly driving the atom with a pulse of on-resonance radiation (a ‘$`\pi `$’ pulse) which will flip it from the ground to the excited state.
With this form of feedback, the obvious coupling of the atom to the system to consider is one proportional to the excited state population operator $`\sigma ^{}\sigma `$. Here $`\sigma =(\sigma _xi\sigma _y)/2`$ is the atomic lowering operator. Specifically,
$$H_{\mathrm{coupling}}=\sigma ^{}\sigma K,$$
(19)
where $`K`$ is an arbitrary Hermitian system operator. When feedback onto the atom in the ground state occurs the upper state population jumps to a value of 1 and then decays away, due to coupling to the continuum of electromagnetic field modes. In other words, $`\sigma ^{}\sigma `$ will tend to follow the photocurrent. Thus there is a strong similarity to simple feedback, if $`K`$ is chosen to be some scalar multiple of $`Z`$.
It is not hard to generalize Eq. (8) to include the TLA ancilla
$`\dot{W}`$ $`=`$ $`i[H_{\mathrm{system}}+\sigma ^{}\sigma K,W]`$ (21)
$`+𝒟[\mathrm{exp}(i{\displaystyle \frac{\pi }{2}}\sigma _x)c]W+\mathrm{\Gamma }𝒟[\sigma ]W,`$
where $`\mathrm{\Gamma }`$ is the damping rate of the atom and $`W`$ is the density matrix for the compound system. The damping rate $`\gamma `$ of the system has been set equal to unity without loss of generality. Of course, the operators are now acting in the joint Hilbert space of the two systems so that $`cc1_{\mathrm{atom}}`$ and $`\sigma 1_{\mathrm{system}}\sigma `$, etc.
The above master equation gives the evolution of the density operator for the compound system. At any time a partial trace of this operator over the atom could be performed to obtain the reduced density matrix for the system alone. However, in general, this cannot be done to the master equation itself in order to obtain a master equation for $`\rho _{\mathrm{system}}(t)`$. The obvious exception to this is the case where $`K=0`$ and the system is unaffected by the atom.
It is logical that a master equation for the system cannot be derived if the atom observables fluctuate, in response to the feedback, on the same time scale as the system observables. The effect of feedback would then depend on the constantly fluctuating state of the atom which, in turn, depends on previous feedback. Removing the atom operators from the master equation without removing information concerning the system is impossible due to the coupling that exists between them. Of course, a non-Markovian expression could be written down for the atom in terms of the system, but this would not lead to a Lindblad master equation without some further approximation.
If the atom reacts very quickly to the feedback and returns to its initial state before more feedback arrives (the next photodetection) then this well defined behavior can be built into a master equation for the system alone. In essence, the atom’s state is approximated by its equilibrium value with respect to the instantaneous state of the system and operators are replaced by their steady state expressions. This procedure is known as adiabatic elimination of the atom.
To proceed with the adiabatic elimination it is noted that the total density matrix can be expanded as
$`W=`$ $`\rho _0||+\rho _1||`$ (22)
$`+\rho _1^{}||+\rho _2||,`$ (23)
where the $`\rho `$s exist in the system subspace. All possible states of the atom have been included ($`|`$ and $`|`$ correspond to the excited and ground state respectively). This approach is particularly appropriate because of the small basis involved. If the above expression for $`W`$ is substituted into the master equation then the atom operators can act on the states of the atom. If the coefficients of the various orthogonal states are equated the following equations for the $`\rho `$s are obtained (the subscript ‘s’ indicates the system):
$`\dot{\rho _0}`$ $`=`$ $`𝒞[H_\mathrm{s}]\rho _0+𝒥[c]\rho _2𝒜[c]\rho _0+\mathrm{\Gamma }\rho _2,`$ (24)
$`\dot{\rho _1}`$ $`=`$ $`𝒞[H_\mathrm{s}]\rho _1+i\rho _1K+𝒥[c]\rho _1^{}𝒜[c]\rho _1{\displaystyle \frac{\mathrm{\Gamma }}{2}}\rho _1,`$ (25)
$`\dot{\rho _2}`$ $`=`$ $`𝒞[H_\mathrm{s}]\rho _2+𝒥[c]\rho _0𝒜[c]\rho _2\mathrm{\Gamma }\rho _2.`$ (26)
By tracing Eq. (23) over the atom the reduced density operator for the system is
$$\rho _\mathrm{s}=\rho _0+\rho _2$$
(27)
and its evolution equation is found to be
$$\dot{\rho }_\mathrm{s}=i[H_\mathrm{s},\rho _\mathrm{s}]+𝒟[c]\rho _\mathrm{s}i[K,\rho _2].$$
(28)
Without some approximation this is as far as the elimination of the atom can be taken. It is not a master equation due to the dependence upon $`\rho _2`$. As discussed previously, the limit in which the atom returns very quickly to the ground state after feedback needs to be considered. Because the probability for photodetection in any infinitesimal time period scales as the size of the period, the atom is in the ground state for almost all time. The approximation that $`\rho _\mathrm{s}\rho _0`$ is therefore made. To obtain a master equation, an expression for $`\rho _2`$ in terms of $`\rho _0`$ is needed. From Eq. (26) it can be seen that if $`\mathrm{\Gamma }`$ is large compared to the other co-efficients of $`\rho _2`$ (except possibly $`K`$) then fluctuations in this operator will be quickly damped out and $`\dot{\rho _2}`$ can then be set to zero. The effect of $`K`$ is to cause rotation of $`\rho _2`$ but not to affect its size. The physical picture already described is consistent with $`\mathrm{\Gamma }`$ being large. Assuming $`K\mathrm{\Gamma }1`$ (where $`K\mathrm{\Gamma }`$ means that the operator $`K`$ scales like $`\mathrm{\Gamma }`$), we find the steady state of $`\rho _2`$ to be
$$\rho _2=(\mathrm{\Gamma }𝒞[K])^1𝒥[c]\rho _0.$$
(29)
When this is substituted into Eq. (28) the master equation for the system alone is obtained. With $`Z=K/\mathrm{\Gamma }`$ it is
$`\dot{\rho }_\mathrm{s}=\{𝒞[H_\mathrm{s}]+𝒟[c]+𝒞[Z](1𝒞[Z])^1𝒥[c]\}\rho _\mathrm{s}.`$ (30)
It is not immediately clear that this master equation is of the Lindblad form . However in appendix A 1 it is shown that it can be written as
$$\dot{\rho }_\mathrm{s}=i[H_\mathrm{s},\rho _\mathrm{s}]+_0^{\mathrm{}}𝑑qe^q𝒟[e^{iqZ}c]\rho _\mathrm{s}.$$
(31)
Some feeling for the nature of the master equation can be obtained by an expansion to third order in $`Z`$ (a small feedback approximation). This gives (subscripts dropped)
$`\dot{\rho }`$ $``$ $`𝒞[H]\rho +𝒟[c]\rho `$ (33)
$`+\left\{𝒞[Z]+(𝒞[Z])^2+(𝒞[Z])^3\right\}𝒥[c]\rho .`$
These terms can be compared to the third order expansion of Eq. (10), with $`\gamma =1`$. The difference in second and higher order terms means that for large feedback the two systems will be significantly different.
### C Electro-optic Feedback via a mode
The more challenging task of adiabatically eliminating an ancilla that has an infinite number of basis states is now considered. Optically, this could correspond to a single-mode cavity. The method of expanding the compound density matrix in terms of the lower number states of the ancilla is not appropriate due to the type of feedback that is utilized. Instead we use Quantum Langevin equations, which place no such restriction on the excitation of the ancilla.
The output field from the system is once again continuously monitored using direct detection (see Fig. 3). We take the feedback to be linear driving of the ancilla cavity. This causes a jump in amplitude of the ancilla cavity when there is a photodetection. It is described by the feedback Hamiltonian
$$H_{\mathrm{fb}}=\frac{ϵ}{2}(ib+ib^{})I(t),$$
(34)
where $`b`$ is the annihilation operator for the cavity, $`ϵ`$ represents the amplitude of the coherent driving field and $`I(t)`$ is the operator for the photocurrent output from the system. Its effect can be determined from the Heisenberg equation of motion for $`b`$,
$$\dot{b}_{\mathrm{fb}}=i[b,H_{\mathrm{fb}}]=\frac{ϵ}{2}I(t).$$
(35)
Since $`I(t)`$ consists of $`\delta `$ functions, it is clear that the cavity field amplitude changes by an amount $`ϵ/2`$ whenever a photodetection occurs. Note that here the implicit equation of motion for $`b`$ is sufficient to determine its evolution because the stochastic term is not dependent upon $`b`$ . To provide a feedback circuit that is classically equivalent to simple feedback in the limit of large damping of the cavity, the following choice of coupling is made:
$$V=\frac{K}{2}(b+b^{}).$$
(36)
The equivalence can be seen if linear damping is included in Eq. (35). The slaved value of $`b`$ (in the limit of large damping $`\dot{b}_{\mathrm{fb}}`$ is set equal to zero) is then substituted into the coupling, which leaves it in the same form as a simple feedback Hamiltonian, given an appropriate choice of $`K`$.
The total master equation is
$`\dot{W}`$ $`=`$ $`i[{\displaystyle \frac{K}{2}}(b+b^{})+H_\mathrm{s},W]`$ (38)
$`+𝒟[e^{ϵ(b+b^{})/2}c]W+\mathrm{\Gamma }𝒟[b]W,`$
where once again $`W`$ is the density matrix describing the compound system and the damping of the system has been set equal to unity. The damping rate of the ancilla cavity is given by $`\mathrm{\Gamma }`$. The quantum Langevin equation that corresponds to this master equation can be found by extending Eq. (17). The result for an arbitrary operator $`r`$ from either sub-system is
$`dr=`$ $`v_{\mathrm{out}}^{}\left[e^{ϵ(bb^{})/2}re^{ϵ(bb^{})/2}r\right]v_{\mathrm{out}}dt`$ (42)
$`+𝒟[c^{}]rdt[dV_{\mathrm{in}}^{}cc^{}dV_{\mathrm{in}},r]`$
$`+\mathrm{\Gamma }𝒟[b^{}]rdt\sqrt{\mathrm{\Gamma }}[dU_{\mathrm{in}}^{}bb^{}dU_{\mathrm{in}},r]`$
$`+i[{\displaystyle \frac{K}{2}}(b+b^{})+H_\mathrm{s},r]dt,`$
where $`dU_{\mathrm{in}}=u_{\mathrm{in}}dt`$. The vacuum field input for the driven cavity, $`u_{\mathrm{in}}`$, has the same properties as $`v_{\mathrm{in}}`$.
To adiabatically eliminate the cavity, in the limit of heavy damping, a QLE will first be determined for a system operator, $`s`$. Eq. (42) is greatly simplified, as $`s`$ commutes with all driven cavity operators, to give
$`ds`$ $`=`$ $`𝒟[c^{}]sdt[dV_{\mathrm{in}}^{}cc^{}dV_{\mathrm{in}},s]`$ (44)
$`+i[{\displaystyle \frac{K}{2}}(b+b^{})+H_\mathrm{s},s]dt.`$
From this it is evident that an expression for $`b`$ is required if a master equation for the system alone is to be derived. The QLE for $`b`$ is
$$\dot{b}=i\frac{K}{2}\frac{\mathrm{\Gamma }}{2}b\sqrt{\mathrm{\Gamma }}u_{\mathrm{in}}+\frac{ϵ}{2}v_{\mathrm{out}}^{}v_{\mathrm{out}}.$$
(45)
For large $`\mathrm{\Gamma }`$ the fluctuations in $`b`$ due to system operators will be quickly damped out. However, the stochastic terms have an infinite bandwidth, so that it is not strictly possible to slave an operator that only responds to a finite bandwidth, $`\mathrm{\Gamma }`$, to these fluctuations. Although this problem can be side-stepped it will prove advantageous to use the following equilibrium value of $`b`$
$`b`$ $`=`$ $`{\displaystyle \frac{iK}{\mathrm{\Gamma }}}{\displaystyle _0^{\mathrm{}}}d\tau e^{\mathrm{\Gamma }\tau /2}[\sqrt{\mathrm{\Gamma }}u_{\mathrm{in}}(t\tau )`$ (47)
$`{\displaystyle \frac{ϵ}{2}}v_{\mathrm{out}}^{}v_{\mathrm{out}}(t\tau )].`$
The integral serves to determine the present contribution to $`b`$ from the stochastic terms at time $`t\tau `$. This contribution falls off at rate $`\mathrm{\Gamma }/2`$, the amplitude decay rate for the ancilla cavity. The term that is not under the integral comes from $`K`$ which is not stochastic and is therefore slowly varying compared to the highly damped cavity operators. Thus, $`b`$ can follow its evolution to a very good approximation.
To simplify matters the Langevin equation for $`s`$ will now be rearranged before substitution so that $`u_{\mathrm{in}}`$ will annihilate the vacuum when the expectation value is taken. This gives
$`ds`$ $`=`$ $`𝒟[c^{}]sdt[dV_{\mathrm{in}}^{}cc^{}dV_{\mathrm{in}},s]`$ (49)
$`+{\displaystyle \frac{i}{2}}\left(b^{}[K,s]+[K,s]b\right)dt+i[H_\mathrm{s},s]dt.`$
This is valid as $`b`$ and $`b^{}`$ commute with system operators. We cannot move the stochastic part $`v_{\mathrm{in}}(t)`$ of $`v_{\mathrm{out}}(t\tau )`$ through the system commutator term to annihilate on the vacuum. However, it is possible to move the photocurrent itself at time $`t\tau `$ as it commutes . If the integrals that will annihilate on the vacuum when the trace over the bath is taken are ignored, then we are left with
$`\dot{s}`$ $`=`$ $`{\displaystyle \frac{iϵ}{2}}[K,s]{\displaystyle _0^{\mathrm{}}}𝑑\tau e^{\mathrm{\Gamma }\tau /2}I(t\tau ){\displaystyle \frac{1}{2\mathrm{\Gamma }}}[K,[K,s]]`$ (51)
$`+𝒟[c^{}]s[v_{\mathrm{in}}^{}cc^{}v_{\mathrm{in}},s]+i[H_\mathrm{s},s].`$
If the limit $`\mathrm{\Gamma }\mathrm{}`$ is taken the integral reduces to $`2I(t)/\mathrm{\Gamma }`$. The resultant equation for $`\dot{s}`$ is an implicit equation as it was derived by idealizing the properties of the cavity and environment . An explicit equation is now required.
The term that needs to be treated in Eq. (51) can be written as
$$\dot{s}_{\mathrm{implicit}}=\frac{ϵI𝒞[K]s}{\mathrm{\Gamma }}.$$
(52)
This gives an explicit increment of the form
$$ds_{\mathrm{explicit}}=dN[\mathrm{exp}(ϵ𝒞[K]/\mathrm{\Gamma })1]s,$$
(53)
where $`dN=Idt=dN^2=v_{\mathrm{out}}^{}v_{\mathrm{out}}dt`$. Remembering that the photocurrent is actually evaluated at a slightly earlier time than the system operators allows $`v_{\mathrm{out}}`$ to be moved to the right of the expression. If we put $`Z=ϵK/\mathrm{\Gamma }`$, in order that our equations can be compared to simple feedback, then the total Langevin equation is
$`ds`$ $`=`$ $`[v_{\mathrm{in}}^{}+c^{}](e^{iZ}se^{iZ}s)[v_{\mathrm{in}}+c]dt`$ (56)
$`{\displaystyle \frac{\mathrm{\Gamma }}{2ϵ^2}}[Z,[Z,s]]dt+𝒟[c^{}]sdt`$
$`[dV_{\mathrm{in}}^{}cc^{}dV_{\mathrm{in}},s]+i[H_\mathrm{s},s]dt.`$
When the expectation value is taken the stochastic part annihilates on the vacuum and the following master equation is obtained
$$\dot{\rho }=i[H_\mathrm{s},\rho ]+𝒟[e^{iZ}c]\rho +\frac{\mathrm{\Gamma }}{ϵ^2}𝒟[Z]\rho .$$
(57)
The only difference from simple feedback is the third term. This is a term of second order in the feedback operator $`Z`$, and represents a type of noise that will tend to smooth over the interesting behavior of the system. Clearly it can be made arbitrarily small if $`ϵ`$ is made large enough. A more detailed discussion of this term is given in Sec. III C
### D All-optical Feedback via an atom
We turn now to coherent, or all-optical feedback. Once again we begin with the simplest possible ancilla, a two-level atom. All-optical feedback via an atom involves the reflection of the output field from the system onto the atom, where the atom is reversibly coupled to the system. Here, the resonant frequencies of the two systems are taken to be equal. It is different from electro-optic feedback as there is no measurement ste; the light is just reflected around a loop with the use of mirrors (see Fig. 4). The theoretical description of such systems was developed largely by Carmichael and Gardiner and has been termed Cascaded Open Systems theory. If linear bath-system couplings are assumed then the compound master equation is
$`\dot{W}=`$ $`i[H_\mathrm{s}+V,W]+𝒟[c]W+\mathrm{\Gamma }𝒟[\sigma ]W`$ (58)
$`+\sqrt{\mathrm{\Gamma }}\left([cW,\sigma ^{}]+[\sigma ,Wc^{}]\right).`$ (59)
The system damping has been set equal to unity as usual and $`\mathrm{\Gamma }`$ is the damping rate of the atom.
In order to investigate the degree to which all-optical feedback can replicate electro-optical simple feedback, a coupling is chosen that is linear in the excited state population of the atom. We expect this operator to follow the output photocurrent from the system. That is, we assume a coupling
$$V=K\sigma ^{}\sigma $$
(60)
identical to that in Sec. II B. Making the expansion of Eq. (23) gives the following for the $`\rho `$s
$`\dot{\rho _0}`$ $`=`$ $`𝒞[H_\mathrm{s}]\rho _0+𝒟[c]\rho _0+\mathrm{\Gamma }\rho _2+\sqrt{\mathrm{\Gamma }}(c\rho _1+\rho _1c^{}),`$ (61)
$`\dot{\rho _1}`$ $`=`$ $`𝒞[H_\mathrm{s}]\rho _1+𝒟[c]\rho _1+\sqrt{\mathrm{\Gamma }}(\rho _2\rho _0)c^{}`$ (63)
$`+i\rho _1K{\displaystyle \frac{\mathrm{\Gamma }}{2}}\rho _1,`$
$`\dot{\rho _2}`$ $`=`$ $`𝒞[H_\mathrm{s}]\rho _2+𝒟[c]\rho _2\sqrt{\mathrm{\Gamma }}(c\rho _1+\rho _1c^{})`$ (65)
$`i[K,\rho _2]\mathrm{\Gamma }\rho _2.`$
The above equations lead to an equation of motion for the system density operator of
$$\dot{\rho }=𝒞[H_\mathrm{s}]\rho +𝒟[c]\rho i[K,\rho _2],$$
(66)
which is the same as Eq. (28). To find an expression for $`\rho _2`$ the normal procedure of taking $`\mathrm{\Gamma }`$ large compared to $`𝒞[H_\mathrm{s}]`$ is performed. Thus, $`\rho _1`$ can be slaved to system operators, $`\rho _0`$ and $`\rho _2`$. Now as we only require a master equation which gives the leading order effect in $`\mathrm{\Gamma }^1`$ of the ancilla on the system, $`\rho _2`$ can be set equal to zero in the $`\rho _1`$ equation, which is the approximation $`\rho _0\rho `$. This is valid as $`\rho _2\rho _0/\mathrm{\Gamma }`$. By substituting the slaved expression for $`\rho _1`$ into that for $`\rho _2`$ we find after simplification
$$\rho _2=\frac{4}{\mathrm{\Gamma }}𝒥\left[\left(1+\frac{2iK}{\mathrm{\Gamma }}\right)^1c\right]\rho _0.$$
(67)
This can now be substituted into Eq. (66) to obtain a master equation. Writing $`Z=4K/\mathrm{\Gamma }`$, we have
$$\dot{\rho }=𝒞[H_\mathrm{s}]\rho +𝒟[c]\rho +𝒞[Z]𝒥\left[\left(1+\frac{Zi}{2}\right)^1c\right]\rho ,$$
(68)
which is the same as the simple feedback Eq. (10) to second order. The third order term is
$$\frac{1}{4}𝒞[Z](𝒥[Z]2𝒜[Z])𝒥[c]\rho .$$
(69)
Again it is not obvious that Eq. (68) is in the Lindblad form, but it is shown in Appendix A 2 that it can be written as
$$\dot{\rho }=i[H_\mathrm{s},\rho ]+𝒟\left[\mathrm{exp}\left(2i\mathrm{arctan}\frac{Z}{2}\right)c\right]\rho .$$
(70)
### E All-optical Feedback via a mode
The final compound system that will be considered involves the output field from a system being reflected onto an optical cavity that is coupled back to the system (see Fig. 5). A Faraday Isolator (comprised of a Faraday Rotator and a Polarization dependent Beam Splitter) prevents reflected light from the cavity returning to the system. The only difference in the total master equation from the previous section is the replacement of the atom lowering operator $`\sigma `$ with the annihilation operator $`b`$. Thus a coupling of the form $`V=Kb^{}b`$ is considered.
The derivation of a master equation for the system alone follows similar lines to that of Sec. II C. The QLE for an arbitrary operator is
$`dr`$ $`=`$ $`+i[H_\mathrm{s}+V,r]dt+𝒟[c^{}]rdt[dV_{\mathrm{in}}^{}cc^{}dV_{\mathrm{in}},r]`$ (73)
$`+\mathrm{\Gamma }𝒟[b^{}]rdt\sqrt{\mathrm{\Gamma }}[dV_{\mathrm{in}}^{}bb^{}dV_{\mathrm{in}},r]`$
$`+\sqrt{\mathrm{\Gamma }}(b^{}rc+c^{}rbrb^{}cc^{}br)dt.`$
For a system operator this becomes
$`ds`$ $`=`$ $`𝒟[c^{}]sdt[dV_{\mathrm{in}}^{}cc^{}dV_{\mathrm{in}},s]`$ (75)
$`+i[H_\mathrm{s}+Kb^{}b,s]dt.`$
The next step is to find an equation for $`b`$. The QLE that governs it is
$$db=(\frac{\mathrm{\Gamma }}{2}b+\sqrt{\mathrm{\Gamma }}v_{\mathrm{in}}+\sqrt{\mathrm{\Gamma }}c+iKb)dt.$$
(76)
This justifies our initial presumption that the cavity photon number would follow the photocurrent. For $`\mathrm{\Gamma }`$ large it is possible to slave $`b`$ to the system operators and to form an integral expression for the contribution from the stochastic term, as in Sec. II C. The result is
$`b`$ $`=`$ $`{\displaystyle \frac{2}{\sqrt{\mathrm{\Gamma }}}}\left(1+{\displaystyle \frac{2iK}{\mathrm{\Gamma }}}\right)^1c`$ (78)
$`+\sqrt{\mathrm{\Gamma }}{\displaystyle _0^{\mathrm{}}}𝑑\tau e^{\mathrm{\Gamma }\tau (1+2iK/\mathrm{\Gamma })/2}v_{\mathrm{in}}(t\tau ).`$
The same trick of rearranging the QLE for the system operator is again used so that, in this case, all of the integral terms annihilate. We put
$$i[Kb^{}b,s]dt=ib^{}[K,s]bdt.$$
(79)
Substituting into this the expression for $`b`$ and $`b^{}`$ gives four terms, only one of which is non-zero when the trace over the bath is taken. This term is
$$\frac{4ic^{}}{\mathrm{\Gamma }}\left(1\frac{2iK}{\mathrm{\Gamma }}\right)^1[K,s]\left(1+\frac{2iK}{\mathrm{\Gamma }}\right)^1c.$$
(80)
In effect, an implicit equation has been derived that has no contribution from stochastic operators, resulting in there being no need for an implicit/explicit distinction. It is now possible to turn the equation for $`ds`$ into a master equation for the system. When this is done we arrive at the same result as Eq. (68). The conclusion is that to first order in $`\mathrm{\Gamma }^1`$, the cavity has the same effect on the system that the atom does, when included in an all-optical feedback loop.
In hindsight, this is what we should have expected, as in the limit of large damping only the lowest number states of the cavity will be occupied with significant probability. One could therefore have expanded the total density matrix analogously to the TLA system to obtain the same equations immediately.
The reason why electro-optic feedback onto an atom and a cavity were not equivalent is due to the more singular nature of the driving of the ancilla. When a detection on the system is made the field amplitude of the cavity jumps, leading to occupation of higher photon number states. These states are, therefore, essential to the description of the compound system. Electro-optic feedback onto the atom cannot replicate this behavior.
## III COMPARISON WITH EXACT RESULTS
We have shown that in principle it is possible to consider a variety of different sorts of feedback in compound quantum systems, and to adiabatically eliminate the ancillary system to arrive at master equations for the system of interest alone. These master equations should be exact in the limit that the ancilla is damped infinitely faster than the system. In practice, this will never be the case, so it is an interesting question to find out under what conditions the equations are valid. This can be done by simulating the full master equation for the compound system and comparing to the results of the master equation for the system alone.
To make such a comparison requires specifying the feedback operator, $`Z`$, and the system Hamiltonian, $`H_\mathrm{s}`$. Once this is done, a comparison can be made by looking at the stationary solutions of the respective master equations. While this could be criticised as not being a complete test, it has the advantages of definiteness and ease of calculation (in some cases at least). Furthermore, we choose a system (a damped optical mode) and Hamiltonians $`H_\mathrm{s}`$ and $`Z`$ such that the stationary solutions have enough structure for the comparison to be interesting. The comparison is both quantitative and qualitative, with the use of the Bures distance as a measure of the difference between the state matrices and the Wigner function to illustrate them.
In the hope of getting some interesting states we take the system to be a damped single mode optical cavity. That is, we choose $`c=a`$, an annihilation operator satisfying $`[a,a^{}]=1`$. We choose a system Hamiltonian (in a rotating frame) of
$$H=\frac{i\lambda }{4}\left[a^2(a^{})^2\right],$$
(81)
This describes a degenerate parametric amplifier (“two photon” driving), which can be realized by driving an intracavity crystal with a $`\chi ^{(2)}`$ non-linearity with light at twice the resonant frequency. For $`\lambda `$ positive, this results in squeezing of the $`X_2`$ quadrature of the field inside the cavity, and stretching of the $`X_1`$ quadrature. The two quadratures are defined in this paper as
$`X_1`$ $`=`$ $`a+a^{}`$ (82)
$`X_2`$ $`=`$ $`i(aa^{}).`$ (83)
Without feedback, the master equation with two-photon driving and damping will have a stationary solution only for $`\lambda <1`$. That is, $`\lambda `$ is the threshold parameter.
The feedback operator is chosen to be
$$Z=\chi a^{}a.$$
(84)
We can get a feel for the effect of this type of feedback by using $`Z`$ in the simple feedback Hamiltonian given in Eq. (7). That is,
$$H_{\mathrm{fb}}(t)=\chi I(t)a^{}a.$$
(85)
This represents a detuning of the system cavity proportional to the photocurrent. It will cause the master equation to have a stationary solution regardless of $`\lambda `$, as will be shown. As the mean photocurrent is equal to the expectation value of the photon number operator for the system, this Hamiltonian is akin to a $`\chi ^{(3)}`$ Kerr non-linearity . In Sec. III F a comparison of feedback to such a nonlinearity is made.
### A Simple Feedback
The master equation for simple feedback is now
$$\dot{\rho }=\frac{\lambda }{4}[a^2(a^{})^2,\rho ]+𝒟[e^{i\chi a^{}a}a]\rho .$$
(86)
To simplify the numerical analysis we choose a single feedback strength for which simulations will be run. To aid this decision the effect of feedback is analyzed. Consider the following quantity:
$`𝒥\left[e^{i\chi a^{}a}\right]\rho .`$ (87)
If this is evaluated in the number basis then we get
$$n|𝒥\left[(e^{i\chi })^{a^{}a}\right]\rho |m=(e^{i\chi })^{nm}\rho _{nm}.$$
(88)
Now this particular system has the property that $`\rho _{nm}=0`$ for $`|nm|`$ odd as the two photon driving is the only source of coherences. These coherences exist between elements with $`|nm|`$ even. Hence, if $`\chi =q\pi `$, with $`q`$ an integer, then the feedback has no effect. Investigation into the states produced with a value of feedback close to this revealed that they are extremely sensitive to any parameter variation. This implies that it is not a suitable regime for the testing of adiabatic elimination. The most obvious alternative is to choose the maximum feedback regime. It is clear that this is achieved with $`\chi =(q+1/2)\pi `$. The states produced are much less sensitive and also have the advantage that, for simple feedback, there is no threshold to the driving strength above which the photon number becomes infinite. For the remainder of the paper we choose $`\chi =\pi /2`$.
The two-photon driving strength $`\lambda `$ was chosen to be as large as possible, given the constraints on the maximum basis size that could be simulated. This amplified the interesting effects of feedback. Not surprisingly, the simulations of the compound systems are the most computationally intensive and provide the upper basis size. It was found that the limit for the system cavity basis size required that photon numbers above 35 had to be truncated. For an accurate simulation this gives a maximum driving strength of about $`\lambda =2.2`$. Where possible, the compound systems were examined in the same regime as simple feedback, but for some the driving threshold of $`\lambda =1`$ remains in force, so $`\lambda =0.97`$ was then chosen.
The numerical simulations were greatly aided by the use of the Quantum Optics toolbox for Matlab . As noted above, we gauged whether the adiabatic elimination is valid by investigating the steady states of the systems. The simple feedback system involved a small enough Liouvillian that matrix inversion methods can be used. The Wigner function of the steady state density matrix for simple feedback, with $`\lambda =2.2`$ and $`\chi =\pi /2`$, is shown in Fig. 6. A plot with $`\lambda =0.97`$ is also included.
### B Electro-optic Feedback via an Atom
Electro-optic feedback via an atom can be compared to the simple feedback just considered if we insert in Eq. (19) $`K=\mathrm{\Gamma }Z=ga^{}a`$, where $`g=\mathrm{\Gamma }\pi /2`$. To test the adiabatic elimination simulations were run for various values of $`\mathrm{\Gamma }`$. It is only for large $`\mathrm{\Gamma }`$ that correspondence between the full dynamics and the adiabatically eliminated master equation is expected. A physical realization of this coupling is a far detuned atom in the standing wave of a single mode cavity . This also introduces a term into the system Hamiltonian of the form $`\delta \sigma ^{}\sigma `$, where $`\delta `$ is the difference in resonant frequency of the atom and system. It is of interest to determine whether the same results are obtained if the adiabatic elimination is done at the same time, rather than after, the large-detuning approximation is made. This is addressed in appendix B, and the answer is affirmative.
The full master equation is of the form
$`\dot{W}`$ $`=`$ $`i[{\displaystyle \frac{i\lambda }{4}}\left\{a^2(a^{})^2\right\}+g\sigma ^{}\sigma a^{}a+\delta \sigma ^{}\sigma ,W]`$ (90)
$`+𝒟[\mathrm{exp}(i{\displaystyle \frac{\pi }{2}}\sigma _x)a]W+\mathrm{\Gamma }𝒟[\sigma ]W.`$
The reduced density matrix for the system at steady state needs to be found. Once again, the Liouvillian is small enough that we can set $`\dot{W}=0`$ and solve the equation $`W=0`$ for the non-trivial solution. Simulations were run for values of $`\mathrm{\Gamma }`$ from 1 to 100, with $`g`$ altered accordingly. Note that the detuning actually has no effect on the system dynamics. The reduced density matrices produced are compared with those found from Eq. (30) with the aid of the Bures distance, which gives a measure of how distinguishable two mixed states ($`\rho _1`$ and $`\rho _2`$) are. The Bures distance is defined as
$$d_{\mathrm{Bures}}(\rho _1,\rho _2)=\sqrt{2}\left(1\mathrm{Tr}\left[\sqrt{\sqrt{\rho _1}\rho _2\sqrt{\rho _1}}\right]\right).$$
(91)
All pairs of density matrices of the same size have a Bures measure that is mapped onto the real numbers between zero and $`\sqrt{2}`$. Fig. 7 shows how the state produced by the compound master equation approaches that produced by the adiabatically eliminated master equation. As $`\mathrm{\Gamma }`$ is increased the Bures distance decreases and the Wigner functions become more similar to the adiabatic state. This shows that the adiabatic elimination is valid in this system for surprisingly small values of $`\mathrm{\Gamma }`$.
A comparison of the stationary Wigner functions produced here with those of simple feedback reveals that there exists vast differences between these feedback schemes. This is not surprising as it is only to first order in $`Z`$ that the equations are the same, and the parameters we have chosen correspond to $`Z`$ quite large. The most obvious visual differences include the presence of a shearing effect and the loss of reflective symmetry in the $`X_2`$ quadrature.
### C Electro-optic Feedback via a Mode
In Sec. II C electro-optic feedback via a mode was considered. In the limit of the ancilla mode being damped on a time scale small compared to those of the system, Eq. (57) was derived. The feedback operator was set as $`Z=ϵK/\mathrm{\Gamma }`$ so that we could make a comparison to simple feedback. It follows that the system coupling operator, $`K`$, is of the same form as the previous section: $`K=ga^{}a`$. The coupling $`V=ga^{}a(b+b^{})/2`$ could be physically achieved via a four wave mixing process in a $`\chi ^{(3)}`$ material . The fourth field would have to have the same frequency as the ancilla cavity for conservation of energy.
Now that $`Z`$ has been specified, the third term in Eq. (57) can be discussed more explicitly. This can be done by considering the evolution of the phase operator, which has an approximate commutation relation with the number operator of $`[\mathrm{\Phi },n]=i`$ . It can then be shown that this term causes phase diffusion at a constant rate, implying that the features of the state which are dependent upon a distinct phase are lost. With the notable exception that the photon number is not directly affected, there are many similarities with damping.
For simulation, parameters are chosen so that $`ϵg/\mathrm{\Gamma }=\pi /2`$, $`\lambda =2.2`$ and $`\mathrm{\Gamma }/2ϵ^2=0.001`$. The last equality maintains the phase diffusion term at a small and constant level. This ensures that the same state is always produced by the adiabatically eliminated master equation.
It is worth mentioning how the full dynamics were simulated. Due to the jump in the field amplitude of the ancilla cavity when a detection on the output of the system is made, the basis size required for an accurate simulation is large. If the amplitude jumps by an amount $`ϵ/2`$ then the photon number will increase by (presuming the initial field was small) $`ϵ^2/4`$. A second detection on the system occurring very soon after the first, would push the photon number even higher. In fact the computational resources available were not sufficient to allow even a quantum trajectory simulation of Eq. (38). The solution was to make a unitary transformation to a frame in which the evolution of the driven cavity due to feedback was separated from that due to quantum noise. That is, the mean amplitude of the field was described classically while the quantum representation of the noise was maintained. The unitary transformation used was
$$U=\mathrm{exp}[ϵf(t)(bb^{})/2],$$
(92)
where $`f(t)`$ is defined by
$$f(t)=_{\mathrm{}}^t𝑑s\mathrm{exp}[\mathrm{\Gamma }(ts)/2]I(s).$$
(93)
Here, $`I(t)`$ is the $`c`$-number stochastic photocurrent. The price of a reduced basis size is a time dependent Liouvillian. When the transformation of Eq. (92) is applied to the implicit master equation (feedback is described by a feedback Hamiltonian instead of the exponentials) an equation is obtained that is already of an explicit form (see appendix C)
$`\dot{W}`$ $`=`$ $`i[ga^{}a\{b+b^{}+ϵf(t)\}{\displaystyle \frac{i\lambda }{4}}\left\{a^2(a^{})^2\right\},W]`$ (95)
$`+𝒟[a]W+\mathrm{\Gamma }𝒟[b]W.`$
It can be seen that $`ϵf(t)`$ represents the amplitude of the driven cavity. Although $`f(t)`$ is stochastic, it is a smoothed (non-singular) version of the photocurrent and can therefore be treated without worrying about the stochastic calculus. Note also that since $`U`$ contains only ancilla operators, the system state matrix $`\rho =\mathrm{Tr}_b[\stackrel{~}{W}]`$ is the same as before, $`\mathrm{Tr}_b[W]`$.
The transformed master equation was simulated using quantum trajectory methods. It is shown in Fig. 8 that as $`\mathrm{\Gamma }`$ becomes large the adiabatically eliminated master equation becomes a very good approximation to the full dynamics. Clearly, though, $`\mathrm{\Gamma }`$ has to be pushed to much higher levels than the TLA damping for this correspondence to hold. One reason for this is that the Wigner functions of the steady state density matrices for electro-optic feedback onto a mode have much greater structure, meaning that a measure such as the Bures distance (which measures the distinguishability of states) will be more sensitive to small differences. It also is likely that the parameter regime chosen is one in which this system varies quickly, with the result that adiabatic elimination will only be valid at very large $`\mathrm{\Gamma }`$.
A comparison of the Wigner functions \[Figs. 8(a) and (c)\] with that produced with simple feedback \[Fig. 6(a)\] shows the expected similarity. It is expected because Eq. (57) only differs from simple feedback due to the presence of the double commutator noise term, which was chosen to be small.
### D All-optical Feedback onto an Atom
The basis size of the TLA ensures that simulating the full dynamics of all-optical feedback \[Eq. (59)\] is relatively easy. However, a threshold driving strength exists ($`\lambda =1`$) for this system which means that the adiabatically eliminated master equation cannot be tested in the same regime as the previous sections. Instead we set $`\lambda =0.97`$ which enabled us to perform an accurate simulation with the computational resources available.
The Hamiltonian of Eq. (59) includes the parametric amplifier driving and also the coupling of Eq. (60). Once again we choose $`K=ga^{}a`$ and set $`4g/\mathrm{\Gamma }=\pi /2`$, while varying $`\mathrm{\Gamma }`$ and $`g`$. The Bures distance between the states produced by Eq. (59) and Eq. (68) is shown in Fig. 9, as are some Wigner functions for the full dynamics and the adiabatic state. It can be seen that the state produced with the full dynamics approaches the adiabatic state at a similar rate, as $`\mathrm{\Gamma }`$ is increased, to electro-optic feedback via a TLA.
There is a large similarity between the state produced via simple feedback in Fig. 6 (c) and that in Fig. 9 (a), with the presence of shearing being the most notable difference. This closer correspondence to simple feedback than that of the electro-optic feedback systems is not surprising given that the adiabatic all-optical master equation was the same as simple feedback to a higher order (second). The smaller driving also contributes to the closeness of the states.
### E All-optical Feedback via a Mode
It was shown in Sec. II E that in the adiabatic limit all-optical feedback onto a mode has the same effect as feeding back onto a TLA. Therefore, the same threshold for the driving strength exists for this system ($`\lambda =1`$).
The basis size required here is not as large as for electro-optic feedback because the photons leak out of the system and into the ancilla cavity, giving a smooth variation of photon number. Despite this, a quantum trajectory simulation was still found to be necessary. The results obtained for $`\lambda =0.97`$ and $`4g/\mathrm{\Gamma }=\pi /2`$ can be found in Fig. 10. The adiabatic state is, of course, the same as for all-optical feedback onto a TLA. There is a notable difference in the speed at which the full dynamics approaches this state. At low damping the Bures distance is already very low. The conclusion is that the ancilla mode has minimal effect on the system when included in an all-optical feedback loop.
### F Comparison with “Reversible Feedback” Generated by a $`\chi ^{(3)}`$ Non-linearity
Finally, we consider the effect of placing a $`\chi ^{(3)}`$ material inside an optical cavity driven by a parametric oscillator. There is no feedback loop involved. The Hamiltonian generated by the $`\chi ^{(3)}`$ non-linearity (a Kerr non-linearity) is given by
$$H_{\mathrm{Kerr}}=\frac{\chi }{2}(a^{})^2a^2.$$
(96)
The Heisenberg equation of motion of the annihilation operator due to this Hamiltonian is found to be
$$\dot{a}=i\chi (a^{}a)a.$$
(97)
Thus, it is clear that the $`\chi ^{(3)}`$ non-linearity causes a detuning proportional to the intensity of the field inside the cavity. In this way the system has a self-awareness that is similar to simple feedback, which is why a comparison is relevant. In fact, it can be shown that the two systems are classically equivalent given the same choice of the parameter $`\chi `$. For large feedback the two systems diverge when treated quantum mechanically. One of the main reasons behind this is that the Kerr effect displays no periodic dependence upon its magnitude, whereas the simple feedback does. This is illustrated in Fig. 11(c), where the Bures distance between the steady states of the two systems is plotted for varying $`\chi `$. The Wigner function of the “reversible feedback” steady state for $`\chi =\pi /2`$ with $`\lambda =2.2`$ and $`\lambda =0.97`$ is given in Figs. 11(a) and (b) respectively. They are seen to be very different from any of the steady states produced by feedback.
## IV Discussion
### A Summary
The description of feedback in compound quantum systems (where the output from the system is used to control the evolution of the ancilla, which is reversibly coupled to the system) is greatly simplified if the ancilla can be adiabatically eliminated. We have shown how this can be done for four generic cases, arising from considering two forms of feedback (all optical or coherent, and electro-optical or incoherent) and two types of ancilla (a two-level atom, and an optical mode). The four resulting master equations for the system alone are given below. They are the most important results of this paper. We also include the perturbative expansions of these master equations to third order in the feedback operator $`Z`$. All of the equations are identical to first order in $`Z`$, but differ in second or third order.
For comparison, we begin with simple feedback (that is, with no ancilla) based on detection of the intensity $`I=b_{\mathrm{out}}^{}b_{\mathrm{out}}`$ of the output field $`b_{\mathrm{out}}=b_{\mathrm{in}}+c`$, and using the feedback Hamiltonian
$$H_{\mathrm{fb}}=I(t)Z.$$
(98)
The master equation for this is
$`\dot{\rho }`$ $`=`$ $`i[H,\rho ]+𝒟[e^{iZ}c]\rho `$ (99)
$``$ $`i[H,\rho ]+𝒟[c]\rho +𝒞[Z]𝒥[c]\rho `$ (101)
$`+\left\{{\displaystyle \frac{1}{2}}\left(𝒞[Z]\right)^2+{\displaystyle \frac{1}{6}}\left(𝒞[Z]\right)^3\right\}𝒥[c]\rho .`$
Here $`𝒞[Z]Bi[Z,B]`$ as before. The remaining master equations result from trying to reproduce this form of feedback via an ancilla.
The first master equation derived using adiabatic elimination is for electro-optic feedback via the inversion of a two-level atom:
$`\dot{\rho }`$ $`=`$ $`i[H,\rho ]+{\displaystyle _0^{\mathrm{}}}𝑑qe^q𝒟[e^{iqZ}c]\rho ,`$ (102)
$``$ $`i[H,\rho ]+𝒟[c]\rho +𝒞[Z]𝒥[c]\rho `$ (104)
$`+\left\{(𝒞[Z])^2+(𝒞[Z])^3\right\}𝒥[c]\rho .`$
This differs from the simple feedback master equation (99) at second order in $`Z`$. The second is for electro-optic feedback via one quadrature of an optical mode:
$`\dot{\rho }`$ $`=`$ $`i[H,\rho ]+𝒟[e^{iZ}c]\rho +{\displaystyle \frac{\mathrm{\Gamma }}{ϵ^2}}𝒟[Z].`$ (105)
The expansion of the above equation can be found from that of the simple feedback. The size of the extra second-order term is determined by $`\mathrm{\Gamma }`$, the damping rate for the ancilla mode, and $`ϵ`$, the strength of driving of the ancilla mode.
Turning now to all-optical feedback, we have found that the same master equation arises regardless of whether the feedback is via the inversion of two-level atom or the intensity of an optical mode. It is:
$`\dot{\rho }`$ $`=`$ $`i[H,\rho ]+𝒟\left[\mathrm{exp}\left(2i\mathrm{arctan}{\displaystyle \frac{Z}{2}}\right)c\right]\rho ,`$ (106)
$``$ $`i[H,\rho ]+𝒟[c]\rho +𝒞[Z]𝒥[c]\rho `$ (108)
$`+\left\{{\displaystyle \frac{1}{2}}\left(𝒞[Z]\right)^2+{\displaystyle \frac{1}{4}}𝒞[Z](𝒥[Z]2𝒜[Z])\right\}𝒥[c]\rho .`$
Unlike Eq. (102), this differs from Eq. (99) only at third order in $`Z`$.
### B Relation to Previous Work
As mentioned in the introduction Slosser and Milburn perform adiabatic elimination of the pump mode of a non-degenerate parametric oscillator. In their system the pump mode is driven by the output photocurrent from the idler mode. The procedure they adopt is similar to that contained in Sec. II B of this paper, in that they expand the density matrix in terms of the lower number states of the pump mode. However, in Sec. II C we have already noted that this is not appropriate when dealing with direct detection feedback onto a mode. Higher number states are essential to the description of the system if the feedback strength is large. For this reason they limit the feedback strength to small and moderate values, with a generalization to larger feedback contained in their appendix. This appendix does not explain the origin of the all-orders feedback term. The techniques of adiabatic elimination using QLE’s that are presented in this paper make it easy to treat their system rigorously to all-orders in the feedback strength. The final result, using their definitions and our superoperators, is (with perfect detection assumed):
$`\dot{\rho }`$ $`=`$ $`ϵ[a^{}b^{}ab,\rho ]+2\mathrm{\Gamma }𝒟[ab]\rho +\gamma _a𝒟[a]\rho `$ (111)
$`+\gamma _b𝒟[\mathrm{exp}\left(\chi ab\chi a^{}b^{}\right)b]\rho .`$
Note that the second term here is the one analogous to the final term in our Eq. (105).
Doherty and co-workers consider a strongly interacting system comprised of an atom inside a cavity . The methods used for adiabatic elimination are similar to those used in this paper. They form QLE’s for operators from any of the three subsystems (center of mass motion, internal state and the cavity mode) and then set the time derivatives of, first, the internal state operator and, second, the cavity operator, to zero. They then substitute into the QLE for the momentum operator $`p_x`$. After a conversion to the explicit form of the QLE, they show that the QLE they derive is compatible with the master equation (using their notation)
$$\dot{\rho }=\frac{i}{\mathrm{}}[\frac{p_x^2}{2m},\rho ]+\frac{\kappa }{2}𝒟\left[\mathrm{exp}\left(2i\mathrm{arctan}\frac{Z}{2}\right)\alpha \right]\rho ,$$
(112)
where
$$Z=\frac{g_0^2\mathrm{cos}^2k_Lx}{\mathrm{\Delta }\kappa }.$$
(113)
Note the similarity with our equation resulting from adiabatic elimination of an optical mode where the coupling is via the intensity, (but of course there is no feedback here so our operator $`c`$ is replaced by the $`c`$-number $`\alpha `$). The derivation of this master equation in Ref. is not completely rigorous in that other master equations would also be compatible with the QLE they derive for $`p_x`$. However, it would be straightforward, using the technique we introduced in Sec. II E, to make it rigorous.
The work done on all-optical feedback in this paper follows on from that done by Wiseman and Milburn . They were able to show that all-optical feedback onto a mode could replicate electro-optic homodyne-detection feedback, but they could only prove equivalence with direct-detection feedback to second order. Here, we have shown that this is because the equivalence only holds to second order. We have done this by finding the master equation to all-orders in the feedback strength, and showing it to be of the Lindblad form.
Showing that all-optical feedback via an ancilla (be it a two-level atom or a mode) cannot replicate electro-optical direct detection feedback, leads naturally to the question of whether a more complicated all-optical feedback scheme can replicate direct electro-optic feedback. Since the feedback is replicated to second order, a fruitful approach would seem to be to make the feedback weak, while multiplying up the number of ancillae to compensate. In appendix D we consider the case of $`N`$ ancillae, with coupling to the system scaling as $`1/N`$, where the output of the system is fed sequentially into all of the ancillae. We show that in the limit $`N\mathrm{}`$, this hypothetical all-optical feedback scheme does indeed produce the simple electro-optic feedback master equation (99).
### C Conclusion
We have shown that it is possible to greatly simplify the description of feedback in compound quantum systems by adiabatically eliminating the ancilla, to give master equations for the system alone. In essence, we have found the first order in $`\mathrm{\Gamma }^1`$ effect of the ancilla upon the system, where $`\mathrm{\Gamma }`$ is the ancilla decay rate. We have done this for a variety of ancillae and forms of feedback, and found good agreement with numerical simulations of the dynamics for the full compound quantum system. The master equations in the various cases are quite different, and their range of validity (that is, how large $`\mathrm{\Gamma }`$ has to be for them to be valid) was also found numerically to differ. For the numerical simulations we of course used a particular system, but the equations we derive are very general.
The primary motivation for this work is the reduction of basis size that is necessary to describe the evolution of the system. It is hoped that the derived equations will prove to be helpful to co-workers. However, we note that numerical testing (to find the regime in which these equations are a good approximation) may be necessary to determine when it is appropriate to use them. Apart from these practical advances, we feel that the previously existing confusion in the literature, as discussed in the introduction, has been resolved, and the procedure of adiabatic elimination in compound quantum systems with feedback is now on stable ground.
###### Acknowledgements.
We would like to acknowledge discussions with W.J. Munro and S.M. Tan. This work was supported in part by the Australian Research Council.
## A Proof of Lindblad Form
### 1 Electro-optic Feedback onto a TLA
To show that Eq. (30) can be written in the Lindblad form the following identity will first be established:
$$(\mathrm{\Gamma }𝒞[K])_0^{\mathrm{}}𝑑x𝒥\left[e^{x(\mathrm{\Gamma }+2iK)/2}\right]\rho =\rho .$$
(A1)
Multiplying the equation through by two arbitrary eigenstates of $`K`$, $`\alpha |`$ and $`|\beta `$, from the left and right respectively, the following is obtained:
$$\rho _{\alpha \beta }[\mathrm{\Gamma }+i(\alpha \beta )]_0^{\mathrm{}}𝑑xe^{x[\mathrm{\Gamma }+i(\alpha \beta )]}=\rho _{\alpha \beta }.$$
(A2)
After the simple integration is performed the identity is proved. Before using this the following rearrangement is made:
$$𝒞[Z](1𝒞[Z])^1=(1𝒞[Z])^11.$$
(A3)
Upon use of the identity with $`Z=K/\mathrm{\Gamma }`$ the master equation Eq. (30) becomes Eq. (31).
### 2 All-optical Feedback onto an Atom
In order to show that the master equation can be written as in Eq. (70) it is sufficient to show that
$`\mathrm{exp}\left[2i\mathrm{arctan}\left(Z/2\right)\right]\rho \mathrm{exp}\left[2i\mathrm{arctan}\left(Z/2\right)\right]\rho `$ (A4)
$`=𝒞[Z]𝒥\left[\left(1+iZ/2\right)^1\right]\rho .`$ (A5)
Note that $`𝒥[c]`$ has been omitted as it is a multiplicative factor on both of the superoperators. Consider the following non-Hermitian operator that can be put into a modulus and argument form:
$$\frac{1+iZ/2}{\sqrt{1+\left(Z/2\right)^2}}=re^{i\theta }.$$
(A6)
That $`r=1`$ can be quickly verified. The argument is given by
$$\theta =\mathrm{arctan}\left(Z/2\right).$$
(A7)
By taking the logarithm of A6 an alternative expression for the argument is obtained
$$\theta =\mathrm{ln}\left(\frac{\sqrt{1+\left(Z/2\right)^2}}{1+iZ/2}\right)^i.$$
(A8)
If the logarithmic form of $`\mathrm{arctan}\left(Z/2\right)`$ is used, then the exponentials of Eq. (A5) disappear. The LHS of that equation becomes
$$\frac{1+\left(Z/2\right)^2}{\left(1+iZ/2\right)^2}\rho \frac{\left(1+iZ/2\right)^2}{1+\left(Z/2\right)^2}\rho .$$
(A9)
It is now noted that $`1+\left(Z/2\right)^2=\left(1+iZ/2\right)\left(1iZ/2\right)`$. The above expression can be re-written as
$`{\displaystyle \frac{1iZ/2}{1+iZ/2}}\rho {\displaystyle \frac{1+iZ/2}{1iZ/2}}{\displaystyle \frac{1+iZ/2}{1+iZ/2}}\rho {\displaystyle \frac{1iZ/2}{1iZ/2}}.`$ (A10)
After algebraic manipulation this can be shown to be equal to the RHS of A5, as required.
## B Electro-optic Feedback via an atom with Jaynes-Cummings coupling and detuning
In this section we take the compound system as being a single mode optical cavity, with electro-optic feedback onto a TLA that is placed in the standing wave of the cavity. The Jaynes-Cummings coupling that will be used is $`V=g(a\sigma ^{}+\sigma a^{})`$, with $`g`$ being a real constant and $`a`$ the annihilation operator for the cavity mode. A detuning of $`\delta \sigma ^{}\sigma `$ is also included. The following hierachy of parameters will be investigated:
$$\delta g\mathrm{\Gamma }𝒞[H_\mathrm{s}].$$
(B1)
Of course, $`𝒞[H_\mathrm{s}]`$ is really a superoperator (containing the system Hamiltonian terms) so here we are only referring to its scalar part.
As will be shown, when the the necessary variables are slaved a Hamiltonian term of the form $`g^2a^{}a/\delta `$ is obtained in the final master equation. With the above scaling, this Hamiltonian is not necessarily small compared to $`\mathrm{\Gamma }`$. This makes the adiabatic elimination of the atom more difficult since the presumption that the atomic relaxation time is much shorter than any system time scale is not necessarily true. To do the elimination of the atom rigorously we therefore transform to an interaction picture defined by $`H_0=g^2(a^{}a+\sigma ^{}\sigma )/\delta `$. This transformation has the additional effect of adding a time dependence into the feedback term of the master equation. To nullify this we will start with a time dependent feedback Hamiltonian whose effect, when moved to the interaction picture, is time independent. The master equation in the Schrödinger picture is thus
$`\dot{W}`$ $`=`$ $`𝒞[H_\mathrm{s}]Wi[\delta \sigma ^{}\sigma +g(a\sigma ^{}+\sigma a^{}),W]+\mathrm{\Gamma }𝒟[\sigma ]W`$ (B3)
$`+𝒟[\left\{\sigma \mathrm{exp}(ig^2t/\delta )+\sigma ^{}\mathrm{exp}(ig^2t/\delta )\right\}a]W.`$
In the interaction picture with respect to $`H_0`$ the master equation is
$`\dot{\stackrel{~}{W}}`$ $`=`$ $`U^{}𝒞[H_\mathrm{s}](U\stackrel{~}{W}U^{})Ui[g^2(a^{}a+\sigma ^{}\sigma )/\delta ,\stackrel{~}{W}]`$ (B6)
$`i[g(a\sigma ^{}+\sigma a^{})+\delta \sigma ^{}\sigma ,\stackrel{~}{W}]+\mathrm{\Gamma }𝒟[\sigma ]\stackrel{~}{W}`$
$`+𝒟[\mathrm{exp}(i\pi \sigma _x/2)a]\stackrel{~}{W}.`$
For simplification we will put $`\mathrm{\Delta }=\delta +g^2/\delta `$.
The expansion of Eq. (23) is made, with the $`\rho `$s now understood to be in the interaction picture. The time rates of change are
$`\dot{\rho _0}`$ $`=`$ $`𝒞[\stackrel{~}{H}_\mathrm{s}]\rho _0+𝒥[a]\rho _2𝒜[a]\rho _0+\mathrm{\Gamma }\rho _2`$ (B8)
$`ig(a^{}\rho _1^{}\rho _1a)i[{\displaystyle \frac{g^2a^{}a}{\delta }},\rho _0],`$
$`\dot{\rho _1}`$ $`=`$ $`𝒞[\stackrel{~}{H}_\mathrm{s}]\rho _1+𝒥[a]\rho _1^{}𝒜[a]\rho _1+ig(\rho _0a^{}a^{}\rho _2)`$ (B10)
$`\left({\displaystyle \frac{\gamma }{2}}i\mathrm{\Delta }\right)\rho _1i[{\displaystyle \frac{g^2a^{}a}{\delta }},\rho _1],`$
$`\dot{\rho _2}`$ $`=`$ $`𝒞[\stackrel{~}{H}_\mathrm{s}]\rho _2+𝒥[a]\rho _0𝒜[a]\rho _2\mathrm{\Gamma }\rho _2`$ (B12)
$`ig(a\rho _1\rho _1^{}a^{})i[{\displaystyle \frac{g^2a^{}a}{\delta }},\rho _2].`$
When $`\rho =\rho _0+\rho _2`$ is used, the above equations give
$`\dot{\rho }`$ $`=`$ $`𝒞[\stackrel{~}{H}_\mathrm{s}]\rho +𝒟[a]\rho i[g^2a^{}a/\delta ,\rho ]`$ (B14)
$`ig[a^{},\rho _1^{}]ig[a,\rho _1].`$
In the limit $`\mathrm{\Gamma }𝒞[\stackrel{~}{H}_\mathrm{s}]`$ the amplitudes of $`\rho _1`$ and $`\rho _2`$ respond to changes in the cavity mode much more quickly than $`\rho _0`$. Their equilibrium values are
$`\rho _1`$ $`=`$ $`\left(\mathrm{\Delta }{\displaystyle \frac{\mathrm{\Gamma }}{2i}}\right)^1\left(1{\displaystyle \frac{ig^2}{\mathrm{\Delta }\delta }}\left(1{\displaystyle \frac{\mathrm{\Gamma }}{2i\mathrm{\Delta }}}\right)^1𝒞[a^{}a]\right)^1`$ (B16)
$`\times (g\rho a^{}+g^2\{a^{},\rho _2\}),`$
$`\rho _2`$ $`=`$ $`{\displaystyle \frac{1}{\mathrm{\Gamma }}}\left(1+{\displaystyle \frac{𝒥[a]}{\mathrm{\Gamma }}}{\displaystyle \frac{g^2}{\mathrm{\Gamma }\delta }}𝒞[a^{}a]\right)^1`$ (B18)
$`\times \left\{𝒥[a]\rho ig(a\rho _1\rho _1^{}a^{})\right\}.`$
These two equations can be rearranged to give $`\rho _2`$ in terms of $`\rho _0`$. In the limit $`\mathrm{\Delta }\mathrm{\Gamma }^2`$ and $`g^2\mathrm{\Gamma }\mathrm{\Delta }`$, we find to first order
$$\rho _2=\frac{1}{\mathrm{\Gamma }}\left(1\frac{2g^2}{\mathrm{\Delta }\mathrm{\Gamma }}𝒞[a^{}a]\right)^1𝒥[a]\rho +\frac{g^2}{\mathrm{\Delta }^2}𝒥[a]\rho .$$
(B19)
Using this in Eq. (B16) allows the following master equation to be derived:
$`\dot{\rho }`$ $`=`$ $`𝒞[\stackrel{~}{H}_\mathrm{s}]\rho +𝒟[a]\rho +{\displaystyle \frac{\mathrm{\Gamma }g^2}{\mathrm{\Delta }^2}}𝒟[a]\rho i[{\displaystyle \frac{g^4(a^{})^2a^2}{\mathrm{\Delta }^3}},\rho ]`$ (B21)
$`+{\displaystyle \frac{2g^2}{\mathrm{\Delta }\mathrm{\Gamma }}}𝒞[a^{}a]\left(1{\displaystyle \frac{2g^2}{\mathrm{\Delta }\mathrm{\Gamma }}}𝒞[a^{}a]\right)^1𝒥[a]\rho .`$
In the limit of $`\mathrm{\Delta }\mathrm{\Gamma }^2`$, while still maintaining $`g^2\mathrm{\Gamma }\mathrm{\Delta }`$, the third and fourth terms drop out, leaving the same master equation derived in Sec. II B, with $`Z=2g^2a^{}a/\mathrm{\Gamma }\mathrm{\Delta }`$ of order unity. This is the same limit in which Walls and Milburn arrive at the effective Hamiltonian used in Eq. (19) . The third and fourth terms correspond to, respectively, an increased damping rate and a $`\chi ^{(3)}`$ nonlinearity for the cavity mode.
Note that the derived Hamiltonian term which threw doubt upon the adiabatic elimination process has been canceled. Of course, when we return to the Schrödinger picture it will reappear, leaving a different master equation from that of Sec. II B. The solution is to start with an extra Hamiltonian term of the form $`g^2a^{}a/\mathrm{\Delta }`$ when using the Jaynes-Cummings coupling. A transformation to the interaction picture is then not required, nor is the time dependence in the feedback.
## C Unitary Transformation of the Total Master Equation for Electro-optic Feedback Onto a Mode
In this appendix the total master equation for electro-optic feedback onto a cavity is unitarily transformed so that the amplitude of the driven cavity may be treated classically, thus reducing the necessary basis size. The implicit form of the master equation will be used as this proves to be more straightforward. That is to say, the photocurrent will be approximated by a slightly smoothed version of $`dN/dt`$. Before transformation the implicit master equation is
$`\dot{W}`$ $`=`$ $`i[ga^{}a(b+b^{})+{\displaystyle \frac{i\lambda }{4}}\{(a^{})^2a^2\},W]+𝒟[a]W`$ (C2)
$`{\displaystyle \frac{iϵ}{2}}[(ib+ib^{})I(t),W]+\mathrm{\Gamma }𝒟[b]W.`$
We now put $`\stackrel{~}{W}=UWU^{}`$ where,
$$U=\mathrm{exp}[ϵf(t)(bb^{})/2]$$
(C3)
and $`f(t)`$ is given in Eq. (93). The unitarily transformed master equation is given by
$$\dot{\stackrel{~}{W}}=\dot{U}U^{}\stackrel{~}{W}+\stackrel{~}{W}U\dot{U}^{}+U(U^{}\stackrel{~}{W}U)U^{},$$
(C4)
where $`\dot{W}=W`$. Now, $`\dot{U}`$ is given by
$$\dot{U}=\frac{ϵ}{2}(bb^{})[\mathrm{\Gamma }f(t)/2+I(t)]U,$$
(C5)
thus the first two terms of Eq. (C4) give
$$\frac{ϵ}{2}\left[\mathrm{\Gamma }f(t)/2+I(t)\right][bb^{},\stackrel{~}{W}].$$
(C6)
The last term of Eq. (C4) will only cause a change to terms that are dependent upon the driven cavity operators. Thus, the non-linear driving and the damping of the system may be ignored for the present. The expression that needs to be simplified contains three terms (damping, coupling, and feedback), which can be evaluated using $`UbU^{}=ϵf(t)/2+b`$. The damping term is
$`\mathrm{\Gamma }U\left\{𝒟[b](U^{}\stackrel{~}{W}U)\right\}U^{}`$ (C7)
$`=\mathrm{\Gamma }\left(𝒟[b]\stackrel{~}{W}+ϵf[bb^{},\stackrel{~}{W}]/4\right),`$ (C8)
the coupling term is
$`iU[ga^{}a(b+b^{}),U^{}\stackrel{~}{W}U]U^{}`$ (C9)
$`=i[ga^{}a(b+b^{}+ϵf),\stackrel{~}{W}],`$ (C10)
and the feedback term is
$`{\displaystyle \frac{iU}{2}}[ϵ(ib+ib^{})I(t),U^{}\stackrel{~}{W}U]U^{}`$ (C11)
$`={\displaystyle \frac{ϵ}{2}}[I(t)(bb^{}),\stackrel{~}{W}].`$ (C12)
Adding up the contributions from the damping, feedback, coupling, Eq. (C6) and also the system Hamiltonian, the following is obtained:
$`\dot{\stackrel{~}{W}}`$ $`=`$ $`i[ga^{}a(b+b^{}+ϵf)+{\displaystyle \frac{i\lambda }{4}}\{(a^{})^2a^2\},\stackrel{~}{W}]`$ (C14)
$`+\mathrm{\Gamma }𝒟[b]\stackrel{~}{W},`$
which is the master equation after transformation. There is now no distinction between the implicit and explicit forms as $`f`$ is a bounded function.
## D All-optical Feedback onto an Infinite Number of Cavities
It is of interest whether all-optical feedback can ever have the same effect on a system as simple electro-optic feedback. In this section we show that this can be achieved with a very large number of ancilla cavities that are coupled back to the system. The basic idea of the all-optical feedback remains the same, in that the output of one cavity becomes the input to the next cavity. The cavities all have the same damping co-efficient, $`\mathrm{\Gamma }`$. Damping of the system is set equal to unity. For large damping, the infinite number of ancilla cavities will be adiabatically eliminated. See Fig. 12.
The form of the coupling of the $`j^{\mathrm{th}}`$ cavity is similar to that for the all-optical feedback via a single mode. It is
$$V_j=\frac{Kb_j^{}b_j}{N},$$
(D1)
where $`N`$ is the total number of driven cavities, $`K`$ is proportional to an Hermitian system operator and $`b_j`$ is the annihilation operator of the $`j^{\mathrm{th}}`$ cavity. The input and output fields to and from the $`j^{\mathrm{th}}`$ cavity are represented by $`u_{\mathrm{in}(j)}`$ and $`u_{\mathrm{out}(j)}`$. This means that $`u_{\mathrm{in}(j+1)}=u_{\mathrm{out}(j)}`$. The output field from the system is given by $`v_{\mathrm{out}}`$.
From Eq. (73) the contribution to the QLE for an arbitrary operator due to the shining of the $`(j1)^{\mathrm{th}}`$ output field onto the $`j^{\mathrm{th}}`$ cavity can be found. The total QLE if there are $`N`$ driven cavities is
$`dr`$ $`=`$ $`𝒟[c^{}]rdt[dV_{\mathrm{in}}^{}cc^{}dV_{\mathrm{in}},r]+\mathrm{\Gamma }{\displaystyle \underset{j=1}{\overset{N}{}}}𝒟[b_j^{}]rdt`$ (D4)
$`\sqrt{\mathrm{\Gamma }}{\displaystyle \underset{j=1}{\overset{N}{}}}[dU_{\mathrm{in}(j)}^{}b_jb_j^{}dU_{\mathrm{in}(j)},r]`$
$`+i[H_\mathrm{s},r]dt+i{\displaystyle \underset{j=1}{\overset{N}{}}}b_j^{}[{\displaystyle \frac{K}{N}},r]b_jdt,`$
where $`dU_{\mathrm{in}(j)}=u_{\mathrm{in}(j)}dt`$. Note that the same idea as in Eq. (79) has been used. The QLE for a system operator is
$`ds`$ $`=`$ $`𝒟[c^{}]sdt[dV_{\mathrm{in}}^{}cc^{}dV_{\mathrm{in}},s]`$ (D6)
$`+i{\displaystyle \underset{j=1}{\overset{N}{}}}b_j^{}[{\displaystyle \frac{K}{N}},s]b_jdt+i[H_\mathrm{s},s]dt.`$
Although the input fields $`u_{\mathrm{in}(j)}`$ obviously depend on system operators, they are evaluated at a slightly earlier time due to the small, but finite, propagation time of the field from the system to the driven cavities. The system operator $`s`$, therefore, commutes with $`dU_j`$.
We now note that the QLE for $`b_j`$ has the same form as Eq. (76)
$$db_j=\left(\frac{\mathrm{\Gamma }b_j}{2}+\sqrt{\mathrm{\Gamma }}u_{\mathrm{in}(j)}+\frac{iKb_j}{N}\right)dt.$$
(D7)
Also,
$$u_{\mathrm{in}(j+1)}=u_{\mathrm{in}(j)}+\sqrt{\mathrm{\Gamma }}b_j.$$
(D8)
To simplify matters only the non-stochastic part of $`b_j`$ and $`u_{\mathrm{in}(j)}`$ will be considered in the derivation of the master equation for the system. This can be justified by mathematical induction. Suppose that $`u_{\mathrm{in}(j)}`$ and $`b_j`$ can both be grouped into stochastic terms linearly dependent upon $`v_{\mathrm{in}}`$ and non-stochastic terms. Then it is clear from Eq. (D8) that $`u_{\mathrm{in}(j+1)}`$ can also be grouped in such a manner. Therefore, in the limit in which Eq. (D7) can be slaved to produce the equivalent of Eq. (78) it can be seen that $`b_{j+1}`$ will consist of non-stochastic terms, arising from the non-stochastic terms of $`u_{\mathrm{in}(j)}`$ and $`b_j`$, as well as stochastic terms linear in $`v_{\mathrm{in}}`$. To complete the mathematical induction, $`b_1`$ and $`u_{\mathrm{in}(1)}`$ can obviously grouped in the manner suggested. Now, terms in $`b_j`$ that go as $`v_{\mathrm{in}}`$ will annihilate onto the vacuum state when the expectation value of Eq. (D6) is taken, thus, the stochastic parts can be ignored as they give a zero contribution.
An expression for non-stochastic part of $`u_{\mathrm{in}(j)}`$ (denoted by $`\overline{u}_{\mathrm{in}(j)}`$) needs to be found in order to evaluate $`\overline{b}_j`$. Using Eq. (D8) and the slaved value $`\overline{b}_{j1}`$ it is found to be
$$\overline{u}_{\mathrm{in}(j)}=\left(\frac{1+2iK/\mathrm{\Gamma }N}{1+2iK/\mathrm{\Gamma }N}\right)^{j1}c.$$
(D9)
Substituting into Eq. (D7) gives $`b_j`$. This is then used in Eq. (D6). Writing $`Z=4K/\mathrm{\Gamma }`$, the summation term is
$`{\displaystyle \frac{i}{N}}𝒥\left[c^{}\left(1iZ/2N\right)^1\right]`$ (D10)
$`\times {\displaystyle \underset{j=0}{\overset{N}{}}}\left({\displaystyle \frac{1iZ/2N}{1iZ/2N}}\right)^j[Z,s]\left({\displaystyle \frac{1+iZ/2N}{1+iZ/2N}}\right)^j.`$ (D11)
Firstly, the quotients are expanded to second order in $`Z/N`$. Then the contributions from the first and second orders are factorized, with the latter expanded to first order in $`j/N^2`$. This gives
$`{\displaystyle \frac{i}{N}}𝒥\left[c^{}\left(1iZ/2N\right)^1\right]`$ (D12)
$`\times {\displaystyle \underset{j=0}{\overset{N}{}}}(1+iZ/N)^j(1{\displaystyle \frac{j}{2}}(Z/N)^2)[Z,s]`$ (D13)
$`\times \left(1iZ/N\right)^j\left(1{\displaystyle \frac{j}{2}}\left(Z/N\right)^2\right).`$ (D14)
It is not difficult to show that the contribution of the $`j/N^2`$ terms is small (of order $`N^1`$). Also, $`(1+iZ/2N)^11`$, so we can now write the summation as
$`i{\displaystyle \underset{j=0}{\overset{N}{}}}𝒥\left[c^{}\left(1+iZ/N\right)^j\right][Z/N,s]`$ (D15)
$``$ $`{\displaystyle \underset{j=0}{\overset{N}{}}}𝒥\left[c^{}\left(1+iZ/N\right)^j\right]\left(𝒥\left[\left(1+iZ/N\right)\right]1\right)s`$ (D16)
$`=`$ $`𝒥[c^{}]{\displaystyle \underset{j=0}{\overset{N}{}}}\left(𝒥\left[\left(1+iZ/N\right)^{j+1}\right]𝒥\left[\left(1+iZ/N\right)^j\right]\right)s`$ (D17)
$``$ $`𝒥[c^{}]\left(𝒥\left[\left(1+iZ/N\right)^N\right]1\right)s`$ (D18)
$``$ $`𝒥[c^{}](𝒥[\mathrm{exp}(iZ)]1)s.`$ (D19)
Terms of order $`N^1`$ have been ignored as the limit $`N\mathrm{}`$ has been taken. Returning to Eq. (D6) gives the same QLE as for simple feedback. In this rather impractical way, all-optical feedback can replicate electro-optic feedback.
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# THE PARTON CONTENT OF THE PHOTON
## 1 Introduction
The photon is the gauge boson of the electromagnetic force and therefore has fundamental direct couplings to all charged particles. However, it can fluctuate into hadronic matter before it takes part in an interaction, thereby giving rise to a wide range of phenomena which can be quantified in observables. Once measured, they can be used to improve our understanding of both perturbative and non-perturbative QCD. In this paper the experimental environment is provided by photon-photon collision events recorded by the four LEP collaborations. Scattering events of this type are commonly characterised by the absolute four momentum squared of the photons participating in the interaction, $`Q^2`$ and $`P^2`$, the invariant mass $`W_{\gamma \gamma }`$ of the resulting hadronic final state and the Bjorken variable $`x=Q^2/(Q^2+W_{\gamma \gamma }^2+P^2)`$. For large enough virtuality $`Q^2(P^2)`$ one (both) electron(s) has a scattering angle large enough to be detected in the experiment. Such events are referred to as single (double) tag events. The observables studied are the structure function of the photon $`F_2^\gamma `$ at low $`x`$ and the production of heavy quarks in untagged and single-tagged photon-photon scattering events.
## 2 The Photon Structure Function $`F_2^\gamma `$
The differential cross section of deep-inelastic electron-photon scattering, as can be obtained from single tag events, can be expressed as
$$\frac{d^2\sigma _{\mathrm{e}\gamma \mathrm{eX}}}{dxdQ^2}=\frac{2\pi \alpha \text{ }^2}{xQ^4}\left[\left(1+(1y)^2\right)F_2^\gamma (x,Q^2)y^2F_\mathrm{L}^\gamma (x,Q^2)\right],$$
(1)
where $`\alpha `$ is the fine structure constant and $`y`$ is the inelasticity. In the region of small $`y`$ under study the contribution from $`F_\mathrm{L}^\gamma (x,Q^2)`$ can be neglected and one obtains the photon structure function $`F_2^\gamma `$. The latest measurements in the low-$`x`$ region come from LEP, namely L3 and OPAL, as shown in Figure 1 $`^\mathrm{?}`$. For mean values of $`Q^2`$ = 1.9 $`\mathrm{GeV}^2`$ values of $`x`$ down to $`210^3`$ have been reached. Particular emphasis is put on the behaviour of $`F_2^\gamma `$ at low $`x`$, since, driven by the hadron-like properties of the interacting photon, one suspects a rise of the photon structure function towards small $`x`$, similar to the rise observed for the proton structure function $`F_2^p`$.
The precision of the measurement at low $`x`$ has been considerably improved and there appears to be mounting evidence for the expected rise of $`F_2^\gamma `$ at low $`x`$ although a final conclusion should not be drawn yet. The GRV LO and SaS1D parameterisations are generally consistent with the data in all the accessible $`x`$ and $`Q^2`$ regions, with the exception of the measurement at the lowest scale, $`Q^2=1.9`$ $`\mathrm{GeV}^2`$, where GRV is too low. In contrast, the naive quark-parton model is not able to describe the data for $`x<0.1`$. These results show that the photon must contain a significant hadron-like component at low $`x`$.
## 3 Heavy Quark Production
Charged $`\mathrm{D}^{}`$ mesons provide a clean tag to study open charm production in photon-photon collisions. The inclusive cross-section for the production of $`\mathrm{D}^{}`$ mesons can be calculated in NLO perturbative QCD (pQCD) and this process is therefore suitable to test the validity of the theory in this region. L3 and OPAL have determined the differential cross-sections $`\mathrm{d}\sigma /\mathrm{d}p_\mathrm{T}^\mathrm{D}^{}`$ in anti-tagged $`\text{e}^+\text{e}^{}`$$``$ $`\text{e}^+\text{e}^{}`$$`\mathrm{D}^{}X`$ events as a function of the transverse momentum $`p_\mathrm{T}^\mathrm{D}^{}`$ as shown in Figure 2$`^{\mathrm{?},\mathrm{?}}`$. Both measurements can been seen to be in agreement with each other within the experimental uncertainties. They are compared to a NLO calculation by Binnewies et al. $`^\mathrm{?}`$ using the massless approach, and by Frixione et al. $`^\mathrm{?}`$ using the massive approach. It is found that despite the low values of $`p_\mathrm{T}^\mathrm{D}^{}`$ studied the massless calculation is in good agreement with the data. The massive calculation agrees with the measured cross-section for $`p_\mathrm{T}^\mathrm{D}^{}>3\mathrm{GeV}`$ but underestimates the data for lower values of $`p_\mathrm{T}^\mathrm{D}^{}`$.
The total cross-section of the process $`\text{e}^+\text{e}^{}`$$``$ $`\text{e}^+\text{e}^{}`$$`\text{c}\overline{\text{c}}`$, where the charm quarks are produced in the collision of two quasi-real photons, can be deduced from measurements like the above by extrapolating to the full phase space using MC model predictions. The results for various experiments are shown in Figure 2$`^{\mathrm{?},\mathrm{?}}`$. All results are in agreement with each other where comparable, and are well described by the NLO calculation of Ref. $`^\mathrm{?}`$ over the whole range of $`\text{e}^+\text{e}^{}`$centre-of-mass energies covered. Also shown in this figure is the total $`\text{b}\overline{\text{b}}`$ cross-section in photon-photon collisions as reported for the first time by L3 $`^\mathrm{?}`$ to be $`\sigma ^{eeeebbX}=9.9\pm 2.9(\mathrm{stat})\pm 3.8(\mathrm{syst})`$ pb. This first sign of $`b`$ production in this process has been obtained by exploiting the more energetic leptons from $`b`$ semi-leptonic decays as opposed to those from charm semi-leptonic decays. The measured value is somewhat above the QCD prediction, although the still sizable experimental uncertainties inhibit a clear conclusion here.
The first measurement of the charm structure function $`F_{2,\mathrm{c}}^\gamma (x,Q^2)`$ of the photon has been performed by OPAL $`^\mathrm{?}`$ based on about 30 $`\mathrm{D}^{}`$ mesons reconstructed in single-tagged events. The value of $`F_{2,\mathrm{c}}^\gamma (x,Q^2)`$ is determined for an average $`Q^2`$ of $`20\mathrm{GeV}^2`$ and in two regions of $`x`$, $`0.0014<x<0.1`$ and $`0.1<x<0.87`$. For $`x>0.1`$, the perturbative NLO calculation of Laenen et al. $`^\mathrm{?}`$ is in good agreement with the measurement. For $`x<0.1`$, the measurement suffers from large uncertainties of the invisible cross-section predicted by the HERWIG and Vermaseren Monte Carlo models, and therefore the result is not very precise. However, the data clearly suggest a non-zero hadron-like component of $`F_{2,\mathrm{c}}^\gamma `$.
## 4 Conclusion
Data is available now which explores the structure function $`F_2^\gamma `$ of the photon in deep-inelastic electron-photon scattering down to values of $`x210^3`$. Thanks to the improved precision of the measurements there appears to be mounting evidence for a rise of $`F_2^\gamma `$ at low $`x`$, but no final conclusion can be drawn yet. Open charm production is studied via the measurement of charged $`\mathrm{D}^{}`$-mesons and found to be in good agreement with NLO perturbative QCD calculations for both differential cross-sections and the total charm cross-section measured. A first sign of $`b`$ production in photon-photon collisions has been reported and is seen to be somewhat higher than the QCD prediction, although experimental uncertainties are still large. A first measurement of the charm structure function has been performed in single tagged events and has been found to be in good agreement with NLO pQCD for $`x>0.1`$. A hadron-like component of $`F_{2,\mathrm{c}}^\gamma `$ at low $`x`$ is clearly favoured by the data.
## Acknowledgements
I would like to thank the CERN SL Division for their successful operation of the accelerator and the four LEP collaboration for providing me with their results. Many thanks also to S. Söldner-Rembold and R. Nisius for providing me with collective figures on the available data and phenomenology.
## References
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# Gravitational Waves from Mesoscopic Dynamics of the Extra Dimensions
## I Mesoscopic Cosmology of the Extra Dimensions
Recently there has been considerable interest in models of spacetimes with relatively large extra dimensions, in which the familiar Standard Model fields are confined to a “brane”, a 3-dimensional defect in the larger space, while gravity propagates in all the dimensions (the “bulk”). In some models, the extra dimensions can be as large as the current direct experimental gravitational probes of the order of a millimeter. The apparent (usual) Planck mass in 3+1D, $`M_{Planck}`$, is related to the true fundamental scale $`M_{}`$ by $`M_{Planck}^2M_{}^2(M_{}^nb^n)`$ where $`b^n`$ is the volume of $`n`$ extra large dimensions. In another class of models (“nonfactorizable geometries”, ), the extra dimensions can be even larger, but have curvature $`k`$ which traps gravitons in a bound state close to a brane. The curvature radius $`bk^1M_{Planck}^2/M_{}^3`$ is again on a mesoscopic scale which may be as large as $`1`$ mm.
The cosmology of these models might well include radical departures from the orderly evolution of the standard model, including episodes of violent mesoscopic geometrical activity at relatively late times. New classical geometrical effects remain important long after inflation, until the Hubble length $`H^1M_{Planck}/T^2`$ is larger than the size or curvature radius $`b`$ of the largest extra dimensions. Of particular interest are two new geometrical degrees of freedom common to many of these models: “radion” modes controlling the size or curvature of the extra dimensions, and new Nambu-Goldstone modes corresponding to inhomogeneous displacements of the brane in the extra dimensions. Previous treatments have focused on microscopic thermal and quantum emission in these modes, but cosmological symmetry breaking can also create large-amplitude, coherent classical excitations on much larger scales (of order $`H^1`$) as the configuration of the extra dimensions and the position of the brane settle into their present state. This highly dynamic geometric activity generically produces a conspicuous detectable relic: an intense, classically generated background of gravitational radiation. In this paper I explore the possibility of a direct detection of a gravitational wave background generated by coherently excited radion and Nambu-Goldstone modes on the threshold of standard cosmology.
## II Hubble Frequency and Energy Equipartition
The main features of the background spectra can be estimated from general scaling considerations, without reference to particular models.
### a Frequency:
The characteristic gravitational (“Hubble”) frequency redshifted to the present day is $`f_{H0}(T)H(T)a(T)/a_0`$ or:
$$f_{H0}=7.65\times 10^5\mathrm{Hz}(T/TeV)g_{}^{1/6}(g_{}/g_S)^{1/3}T_{2.728}=9.37\times 10^5\mathrm{Hz}(H\times 1\mathrm{m}\mathrm{m})^{1/2}g_{}^{1/12}(g_{}/g_S)^{1/3}T_{2.728}$$
(1)
The estimate is valid back to the threshold of the extradimensional dynamics, $`H^1b`$. The mapping between observed frequency, temperature $`T`$ and Hubble length is shown by axis labels in figure 1 for frequencies in the LISA band.<sup>*</sup><sup>*</sup>*Note the weak dependence on the particle-physics uncertainties encapsulated in the number of effective relativistic degrees of freedom, defined in the usual way $`g_{}_{bosons}(T_i/T)^4g_i+(7/8)_{fermions}(T_i/T)^4g_i`$ and $`g_S_{bosons}(T_i/T)^3g_i+(7/8)_{fermions}(T_i/T)^3g_i`$). Even above 100 GeV we expect standard relativistic cosmology derived from relativity and thermodynamics to hold(e.g.), $`H=2.07\times 10^5\mathrm{Hz}g_{}^{1/2}(T/GeV)^2`$ and $`a(T)/a_0=3.70\times 10^{13}T_{2.728}g_S^{1/3}(GeV/T)`$. Extra dimensions between $`b1`$ mm and $`10^6`$ mm produce backgrounds peaked in the LISA band ($`10^1`$ to $`10^4`$ Hz); observations with LIGO (up to $`1000`$Hz) could detectThis requires high enough sensitivity, $`h_{rms}10^{22}`$, to detect waves with $`\mathrm{\Omega }_{GW}\mathrm{\Omega }_{rel}`$, which may be possible with future instrumentation. activity from dimensions down to $`b10^{14}`$mm.
### b Amplitude:
The mechanisms considered here excite geometrical degrees of freedom in approximate equipartition of energy density with the thermal relativistic matter. A stochastic background of gravitational radiation with rms metric perturbations $`h_{rms}(f)`$ over bandwidth $`\mathrm{\Delta }f`$ contributes a fraction of the critical density $`\mathrm{\Omega }_{GW}(f,\mathrm{\Delta }f=f)(f/H)^2h_{rms}^2(f,\mathrm{\Delta }f=f)`$. In an experiment such as LISA a stochastic background can only be distinguished from other sources of noise and astrophysical wave sources by resolving the background in frequency (and to some extent, direction and polarization). The rms strain produced in a LISA frequency resolution element is
$$h_{rms}(f,\mathrm{\Delta }f)=4.74\times 10^{20}f_{mHz}^{3/2}T_{2.728}^2h_{70}^0(\mathrm{\Delta }f/3\times 10^8\mathrm{Hz})^{1/2}[\mathrm{\Omega }_{GW}(\mathrm{\Delta }f=f)/\mathrm{\Omega }_{rel}]^{1/2}.$$
(2)
The reference density is set by the mean energy density in all relativistic species (photons and three massless neutrinos), $`\mathrm{\Omega }_{rel}=8.51\times 10^5h_{70}^2T_{2.728}^4`$, where $`h_{70}`$ refers to the Hubble constant. Since the energy density of gravitational waves redshifts like relativistic matter, this gives a maximal bound for primordial broad-band backgrounds, shown in figure 1 along with projected LISA sensitivity.A chaotic cosmology is almost unconstrained by other data if it happens early enough. Scalar inhomogeneities (including entropy perturbations) from subhorizon events earlier than about 100 GeV are erased by neutrino and nucleon diffusion before nucleosynthesis begins, so gravitational waves may be the only direct surviving relic. Overproduction of black holes or other stable relics must be avoided but these are very model-dependent. Existing data constrain gravitational waves only at much lower frequencies; above $`f_{msp}=4.4\times 10^9`$Hz (corresponding to $`f_{H0}(T60\mathrm{M}\mathrm{e}\mathrm{V},H^1453\mathrm{k}\mathrm{m}))`$, millisecond pulsar timing gives a limit at 90% confidence, $`\mathrm{\Omega }_{GW}<1.14\times 10^4(f/f_{msp})^2T_{2.728}^4h^0\mathrm{\Omega }_{rel}.`$
## III Turbulent Flow from a First-Order Radion Transition
The radion can be a significant source of mesoscopic activity if its potential has a first-order phase transition. The radion is stuck initially in a metastable state with some initial value $`b_i`$ and thermal or quantum nucleation leads to randomly nucleated regions corresponding to the final value $`b_0`$, accompanied by a release of internal energy. The gravitational waves from this mode are easiest to describe if for the largest extra dimension, $`b_0H^1(T=M_{})`$: the final dimensional stabilization then happens within the context of an approximately 3+1-dimensional cosmology.
On our brane this process resembles nucleation of bubbles or vacuum domains of Higgs scalars, with the extra complication of modifications in gravity. In an expanding universe the bubbles collide and overlap, creating flows of energy with velocities of the order of unity. The coherence scale $`R`$ of the flows is determined by nucleation dynamics which generally yields $`R\mathrm{log}(HT)/H10^2H^1`$. A number of models have been used to estimate the gravitational radiation in similar situations, from colliding bubbles and the resulting energy flows in the context of QCD and electroweak phase transitions. In the absence of a definite model of radion bubbles, we estimate the maximal gravitational wave spectrum, based on establishing a sustained relativistic turbulent cascade up to the scale $`R`$ for a time $`H^1`$.
The power output of a system in gravitational waves, $`L_{GW}0.1L_{internal}^2/L_0`$, is determined by the (changing quadrupolar) flow of mass-energy $`L_{internal}`$, where $`L_0=c^5/G=M_{Planck}^2`$ and the numerical factor $`0.1`$ is typical of simple asymmetric geometries. Flows of mass-energy on the Hubble scale produce a gravitational wave power per volume close to $`10^1H^3L_0`$, and integrated over time $`H^1`$ produce broad-band $`\mathrm{\Omega }_{GW}10^1\mathrm{\Omega }_{rel}`$ at characteristic frequency $`H`$ (now shifted to $`f_{H0}`$). Flows on smaller scale $`R`$ create a spectrum peaked at higher frequency $`f_{peak}f_{H0}(RH)^1`$ and with a smaller amplitude. Relativistic motions in a volume $`R^3`$ with density perturbations of order unity, $`MR^3`$ create a mass-energy flow $`L_{internal}(M/R)R^2`$ which if sustained for time $`H^1`$ gives $`\mathrm{\Omega }(\mathrm{\Delta }f=f)Rf_{peak}^1`$. In a narrow band this translates to amplitude $`h_{rms}(f,\mathrm{\Delta }f=3\times 10^8\mathrm{Hz})\mathrm{\Omega }(\mathrm{\Delta }f=f)^{1/2}(\mathrm{\Delta }f/f)^{1/2}(H/f)f_{peak}^2`$. To estimate the low-frequency spectrum consider motions of smaller velocity on the same length scale, $`vRf<1`$. The mass-energy flow is $`L_{internal}Mv^2(v/R)v^3`$, hence the gravitational wave power $`f^6`$. The maximal tail of low-frequency waves then has $`\mathrm{\Omega }(\mathrm{\Delta }f=f)f^6`$, hence $`h_{rms}(f,\mathrm{\Delta }f=3\times 10^8\mathrm{Hz})f^{3/2}`$. These estimates lead to the maximal spectrum shown in figure 1 for $`R,T=0.1H^1,10`$TeV and $`R,T=0.01H^1,100`$GeV. (In the quieter case where turbulence is not sustained, the amplitude is less than this estimate by $`(RH)^{1/2}`$.)
## IV Gravitational Waves from Brane Displacement
In many scenarios there is another degree of freedom, the position of the brane. Spatial inhomogeneities in this displacement correspond to nearly-massless Nambu-Goldstone modes. In the cosmological formation of the brane/defect, the position is in general a random variable uncorrelated on large scales, and large-scale modes are excited by the Kibble mechanism. Since the modes also represent curvature of the brane, they are efficiently coupled into gravitational waves as viewed on the brane. This mode is generally excited if for the largest extra dimension, $`b_0H^1(T=M_{})`$: the 3-brane condenses as a defect before cosmology enters the 3+1-D era.
Suppose our brane forms at some early time as a defect in $`3+n+1`$ space, spontaneously breaking the Poincaré symmetry of the full theory. The position of the brane in each higher dimension $`j`$ can be described by a field $`y_j(x_i)`$ which depends on the coordinate $`x_i`$ on the brane.<sup>§</sup><sup>§</sup>§The $`y_j`$’s are analogous to group-transformation angles $`\theta `$ in the Goldstone description of the pion or axion, but have dimensions of length. They are related to the canonically-normalized physical fields $`\pi `$ by $`y=\pi /f_b^2`$ (rather than $`\theta =\pi /f_b`$) where $`f_b^4`$ is the brane tension. When a $`y`$ is “eaten” it gets a mass $`m_yf_b^2/M_{Planck}`$ or $`(1\mathrm{m}\mathrm{m})^1`$ for $`f_b`$TeV. These fields represent propagating Nambu-Goldstone bosons, new modes in addition to the Standard Model fields on the brane and gravitation propagating in the bulk.
In a cosmological setting, the initial values of the $`y_j`$’s are not generally correlated on large scales but are random for points with large separation, since the position of the brane, when it condenses as a defect, is determined by local accidents. A topologically stable 3-wall forms for example if vacua in a 4-bulk fall into two discrete degenerate minima; the spontaneous choice of one minimum or the other is a random variable at large separation, so the wall is not formed initially in its (flat) ground state, but in an excited (wrinkled) state. This is the same “Kibble mechanism” which leads in other contexts to formation of cosmic strings, axion miniclusters, or Goldstone excitations of scalar fields (which are also gravitational wave sources, e.g. ).
If the cosmology thereby excites these modes, the brane is wrinkled on all scales up to the dimensional size or curvature radius $`b_0`$ (which we take to be fixed here), curving through the extra dimensions. Ignoring the subtleties of the cosmological metric and imagining the brane in a flat spacetime, displacements directly induce an effective 4-metric on the brane $`g_{\mu \nu }=\eta _{\mu \nu }+h_{\mu \nu }`$, where the perturbed 4-metric is $`h_{\mu \nu }=_{j=1}^n_\mu y_j_\nu y_j`$. Although these initial perturbations are scalars, no symmetry prevents transverse-traceless tensor components from being excited dynamically for waves with $`fH`$, creating classical gravitational waves. Kibble excitations corresponding to perturbations of magnitude $`\delta y`$ on scale $`f^1`$ thereby lead to wave amplitudes of the order of $`h_{rms}(f)|h_{jk}^{TT}||\delta yf|^2.`$
In the cosmological context, the Kibble excitation is regulated by the expansion. The brane location is initially uncorrelated on large scales, so on the Hubble scale variations in $`y`$ are of order $`H^1`$, producing waves with amplitude $`h_{rms}H^2\delta y^2O(1)`$, comparable in total energy with the other forms of energy on the brane ($`\rho _{rel}H^2M_{Planck}^2`$) which now appear in $`\mathrm{\Omega }_{rel}`$. These waves show up today as gravitational waves at the redshifted Hubble frequency $`f_{H0}(H)`$. The maximal spectrum corresponds to a background of the order of the energy density in brane fields (that is, $`h=O(1)`$ on each scale as it enters the horizon), leading to constant $`\mathrm{\Omega }(GW)`$ over a range of higher frequencies. The amplitude of excitation is reduced after the universe enters the classical 3+1-D era, when $`H^1b_0`$ (or after $`Hm_y`$); variations in $`y`$ are at most of order $`b_0`$ (or $`m_y`$), leading to $`h_{rms}(Hb_0)^2`$ and to a spectrum with a corresponding $`hf^{5/2}`$ rolloff at lower frequencies. The damping scale and low-frequency spectrum contain information on the scale of extra dimensions and/or the stabilization of the brane. The spectrum at higher frequencies is a probe of the cosmological model during the multidimensional/preclassical period, including the excitation mechanism.
The predicted backgrounds can be distinguished from other expected signals by their isotropy, and by the distinctive spectra with rolloffs at both low and high frequencies . Under some circumstances these backgrounds could be stronger than previously contemplated sources.
###### Acknowledgements.
I am grateful for useful discussions with E. Adelberger, P. Bender, D. Kaplan, A. Nelson, S. Sigurdsson, E. Witten, and especially G. Fuller. I thank the Aspen Center for Physics, the Max-Planck-Institut für Astrophysik, the Isaac Newton Institute for Mathematical Sciences and the Ettore Majorana Centre for Scientific Culture for hospitality. This work was supported at the University of Washington by NASA, and at the Max-Planck-Institute für Astrophysik by a Humboldt Research Award.
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# The 19-Vertex Model at critical regime |𝑞|=1
## 1 Introduction
In this paper we shall study the 19 vertex model associated with the quantum affine symmetry $`U_q(\widehat{sl_2})`$. The 19-vertex model is a higher spin generalization of the celebrated 6-vertex model, whose Boltzmann weights are given in (2.11). For the massive parameter case $`|q|<1`$, M.Idzumi derived the integral representations of correlation functions from the viewpoint of the representation thory of the quantum group $`U_q(\widehat{sl_2})`$. In this paper we shall consider the problem at the critical regime $`|q|=1`$, where the representation theory of the quantum group $`U_q(\widehat{sl_2})`$ cannot be used. M.Jimbo, H.Konno and T.Miwa studied the 6-vertex model at critical regime $`|q|=1`$. They presented the free boson realizations of the vertex operators and gave the trace constructions of the solutions of the quantum Knizhnik-Zamolodchikov equations, which represent the correlation functions. In this paper we shall give the integral representations for the correlation functions of the 19-vertex model at critical regime $`|q|=1`$. In order to give trace construction of the correlation functions we need the free field realization of the vertex operators. In this paper we present the free field realizations of the vertex operators in terms of free bosons and free fermions.
The critical 19 vertex model is a limiting case of the fusion 8-vertex model. The latter is massive, and the corner transfer method can be applied . Therefore we can conclude that the correlation functions of the critical 19-vertex model are governed by the following system of difference equations.
$`G_{2N}(\mathrm{},\beta _{j+1},\beta _j,\mathrm{})_{\mathrm{}ϵ_{j+1}ϵ_j\mathrm{}}`$
$`=`$ $`{\displaystyle \underset{ϵ_j^{}ϵ_{j+1}^{}=0,1,2}{}}R_{ϵ_jϵ_{j+1}}^{ϵ_j^{}ϵ_{j+1}^{}}(\beta _j\beta _{j+1})G_{2N}(\mathrm{},\beta _j,\beta _{j+1},\mathrm{})_{\mathrm{}ϵ_j^{}ϵ_{j+1}^{}\mathrm{}},`$
$`G_{2N}(\beta _1\mathrm{},\beta _{2N1},\beta _{2N}i\lambda )_{ϵ_1\mathrm{}ϵ_{2N}}=G_{2N}(\beta _{2N},\beta _1\mathrm{},\beta _{2N1})_{ϵ_{2N}ϵ_1\mathrm{}ϵ_{2N1}},`$
where we the R-matrix is given in (2.11). In this paper we set the deformation parameter $`q`$ as following.
$`q=\mathrm{exp}\left({\displaystyle \frac{\pi i}{\xi }}\right),\xi >2.`$ (1.3)
We note that the above equations imply in particular the quantum Knizhnik-Zamolodchikov equations. In this paper we give the integral representations of the above system of diffence equations. The correlation functions are obtained by taking the shift parameter $`\lambda =2\pi `$.
In this connection we should mention about the work , in which S.Lukyanov give the integral representations of the form factors of sine-Gordon field theory. We should mention about the work , in which the authors gave the integral representations of the correlation functions of the $`U_q(\widehat{sl_n})`$ analogue of the 6-vertex model at critical regime $`|q|=1`$. The special form factors of the $`U_q(\widehat{sl_n})`$ analogue of the 6-vertex model at critical regime $`|q|=1`$ are given by T.Miwa and Y.Takeyama .
Now a few words about the organization of the paper. In section 2 we formulate our problem. In section 3 we give the free field realizations of the vertex operators. In section 4 we give proofs of properties of the free field realizations. In section 5 we give the integral representations of the correlation functions. In Appendix we summarize the multi-Gamma functions.
## 2 Problem
The purpose of this section is to formulate our problem. At first we introduce two dimensional solvable lattice model, so called 19-vertex model at critical regime $`|q|=1`$. Consider an infinite square lattice consisting of oriented lines, each carrying a spectral parameter varying fromline to line. The orientation of each line will be shown by an arrow on it. A vertex is a crossing of two lines with spectral parameters, say $`\beta _1`$ and $`\beta _2`$, together with the adjacent 4 edges belonging to the cross lines, as shown in Figure 1. The edges are assigned spin-state variables : $`j_1,j_2,k_1,k_2`$. In the 19-vertex model, each spin-state can take one of three different values $`0,1,2`$. A spin configuration around the vertex is an assignment of $`0,1,2`$ on the four edges. There are 81 possible vertex configurations. We assign each configuration a Boltzmann weight. The set of all Boltzmann weight form the elements of the $`R`$-matrix:
The matrix $`R(\beta )`$ acts on $`^3^3`$ via the natural basis $`\{v_0,v_1,v_2\}`$ of $`^3`$ as following.
$`R(\beta )v_{k_1}v_{k_2}={\displaystyle \underset{j_1,j_2=0,1,2}{}}v_{j_1}v_{j_2}R_{j_1j_2}^{k_1k_2}(\beta ).`$ (2.1)
The Boltzmann weights are given explicitly by
$`R(\beta )={\displaystyle \frac{1}{\kappa (\beta )}}\left(\begin{array}{ccccccccc}1& & & & & & & & \\ & p& & e_2& & & & & \\ & & b& & g_2& & c_2& & \\ & e_1& & p& & & & & \\ & & h_1& & o& & h_2& & \\ & & & & & p& & e_2& \\ & & c_1& & g_1& & b& & \\ & & & & & e_1& & p& \\ & & & & & & & & 1\end{array}\right).`$ (2.11)
Here we have set the normarized partition as
$`\kappa (\beta )={\displaystyle \frac{\mathrm{sh}\left({\displaystyle \frac{1}{\xi }}(\beta +\pi i)\right)}{\mathrm{sh}\left({\displaystyle \frac{1}{\xi }}(\beta \pi i)\right)}}.`$ (2.12)
The nonzero entries are given by
$`p(\beta )`$ $`=`$ $`{\displaystyle \frac{\mathrm{sh}\left({\displaystyle \frac{1}{\xi }}\beta \right)}{\mathrm{sh}\left({\displaystyle \frac{1}{\xi }}(\beta 2\pi i)\right)}},`$ (2.13)
$`b(\beta )`$ $`=`$ $`{\displaystyle \frac{\mathrm{sh}\left({\displaystyle \frac{1}{\xi }}\beta \right)\mathrm{sh}\left({\displaystyle \frac{1}{\xi }}(\beta +\pi i)\right)}{\mathrm{sh}\left({\displaystyle \frac{1}{\xi }}(\beta \pi i)\right)\mathrm{sh}\left({\displaystyle \frac{1}{\xi }}(\beta 2\pi i)\right)}},`$ (2.14)
$`o(\beta )`$ $`=`$ $`{\displaystyle \frac{\mathrm{sh}\left({\displaystyle \frac{\pi i}{\xi }}\right)\mathrm{sh}\left({\displaystyle \frac{2\pi i}{\xi }}\right)+\mathrm{sh}\left({\displaystyle \frac{1}{\xi }}\beta \right)\mathrm{sh}\left({\displaystyle \frac{1}{\xi }}(\beta \pi i)\right)}{\mathrm{sh}\left({\displaystyle \frac{1}{\xi }}(\beta \pi i)\right)\mathrm{sh}\left({\displaystyle \frac{1}{\xi }}(\beta 2\pi i)\right)}},`$ (2.15)
and
$`c_1(\beta )`$ $`=`$ $`{\displaystyle \frac{\mathrm{exp}\left({\displaystyle \frac{2}{\xi }}\beta \right)\mathrm{sh}\left({\displaystyle \frac{\pi i}{\xi }}\right)\mathrm{sh}\left({\displaystyle \frac{2\pi i}{\xi }}\right)}{\mathrm{sh}\left({\displaystyle \frac{1}{\xi }}(\beta \pi i)\right)\mathrm{sh}\left({\displaystyle \frac{1}{\xi }}(\beta 2\pi i)\right)}},`$ (2.16)
$`e_1(\beta )`$ $`=`$ $`{\displaystyle \frac{\mathrm{exp}\left({\displaystyle \frac{1}{\xi }}\beta \right)\mathrm{sh}\left({\displaystyle \frac{2\pi i}{\xi }}\right)}{\mathrm{sh}\left({\displaystyle \frac{1}{\xi }}(\beta 2\pi i)\right)}},`$ (2.17)
$`c_2(\beta )`$ $`=`$ $`\mathrm{exp}\left({\displaystyle \frac{4}{\xi }}\beta \right)c_1(\beta ),e_2(\beta )=\mathrm{exp}\left({\displaystyle \frac{2}{\xi }}\beta \right),`$ (2.18)
$`g_1(\beta )`$ $`=`$ $`{\displaystyle \frac{\mathrm{exp}\left({\displaystyle \frac{1}{\xi }}(\beta +\pi i)\right)\mathrm{sh}\left({\displaystyle \frac{\pi i}{\xi }}\right)\mathrm{sh}\left({\displaystyle \frac{1}{\xi }}\beta \right)}{\mathrm{sh}\left({\displaystyle \frac{1}{\xi }}(\beta \pi i)\right)\mathrm{sh}\left({\displaystyle \frac{1}{\xi }}(\beta 2\pi i)\right)}},`$ (2.19)
$`h_1(\beta )`$ $`=`$ $`{\displaystyle \frac{2\mathrm{exp}\left({\displaystyle \frac{1}{\xi }}(\beta \pi i)\right)\mathrm{sh}\left({\displaystyle \frac{1}{\xi }}\beta \right)\mathrm{sh}\left({\displaystyle \frac{2\pi i}{\xi }}\right)\mathrm{ch}\left({\displaystyle \frac{\pi i}{\xi }}\right)}{\mathrm{sh}\left({\displaystyle \frac{1}{\xi }}(\beta \pi i)\right)\mathrm{sh}\left({\displaystyle \frac{1}{\xi }}(\beta 2\pi i)\right)}},`$ (2.20)
$`g_2(\beta )`$ $`=`$ $`\mathrm{exp}\left({\displaystyle \frac{2}{\xi }}\beta \right)g_1(\beta ),h_2(\beta )=\mathrm{exp}\left({\displaystyle \frac{2}{\xi }}\beta \right)h_1(\beta ).`$ (2.21)
The $`R`$-matrix satisfies the Yang-Baxter equation :
$`R_{12}(\beta _1\beta _2)R_{13}(\beta _1\beta _3)R_{23}(\beta _2\beta _3)=R_{23}(\beta _2\beta _3)R_{13}(\beta _1\beta _3)R_{12}(\beta _1\beta _2).`$ (2.22)
The 19-vertex model is a limiting case of the fusion 8-vertex model . The latter is massive, and the corner transfer matrix method can be applied . Now we can conclude that the correlation functions of the critical 19-vertex model are governed by the following system of difference equations.
$`R`$-matrix Symmetry.
$`G_{2N}(\mathrm{},\beta _{j+1},\beta _j,\mathrm{})_{\mathrm{}ϵ_{j+1}ϵ_j\mathrm{}}`$ (2.23)
$`=`$ $`{\displaystyle \underset{ϵ_j^{}ϵ_{j+1}^{}=0,1,2}{}}R_{ϵ_jϵ_{j+1}}^{ϵ_j^{}ϵ_{j+1}^{}}(\beta _j\beta _{j+1})G_{2N}(\mathrm{},\beta _j,\beta _{j+1},\mathrm{})_{\mathrm{}ϵ_j^{}ϵ_{j+1}^{}\mathrm{}}.`$
Cyclicity Condition.
$`G_{2N}(\beta _1\mathrm{},\beta _{2N1},\beta _{2N}i\lambda )_{ϵ_1\mathrm{}ϵ_{2N}}=G_{2N}(\beta _{2N},\beta _1\mathrm{},\beta _{2N1})_{ϵ_{2N}ϵ_1\mathrm{}ϵ_{2N1}}.`$ (2.24)
The correlation functions :
$`G_N(\beta _N^{}+\pi i,\mathrm{},\beta _1^{}+\pi i,\beta _1,\mathrm{},\beta _N)_{2j_N,\mathrm{},2j_1,j_1,\mathrm{},j_N}.`$ (2.25)
represents the configuration functions in Figure 2, up to constant factors.
Now we can translate the problem to the following.
Find out the realizations of the vertex operators, which satisfy the following conditions.
$`R`$-matrix Symmetry.
$`\mathrm{\Phi }_{j_2}(\beta _2)\mathrm{\Phi }_{j_1}(\beta _1)={\displaystyle \underset{k_1,k_2=0,1,2}{}}R_{j_1j_2}^{k_1k_2}(\beta _1\beta _2)\mathrm{\Phi }_{k_1}(\beta _1)\mathrm{\Phi }_{k_2}(\beta _2).`$ (2.26)
Homogeneity Condition.
$`e^{\lambda D}\mathrm{\Phi }_j(\beta )e^{\lambda D}=\mathrm{\Phi }_j(\beta +i\lambda ).`$ (2.27)
Using the above vertex operators and the degree operator, we can construct the solutions of the system of difference equations as following.
$`G_{2N}(\beta _1,\mathrm{},\beta _{2N})_{j_1\mathrm{}j_{2N}}={\displaystyle \frac{\mathrm{tr}_{}\left(e^{\lambda D}\mathrm{\Phi }_{j_1}(\beta _1)\mathrm{}\mathrm{\Phi }_{j_{2N}}(\beta _{2N})\right)}{\mathrm{tr}_{}\left(e^{\lambda D}\right)}}.`$ (2.28)
The $`R`$-matrix symmetry (2.23) follows from the condition (2.26). The cyclicity condition (2.23) follows from the homogeneity condition (2.27). In section 3 we give the free field realization of the vertex operators $`\mathrm{\Phi }_j(\beta )`$. In section 5 we give the degree operator $`D`$ and the space $``$, on which the trace is evaluated.
## 3 Free field realizations
The purpose of this section is to give the free field realization of the vertex operators.
Let us introduce the bose-fields by
$`[b(t),b(t^{})]={\displaystyle \frac{1}{t}}{\displaystyle \frac{\mathrm{sh}\left({\displaystyle \frac{\pi }{2}}(\xi 2)t\right)}{\mathrm{sh}\left({\displaystyle \frac{\pi }{2}}\xi t\right)}}\delta (t+t^{}).`$ (3.1)
Let us set the basic operators by
$`U_0(\beta )=:\mathrm{exp}\left({\displaystyle _{\mathrm{}}^{\mathrm{}}}b(t)e^{i\beta t}dt\right):,`$ (3.2)
$`U_1(\beta )=:\mathrm{exp}({\displaystyle _{\mathrm{}}^{\mathrm{}}}b(t)e^{i\beta t}dt):`$ (3.3)
Let us set the fermion-fields by
$`[\psi (t),\psi (t^{})]_+=2\mathrm{c}\mathrm{h}(\pi t)\delta (t+t^{}).`$ (3.4)
Fourier transformation of the fermion-field is given by
$`\widehat{\psi }(\beta )={\displaystyle \frac{1}{\sqrt{2\pi }}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}\psi (t)e^{it\beta }𝑑t.`$ (3.5)
The free-field realizations of the vertex operators are given by
$`\mathrm{\Phi }_2(\beta )`$ $`=`$ $`U_0(\beta ),`$ (3.6)
$`\mathrm{\Phi }_1(\beta )`$ $`=`$ $`\left(e^{\frac{\pi i}{\xi }}+e^{\frac{\pi i}{\xi }}\right){\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\alpha {\displaystyle \frac{\mathrm{exp}\left(\frac{1}{\xi }(\alpha \beta )\right)}{\mathrm{sh}\left({\displaystyle \frac{1}{\xi }}(\alpha \beta +\pi i)\right)}}`$ (3.7)
$`\times U_0(\beta )U_1(\alpha )\widehat{\psi }(\alpha ),`$
$`\mathrm{\Phi }_0(\beta )`$ $`=`$ $`e^{\frac{\pi i}{\xi }}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\alpha _1{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\alpha _2{\displaystyle \underset{k=1}{\overset{2}{}}}{\displaystyle \frac{\mathrm{exp}\left(\frac{1}{\xi }(\alpha _k\beta )\right)}{\mathrm{sh}\left({\displaystyle \frac{1}{\xi }}(\alpha _k\beta +\pi i)\right)}}`$ (3.8)
$`\times U_0(\beta )U_1(\alpha _1)U_1(\alpha _2)\widehat{\psi }(\alpha _1)\widehat{\psi }(\alpha _2)`$
## 4 Proof
The purpose of this section is to give proofs of properties of vertex operators. At first we explain the formulas of the form
$`X(\beta _1)Y(\beta _2)=C_{XY}(\beta _1\beta _2):X(\beta _1)X(\beta _2):,`$ (4.1)
where $`X,Y=U_j`$, and $`C_{XY}(\beta )`$ is a meromorphic function on $``$. These formulae follow from the commutation relation of the free bosons. When we compute the contraction of the basic operators, we often encounter an integral
$`{\displaystyle _0^{\mathrm{}}}F(t)𝑑t,`$ (4.2)
which is divergent at $`t=0`$. Here we adopt the following prescription for regularization : it should be understood as the countour integral,
$`{\displaystyle _C}F(t){\displaystyle \frac{\mathrm{log}(t)}{2\pi i}}𝑑t,`$ (4.3)
where the countour $`C`$ is given by
The contractions of the basic operators have the following forms.
$`U_j(\beta _1)U_j(\beta _2)=:U_j(\beta _1)U_j(\beta _2):{\displaystyle \frac{\mathrm{\Gamma }\left(i{\displaystyle \frac{\beta _2\beta _1}{\pi \xi }}+1{\displaystyle \frac{1}{\xi }}\right)}{\mathrm{\Gamma }\left(i{\displaystyle \frac{\beta _2\beta _1}{\pi \xi }}+{\displaystyle \frac{1}{\xi }}\right)}}\mathrm{exp}\left({\displaystyle \frac{\xi 2}{\xi }}(\gamma +\mathrm{log}(\pi \xi ))\right).`$ (4.4)
Please see the Appendix. The commutation relations of the basic operators are given by
$`U_j(\beta _1)U_j(\beta _2)`$ $`=`$ $`U_j(\beta _2)U_j(\beta _1){\displaystyle \frac{\mathrm{sh}\left({\displaystyle \frac{1}{\xi }}(\beta _1\beta _2+\pi i)\right)}{\mathrm{sh}\left({\displaystyle \frac{1}{\xi }}(\beta _2\beta _1+\pi i)\right)}},(j=0,1),`$ (4.5)
$`U_0(\beta _1)U_1(\beta _2)`$ $`=`$ $`U_1(\beta _2)U_0(\beta _1){\displaystyle \frac{\mathrm{sh}\left({\displaystyle \frac{1}{\xi }}(\beta _2\beta _1+\pi i)\right)}{\mathrm{sh}\left({\displaystyle \frac{1}{\xi }}(\beta _1\beta _2+\pi i)\right)}}.`$ (4.6)
The anti-commutation relation becomes
$`[\widehat{\psi }(\beta _1),\widehat{\psi }(\beta _2)]_+=\delta (\beta _1\beta _2+\pi i)+\delta (\beta _2\beta _1+\pi i).`$ (4.7)
Now let us start to prove the properties of vertex operators.
Proof of $`R`$-matrix symmetery
Let us prove the equation :
$`R_{21}^{12}(\beta _1\beta _2)\mathrm{\Phi }_1(\beta _1)\mathrm{\Phi }_2(\beta _2)+R_{21}^{21}(\beta _1\beta _2)\mathrm{\Phi }_2(\beta _1)\mathrm{\Phi }_1(\beta _2)=\mathrm{\Phi }_1(\beta _2)\mathrm{\Phi }_2(\beta _1).`$ (4.8)
By using the commutation relations of basic operators, we can rearrange the operator part as
$`U_0(\beta _1)U_0(\beta _2)U_1(\alpha )\widehat{\psi }(\alpha ).`$ (4.9)
The equation (4.8) follows from the integrand identity :
$`e_1(\beta _1\beta _2)e^{\frac{1}{\xi }\beta _1}\mathrm{sh}\left({\displaystyle \frac{1}{\xi }}(\beta _2\alpha +\pi i)\right)+p(\beta _1\beta _2)e^{\frac{1}{\xi }\beta _2}\mathrm{sh}\left({\displaystyle \frac{1}{\xi }}(\beta _1\alpha \pi i)\right)`$
$`=e^{\frac{1}{\xi }\beta _2}\mathrm{sh}\left({\displaystyle \frac{1}{\xi }}(\beta _1\alpha +\pi i)\right).`$ (4.10)
Let us prove the equation :
$`R_{11}^{02}(\beta _1\beta _2)\mathrm{\Phi }_0(\beta _1)\mathrm{\Phi }_2(\beta _2)+R_{11}^{11}(\beta _1\beta _2)\mathrm{\Phi }_1(\beta _1)\mathrm{\Phi }_1(\beta _2)`$
$`+R_{11}^{20}(\beta _1\beta _2)\mathrm{\Phi }_2(\beta _1)\mathrm{\Phi }_0(\beta _2)=\mathrm{\Phi }_1(\beta _2)\mathrm{\Phi }_1(\beta _1).`$ (4.11)
By using the commutation relations of basic operators, we can rearrange the operator part as
$`U_0(\beta _1)U_0(\beta _2)U_1(\alpha _1)U_1(\alpha _2)\widehat{\psi }(\alpha _1)\widehat{\psi }(\alpha _2).`$ (4.12)
Let us set
$`H(\alpha )={\displaystyle \frac{\mathrm{sh}\left({\displaystyle \frac{1}{\xi }}(\alpha +\pi i)\right)}{\mathrm{sh}\left({\displaystyle \frac{1}{\xi }}(\alpha +\pi i)\right)}}.`$ (4.13)
Consider the integral of the form :
$`{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\alpha _1{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\alpha _2F(\alpha _1,\alpha _2)U_1(\alpha _1)U_1(\alpha _2)\widehat{\psi }(\alpha _1)\widehat{\psi }(\alpha _2).`$ (4.14)
Due to the commutation relation of $`U_1(\alpha )`$ and the anti-commutation relation of $`\widehat{\psi }(\alpha )`$, the above integral equals to
$`{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\alpha _1{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\alpha _2H(\alpha _2\alpha _1)F(\alpha _1,\alpha _2)U_1(\alpha _1)U_1(\alpha _2)\widehat{\psi }(\alpha _1)\widehat{\psi }(\alpha _2)`$
$`H(\pi i){\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\alpha F(\alpha \pi i,\alpha )U_1(\alpha )U_1(\alpha \pi i),`$ (4.15)
where we have used the relation: $`H(\pi i)=0`$. Note that the part $`H(\pi i)U_1(\alpha )U_1(\alpha \pi i)`$ is convergent.
Observing this we define ’weak equality’ in the following sense. We say that the function $`G_1(\alpha _1,\alpha _2)`$ and $`G_2(\alpha _1,\alpha _2)`$ are equal in weak sense if
$`G_1(\alpha _1,\alpha _2)H(\alpha _2\alpha _1)G_1(\alpha _2,\alpha _1)=G_2(\alpha _1,\alpha _2)H(\alpha _2\alpha _1)G_2(\alpha _2,\alpha _1).`$ (4.16)
We write
$`G_1(\alpha _1,\alpha _2)G_2(\alpha _1,\alpha _2).`$ (4.17)
Note that the equation
$`G_1(\alpha \pi i,\alpha )=G_2(\alpha \pi i,\alpha ),`$ (4.18)
is a special case of weakly equality. In oreder to prove the equation (4.11) it is enough to prove the equality of the integrand part in weakly sense. The equation (4.11) follows from the following weakly sense identity.
$`h_1(\beta _1\beta _2)e^{\frac{\pi i}{\xi }}\mathrm{sh}\left({\displaystyle \frac{1}{\xi }}(\beta _2\alpha _1+\pi i)\right)\mathrm{sh}\left({\displaystyle \frac{1}{\xi }}(\beta _2\alpha _2+\pi i)\right)`$
$`+o(\beta _1\beta _2)(1+e^{\frac{2\pi i}{\xi }})^2\mathrm{sh}\left({\displaystyle \frac{1}{\xi }}(\beta _2\alpha _1+\pi i)\right)\mathrm{sh}\left({\displaystyle \frac{1}{\xi }}(\alpha _2\beta _1+\pi i)\right)`$
$`+h_2(\beta _1\beta _2)e^{\frac{\pi i}{\xi }}\mathrm{sh}\left({\displaystyle \frac{1}{\xi }}(\alpha _1\beta _1+\pi i)\right)\mathrm{sh}\left({\displaystyle \frac{1}{\xi }}(\alpha _2\beta _1+\pi i)\right)`$
$`(1+e^{\frac{2\pi i}{\xi }})^2\mathrm{sh}\left({\displaystyle \frac{1}{\xi }}(\beta _1\alpha _1+\pi i)\right)\mathrm{sh}\left({\displaystyle \frac{1}{\xi }}(\alpha _2\beta _2+\pi i)\right).`$ (4.19)
As the same arguments as the above we obtain the commutation relation of the vertex operators (2.26).
## 5 Correlation functions
In this section we derive a solution of the system of difference equations (2.23) and (2.24), algebraically, and obtain an integral representation of it.
Let us introduce the Fock space $`^b`$ generated by $`|vac_b`$ which satisfies
$`b(t)|vac_b=0,\mathrm{if}t>0.`$ (5.1)
Let us introduce the Fock space $`^\psi `$ generated by $`|vac_\psi `$ which satisfies
$`\psi (t)|vac_\psi =0,\mathrm{if}t>0.`$ (5.2)
Set the space $``$ by
$`=^b^\psi .`$ (5.3)
Let us introduce the degree operators $`D^b`$ and $`D^\psi `$ by
$`D^bb(t)|vac_b=tb(t)|vac_b,D^\psi \psi (t)|vac_\psi =t\psi (t)|vac_\psi ,t>0.`$ (5.4)
Set the degree operator $`D`$ on $``$ by
$`D=D^bid+idD^\psi .`$ (5.5)
We have
$`e^{\lambda D}U_j(\beta )e^{\lambda D}=U_j(\beta +i\lambda ),e^{\lambda D}\widehat{\psi }(\beta )e^{\lambda D}=\widehat{\psi }(\beta +i\lambda ).`$ (5.6)
Therefore the vertex operators satisfy the homogeneity condition.
$`e^{\lambda D}\mathrm{\Phi }_j(\beta )e^{\lambda D}=\mathrm{\Phi }_j(\beta +i\lambda ).`$ (5.7)
Now let us consider the trace function for $`\lambda >0`$ defined by
$`G_{2N}(\beta _1,\mathrm{},\beta _{2N})_{j_1\mathrm{}j_{2N}}={\displaystyle \frac{\mathrm{tr}_{}\left(e^{\lambda D}\mathrm{\Phi }_{j_1}(\beta _1)\mathrm{}\mathrm{\Phi }_{j_{2N}}(\beta _{2N})\right)}{\mathrm{tr}_{}\left(e^{\lambda D}\right)}}.`$ (5.8)
The trace of the bosonic parts is evaluated as followings.
$`{\displaystyle \frac{\mathrm{tr}_^b\left(e^{\lambda D^b}b(t)b(t^{})\right)}{\mathrm{tr}_^b\left(e^{\lambda D^b}\right)}}={\displaystyle \frac{e^{\lambda t}}{e^{\lambda t}1}}[b(t),b(t^{})].`$ (5.9)
The product of the basic operator $`U_j(\beta ),(j=0,1)`$ is evaluated as
$`{\displaystyle \frac{\mathrm{tr}_^b\left(e^{\lambda D^b}U_j(\beta _1)U_j(\beta _2)\right)}{\mathrm{tr}_^b\left(e^{\lambda D^b}\right)}}=Const.\phi _1(\beta _1\beta _2)`$ (5.10)
$`\times `$ $`\mathrm{sh}\left({\displaystyle \frac{\pi }{\lambda }}(\beta _1\beta _2+\pi i)\right)\mathrm{sh}\left({\displaystyle \frac{\pi }{\lambda }}(\beta _1\beta _2\pi i)\right)\mathrm{sh}\left({\displaystyle \frac{1}{\xi }}(\beta _1\beta _2+\pi i)\right),`$
where we set the kernel function as
$`\phi _1(\beta )={\displaystyle \frac{1}{S_2(i\beta \pi |\pi \xi ,\lambda )S_2(i\beta \pi |\pi \xi ,\lambda )}}.`$ (5.11)
The product of the basic operators $`U_0(\beta )`$ and $`U_1(\alpha )`$ is evaluated as
$`{\displaystyle \frac{\mathrm{tr}_^b\left(e^{\lambda D^b}U_0(\beta )U_1(\alpha )\right)}{\mathrm{tr}_^b\left(e^{\lambda D^b}\right)}}=Const.\phi _2(\beta \alpha )\mathrm{sh}\left({\displaystyle \frac{1}{\xi }}(\beta \alpha \pi i)\right),`$ (5.12)
where we set the kernel function as
$`\phi _2(\alpha )={\displaystyle \frac{1}{S_2(i\alpha +\pi |\pi \xi ,\lambda )S_2(i\alpha +\pi |\pi \xi ,\lambda )}}.`$ (5.13)
The trace of the fermionic parts is evaluated as followings.
$`{\displaystyle \frac{\mathrm{tr}_^\psi \left(e^{\lambda D^\psi }\psi (t)\psi (t^{})\right)}{\mathrm{tr}_^\psi \left(e^{\lambda D^\psi }\right)}}={\displaystyle \frac{e^{\lambda t}}{e^{\lambda t}+1}}[\psi (t),\psi (t^{})]_+.`$ (5.14)
Let us set the auxiliary function $`𝒥(\alpha )`$ by
$`𝒥(\alpha )={\displaystyle \frac{1}{2\pi }}{\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle \frac{e^{i\alpha t}}{1+e^{\lambda t}}}𝑑t.`$ (5.15)
We then have
$`{\displaystyle \frac{\mathrm{tr}_^\psi \left(e^{\lambda D^\psi }\widehat{\psi }(\alpha _1)\widehat{\psi }(\alpha _2)\right)}{\mathrm{tr}_^\psi \left(e^{\lambda D^\psi }\right)}}=𝒥(\alpha _1\alpha _2+\pi i)+𝒥(\alpha _1\alpha _2\pi i).`$ (5.16)
The trace of the vertex operators is evaluated by applying the Wick’s theorem.
The one-point correlation functions are evaluated as follows.
$`G_2(\beta _1,\beta _2)_{2,0}`$ $`=`$ $`e^{\frac{2}{\xi }\beta _2}{\displaystyle \frac{S_2(i(\beta _2\beta _1)+\pi |\pi \xi ,\lambda )}{S_2(i(\beta _2\beta _1)+\pi \xi \pi |\pi \xi ,\lambda )}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\alpha _1{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\alpha _2`$
$`\times `$ $`\left\{{\displaystyle \underset{j,k=1}{\overset{2}{}}}\phi _2(\beta _j\alpha _k)\right\}\phi _1(\alpha _1\alpha _2)(𝒥(\alpha _1\alpha _2+\pi i)+𝒥(\alpha _1\alpha _2\pi i))`$
$`\times `$ $`e^{\frac{1}{\xi }(\alpha _1+\alpha _2)}\mathrm{sh}\left({\displaystyle \frac{\pi }{\lambda }}(\alpha _1\alpha _2+\pi i)\right)\mathrm{sh}\left({\displaystyle \frac{\pi }{\lambda }}(\alpha _2\alpha _1+\pi i)\right)`$
$`\times `$ $`\mathrm{sh}\left({\displaystyle \frac{1}{\xi }}(\alpha _1\alpha _2+\pi i)\right)\mathrm{sh}\left({\displaystyle \frac{1}{\xi }}(\alpha _1\beta _1+\pi i)\right)\mathrm{sh}\left({\displaystyle \frac{1}{\xi }}(\alpha _2\beta _1+\pi i)\right),`$
$`G_2(\beta _1,\beta _2)_{2,0}`$ $`=`$ $`e^{\frac{2}{\xi }\beta _1}{\displaystyle \frac{S_2(i(\beta _2\beta _1)+\pi |\pi \xi ,\lambda )}{S_2(i(\beta _2\beta _1)+\pi \xi \pi |\pi \xi ,\lambda )}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\alpha _1{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\alpha _2`$
$`\times `$ $`\left\{{\displaystyle \underset{j,k=1}{\overset{2}{}}}\phi _2(\beta _j\alpha _k)\right\}\phi _1(\alpha _1\alpha _2)(𝒥(\alpha _1\alpha _2+\pi i)+𝒥(\alpha _1\alpha _2\pi i))`$
$`\times `$ $`e^{\frac{1}{\xi }(\alpha _1+\alpha _2)}\mathrm{sh}\left({\displaystyle \frac{\pi }{\lambda }}(\alpha _1\alpha _2+\pi i)\right)\mathrm{sh}\left({\displaystyle \frac{\pi }{\lambda }}(\alpha _2\alpha _1+\pi i)\right)`$
$`\times `$ $`\mathrm{sh}\left({\displaystyle \frac{1}{\xi }}(\alpha _1\alpha _2+\pi i)\right)\mathrm{sh}\left({\displaystyle \frac{1}{\xi }}(\alpha _1\beta _2\pi i)\right)\mathrm{sh}\left({\displaystyle \frac{1}{\xi }}(\alpha _2\beta _2\pi i)\right),`$
and
$`G_2(\beta _1,\beta _2)_{1,1}`$ $`=`$ $`e^{\frac{1}{\xi }(\beta _1+\beta _2)}{\displaystyle \frac{S_2(i(\beta _2\beta _1)+\pi |\pi \xi ,\lambda )}{S_2(i(\beta _2\beta _1)+\pi \xi \pi |\pi \xi ,\lambda )}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\alpha _1{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\alpha _2`$
$`\times `$ $`\left\{{\displaystyle \underset{j,k=1}{\overset{2}{}}}\phi _2(\beta _j\alpha _k)\right\}\phi _1(\alpha _1\alpha _2)(𝒥(\alpha _1\alpha _2+\pi i)+𝒥(\alpha _1\alpha _2\pi i))`$
$`\times `$ $`e^{\frac{1}{\xi }(\alpha _1+\alpha _2)}\mathrm{sh}\left({\displaystyle \frac{\pi }{\lambda }}(\alpha _1\alpha _2+\pi i)\right)\mathrm{sh}\left({\displaystyle \frac{\pi }{\lambda }}(\alpha _2\alpha _1+\pi i)\right)`$
$`\times `$ $`\mathrm{sh}\left({\displaystyle \frac{1}{\xi }}(\alpha _1\alpha _2+\pi i)\right)\mathrm{sh}\left({\displaystyle \frac{1}{\xi }}(\alpha _1\beta _2\pi i)\right)\mathrm{sh}\left({\displaystyle \frac{1}{\xi }}(\alpha _2\beta _1+\pi i)\right).`$
Here we omit an irrelevant constant factor.
By applying Wick’s theorem we obtain the $`N`$-point correlation functions. We consider the special components :
$`j_1,\mathrm{},j_L=2,j_{L+1},\mathrm{},j_{L+2M}=1,j_{L+2M+1},\mathrm{},j_{2(L+M)}=0.`$ (5.20)
By using the $`R`$-matrix symmetry (2.23) we obtain every components from this component. The $`N`$-point function is evaluated as following.
$`G_{2(L+M)}(\beta _1\mathrm{}\beta _L|\beta _{L+1},\mathrm{},\beta _{L+2M}|\beta _{L+2M+1}\mathrm{}\beta _{2(L+M)})_{2,\mathrm{},2,1,\mathrm{},1,0,\mathrm{},0}`$
$`=E(\{\beta \}){\displaystyle 𝑑\alpha \mathrm{\Psi }(\{\alpha \}|\{\beta \})Pf(\{\alpha \})I_\lambda (\{\alpha \})I_\xi (\{\alpha \}|\{\beta \})},`$ (5.21)
where the integral $`𝑑\alpha `$ represents
$`{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\alpha _{L+1}\mathrm{}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\alpha _{2(L+M)}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\alpha _{L+2M+1}^{}\mathrm{}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\alpha _{2(L+M)}^{}.`$ (5.22)
Here we have set
$`E(\{\beta \})`$ $`=`$ $`e^{\frac{1}{\xi }(\beta _{L+1}+\mathrm{}+\beta _{L+2M})}e^{\frac{2}{\xi }(\beta _{L+2M+1}^{}+\mathrm{}+\beta _{2(L+M)}^{})}`$ (5.23)
$`\times `$ $`{\displaystyle \underset{1j<k2(L+M)}{}}{\displaystyle \frac{S_2(i(\beta _k\beta _j)+\pi |\pi \xi ,\lambda )}{S_2(i(\beta _k\beta _j)+\pi \xi \pi |\pi \xi ,\lambda )}}.`$
The integral kernel is given by
$`\mathrm{\Psi }(\{\alpha \}|\{\beta \})={\displaystyle \underset{L+1j<k2(L+M)}{}}\phi _1(\alpha _j\alpha _k){\displaystyle \underset{L+2M+1j<k2(L+M)}{}}\phi _1(\alpha _j^{}\alpha _k^{})`$
$`\times `$ $`{\displaystyle \underset{j=L+1}{\overset{2(L+M)}{}}}{\displaystyle \underset{k=L+2M+1}{\overset{2(L+M)}{}}}\phi _1(\alpha _j\alpha _k^{}){\displaystyle \underset{j=1}{\overset{2(L+M)}{}}}{\displaystyle \underset{k=L+1}{\overset{2(L+M)}{}}}\phi _2(\beta _j\alpha _k){\displaystyle \underset{j=1}{\overset{2(L+M)}{}}}{\displaystyle \underset{k=L+2M+1}{\overset{2(L+M)}{}}}\phi _2(\beta _j\alpha _k^{}).`$
The $`Pf(\{\alpha \})`$ represents a Pfaffian of $`2(L+M)\times 2(L+M)`$ anti-symmmetry matrix whose entries are given by $`𝒥(\alpha \alpha ^{}+\pi i)+𝒥(\alpha \alpha ^{}\pi i)`$.
The integrand functions are given by
$`I_\lambda (\{\alpha \})={\displaystyle \underset{j=L+1}{\overset{2(L+M)}{}}}{\displaystyle \underset{k=L+1}{\overset{2(L+M)}{}}}\mathrm{sh}\left({\displaystyle \frac{\pi }{\lambda }}(\alpha _j\alpha _k\pi i)\right)`$ (5.25)
$`\times `$ $`{\displaystyle \underset{j=L+1}{\overset{2(L+M)}{}}}{\displaystyle \underset{k=L+2M+1}{\overset{2(L+M)}{}}}\mathrm{sh}\left({\displaystyle \frac{\pi }{\lambda }}(\alpha _j\alpha _k^{}\pi i)\right)\mathrm{sh}\left({\displaystyle \frac{\pi }{\lambda }}(\alpha _j\alpha _k^{}+\pi i)\right)`$
and
$`I_\xi (\{\alpha \}|\{\beta \})=e^{\frac{1}{\xi }(\alpha _{L+1}+\mathrm{}\alpha _{2(L+M)}+\alpha _{l+2M+1}^{}+\mathrm{}+\alpha _{2(L+M)}^{})}`$
$`\times `$ $`{\displaystyle \underset{L+1j<k2(L+M)}{}}\mathrm{sh}\left({\displaystyle \frac{1}{\xi }}(\alpha _j\alpha _k+\pi i)\right){\displaystyle \underset{j=L+1}{\overset{L+2M}{}}}{\displaystyle \underset{k=L+2M+1}{\overset{2(L+M)}{}}}\mathrm{sh}\left({\displaystyle \frac{1}{\xi }}(\alpha _j\alpha _k^{}+\pi i)\right)`$
$`\times `$ $`{\displaystyle \underset{j=L+2M+1}{\overset{2(L+M)}{}}}\mathrm{sh}\left({\displaystyle \frac{1}{\xi }}(\alpha _j\alpha _k^{}\mathrm{sgn}(jk)\pi i)\right)`$
$`\times `$ $`{\displaystyle \underset{j=1}{\overset{L}{}}}\left\{{\displaystyle \underset{k=L+1}{\overset{2(L+M)}{}}}\mathrm{sh}\left({\displaystyle \frac{1}{\xi }}(\beta _j\alpha _k\pi i)\right){\displaystyle \underset{k=L+2M+1}{\overset{2(L+M)}{}}}\mathrm{sh}\left({\displaystyle \frac{1}{\xi }}(\beta _j\alpha _k^{}\pi i)\right)\right\}`$
$`\times `$ $`{\displaystyle \underset{j=L+1}{\overset{L+2M}{}}}\left\{{\displaystyle \underset{k=L+2M+1}{\overset{2(L+M)}{}}}\mathrm{sh}\left({\displaystyle \frac{1}{\xi }}(\beta _j\alpha _k\pi i)\right){\displaystyle \underset{k=L+2M+1}{\overset{2(L+M)}{}}}\mathrm{sh}\left({\displaystyle \frac{1}{\xi }}(\beta _j\alpha _k^{}\pi i)\right)\right\}`$
$`\times `$ $`{\displaystyle \underset{\genfrac{}{}{0pt}{}{j,k=L+1}{jk}}{\overset{L+2M}{}}}\mathrm{sh}\left({\displaystyle \frac{1}{\xi }}(\beta _j\alpha _k+\mathrm{sgn}(jk)\pi i)\right)`$
$`\times `$ $`{\displaystyle \underset{\genfrac{}{}{0pt}{}{j,k=L+2M+1}{jk}}{\overset{2(L+M)}{}}}\left\{\mathrm{sh}\left({\displaystyle \frac{1}{\xi }}(\beta _j\alpha _k+\mathrm{sgn}(jk)\pi i)\right)\mathrm{sh}\left({\displaystyle \frac{1}{\xi }}(\beta _j\alpha _k^{}+\mathrm{sgn}(jk)\pi i)\right)\right\}`$
Here we omit an irrelevant constant factor.
Next we consider the special case $`\lambda =2\pi `$, in which the trace function will become the correlation functions of our original solvable lattice problem. Note that the special case $`\lambda =2\pi `$ the kernel functions simplify.
$`𝒥(\alpha +\pi i)+𝒥(\alpha \pi i)=\delta (\alpha \pi i),`$ (5.28)
and
$`{\displaystyle \frac{\mathrm{tr}_^b\left(e^{\lambda D^b}U_j(\beta _1)U_j(\beta _2)\right)}{\mathrm{tr}_^b\left(e^{\lambda D^b}\right)}}=Const.{\displaystyle \frac{\mathrm{ch}\left({\displaystyle \frac{1}{2}}(\beta _1\beta _2)\right)}{\mathrm{sh}\left({\displaystyle \frac{1}{\xi }}(\beta _1\beta _2\pi i)\right)}}.`$ (5.29)
We have seen that the kernel function of trace formulae gets simplified when specialized at $`\lambda =2\pi `$. Here we summarize the simplified formulae for the one-point correlation functions at $`\lambda =2\pi `$. In this simplified formulae the number of the contour integrals reduces to one. This is due to a property of the fermion two-points function which becomes the delta function.
$`G_2(\beta _1,\beta _2)_{2,0}=e^{\frac{2}{\xi }\beta _2}{\displaystyle \frac{\mathrm{ch}\left({\displaystyle \frac{1}{2}}(\beta _1\beta _2)\right)}{\mathrm{sh}\left({\displaystyle \frac{1}{\xi }}(\beta _1\beta _2\pi i)\right)}}`$ (5.30)
$`\times `$ $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\alpha {\displaystyle \frac{e^{\frac{2}{\xi }\alpha }}{_{k=1}^2\mathrm{sh}\left(\beta _k\alpha \right)}}\mathrm{sh}\left({\displaystyle \frac{1}{\xi }}(\beta _1\alpha \pi i)\right)\mathrm{sh}\left({\displaystyle \frac{1}{\xi }}(\beta _1\alpha 2\pi i)\right).`$
$`G_2(\beta _1,\beta _2)_{0,2}=e^{\frac{2}{\xi }\beta _1}{\displaystyle \frac{\mathrm{ch}\left({\displaystyle \frac{1}{2}}(\beta _1\beta _2)\right)}{\mathrm{sh}\left({\displaystyle \frac{1}{\xi }}(\beta _1\beta _2\pi i)\right)}}`$ (5.31)
$`\times `$ $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\alpha {\displaystyle \frac{e^{\frac{2}{\xi }\alpha }}{_{k=1}^2\mathrm{sh}\left(\beta _k\alpha \right)}}\mathrm{sh}\left({\displaystyle \frac{1}{\xi }}(\beta _2\alpha )\right)\mathrm{sh}\left({\displaystyle \frac{1}{\xi }}(\beta _2\alpha +\pi i)\right).`$
$`G_2(\beta _1,\beta _2)_{1,1}=e^{\frac{1}{\xi }(\beta _1+\beta _2)}{\displaystyle \frac{\mathrm{ch}\left({\displaystyle \frac{1}{2}}(\beta _1\beta _2)\right)}{\mathrm{sh}\left({\displaystyle \frac{1}{\xi }}(\beta _1\beta _2\pi i)\right)}}`$ (5.32)
$`\times `$ $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\alpha {\displaystyle \frac{e^{\frac{2}{\xi }\alpha }}{_{k=1}^2\mathrm{sh}\left(\beta _k\alpha \right)}}\mathrm{sh}\left({\displaystyle \frac{1}{\xi }}(\beta _1\alpha \pi i)\right)\mathrm{sh}\left({\displaystyle \frac{1}{\xi }}(\beta _2\alpha )\right).`$
Here we omit an irrelevant constant factor. The $`R`$-matrix symmetry (2.23) of the above integral representations can be reduced to the special case of the identity (4.19): $`\alpha _1=\alpha +\pi i,\alpha _2=\alpha `$. The cyclicity condition (2.24) can be checked by straightforward calculation. We have seen that the formulae of the one-point functions get simplified when we set $`\lambda =2\pi `$. This feature holds for the $`N`$-point correlation functions, too. The number of the contour integrals reduces to only $`N`$.
Acknowledgements. This work was partly supported by Grant-in-Aid for Encouragements for Young Scientists (A) from Japan Society for the Promotion of Science. (11740099)
## Appendix A Multi Gamma functions
Here we summarize the multiple gamma and the multiple sine functions, following N.Kurokawa .
Let us set the functions $`\mathrm{\Gamma }_1(x|\omega )`$ and $`\mathrm{\Gamma }_2(x|\omega _1,\omega _2)`$ by
$`\mathrm{log}\mathrm{\Gamma }_1(x|\omega )+\gamma B_{11}(x|\omega )`$ $`=`$ $`{\displaystyle _C}{\displaystyle \frac{dt}{2\pi it}}e^{xt}{\displaystyle \frac{\mathrm{log}(t)}{1e^{\omega t}}},`$ (A.1)
$`\mathrm{log}\mathrm{\Gamma }_2(x|\omega _1,\omega _2){\displaystyle \frac{\gamma }{2}}B_{22}(x|\omega _1,\omega _2)`$ $`=`$ $`{\displaystyle _C}{\displaystyle \frac{dt}{2\pi it}}e^{xt}{\displaystyle \frac{\mathrm{log}(t)}{(1e^{\omega _1t})(1e^{\omega _2t})}},`$ (A.2)
where the functions $`B_{jj}(x)`$ are the multiple Bernoulli polynomials defined by
$`{\displaystyle \frac{t^re^{xt}}{_{j=1}^r(e^{\omega _jt}1)}}={\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{t^n}{n!}}B_{r,n}(x|\omega _1\mathrm{}\omega _r),`$ (A.3)
more explicitly
$`B_{11}(x|\omega )`$ $`=`$ $`{\displaystyle \frac{x}{\omega }}{\displaystyle \frac{1}{2}},`$ (A.4)
$`B_{22}(x|\omega )`$ $`=`$ $`{\displaystyle \frac{x^2}{\omega _1\omega _2}}\left({\displaystyle \frac{1}{\omega _1}}+{\displaystyle \frac{1}{\omega _2}}\right)x+{\displaystyle \frac{1}{2}}+{\displaystyle \frac{1}{6}}\left({\displaystyle \frac{\omega _1}{\omega _2}}+{\displaystyle \frac{\omega _2}{\omega _1}}\right).`$ (A.5)
Here $`\gamma `$ is Euler’s constant, $`\gamma =lim_n\mathrm{}(1+\frac{1}{2}+\frac{1}{3}+\mathrm{}+\frac{1}{n}\mathrm{log}n)`$.
Here the contor of integral is given by
Let us set
$`S_1(x|\omega )`$ $`=`$ $`{\displaystyle \frac{1}{\mathrm{\Gamma }_1(\omega x|\omega )\mathrm{\Gamma }_1(x|\omega )}},`$ (A.6)
$`S_2(x|\omega _1,\omega _2)`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Gamma }_2(\omega _1+\omega _2x|\omega _1,\omega _2)}{\mathrm{\Gamma }_2(x|\omega _1,\omega _2)}}.`$ (A.7)
We have
$`\mathrm{\Gamma }_1(x|\omega )=e^{(\frac{x}{\omega }\frac{1}{2})\mathrm{log}\omega }{\displaystyle \frac{\mathrm{\Gamma }(x/\omega )}{\sqrt{2\pi }}},S_1(x|\omega )=2\mathrm{s}\mathrm{i}\mathrm{n}(\pi x/\omega ),`$ (A.8)
$`{\displaystyle \frac{\mathrm{\Gamma }_2(x+\omega _1|\omega _1,\omega _2)}{\mathrm{\Gamma }_2(x|\omega _1,\omega _2)}}={\displaystyle \frac{1}{\mathrm{\Gamma }_1(x|\omega _2)}},{\displaystyle \frac{S_2(x+\omega _1|\omega _1,\omega _2)}{S_2(x|\omega _1,\omega _2)}}={\displaystyle \frac{1}{S_1(x|\omega _2)}},{\displaystyle \frac{\mathrm{\Gamma }_1(x+\omega |\omega )}{\mathrm{\Gamma }_1(x|\omega )}}=x.`$ (A.9)
$`\mathrm{log}S_2(x|\omega _1\omega _2)={\displaystyle _C}{\displaystyle \frac{\mathrm{sh}(x\frac{\omega _1+\omega _2}{2})t}{2\mathrm{s}\mathrm{h}\frac{\omega _1t}{2}\mathrm{sh}\frac{\omega _2t}{2}}}\mathrm{log}(t){\displaystyle \frac{dt}{2\pi it}},(0<\mathrm{Re}x<\omega _1+\omega _2).`$ (A.10)
$`S_2(x|\omega _1\omega _2)={\displaystyle \frac{2\pi }{\sqrt{\omega _1\omega _2}}}x+O(x^2),(x0).`$ (A.11)
$`S_2(x|\omega _1\omega _2)S_2(x|\omega _1\omega _2)=4\mathrm{s}\mathrm{i}\mathrm{n}{\displaystyle \frac{\pi x}{\omega _1}}\mathrm{sin}{\displaystyle \frac{\pi x}{\omega _2}}.`$ (A.12)
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# 1 Introduction.
## 1 Introduction.
Fractional Brownian motion (fBm) was first introduced by Kolmogorov in 1940 \[Kl\] and later studied by Lévy and Mandelbrot \[Lv, Mn\].
Let $`(\mathrm{\Omega },,P)`$ be a probability space, and $`\alpha \mathrm{I}\mathrm{R}`$, $`|\alpha |<1`$ be a parameter. FBm with exponent $`\alpha `$ is a self-similar, centered Gaussian random process $`\xi _\alpha (t,\omega )`$, $`(t,\omega )[0,\mathrm{})\times \mathrm{\Omega }`$ (often abbreviated as $`\xi _\alpha (t)`$) with stationary increments and the correlation function
$$E(\xi _\alpha (s),\xi _\alpha (t))=C(s^{1+\alpha }+t^{1+\alpha }|st|^{1+\alpha }),C=\frac{\mathrm{\Gamma }(1\alpha )}{\alpha }\frac{\mathrm{cos}\frac{1+\alpha }{2}\pi }{\frac{1+\alpha }{2}\pi }.$$
For $`\alpha =0`$ we recover the ordinary Brownian motion (oBm). From the form of the correlation function of the increments
$$E(d\xi _\alpha (s),d\xi _\alpha (t))=C\alpha (\alpha +1)\frac{dsdt}{|st|^{1\alpha }},$$
one can see that fBm is not Markovian for all $`\alpha 0`$. The trajectories of fBm are almost surely Hölder-continuous with exponent less than $`\frac{1+\alpha }{2}`$, and not Hölder-continuous with exponent greater or equal to $`\frac{1+\alpha }{2}`$. The finite-dimensional distributions of fBm are scale-invariant:
$$r,t>0,r^{\frac{1+\alpha }{2}}\xi _\alpha (rt)\stackrel{\mathrm{dist}}{=}\xi _\alpha (t).$$
The scale-invariance and long-range correlations make fBm important in applications.
In recent years there has been much interest in fBm, see for example \[G, K1, K2, Lb, Ml, Sn, T1, T2, W, X, Yr\]. The problem of the construction of stochastic calculus with respect to fBm has been considered in \[Ln, DH, DU\]. The main difficulty is that fBm fails to be a semi-martingale for all $`0<\alpha <1`$ \[Ln\]. Our goal is to construct the stochastic integrals with respect to fBm and to prove the general existence-uniqueness theorem for solutions of stochastic differential equations (SDE’s) with fBm. Due to the strongly non-Markovian nature of fBm, the most natural way to define the stochastic integral is to do it pathwise, for a.e. $`\omega `$. This brings us to the question of existence of Stieltjes integrals with respect to Hölder continuous functions and of the existence and uniqueness of solutions of ODE’s with Hölder continuous forcing (Section 2). The existence of Stieltjes integrals with respect to Hölder continuous functions was established in \[Yg, Kn\]. Below we use the methods of Renormalization group to prove a formula (Theorem 2) which allows us to estimate the $`L^{\mathrm{}}`$-norm of the Stieltjes integral. We then use this estimate to prove the general existence-uniqueness theorem for solutions of ODE’s with Hölder continuous forcing. In section 3 we apply these results to construct the stochastic integrals with respect to fBm and show the existence and uniqueness of solutions of stochastic differential equations with fBm. Previous results in this direction were obtained by \[Ln, DH\]. In section 4 we use Theorem 2 to get estimates on the tails of the stochastic integral with respect to fBm. We also consider the question of existence of fBm stochastic integral of deterministic functions and derive the probability distribution of its maximum. Throughout the paper, $`𝒞^\beta (I)`$ denotes the space of Hölder-continuous functions on the interval $`I`$ with exponent $`\beta `$. We will have many occasions to partition an interval into $`2^n`$ sub-intervals of equal size; the $`i`$-th partition point of the interval under discussion is denoted by $`s_i^n`$, and we put $`\mathrm{\Delta }f(s+s_i^n)=f(s+s_{i+1}^n)f(s+s_i^n)`$ if $`f`$ is a function defined on the interval.
## 2 Main Theorems.
In this section we consider the question of existence of a Stieltjes integral for functions of unbounded variation, give a formula which allows to estimate its upper bound and prove the existence-uniqueness theorem for ordinary differential equations with Hölder continuous forcing.
In applications (such as the construction of the stochastic calculus for fBm), it is often interesting to consider Stieltjes integrals $`f𝑑g`$ for functions of unbounded variation. The difficulty in constructing the integral is that the upper bounds on Riemann sums $`|f||\mathrm{\Delta }g|`$ diverge. However this problem can be solved on certain classes of $`f,g,`$ since, if $`g`$ oscillates fast enough, the nearby terms in the Riemann sum $`f\mathrm{\Delta }g`$ may cancel. It is shown in \[Yg, Kn\] that the Stieltjes integral exists on certain classes of Hölder continuous functions.
###### Theorem 1
(Young-Kondurar) Let $`f𝒞^\beta (\mathrm{I}\mathrm{R})`$, $`g𝒞^\gamma (\mathrm{I}\mathrm{R})`$. If $`\beta +\gamma >1,`$ then $`_0^tf𝑑g`$ exists as a Stieltjes integral for all $`t>0.`$
Generalizations of Theorem 1 can be found in \[Dy\].
The next formula is useful in estimating the $`L^{\mathrm{}}`$-norm of the Stieltjes integral and is the main tool used in this paper.
###### Theorem 2
Let $`f`$, $`g`$, $`\beta `$, and $`\gamma `$ be as in Statement 1. Then
$$_s^tf(\tau )𝑑g(\tau )=f(s)(g(t)g(s))+\underset{k=1}{\overset{\mathrm{}}{}}\underset{i=0}{\overset{2^{k1}1}{}}\mathrm{\Delta }f(s+s_{2i}^k)\mathrm{\Delta }g(s+s_{2i+1}^k).$$
(1)
As far as we know this form of the Stieltjes integral has not appeared before. The idea behind the proof of Theorem 2 is to write a recursion between the Riemann sums on finer partitions of the interval and the Riemann sum on coarser partitions of the interval, very much like in Renormalization group. The same idea is used in \[Ru\] to give a new proof of Theorem 1, which amounts to showing that the right hand side of the equation (1) does not depend on the sequence of partitions we choose. Below we will need the change of variables formula which is a deterministic analog of the Itô’s formula for Brownian motion.
###### Lemma 2.1
Let $`u:[0,\mathrm{})\times \mathrm{I}\mathrm{R}\mathrm{I}\mathrm{R}`$, $`e,f:[0,\mathrm{})\mathrm{I}\mathrm{R}`$, $`1>\gamma >\frac{1}{2}`$, $`\beta >1\gamma `$, $`T>0`$. Suppose $`fC^\beta ([0,T])`$, $`gC^\gamma ([0,T])`$, $`e`$ is continuous, $`u`$ is differentiable in $`t`$ with continuous $`u/t`$ and twice differentiable in $`x`$.For $`0tT`$, consider the Stieltjes integral
$`\eta (t)=\eta (0)+{\displaystyle _0^t}e(s)𝑑s+{\displaystyle _0^t}f(s)𝑑g(s).`$
Then for all $`0tT`$, $`v(t)=u(t,\eta (t))`$ is also a Stieltjes integral whose differential is
$$dv=\frac{u}{t}dt+\frac{u}{\eta }d\eta .$$
(2)
Note that for $`\gamma >\frac{1}{2}`$, the change of variable formula is the same as in the case of ordinary calculus. This follows from the fact that the quadratic variation of $`g`$ is zero and therefore the terms of order $`d\eta ^2`$ are negligible.
Theorems 1 and 2 can be used to prove the following general theorem on the existence and uniqueness of solutions of ordinary differential equations with Hölder continuous forcing.
###### Theorem 3
Let $`b,\sigma :[0,\mathrm{})\times \mathrm{I}\mathrm{R}\mathrm{I}\mathrm{R}`$, $`g𝒞^\gamma (\mathrm{I}\mathrm{R})`$ and $`1/2<\gamma 1`$. Suppose $`b`$ is globally Lipschitz in $`t`$ and $`x`$, and $`\sigma 𝒞^1(\mathrm{I}\mathrm{R})`$ with $`\sigma `$, $`\sigma _t^{}`$, $`\sigma _x^{}`$ globally Lipschitz in $`t`$ and $`x`$. Then for every $`T>0`$ and $`\gamma >\beta >1\gamma `$, the ordinary differential equation
$$dx(t)=b(t,x(t))dt+\sigma (t,x(t))dg(t),x(0)=x_0$$
(3)
has a unique solution in $`C^\beta ([0,T])`$.
## 3 Applications to fractional Brownian motion.
The natural way of constructing the fBm stochastic integral $`_0^tf(\tau ,\omega )𝑑\xi _\alpha (\tau ,\omega )`$ is to define it as a Stieltjes integral for a.e. $`\omega `$. Since $`\xi _\alpha (\tau ,\omega )𝒞^\gamma (\mathrm{I}\mathrm{R})`$ for $`\gamma \frac{1+\alpha }{2}`$ with probability one, Theorem 1 implies that the fBm stochastic integral $`_0^tf(\tau ,\omega )𝑑\xi _\alpha (\tau ,\omega )`$ exists for all $`f𝒞^\beta (\mathrm{I}\mathrm{R})`$ with $`\beta >\frac{1\alpha }{2}`$. The paper \[Ln\] contains a special case of this result for functions $`f(\xi _\alpha (,\omega ))𝒞^1(\mathrm{I}\mathrm{R})`$ for a.e. $`\omega `$, derived by expanding $`f`$ in Taylor series in $`\xi _\alpha `$ and using the fact that the quadratic variation of $`\xi _\alpha `$ is zero.
Similarly, Theorem 2 holds with probability one for $`_0^tf(\tau ,\omega )𝑑\xi _\alpha (\tau ,\omega )`$ for all $`f𝒞^\beta (\mathrm{I}\mathrm{R})`$ with $`\beta >\frac{1\alpha }{2}`$.
Itô’s formula for fBm can be stated pathwise, as a corollary of Lemma 2.1, however it can be stated also under weaker assumptions. Itô’s formula for fBm has been established under very different assumptions in \[DH\] for $`0<\alpha <1,`$ and in \[DU\] for $`1<\alpha <1`$.
###### Lemma 3.1
Let $`u:[0,\mathrm{})\times \mathrm{I}\mathrm{R}\mathrm{I}\mathrm{R}`$, $`e,f:[0,\mathrm{})\times \mathrm{\Omega }\mathrm{I}\mathrm{R}`$, $`0<\alpha <1`$, $`\beta >\frac{1\alpha }{2}`$, $`T>0`$. Suppose $`fC^\beta ([0,T])`$, $`e`$ is continuous, $`u`$ is differentiable with continuous $`u/t`$ and $`u/xC^\gamma ([0,T])`$. For $`0tT`$, consider the stochastic integral
$`\eta (t,\omega )=\eta (0,\omega )+{\displaystyle _0^t}e(s,\omega )𝑑s+{\displaystyle _0^t}f(s,\omega )𝑑\xi _\alpha (s,\omega ).`$
Suppose $`\underset{0tT}{sup}E\left(\frac{^2u}{(\eta )^2}(t,\eta (t,\omega ))\right)^2<\mathrm{}`$. Then for all $`0tT`$, $`v(t)=u(t,\eta (t,\omega ))`$ is also a stochastic integral whose differential is
$`dv={\displaystyle \frac{u}{t}}dt+{\displaystyle \frac{u}{\eta }}d\eta .`$
Note that for $`0<\alpha <1`$, Itô’s formula for fBm is the same as in the deterministic case. This follows from the fact that the quadratic variation of fBm is zero. Theorem 3 implies the following existence-uniqueness theorem for SDE’s with fBm:
###### Theorem 4
Let $`b,\sigma :[0,\mathrm{})\times \mathrm{I}\mathrm{R}\mathrm{I}\mathrm{R}`$, $`Z:𝛀\mathrm{I}\mathrm{R}`$ and $`0<\alpha <1`$, $`\frac{1\alpha }{2}<\beta <\frac{1+\alpha }{2}`$. Suppose $`b`$ is globally Lipshitz in $`t`$ and $`x`$, and $`\sigma 𝒞^1(\mathrm{I}\mathrm{R})`$ with $`\sigma `$, $`\sigma _t^{}`$, $`\sigma _x^{}`$ globally Lipshitz in $`t`$ and $`x`$. Then for every $`T>0`$ the SDE
$$dX(t,\omega )=b(t,X(t,\omega ))dt+\sigma (t,X(t,\omega ))d\xi _\alpha (t,\omega ),X(0,\omega )=Z(\omega )$$
(4)
has a unique solution in $`C^\beta ([0,T])`$ with probability 1.
In \[Ln, DH\], the existence and uniqueness theorem for solutions of stochastic differential equations was proved when the diffusion coefficient is a function of time $`t`$ only; in \[Ln\], the existence theorem was proved also when the drift and diffusion coefficients are functions of $`X`$ only. Both papers adapt the methods used for oBm. The new idea in this paper is to use the formula in Theorem 2, which allows us to prove the existence and uniqueness theorem in the general case, when the drift and diffusion are functions of both $`t`$ and $`X`$.
## 4 Additional results for fractional Brownian motion.
When the integrand is of the form $`f(\tau ,\xi _\alpha (\tau ,\omega ))`$, we can obtain estimates of the tail of the maximum of stochastic integrals from Theorem 2 and from tail estimates of the Hölder coefficient of fBm.
###### Theorem 5
Let $`f:[0,1]\times \mathrm{I}\mathrm{R}\mathrm{I}\mathrm{R}`$ and $`\frac{1}{2}<\gamma <\frac{1+\alpha }{2}`$, $`\delta =\frac{1+\alpha }{2}\gamma `$. Suppose $`f`$ is differentiable with bounded $`|f_t^{}(t,x)|`$, $`|f_x^{}(t,x)|`$. Write $`|f_t^{}|=\underset{[0,1],\mathrm{I}\mathrm{R}}{sup}|f_t^{}(t,x)|`$ and $`|f_x^{}|=\underset{[0,1],\mathrm{I}\mathrm{R}}{sup}|f_x^{}(t,x)|`$. Then
$`P\left(\underset{0t1}{\mathrm{max}}{\displaystyle _0^t}f(\tau ,\xi _\alpha (\tau ,\omega ))𝑑\xi _\alpha (\tau ,\omega )>\lambda \right){\displaystyle \frac{2^\gamma +1}{2^\gamma 1}}\sqrt{{\displaystyle \frac{2}{\pi }}}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{2^{(1\delta )n}}{\nu }}\mathrm{exp}\{{\displaystyle \frac{2^\gamma 1}{2^\gamma +1}}\nu ^22^{2n\delta 1}\}.`$ (5)
where
$`\nu ={\displaystyle \frac{2^{2\gamma 1}1}{|f_x^{}|}}\left[\sqrt{(|f(0,0)|+{\displaystyle \frac{|f_t^{}|}{2^{\gamma +1}2}})^2+{\displaystyle \frac{4}{2^{2\gamma }2}}|f_x^{}|\lambda }|f(0,0)|{\displaystyle \frac{|f_t^{}|}{2^{\gamma +1}2}}\right].`$ (6)
More information on the stochastic integral is available when $`f`$ is a function of $`t`$ only, since in this case the techniques used for oBm can be applied.
###### Statement 4.1
Let $`0<\alpha <1`$, and let $`f(t,\omega )=f(t)`$ be a function of $`t`$ only. If $`fL^{\frac{2}{1+\alpha }}([0,\mathrm{}))`$, then the fBm stochastic integral $`_0^tf(\tau )𝑑\xi _\alpha (\tau ,\omega )`$ exists in $`L^2([0,\mathrm{})\times \mathrm{\Omega })`$ for all $`t[0,\mathrm{})`$.
The proof is based on the Hardy-Littlewood-Sobolev inequality (see \[LL\]).
Since $`_0^tf(\tau )𝑑\xi _\alpha (\tau ,\omega )`$ is a Gaussian process, we can show the following properties:
###### Statement 4.2
Let $`\alpha `$ and $`f`$ be as in Theorem 4.1, and let $`0<\beta <\alpha `$. Write $`q_f(s,t)=_s^t_s^tf(u)f(v)\frac{dudv}{|uv|^{1\alpha }}`$. If $`fL^{\frac{2}{1+\beta }}([0,1])`$, then
1. for a.e. $`\omega `$, $`_0^tf(\tau )𝑑\xi _\alpha (\tau ,\omega )`$ has a $`t`$-continuous version for all $`t[0,1]`$;
2. for every real $`r`$, $`P\left(\underset{0t1}{\mathrm{max}}_0^tf(\tau )𝑑\xi _\alpha (\tau ,\omega )>\lambda \right)`$ is bounded from above by
$$_{\lambda r/\sqrt{q_{f_+}(0,1)}}^{\mathrm{}}+_{\lambda (1r)/\sqrt{q_f_{}(0,1)}}^{\mathrm{}}\sqrt{\frac{2}{\pi }}e^{x^2/2}𝑑x,\mathrm{where}f_\pm =\frac{|f|\pm f}{2};$$
3. for every integer $`m2`$ and real $`\lambda \sqrt{1+\mathrm{log}m^4}`$, $`P\left(\underset{0t1}{\mathrm{max}}\left|_0^tf(\tau )𝑑\xi _\alpha (\tau ,\omega )\right|>\lambda \right)`$ is bounded from above by
$$_{\lambda /c}^{\mathrm{}}\frac{5}{2}m^2e^{x^2/2}𝑑x,\mathrm{where}c=\underset{0s,t1}{sup}\sqrt{q_f(s,t)}+(2+\sqrt{2})_1^{\mathrm{}}\underset{|st|<m^{x^2}}{sup}\sqrt{q_f(s,t)}dx<\mathrm{}.$$
(1) follows from Kolmogorov’s continuity criterion, the bound (2) follows from Slepian’s lemma \[Sl\] (the interesting case is $`0<r<1`$) , and (3) follows from Fernique’s inequality \[F\].
### 4.1 Acknowledgments.
I would like to thank my advisor Yakov Sinai for suggesting the problem and for all his help during the course of the work, Almut Burchard for suggesting a way to simplify the argument of section 5.3, the referee for his helpful suggestions in restructuring the paper and especially for an elegant way to simplify the argument of section 5.4, Tadashi Tokieda for his help and comments on the text and Luc Rey-Bellet for his comments on the text.
## 5 Proofs of theorems.
### 5.1 Proof of Theorem 2.
Proof: The idea of the proof is to write the Riemann sums on the smaller scales in terms of the Riemann sums on the larger scales as in Renormalization Group. Denote by $`S^n(f)`$ the Riemann sum of $`f`$ corresponding to the partition of $`[0,1]`$ into $`2^n`$ equal sub-intervals. We have
$`S^n(f)`$ $`=`$ $`S^{n1}(f)+{\displaystyle \underset{i=0}{\overset{2^{n1}1}{}}}\mathrm{\Delta }f(s_{2i}^n)\mathrm{\Delta }g(s_{2i+1}^n)`$
$`=`$ $`\mathrm{}`$
$`=`$ $`S^0(f)+{\displaystyle \underset{k=1}{\overset{n}{}}}{\displaystyle \underset{i=0}{\overset{2^{k1}1}{}}}\mathrm{\Delta }f(s_{2i}^k)\mathrm{\Delta }g(s_{2i+1}^k).`$
As $`n\mathrm{}`$, $`S^n(f)`$ converges to $`_0^tf(\tau )𝑑g(\tau )`$ by Theorem 1. The right hand side converges provided $`\beta +\gamma >1`$.
### 5.2 Proof of the change of variables formula.
The proof given here follows \[Mk\] for the most part.
Proof: We can write the integral version of equation (2):
$$v(t)v(0)=_0^t\frac{u}{s}𝑑s+_0^t\frac{u}{\eta }𝑑\eta $$
From Theorem 2 and from continuity of $`e`$ it follows that:
$$\mathrm{\Delta }\eta (t)=f(t)\mathrm{\Delta }g(t)+e(t)\mathrm{\Delta }t+O(\mathrm{\Delta }t)^{\beta +\gamma }+o(\mathrm{\Delta }t).$$
$`v(t,g(t))v(0)={\displaystyle \underset{k[2^nt]}{}}\{u({\displaystyle \frac{k}{2^n}},\eta ({\displaystyle \frac{k}{2^n}}))u({\displaystyle \frac{k1}{2^n}},\eta ({\displaystyle \frac{k}{2^n}}))\}`$
$`+{\displaystyle \underset{k[2^nt]}{}}\{u({\displaystyle \frac{k1}{2^n}},\eta ({\displaystyle \frac{k}{2^n}}))u({\displaystyle \frac{k1}{2^n}},\eta ({\displaystyle \frac{k1}{2^n}}))\}+u(t,\eta (t))u({\displaystyle \frac{[2^nt]}{2^n}},\eta ({\displaystyle \frac{[2^nt]}{2^n}}))`$
$`={\displaystyle \underset{k[2^nt]}{}}\{{\displaystyle \frac{u}{t}}({\displaystyle \frac{k1}{2^n}},\eta ({\displaystyle \frac{k}{2^n}})){\displaystyle \frac{1}{2^n}}+o({\displaystyle \frac{t}{2^n}})\}`$
$`+{\displaystyle \underset{k[2^nt]}{}}\{{\displaystyle \frac{u}{\eta }}[{\displaystyle \frac{k1}{2^n}},\eta ({\displaystyle \frac{k1}{2^n}})](\eta ({\displaystyle \frac{k}{2^n}})\eta ({\displaystyle \frac{k1}{2^n}}))`$
$`+{\displaystyle \frac{1}{2}}{\displaystyle \frac{^2u}{\eta ^2}}({\displaystyle \frac{k1}{2^n}},\eta ({\displaystyle \frac{k1}{2^n}}))(\eta ({\displaystyle \frac{k}{2^n}})\eta ({\displaystyle \frac{k1}{2^n}}))^2+o(\mathrm{\Delta }\eta ^2)\}`$
$`={\displaystyle \underset{k[2^nt]}{}}\{({\displaystyle \frac{u}{t}}({\displaystyle \frac{k1}{2^n}},\eta ({\displaystyle \frac{k}{2^n}}))+{\displaystyle \frac{u}{\eta }}({\displaystyle \frac{k1}{2^n}},\eta ({\displaystyle \frac{k1}{2^n}}))e[{\displaystyle \frac{k1}{2^n}}]){\displaystyle \frac{1}{2^n}}+o({\displaystyle \frac{t}{2^n}})\}`$
$`+{\displaystyle \underset{k[2^nt]}{}}\{{\displaystyle \frac{u}{\eta }}({\displaystyle \frac{k1}{2^n}},\eta ({\displaystyle \frac{k1}{2^n}}))f({\displaystyle \frac{k1}{2^n}})(g({\displaystyle \frac{k}{2^n}})g({\displaystyle \frac{k1}{2^n}}))`$
$`+{\displaystyle \frac{1}{2}}{\displaystyle \frac{^2u}{\eta ^2}}({\displaystyle \frac{k1}{2^n}},\eta ({\displaystyle \frac{k1}{2^n}}))f({\displaystyle \frac{k1}{2^n}})^2(g({\displaystyle \frac{k}{2^n}})g({\displaystyle \frac{k1}{2^n}}))^2+o(\mathrm{\Delta }g^2)\}`$
$`={\displaystyle _0^t}\left({\displaystyle \frac{u}{s}}+{\displaystyle \frac{u}{\eta }}e(s)\right)𝑑s+o(1)+{\displaystyle _0^t}{\displaystyle \frac{u}{\eta }}f(s)𝑑g(s)+{\displaystyle \frac{1}{2}}{\displaystyle \underset{k[2^nt]}{}}u_{\eta \eta }f^2\mathrm{\Delta }g^2+o\left({\displaystyle \mathrm{\Delta }g^2}\right).`$
Since $`u_{\eta \eta }f^2\mathrm{\Delta }g^2O\left(\frac{t}{2^n}\right)^{2\gamma },`$ and $`\gamma >\frac{1}{2}`$, $`_{k[2^nt]}u_{\eta \eta }f^2\mathrm{\Delta }g^20`$ as $`n\mathrm{}`$. Similarly, $`o\left(\mathrm{\Delta }g^2\right)0`$ as $`n\mathrm{}`$.
### 5.3 Proof of Theorem 3: local existence and uniqueness.
In this section we will prove the local existence and uniqueness result. We will derive the global result in Theorem 3 by showing that we can apply the local existence and uniqueness result repeatedly, taking as initial condition the value of the solution at the end of the previous interval.
Fix $`T>0`$, $`s[0,T]`$ and $`a\mathrm{I}\mathrm{R}`$. Define an integral operator
$$F_tX=_s^tb(\tau ,X(\tau ))𝑑\tau +_s^t\sigma (\tau ,X(\tau ))𝑑g(\tau )+a.$$
For $`s=0`$ and $`a=x_0`$, solutions of the ODE (3) are exactly fixed points of $`F`$.
For $`0<\beta <1`$, and for $`\epsilon >0`$, consider the Banach space $`𝒞^\beta ([s,s+\epsilon ])`$ with the norm $`f_\beta =\underset{t}{\mathrm{max}}|f(t)|+\underset{t_1t_2}{\mathrm{max}}\frac{|f(t_1)f(t_2)|}{|t_1t_2|^\beta }`$. For $`K>0`$, consider the closed subset $`𝒞_K^\beta ([s,s+\epsilon ])=\{f:|f(t_2)f(t_1)|K|t_2t_1|^\beta ,t_1,t_2[s,s+\epsilon ]\}`$ of $`𝒞^\beta ([s,s+\epsilon ])`$. Finally, consider the closed subset $`𝒞_K^\beta ([s,s+\epsilon ],a)=\{f𝒞_K^\beta ([s,s+\epsilon ]):f(s)=a\}`$ of $`𝒞^\beta ([s,s+\epsilon ])`$. For a given $`g𝒞^\gamma ([0,T])`$, there is $`L>0`$ such that $`g𝒞_L^\gamma ([0,T])`$. We will show by the contraction mapping theorem that for given $`T`$, $`s`$, $`a`$, $`K`$, $`L`$, there exists $`\epsilon >0`$ such that $`F_t`$ has a unique fixed point on $`𝒞_K^\beta ([s,s+\epsilon ],a)`$. We need to establish that $`F_t`$ maps $`𝒞_K^\beta ([s,s+\epsilon ],a)`$ into itself and that it is a contraction. This will be done in Lemma 5.3. The necessary preliminary estimates are obtained in Corollaries 5.1 and 5.2 of Theorem 2. In what follows we will consider $`T,K,L>0`$ to be fixed.
###### Corollary 5.1
Let $`\beta `$ and $`\gamma `$ be as in Theorem 1, $`T>0`$ and let $`g𝒞_L^\gamma ([0,T])`$. Then for every $`s,t[0,T]`$,
$$_s^tX(\tau )𝑑g(\tau )_{\mathrm{}}=LX_\beta (ts)^\gamma (1+\frac{(ts)^\beta }{2^{\beta +\gamma }2}).$$
(7)
Proof: By Theorem 2
$$|_s^tX(\tau )𝑑g(\tau )||X(s)||g(t)g(s)|+\underset{k=1}{\overset{\mathrm{}}{}}\underset{i=0}{\overset{2^{k1}1}{}}|\mathrm{\Delta }X(s+s_{2i}^k)||\mathrm{\Delta }g(s+s_{2i+1}^k)|.$$
(8)
Since $`|X(s)|X_\beta `$, $`|\mathrm{\Delta }X(s+s_{2i}^k)|X_\beta (\frac{ts}{2^k})^\beta `$ and $`g𝒞_L^\gamma ([0,T])`$, we obtain after performing the sum in (8),
$$|_s^tX(\tau )𝑑g(\tau )|LX_\beta (ts)^\gamma +LX_\beta \frac{(ts)^{\beta +\gamma }}{2^{\beta +\gamma }2}.$$
###### Corollary 5.2
Let $`\epsilon >0`$ and let $`\beta `$, $`\gamma `$, $`g`$ and $`T`$ be as in Corollary 5.1. Then, for every $`s[0,T]`$ and $`t[s,s+\epsilon ]`$,
$$_s^tX(\tau )𝑑\tau _\beta X_{\mathrm{}}\epsilon ^{1\beta }(1+\epsilon ^\beta ),$$
(9)
and
$$_s^tX(\tau )𝑑g(\tau )_\beta LX_\beta \epsilon ^{\gamma \beta }(1+\epsilon ^\beta )(1+\frac{\epsilon ^\beta }{2^{\beta +\gamma }2}),$$
(10)
Proof:
$`{\displaystyle _s^t}X(\tau )𝑑\tau _\beta `$ $`=`$ $`\underset{t[s,s+\epsilon ]}{\mathrm{max}}|{\displaystyle _s^t}X(\tau )𝑑\tau |+\underset{t_1t_2[s,s+\epsilon ]}{\mathrm{max}}{\displaystyle \frac{|_{t_1}^{t_2}X(\tau )𝑑\tau |}{|t_2t_1|^\beta }}`$
$``$ $`X_{\mathrm{}}(\epsilon +\epsilon ^{1\beta }).`$
From Corollary 5.1 we obtain
$$\underset{t[s,s+\epsilon ]}{\mathrm{max}}|_s^tX(\tau )𝑑g(\tau )|LX_\beta \epsilon ^\gamma (1+\frac{\epsilon ^\beta }{2^{\beta +\gamma }2}),$$
(11)
similarly we obtain
$$\underset{t_1t_2[s,s+\epsilon ]}{\mathrm{max}}\frac{|_{t_1}^{t_2}X(\tau )𝑑g(\tau )|}{|t_1t_2|^\beta }LX_\beta \epsilon ^{\gamma \beta }(1+\frac{\epsilon ^\beta }{2^{\beta +\gamma }2}).$$
(12)
The result (10) follows from (11) and (12).
###### Lemma 5.3
Let $`b,\sigma :[0,\mathrm{})\times \mathrm{I}\mathrm{R}\mathrm{I}\mathrm{R}`$, $`1\gamma >1/2`$, $`T,K,L>0`$. Suppose $`b`$ and $`\sigma `$ are Lipschitz in $`x`$ and $`t`$. Then there exists $`\epsilon _1>0`$, such that for all $`\gamma >\beta >1\gamma `$ and $`t[s,s+\epsilon _1],`$ the operator $`F_t`$ maps $`𝒞_K^\beta ([s,s+\epsilon _1],a)`$ into itself.
Notation. We shall denote the Lipschitz coefficients of $`b`$ and $`\sigma `$ by $`𝐁`$ and $`𝐒`$ respectively.
Proof: Let $`\epsilon >0`$ and $`t_1,t_2[s,s+\epsilon ]`$. It is sufficient to demonstrate that $`F_tX_\beta K`$ for a sufficiently small $`\epsilon >0`$.
By the triangle inequality and by Corollary 5.2,
$$F_tX_\beta \epsilon ^{\gamma \beta }(1+\epsilon ^\beta )(b(t,X(t))_{\mathrm{}}\epsilon ^{1\gamma }+L\sigma (t,X(t))_\beta (1+\frac{\epsilon ^\beta }{2^{\beta +\gamma }2})).$$
(13)
Since $`b`$ and $`\sigma `$ are Lipschitz and $`X𝒞_K^\beta ([s,s+\epsilon ],a)`$, it is easy to see that
$$b(t,X(t))_{\mathrm{}}|b(s,a)|+𝐁(\epsilon +K\epsilon ^\beta ),$$
(14)
and
$$\sigma (t,X(t))_\beta |\sigma (s,a)|+𝐒(\epsilon ^{1\beta }+K)(1+\epsilon ^\beta ).$$
(15)
Substituting (14) and (15) into (13), we obtain
$`F_tX_\beta `$ $``$ $`\epsilon ^{\gamma \beta }(1+\epsilon ^\beta )[|b(s,a)|+𝐁(\epsilon +K\epsilon ^\beta ))\epsilon ^{1\gamma }`$ (17)
$`+L(1+{\displaystyle \frac{\epsilon ^\beta }{2^{\beta +\gamma }2}})(|\sigma (s,a)|+𝐒(1+\epsilon ^\beta )(\epsilon ^{1\beta }+K))].`$
Since the right hand side is an increasing continuous function of $`\epsilon `$ and is $`0`$ at $`\epsilon =0`$, it equals $`K`$ at some $`\epsilon _1`$. For this choice of $`\epsilon _1`$ (or any smaller $`\epsilon _1`$), $`F_t`$ maps $`𝒞_K^\beta ([s,s+\epsilon _1],a)`$ into itself.
###### Lemma 5.4
Assume the same hypothesis as in Lemma 5.3. Suppose $`b`$ is Lipschitz in $`t`$ and $`x`$ and $`\sigma 𝒞^1([0,\mathrm{})\times \mathrm{I}\mathrm{R})`$ with $`\sigma _t^{}(t,x)`$, $`\sigma _x^{}(t,x)`$ Lipschitz in $`x`$. Then there exists $`\epsilon _2>0`$ such that for all $`\gamma >\beta >1\gamma `$ and $`t[s,s+\epsilon _2]`$, the operator $`F_t`$ is a contraction on $`𝒞_K^\beta ([s,s+\epsilon _2])`$.
Notation. We shall denote the Lipschitz coefficient of $`b`$ by $`𝐁`$ and the Lipschitz coefficient of $`\sigma `$, $`\sigma _t^{}`$, $`\sigma _x^{}`$ by $`𝐒`$.
Proof: We need to show that there exist $`\epsilon _2>0`$ and $`\lambda <1`$ such that for all $`t[s,s+\epsilon _2]`$ and all $`X,Y𝒞_K^\beta ([s,s+\epsilon _2])`$,
$`F_tXF_tY_\beta \lambda XY_\beta .`$
By the triangle inequality and by Corollary 5.2,
$`F_tXF_tY_\beta `$ $``$ $`\epsilon ^{\gamma \beta }(1+\epsilon ^\beta )[||b(t,X(t))b(t,Y(t))||_{\mathrm{}}\epsilon ^{1\gamma }`$ (19)
$`+L||\sigma (t,X(t))\sigma (t,Y(t))||_\beta (1+{\displaystyle \frac{\epsilon ^\beta }{2^{\beta +\gamma }2}})].`$
We estimate the two terms in (19) separately. Since $`b`$ is Lipschitz,
$$b(t,X(t))b(t,Y(t))_{\mathrm{}}𝐁|X(t)Y(t)|𝐁XY_\beta .$$
(20)
Now we will estimate $`\sigma (t,X(t))\sigma (t,Y(t))_\beta `$. Since $`\sigma `$ is differentiable,
$$\underset{t}{\mathrm{max}}|\sigma (t,X(t))\sigma (t,Y(t))|𝐒XY_\beta .$$
(21)
By the fundamental theorem of calculus,
$$\sigma (t,X(t))\sigma (t,Y(t))=(X(t)Y(t))_0^1\sigma _x^{}(t,\nu X(t)+(1\nu )Y(t))𝑑\nu .$$
Therefore
$`|\sigma (t_1,X(t_1))\sigma (t_1,Y(t_1))\sigma (t_2,X(t_2))+\sigma (t_2,Y(t_2))|=`$
$`|(X(t_1)Y(t_1)X(t_2)+Y(t_2)){\displaystyle _0^1}\sigma _x^{}(t_1,\nu X(t_1)+(1\nu )Y(t_1))d\nu `$
$`+(X(t_2)Y(t_2)){\displaystyle _0^1}(\sigma _x^{}(t_1,\nu X(t_1)+(1\nu )Y(t_1))\sigma _x^{}(t_2,\nu X(t_2)+(1\nu )Y(t_2)))d\nu |`$
$`XY_\beta (t_2t_1)^\beta 𝐒+XY_\beta \left(𝐒(t_2t_1)+𝐒(\nu |X(t_1)X(t_2)|+(1\nu )|Y(t_1)Y(t_2)|)\right)`$
$`XY_\beta (t_2t_1)^\beta 𝐒(1+(t_2t_1)^{1\beta }+K).`$
Consequently
$$\sigma (t,X(t))\sigma (t,Y(t))_\beta XY_\beta 𝐒(2+\epsilon ^{1\beta }+K).$$
(22)
Substituting (21) and (22) in (19), we obtain
$$F_tXF_tY_\beta \epsilon ^{\gamma \beta }(1+\epsilon ^\beta )\left[𝐁\epsilon ^{1\gamma }+L𝐒(2+\epsilon ^{1\beta }+K)(1+\frac{\epsilon ^\beta }{2^{\beta +\gamma }2})\right]XY_\beta .$$
(23)
The coefficient of $`XY_\beta `$ is an increasing function of $`\epsilon `$ and is $`0`$ at $`\epsilon =0`$. Choose $`\epsilon _2>0`$ small enough so that this coefficient is less than 1 and $`\epsilon _2\epsilon _1`$. Then $`F`$ is a contraction on $`𝒞_K^\beta ([s,s+\epsilon _2]).`$
Combining Lemmas 5.3 and 5.4, we obtain a local existence and uniqueness result:
###### Corollary 5.5
Assume the same hypothesis as in Lemma 5.4. Then there exists $`\epsilon >0`$ depending on $`s`$ such that for all $`t[s,s+\epsilon ]`$, ODE (3) has a unique solution in $`𝒞_K^\beta ([s,s+\epsilon ],a)`$. In particular, if ODE (3) has a solution $`X`$ on $`[0,s]`$, then there exists $`\epsilon >0`$ depending on $`s`$ such that for all $`t[s,s+\epsilon ]`$, ODE (3) has a unique solution $`Y`$ in $`𝒞_K^\beta ([s,s+\epsilon ],X(s))`$.
Proof: Take $`\epsilon `$ to be $`\epsilon _2`$ of Lemma 5.4. Since $`F_t`$ is a contraction on the closed subset $`𝒞_K^\beta ([s,s+\epsilon ],a)`$ of the complete metric space $`𝒞^\beta ([s,s+\epsilon ])`$, it has a unique fixed point $`X`$ in $`𝒞_K^\beta ([s,s+\epsilon ],a)`$. From the definition of $`F_t`$ it follows that $`X`$ is a unique solution of ODE (3) on $`[s,s+\epsilon ]`$ in $`𝒞_K^\beta ([s,s+\epsilon ],a)`$.
The sufficient conditions on $`\epsilon `$ in Corollary 5.5 are given by inequalities (17) and (23) with $`a`$ replaced by $`X(s)`$:
$`\epsilon ^{\gamma \beta }(1+\epsilon ^\beta )[|b(s,X(s))|+𝐁(\epsilon +K\epsilon ^\beta ))\epsilon ^{1\gamma }`$ (24)
$`+L(1+{\displaystyle \frac{\epsilon ^\beta }{2^{\beta +\gamma }2}})(|\sigma (s,X(s))|+𝐒(1+\epsilon ^\beta )(\epsilon ^{1\beta }+K))]K.`$ (25)
$$\epsilon ^{\gamma \beta }(1+\epsilon ^\beta )(𝐁\epsilon ^{1\gamma }+L𝐒(2+\epsilon ^{1\beta }+K)(1+\frac{\epsilon ^\beta }{2^{\beta +\gamma }2}))<1.$$
(26)
Inequality (26) does not depend on $`X(s)`$.
### 5.4 Proof of Theorem 3: global existence and uniqueness.
By Corollary 5.5 with $`s=0`$ and $`a=x_0`$, ODE (3) has a unique solution on $`[0,\epsilon _0]`$, where $`\epsilon _0`$ satisfies (24) and (26) with $`s=0`$. By using Corollary 5.5 $`n`$ times we obtain that the solution exists on $`[0,\epsilon _0+\mathrm{}+\epsilon _{n1}]`$. ODE (3) has a solution on $`[0,T]`$ if there exists $`m>0`$ such that $`_{i=0}^m\epsilon _iT`$. This is true if $`b`$ and $`\sigma `$ are globally bounded, since in this case we can choose $`\epsilon _i=\epsilon _0`$ (to see this we can substitute the global bounds on $`b`$, $`\sigma `$, $`𝐁`$, $`𝐒`$ into (24) and (26)). In the case when $`b`$ and $`\sigma `$ grow at most linearly, we will use a change of variables to reduce it to the case of globally bounded $`b`$ and $`\sigma `$.
Proof: Existence: Suppose that $`b`$ and $`\sigma `$ are bounded on $`[0,\mathrm{})\times \mathrm{I}\mathrm{R}`$, then taking the upper bound on $`b`$ and $`\sigma `$ in (24) we get that $`\epsilon _i`$ satisfying (24) and (26) does not depend on $`i`$. In this case the global existence is established. Now suppose that $`b`$ and $`\sigma `$ satisfy the assumptions of Theorem 3. Consider the ODE
$$dy(t)=\frac{b(t,\mathrm{tan}y(t))}{1+(\mathrm{tan}y(t))^2}dt+\frac{\sigma (t,\mathrm{tan}y(t))}{1+(\mathrm{tan}y(t))^2}dg(t).$$
(27)
This ODE has globally bounded coefficients satisfying the assumptions of Theorem 3, and thus has a global solution on $`[0,T]`$. Now, $`x(t)=\mathrm{tan}y(t)`$ satisfies equation (3) (by Lemma 2.1). Thus equation (3) has a global solution on $`[0,T]`$.
Uniqueness: Let $`Y_1`$ and $`Y_2`$ be two solutions in $`𝒞^\beta ([0,T])`$. Then there exist $`K_1`$ and $`K_2`$ such that $`Y_1𝒞_{K_1}^\beta ([0,T])`$ and $`Y_2𝒞_{K_2}^\beta ([0,T])`$, so $`Y_1`$ and $`Y_2`$ are in $`𝒞_{\mathrm{max}\{K_1,K_2\}}^\beta ([0,T])`$. $`Y_1`$ and $`Y_2`$ coincide at the initial point $`t=0`$. Let $`t_{sup}`$ be the supremum of the set on which they coincide. Since both solutions are continuous, they coincide at $`t_{sup}`$ as well. $`t_{sup}`$ must equal $`T`$, for otherwise we can make $`Y_1`$ and $`Y_2`$ coincide past $`t_{sup}`$ by Corollary 5.5.
### 5.5 Proof of Itô’s formula for fBm (Lemma 3.1).
Proof: By analogy with the proof of Lemma 2.1, we obtain
$`v(t,\xi (t))v(0)={\displaystyle _0^t}\left({\displaystyle \frac{u}{s}}+{\displaystyle \frac{u}{\eta }}e(s)\right)𝑑s+o(1)+{\displaystyle _0^t}{\displaystyle \frac{u}{\eta }}f(s)𝑑\xi _\alpha +{\displaystyle \frac{1}{2}}{\displaystyle \underset{k[2^nt]}{}}u_{\eta \eta }f^2\mathrm{\Delta }\xi _\alpha ^2+o\left({\displaystyle \mathrm{\Delta }\xi _\alpha ^2}\right).`$
Since $`E(u_{\eta \eta }f^2\mathrm{\Delta }\xi _\alpha ^2)f_{\mathrm{}}\sqrt{E(u_{\eta \eta })^2}\sqrt{E(\mathrm{\Delta }\xi _\alpha ^4)}=O\left(\frac{t}{2^n}\right)^{2\gamma },`$ by Chebyshev inequality
$$P(\underset{k[2^nt]}{}u_{\eta \eta }f^2\mathrm{\Delta }\xi _\alpha ^2\frac{1}{n})\mathrm{const}\frac{n}{2^{n(2\gamma 1)}},$$
and by the Borel-Cantelli lemma $`_{k[2^nt]}u_{\eta \eta }f^2\mathrm{\Delta }\xi _\alpha ^20`$ as $`n\mathrm{}`$ a.e. An analogous argument shows that $`o\left(\mathrm{\Delta }\xi _\alpha ^2\right)0`$ as $`n\mathrm{}`$ a.e. Thus Itô’s formula holds.
### 5.6 Proof of Theorem 5
Proof: Since $`f`$ is differentiable, and since $`\xi _\alpha (,\omega )C^\gamma ([0,1])`$ with probability 1, for a.e. $`\omega `$ there exists $`L(\omega )>0`$ such that
$$|\mathrm{\Delta }f(s_{2i}^k,\xi _\alpha (s_{2i}^k))||f_t^{}|\frac{t}{2^k}+|f_x^{}|\frac{L(\omega )t^\gamma }{2^{k\gamma }}$$
holds. From Theorem 2 we get
$`\left|{\displaystyle _0^t}f(\tau ,\xi _\alpha (\tau ))𝑑\xi _\alpha (\tau )\right||f(0,0)|L(\omega )t^\gamma +|f_t^{}|{\displaystyle \frac{L(\omega )t^{\gamma +1}}{2^{\gamma +1}2}}+|f_x^{}|{\displaystyle \frac{L(\omega )^2t^{2\gamma }}{2^{2\gamma }2}},`$
and so,
$`\underset{0t1}{\mathrm{max}}\left|{\displaystyle _0^t}f(\tau ,\xi _\alpha (\tau ))𝑑\xi _\alpha (\tau )\right||f(0,0)|L(\omega )+|f_t^{}|{\displaystyle \frac{L(\omega )}{2^{\gamma +1}2}}+|f_x^{}|{\displaystyle \frac{L(\omega )^2}{2^{2\gamma }2}}.`$
Therefore
$`P\{\omega :\underset{0t1}{\mathrm{max}}\left|{\displaystyle _0^t}f(\tau ,\xi _\alpha (\tau ))𝑑\xi _\alpha (\tau )\right|>\lambda \}`$
$`P\{\omega :|f(0,0)|L(\omega )+|f_t^{}|{\displaystyle \frac{L(\omega )}{2^{\gamma +1}2}}+|f_x^{}|{\displaystyle \frac{L(\omega )^2}{2^{2\gamma }2}}>\lambda \}`$
$`P\{\omega :L(\omega )>\nu \}`$
$`P\{\omega :t_1,t_2[0,1]\mathrm{s}.\mathrm{t}.|\xi _\alpha (t_2)\xi _\alpha (t_1)|>\nu |t_2t_1|^\gamma \},`$
where $`\nu `$ is given by (6).
It is easy to see that if
$`|\xi _\alpha ({\displaystyle \frac{k+1}{2^n}})\xi _\alpha ({\displaystyle \frac{k}{2^n}})|{\displaystyle \frac{2^\gamma 1}{2^\gamma +1}}{\displaystyle \frac{L(\omega )}{2^{n\gamma }}},n>0,0k2^n1,`$
then
$`|\xi _\alpha (t_2)\xi _\alpha (t_1)|L(\omega )|t_2t_1|^\gamma ,t_1,t_2[0,1].`$
Therefore
$`P\{\omega :t_1,t_2[0,1]\mathrm{s}.\mathrm{t}.|\xi _\alpha (t_2)\xi _\alpha (t_1)|>\nu |t_2t_1|^\gamma \}`$
$``$ $`P\{\omega :n>0,0k2^n1,\mathrm{s}.\mathrm{t}.|\xi _\alpha ({\displaystyle \frac{k+1}{2^n}})\xi _\alpha ({\displaystyle \frac{k}{2^n}})|>{\displaystyle \frac{2^\gamma 1}{2^\gamma +1}}{\displaystyle \frac{\nu }{2^{n\gamma }}}\}`$
$``$ $`{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \underset{k=0}{\overset{2^n1}{}}}P\{\omega :|\xi _\alpha ({\displaystyle \frac{k+1}{2^n}})\xi _\alpha ({\displaystyle \frac{k}{2^n}})|>{\displaystyle \frac{2^\gamma 1}{2^\gamma +1}}{\displaystyle \frac{\nu }{2^{n\gamma }}}\}`$
$``$ $`{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}2^n\sqrt{{\displaystyle \frac{2}{\pi }}}{\displaystyle _{\frac{2^\gamma 1}{2^\gamma +1}\nu 2^{n\delta }}^{\mathrm{}}}e^{x^2/2}𝑑x.`$
Using the estimate $`_c^{\mathrm{}}e^{x^2/2}𝑑x\frac{e^{c^2/2}}{c}`$, we obtain the desired result.
### 5.7 Proof of Statement 4.1.
Proof:
The proof will be reached via the step-by-step procedure used for oBm. Fix $`t>0`$.
Step 1. Let $`\varphi :[0,\mathrm{})\mathrm{I}\mathrm{R}`$ be a simple function of the form $`_{i=0}^{2^n1}\varphi (s_i^n)\chi _{[s_i^n,s_{i+1}^n]}`$, where $`\chi `$ is an indicator function and $`s_i^n=\frac{i}{2^n}t`$. Define
$$_0^t\varphi (\tau )𝑑\xi _\alpha (\tau )=\underset{i=0}{\overset{2^n1}{}}\varphi (s_i^n)\mathrm{\Delta }\xi _\alpha (s_i^n).$$
Step 2. Let $`g𝒞^{\mathrm{}}([0,\mathrm{}))`$. Approximate $`g`$ by a sequence of simple functions: $`\varphi _n(\tau )=_{i=0}^{2^n1}g(s_i^n)\chi _{[s_i^n,s_{i+1}^n]}(\tau )`$. Then $`\varphi _ng`$ uniformly on $`[0,t]`$ and $`_0^t|g(\tau )\varphi _n(\tau )|^{\frac{2}{1+\alpha }}𝑑\tau 0`$. Therefore the sequence $`\varphi _n`$ is Cauchy in $`L^{\frac{2}{1+\alpha }}([0,\mathrm{}))`$. To show that the sequence $`_0^t\varphi _n(\tau )𝑑\xi _\alpha (\tau )`$ is Cauchy in $`L^2([0,\mathrm{})\times \mathrm{\Omega })`$, we use the Hardy-Littlewood-Sobolev inequality:
$`E\left({\displaystyle _0^t}\varphi _m(\tau )𝑑\xi _\alpha (\tau ){\displaystyle _0^t}\varphi _n(\tau )𝑑\xi _\alpha (\tau )\right)^2`$
$`=C\alpha (\alpha +1){\displaystyle _0^t}{\displaystyle _0^t}{\displaystyle \frac{(\varphi _m(u)\varphi _n(u))(\varphi _m(v)\varphi _n(v))}{|uv|^{1\alpha }}}𝑑u𝑑v`$
$`\mathrm{const}\varphi _m\varphi _n_{\frac{2}{1+\alpha }}^20\mathrm{as}m,n\mathrm{}.`$
Thus the integral $`_0^tg(\tau )𝑑\xi _\alpha (\tau )`$ exists as the $`L^2`$-limit of $`_0^t\varphi _n(\tau )𝑑\xi _\alpha (\tau )`$.
Step 3. Let $`fL^{\frac{2}{1+\alpha }}([0,\mathrm{}))`$. Let $`j𝒞_c^{\mathrm{}}([0,\mathrm{}))`$ with $`_{[0,\mathrm{})}j=1`$. Define $`j_n(\tau )=\frac{1}{n}j(n\tau )`$ and $`g_n=j_nf`$. Then $`g_n𝒞^{\mathrm{}}([0,\mathrm{}))`$ and $`_0^t|g_n(\tau )f(\tau )|^{\frac{2}{1+\alpha }}𝑑\tau 0`$. In particular, the sequence $`g_n`$ is Cauchy in $`L^{\frac{2}{1+\alpha }}([0,\mathrm{}))`$. The Hardy-Littlewood-Sobolev inequality gives
$`E\left({\displaystyle _0^t}g_m(\tau )𝑑\xi _\alpha (\tau ){\displaystyle _0^t}g_n(\tau )𝑑\xi _\alpha (\tau )\right)^2`$ $``$ $`\mathrm{const}g_mg_n_{\frac{2}{1+\alpha }}^20\mathrm{as}m,n\mathrm{}.`$
Thus the integral $`_0^tf(\tau )𝑑\xi _\alpha (\tau )`$ exists as an $`L^2`$-limit of $`_0^tg_n(\tau )𝑑\xi _\alpha (\tau )`$.
Thus, for all $`fL^{\frac{2}{1+\alpha }}([0,1])`$, we can choose simple functions $`\varphi _n`$ converging in $`L^{\frac{2}{1+\alpha }}`$ to $`f`$ such that the $`L^2`$-limit of $`_0^t\varphi _n(\tau )𝑑\xi _\alpha (\tau )`$ exists.
### 5.8 Proof of Statement 4.2.
### 5.9 Lemmas
We begin with three lemmas. The first is Slepian’s lemma \[Sl, Kl\].
###### Lemma 5.6
(Slepian) Let $`\mathrm{\Gamma }`$ be a countable set, and let $`X(t)`$, $`Y(t)`$ be two real Gaussian processes indexed by $`t\mathrm{\Gamma }`$. Suppose $`EX^2(t)=EY^2(t)`$ and $`EX(s)X(t)EY(s)Y(t)`$ for all $`s,t\mathrm{\Gamma }`$. Then, for all real $`\lambda `$, $`P\left(\underset{t\mathrm{\Gamma }}{\mathrm{max}}X(t)\lambda \right)P\left(\underset{t\mathrm{\Gamma }}{\mathrm{max}}Y(t)\lambda \right).`$
Consequently, if $`X`$ and $`Y`$ have continuous versions, Lemma 5.6 holds when the index set $`\mathrm{\Gamma }`$ is $`[0,1]`$.
The next lemma gives us Markov property.
###### Lemma 5.7
Let $`\alpha `$, $`f`$ be as in Theorem 4.1, and let $`0<\beta <\alpha `$. Let $`Y(t)`$ be a Gaussian process such that $`EY(t)=0`$ and $`EY(s)Y(t)=q_f(0,s)=_0^s_0^sf(u)f(v)\frac{dudv}{|uv|^{1\alpha }}`$ does not depend on $`t`$ whenever $`st`$. If $`fL^{\frac{2}{1+\beta }}([0,1])`$, then $`Y(t)`$ is Markov and has a continuous version.
Proof:
To show that a process is Markov, it is sufficient to show that its non-overlapping increments are independent. For a Gaussian process this amounts to checking that any two non-overlapping increments are uncorrelated: for $`s_1<t_1<s_2<t_2`$,
$`E(Y(t_1)Y(s_1))(Y(t_2)Y(s_2))`$
$`=EY(t_1)Y(t_2)EY(t_1)Y(s_2)EY(s_1)Y(t_2)+EY(s_1)Y(s_2)`$
$`=q_f(0,t_1)q_f(0,t_1)q_f(0,s_1)+q_f(0,s_1)=0.`$
Next we show that $`Y(t)`$ has a continuous version.
$`E(Y(t)Y(s))^2=q_f(0,t)2q_f(0,s)+q_f(0,s)`$
$`=C\alpha (\alpha +1){\displaystyle _s^t}{\displaystyle _s^t}f(u)f(v){\displaystyle \frac{dudv}{|uv|^{1\alpha }}}+2C\alpha (\alpha +1){\displaystyle _0^s}{\displaystyle _s^t}f(u)f(v){\displaystyle \frac{dudv}{|uv|^{1\alpha }}}`$
$`=I_1+I_2.`$
$`I_1`$ can be estimated by the Hardy-Littlewood-Sobolev inequality:
$$I_1=q_f(s,t)C\alpha (\alpha +1)|st|^{\alpha \beta }_s^t_s^t\frac{f(u)f(v)}{|uv|^{1\beta }}𝑑u𝑑v\mathrm{const}|st|^{\alpha \beta }f_{\frac{2}{1+\alpha }}^2.$$
(28)
$`I_2`$ can be estimated by the Hardy-Littlewood-Sobolev and Hölder inequalities:
$`I_2`$ $``$ $`\mathrm{const}f\chi _{[0,s]}_{\frac{2}{1+\alpha }}f\chi _{[s,t]}_{\frac{2}{1+\alpha }}`$
$``$ $`\mathrm{const}f\chi _{[0,s]}_{\frac{2}{1+\alpha }}|st|^{\frac{\alpha \beta }{\alpha +\beta }}\left({\displaystyle _s^t}f(\tau )^{\frac{2}{1+\beta }}𝑑\tau \right)^{\frac{1+\beta }{1+\alpha }}.`$
Choosing an integer $`m`$ such that $`(\alpha \beta )m>1`$ and $`\frac{\alpha \beta }{\alpha +\beta }m>1`$, we can ensure
$`E\left(Y(t)Y(s)\right)^{2m}\mathrm{const}|st|^\gamma ,\mathrm{where}\gamma >1.`$
$`Y`$ has a continuous version by Kolmogorov’s continuity criterion.
Finally, we have a reflection principle:
###### Lemma 5.8
Let $`Y(t)`$ be a centered Gaussian Markov process with continuous paths. Then for all $`\lambda >0`$ and $`T0`$, $`P(\underset{0tT}{\mathrm{max}}Y(t)\lambda )=2P(Y(T)\lambda )`$.
The proof is exactly analogous to the oBm case.
Now we are ready to prove Theorem 4.2.
#### 5.9.1 Proof of Part (1).
Proof:
Choose an integer $`m`$ such that $`(\alpha \beta )m>1`$. Then
$`E\left({\displaystyle _0^t}f(\tau )𝑑\xi _\alpha (\tau ){\displaystyle _0^s}f(\tau )𝑑\xi _\alpha (\tau )\right)^{2m}`$ $``$ $`\mathrm{const}\left(E\left({\displaystyle _0^t}f(\tau )𝑑\xi _\alpha (\tau ){\displaystyle _0^s}f(\tau )𝑑\xi _\alpha (\tau )\right)^2\right)^m`$
$`=`$ $`\left(C\alpha (\alpha +1){\displaystyle _s^t}{\displaystyle _s^t}{\displaystyle \frac{f(u)f(v)}{|uv|^{1\alpha }}}𝑑u𝑑v\right)^m`$
$``$ $`\mathrm{const}|st|^{(\alpha \beta )m},`$
where the first inequality holds because $`_0^tf(\tau )𝑑\xi _\alpha (\tau )`$ is a Gaussian random variable, and the second by (28). By Kolmogorov’s criterion the process $`_0^tf(\tau )𝑑\xi _\alpha (\tau )`$ admits a continuous version.
#### 5.9.2 Proof of Part (2).
Proof:
Let $`X(t)=_0^tf(\tau )𝑑\xi _\alpha (\tau )`$. $`X(t)`$ is a Gaussian process with $`EX(t)=0`$ and $`EX(s)X(t)=_0^s_0^tf(u)f(v)\frac{dudv}{|uv|^{1\alpha }}`$. Define $`Y(t)`$ to be a Gaussian process with $`EY(t)=0`$ and $`EY(s)Y(t)=q_f(0,s)`$ for $`st`$. Clearly $`EX(t)^2=EY(t)^2`$. It can be shown that the process $`Y(t)`$ is well-defined for all $`t`$.
Suppose $`f0`$. Then $`EX(s)X(t)EY(s)Y(t)`$ for $`s,t[0,1]`$. Therefore the processes $`X(t)`$ and $`Y(t)`$ satisfy the assumptions of Slepian’s lemma (Lemma 5.6). By Lemma 5.7, $`Y(t)`$ is a Markov process with continuous paths, and by Lemma 5.8, $`Y`$ obeys the reflection principle:
$`P(\underset{0t1}{\mathrm{max}}Y(t)\lambda )=2P(Y(1)\lambda )={\displaystyle _{\lambda /\sqrt{q_f(0,1)}}^{\mathrm{}}}\sqrt{{\displaystyle \frac{2}{\pi }}}e^{x^2/2}𝑑x.`$
Hence for $`f0`$,
$`P(\underset{0t1}{\mathrm{max}}X(t)\lambda ){\displaystyle _{\lambda /\sqrt{q_f(0,1)}}^{\mathrm{}}}\sqrt{{\displaystyle \frac{2}{\pi }}}e^{x^2/2}𝑑x.`$
Let $`fL^{\frac{2}{1+\beta }}([0,\mathrm{}))`$. Write $`X_\pm =_0^tf_\pm (\tau )𝑑\xi _\alpha (\tau )`$. Define processes $`Y_\pm (t)`$ by replacing $`f`$ by $`f_\pm `$ in the definition of $`Y(t)`$. Since Slepian’s lemma applies also to $`X_{}(t)`$ and $`Y_{}(t)`$, we have
$`P(\underset{0t1}{\mathrm{max}}\pm X_\pm (t)\lambda ){\displaystyle _{\lambda /\sqrt{q_{f_\pm }(0,1)}}^{\mathrm{}}}\sqrt{{\displaystyle \frac{2}{\pi }}}e^{x^2/2}𝑑x.`$
Therefore
$`P(\underset{0t1}{\mathrm{max}}X(t)\lambda )`$ $`=`$ $`P(\underset{0t1}{\mathrm{max}}X_+(t)+\underset{0t1}{\mathrm{max}}(X_{}(t))\lambda )`$
$``$ $`P(\underset{0t1}{\mathrm{max}}X_+(t)\lambda r)+P(\underset{0t1}{\mathrm{max}}(X_{}(t))\lambda (1r))`$
$``$ $`{\displaystyle _{\lambda r/\sqrt{q_{f_+}(0,1)}}^{\mathrm{}}}+{\displaystyle _{\lambda (1r)/\sqrt{q_f_{}(0,1)}}^{\mathrm{}}}\sqrt{{\displaystyle \frac{2}{\pi }}}e^{x^2/2}𝑑x.`$
#### 5.9.3 Proof of Part (3).
Proof:
Inequality (28) shows that $`\sqrt{q_f(s,t)}\mathrm{const}|st|^{\frac{\alpha \beta }{2}}`$, so
$$_1^{\mathrm{}}\underset{|st|<m^{x^2}}{sup}\sqrt{q_f(s,t)}dx<\mathrm{}.$$
This is the condition of applicability of Fernique’s inequality \[F\], of which the claim is a direct consequence.
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# 1 Introduction
## 1 Introduction
The supersymmetric sinh-Gordon (SShG) model - is one of the simplest examples of a $`1+1`$ dimensional integrable quantum field theory with $`N=1`$ supersymmetry. Indeed, the particle spectrum consists of one Boson and one Fermion which have equal mass and which enjoy factorized scattering . As such, SShG is a valuable toy model.
In this article, we consider the boundary SShG model, with boundary conditions that preserve the bulk integrability, but not necessarily the bulk supersymmetry -. In addition to its usefulness as a simple prototype, we expect that this model may also have applications to quantum impurity problems .
An interesting feature of the boundary SShG model is that the $`S`$ matrices which have been conjectured for both bulk and boundary scattering are not diagonal. Our main objective is to perform a thermodynamic Bethe Ansatz (TBA) analysis , - for this model, using these conjectured $`S`$ matrices as inputs. Such analysis can provide checks on the input $`S`$ matrix data, as well as information about the underlying boundary conformal field theory -.
Conventional wisdom suggests that the problem of determining the necessary Bethe Ansatz equations is intractable, due to the fact that the boundary $`S`$ matrix is not diagonal. Nevertheless, we succeed to determine the Bethe Ansatz equations and carry out the TBA analysis for the boundary SShG model. This is the first example of a model with non-diagonal boundary $`S`$ matrix which is exactly solved.
We also obtain an expression for the boundary entropy , for the boundary SShG model. Moreover, we find a rich pattern of boundary roaming trajectories corresponding to $`c<3/2`$ superconformal models , , thereby generalizing previous work on bulk - and boundary , roaming.
The outline of this article is as follows. In Sec. 2 we review the scattering theory of the boundary SShG model, which serves as our input. Here we also show that the strong-weak duality symmetry of the bulk model (see, e.g., ) extends also to the model with boundary. Moreover, we introduce the notations and conventions which are used throughout the paper. In Sec. 3 we formulate the so-called Yang matrix and relate it to a commuting transfer matrix, which is the true starting point of any TBA analysis. For the problem at hand, we require a boundary version of the Yang matrix , , which presents an interesting complication with respect to the more familiar case of periodic boundary conditions. In Sec. 4 we use the open-chain fusion formula to derive an exact inversion identity, using which we obtain the eigenvalues of the transfer matrix in terms of roots of certain Bethe Ansatz equations. That such an inversion identity exists is presumably due to the fact that the bulk $`S`$ matrix satisfies the so-called free Fermion condition , . In Sec. 5 we use these results to derive the TBA equations. Certain remarkable identities lead to very simple formulas, in particular for the boundary entropy. In Sec. 6 we use our result for the boundary entropy to obtain boundary roaming trajectories. Finally, in Sec. 7 we discuss our results and describe some possible generalizations.
## 2 Review of boundary SShG scattering theory
In this Section, we review the bulk and boundary $`S`$ matrices which have been proposed for the boundary supersymmetric sinh-Gordon model. As mentioned in the Introduction, these $`S`$ matrices will be used as inputs in the calculations that follow. We also show that the strong-weak duality symmetry of the bulk model extends also to the model with boundary.
### 2.1 Bulk
In order to understand the SShG scattering theory, it is essential to first consider a related model, namely, the supersymmetric sine-Gordon (SSG) model, whose Euclidean-space Lagrangian density is given by
$`_{bulk}^{SSG}={\displaystyle \frac{1}{2}}\left(_z\varphi _{\overline{z}}\varphi +\overline{\psi }_z\overline{\psi }+\psi _{\overline{z}}\psi \right)+iM\overline{\psi }\psi \mathrm{cos}\beta \varphi +{\displaystyle \frac{M^2}{2\beta ^2}}\mathrm{sin}^2\beta \varphi ,`$ (2.1)
where $`\varphi (z,\overline{z})`$ is a real scalar field, $`\psi (z,\overline{z})`$ and $`\overline{\psi }(z,\overline{z})`$ are the components of a Majorana spinor field, and $`\beta `$ is the dimensionless coupling constant. The Lagrangian density for the supersymmetric sinh-Gordon (SShG) model is obtained by analytic continuation to imaginary coupling, i.e. setting $`\beta =i\widehat{\beta }`$ with $`\widehat{\beta }`$ real:
$`_{bulk}^{SShG}={\displaystyle \frac{1}{2}}\left(_z\varphi _{\overline{z}}\varphi +\overline{\psi }_z\overline{\psi }+\psi _{\overline{z}}\psi \right)+iM\overline{\psi }\psi \mathrm{cosh}\widehat{\beta }\varphi +{\displaystyle \frac{M^2}{2\widehat{\beta }^2}}\mathrm{sinh}^2\widehat{\beta }\varphi .`$ (2.2)
Both of these models have $`N=1`$ supersymmetry (without topological charge) , and are integrable , . <sup>1</sup><sup>1</sup>1Of course, a similar relation exists between the usual (non-supersymmetric) sine-Gordon (SG) and sinh-Gordon (ShG) models. It is more straightforward to infer the scattering theory for the trigonometric (SG, SSG) models than for the hyperbolic (ShG, SShG) models, because the former have kinks with topological charge, whose non-diagonal $`S`$ matrices must satisfy highly restrictive constraints ,,. The $`S`$ matrices for the hyperbolic models are inferred by analytic continuation of the corresponding breather $`S`$ matrices, as is explained in more detail below.
Observe that SSG has a periodic potential, which admits classical soliton solutions that interpolate between neighboring minima. Correspondingly, it has been proposed - that the SSG quantum spectrum consists of supersymmetric multiplets of kinks of mass $`m`$ and breathers (bound states of kinks) of mass $`m_n=2m\mathrm{sin}(n\alpha \pi )`$, $`n=1,2,\mathrm{},[\frac{1}{2\alpha }]`$, where
$`\alpha ={\displaystyle \frac{\frac{\beta ^2}{4\pi }}{1\frac{\beta ^2}{4\pi }}},`$ (2.3)
and $`[x]`$ denotes integer part of $`x`$. Hence, breathers can be present only if $`0<\alpha <\frac{1}{2}`$. The lightest ($`n=1`$) breathers are identified as the elementary particles (Boson, Fermion) corresponding to the fields in the Lagrangian density (2.1).
Upon making the analytic continuation to SShG (which is the model of primary interest), we see that the potential is no longer periodic, and hence, there are no longer any classical soliton solutions. Thus, the SShG quantum spectrum does not contain kinks; it consists only of the elementary particles of some mass $`m`$ corresponding to the fields in the Lagrangian density (2.2), i.e. corresponding to the $`n=1`$ SSG breather. Setting $`\beta =i\widehat{\beta }`$ in Eq. (2.3), we obtain
$`\alpha ={\displaystyle \frac{\frac{\widehat{\beta }^2}{4\pi }}{1+\frac{\widehat{\beta }^2}{4\pi }}}B,`$ (2.4)
where we have introduced the SShG parameter $`B`$.
Since the SShG spectrum corresponds to the $`n=1`$ SSG breather, we infer that the SShG $`S`$ matrix $`S(\theta )`$ for two particles of rapidities $`\theta _1,\theta _2`$ (and corresponding energy $`E_i=m\mathrm{cosh}\theta _i`$ and momentum $`P_i=m\mathrm{sinh}\theta _i`$, $`i=1,2`$) is given by the analytic continuation of the $`n=1`$ SSG breather $`S`$ matrix,
$`S(\theta )=S_{ShG}(\theta )S_{SUSY}(\theta ),`$ (2.5)
where $`\theta =\theta _1\theta _2`$. The scalar factor $`S_{ShG}(\theta )`$ is given by
$`S_{ShG}(\theta )={\displaystyle \frac{\mathrm{sinh}\theta i\mathrm{sin}(2B\pi )}{\mathrm{sinh}\theta +i\mathrm{sin}(2B\pi )}}.`$ (2.6)
This is the $`S`$ matrix of the usual (non-supersymmetric) sinh-Gordon model ,, which is the analytic continuation of the $`n=1`$ SG breather $`S`$ matrix , but with a different dependence on the coupling constant. It satisfies
$`S_{ShG}(\theta )S_{ShG}(\theta )=1,S_{ShG}(\theta )=S_{ShG}(i\pi \theta ).`$ (2.7)
The factor $`S_{SUSY}(\theta )`$ is given by <sup>2</sup><sup>2</sup>2The matrix $`S_{SUSY}(\theta )`$ for SSG was first obtained in terms of an unknown parameter $`\mathrm{\Delta }`$ by solving the constraints coming from supersymmetry and factorization. The identification of $`\mathrm{\Delta }`$ in terms of $`\alpha `$ was made in .
$`S_{SUSY}(\theta )=Y(\theta )R(\theta ),`$ (2.8)
where $`R(\theta )`$ is a $`4\times 4`$ matrix acting on the tensor product space $`VV`$, where $`V`$ is the 2-dimensional vector space of 1-particle states. We choose $`\{|b(\theta ),|f(\theta )\}`$ to be the basis of $`V`$ (corresponding to a Boson, Fermion with rapidity $`\theta `$, respectively); and hence, the basis of $`VV`$ is given by $`\{|b_1,b_2,|b_1,f_2,|f_1,b_2,|f_1,f_2\}`$, where $`|b_1,b_2|b(\theta _1),b(\theta _2)`$ , etc. In this basis, $`R(\theta )`$ is given by
$`R(\theta )=\left(\begin{array}{cccc}a_+(\theta )& 0& 0& d(\theta )\\ 0& b& c(\theta )& 0\\ 0& c(\theta )& b& 0\\ d(\theta )& 0& 0& a_{}(\theta )\end{array}\right),`$ (2.13)
with
$`a_\pm (\theta )=\pm 1{\displaystyle \frac{2i\mathrm{sin}B\pi }{\mathrm{sinh}\theta }},b=1,c={\displaystyle \frac{i\mathrm{sin}B\pi }{\mathrm{sinh}\frac{\theta }{2}}},d={\displaystyle \frac{\mathrm{sin}B\pi }{\mathrm{cosh}\frac{\theta }{2}}}.`$ (2.14)
It is important to note that the matrix elements of $`R(\theta )`$ satisfy the “free Fermion” condition ,
$`a_+a_{}+b^2=c^2+d^2.`$ (2.15)
The matrix $`R(\theta )`$ is a solution of the Yang-Baxter equation <sup>3</sup><sup>3</sup>3We use the very useful convention (which is standard in the spin-chain literature \- , but unfortunately not in the field-theory literature), whereby $`R_{ij}(\theta )`$ acts nontrivially on the $`i^{th}`$ and $`j^{th}`$ vector spaces. For instance, in the Yang-Baxter equation, the $`R`$ matrices act on $`V^3`$, and therefore $`R_{12}(\theta )=R(\theta )𝕀,R_{23}(\theta )=𝕀R(\theta )`$, etc., where $`𝕀`$ is the unit matrix.
$`R_{12}(\theta _1\theta _2)R_{13}(\theta _1\theta _3)R_{23}(\theta _2\theta _3)=R_{23}(\theta _2\theta _3)R_{13}(\theta _1\theta _3)R_{12}(\theta _1\theta _2).`$ (2.16)
This matrix is both $`𝖯`$ and $`𝖳`$ invariant,
$`𝒫_{12}R_{12}(\theta )𝒫_{12}=R_{12}(\theta ),R_{12}(\theta )^{t_1t_2}=R_{12}(\theta ),`$ (2.17)
where $`t_i`$ denotes transposition in the $`i^{th}`$ space, and $`𝒫`$ is the permutation matrix
$`𝒫=\left(\begin{array}{cccc}1& 0& 0& 0\\ 0& 0& 1& 0\\ 0& 1& 0& 0\\ 0& 0& 0& 1\end{array}\right).`$ (2.22)
Moreover, $`Y(\theta )`$ is a scalar factor given by
$`Y(\theta )={\displaystyle \frac{\mathrm{sinh}\frac{\theta }{2}}{\mathrm{sinh}\frac{\theta }{2}i\mathrm{sin}B\pi }}\mathrm{exp}\left({\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{dt}{t}}{\displaystyle \frac{\mathrm{sinh}(it\theta /\pi )\mathrm{sinh}(t(1+B))\mathrm{sinh}(tB)}{\mathrm{cosh}t\mathrm{cosh}^2\frac{t}{2}}}\right),`$ (2.23)
which is a solution of the unitarity and crossing constraints
$`Y(\theta )Y(\theta )={\displaystyle \frac{\mathrm{sinh}^2\frac{\theta }{2}}{\mathrm{sinh}^2\frac{\theta }{2}+\mathrm{sin}^2B\pi }},Y(\theta )=Y(i\pi \theta ).`$ (2.24)
Let us denote the total scalar factor by $`Z(\theta )`$
$`Z(\theta )=S_{ShG}(\theta )Y(\theta ).`$ (2.25)
One can show that $`Z(\theta )`$ has the integral representation
$`Z(\theta )={\displaystyle \frac{\mathrm{sinh}\frac{\theta }{2}}{\mathrm{sinh}\frac{\theta }{2}+i\mathrm{sin}B\pi }}\mathrm{exp}\left({\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{dt}{t}}{\displaystyle \frac{\mathrm{sinh}(it\theta /\pi )\mathrm{sinh}(t(1B))\mathrm{sinh}(tB)}{\mathrm{cosh}t\mathrm{cosh}^2\frac{t}{2}}}\right),`$ (2.26)
which is the same as the expression in Eq. (2.23), except with $`BB`$. It has no poles <sup>4</sup><sup>4</sup>4Although $`Y(\theta )`$ has a pole at $`\theta =i2B\pi `$, it is canceled by a corresponding zero of $`S_{ShG}(\theta )`$. in the physical strip ($`0<Im\theta <\pi `$), provided $`B`$ lies in the range
$`0<B<1,`$ (2.27)
which corresponds to $`0<\widehat{\beta }^2<\mathrm{}`$.
In short, the proposed SShG bulk $`S`$ matrix $`S(\theta )`$ is given by
$`S(\theta )=Z(\theta )R(\theta ),`$ (2.28)
where the scalar factor $`Z(\theta )`$ is given by Eq. (2.26) and the matrix $`R(\theta )`$ is given by Eqs. (2.13), (2.14).
It is known (see, e.g., ) that the SShG bulk $`S`$ matrix is invariant under the strong-weak duality transformation $`\widehat{\beta }4\pi /\widehat{\beta }`$, which implies
$`B1B.`$ (2.29)
Indeed, this invariance can be checked by inspection of the matrix elements (2.14) of $`R(\theta )`$ and the expression (2.26) for $`Z(\theta )`$. (The factors $`S_{ShG}(\theta )`$ and $`Y(\theta )`$ are not separately invariant.) Note that this transformation maps the range (2.27) into itself.
### 2.2 Boundary
We turn now to boundary conditions and boundary scattering, following the framework developed by Ghoshal and Zamolodchikov . An investigation of which boundary terms can be added to the bulk SShG model (2.2) without spoiling (classical) integrability has led to the following results : the boundary Lagrangian
$`L_{boundary}^{SShG}=\mathrm{\Lambda }\mathrm{cosh}\widehat{\beta }(\varphi \varphi _0)+𝖬\overline{\psi }\psi +ϵ\psi +\overline{ϵ}\overline{\psi },𝖬\pm 1,`$ (2.30)
breaks supersymmetry but preserves integrability; and
$`L_{boundary}^{SShG}={\displaystyle \frac{M}{\widehat{\beta }^2}}\mathrm{cosh}\widehat{\beta }\varphi \pm \overline{\psi }\psi `$ (2.31)
preserves both supersymmetry and integrability. Notice that the boundary terms (2.30) involve a total of 5 boundary parameters $`\mathrm{\Lambda },\varphi _0,𝖬,ϵ,\overline{ϵ}`$. If $`ϵ,\overline{ϵ}`$ are nonzero, then Fermion number is not conserved.
The proposed boundary $`S`$ matrix $`𝖲(\theta )`$ for a particle of rapidity $`\theta `$ is given by (compare with Eq. (2.5)) <sup>5</sup><sup>5</sup>5We make an effort to distinguish boundary quantities from the corresponding bulk quantities by using sans serif letters to denote the former, and Roman letters to denote the latter.,<sup>6</sup><sup>6</sup>6The boundary $`S`$ matrix is obtained in terms of a set of boundary parameters $`(\eta ,\vartheta ,\phi ,\epsilon )`$ by solving the boundary Yang-Baxter equation. The relation of these parameters to those in $`L_{boundary}^{SShG}`$ (2.30) is not known.
$`𝖲(\theta )=𝖲_{ShG}(\theta ;\eta ,\vartheta )𝖲_{SUSY}^{(\epsilon )}(\theta ;\phi ).`$ (2.32)
The scalar factor $`𝖲_{ShG}(\theta ;\eta ,\vartheta )`$, which depends on two boundary parameters $`\eta ,\vartheta `$, is given by
$`𝖲_{ShG}(\theta ;\eta ,\vartheta )=𝖷_0(\theta )𝖷_1(\theta ;{\displaystyle \frac{4\eta B}{\pi }})𝖷_1(\theta ;{\displaystyle \frac{4i\vartheta B}{\pi }}),`$ (2.33)
where
$`𝖷_0(\theta )=(1)(1+2B)(22B),𝖷_1(\theta ;F)={\displaystyle \frac{1}{(1F)(1+F)}},`$ (2.34)
with
$`(x){\displaystyle \frac{\mathrm{sinh}(\frac{\theta }{2}+\frac{i\pi x}{4})}{\mathrm{sinh}(\frac{\theta }{2}\frac{i\pi x}{4})}}.`$ (2.35)
This is the boundary $`S`$ matrix of the usual (non-supersymmetric) boundary sinh-Gordon model, which is the analytic continuation of the $`n=1`$ boundary sine-Gordon breather $`S`$ matrix . It satisfies
$`𝖲_{ShG}(\theta ;\eta ,\vartheta )𝖲_{ShG}(\theta ;\eta ,\vartheta )=1,𝖲_{ShG}({\displaystyle \frac{i\pi }{2}}+\theta ;\eta ,\vartheta )S_{ShG}(2\theta )=𝖲_{ShG}({\displaystyle \frac{i\pi }{2}}\theta ;\eta ,\vartheta ).`$
(2.36)
The factor $`𝖲_{SUSY}^{(\epsilon )}(\theta ;\phi )`$ is given by
$`𝖲_{SUSY}^{(\epsilon )}(\theta ;\phi )=𝖸^{(\epsilon )}(\theta ;\phi )𝖱^{(\epsilon )}(\theta ;\phi ),`$ (2.37)
where $`\epsilon `$ is a discrete parameter which can be either $`+1`$ or $`1`$, and $`\phi `$ is a continuous boundary parameter. $`𝖱^{(\epsilon )}(\theta ;\phi )`$ is a $`2\times 2`$ matrix acting on the vector space $`V`$ of 1-particle states, which is given by
$`𝖱^{(\epsilon )}(\theta ;\phi )=\left(\begin{array}{cc}\mathrm{cosh}\frac{\theta }{2}G_+^{(\epsilon )}+i\mathrm{sinh}\frac{\theta }{2}G_{}^{(\epsilon )}& \epsilon i\mathrm{sinh}\theta \\ i\mathrm{sinh}\theta & \mathrm{cosh}\frac{\theta }{2}G_+^{(\epsilon )}i\mathrm{sinh}\frac{\theta }{2}G_{}^{(\epsilon )}\end{array}\right),`$ (2.40)
where
$`G_\epsilon ^{}^{(\epsilon )}=\{\begin{array}{cc}\hfill r\left(\mathrm{cosh}\phi +e^{\epsilon \phi }\frac{\mathrm{sinh}^2\frac{\theta }{2}}{1+\epsilon \mathrm{sin}B\pi }\right)\text{if}& \epsilon ^{}=\epsilon \hfill \\ \hfill r\left(\mathrm{sinh}\phi +\epsilon e^{\epsilon \phi }\frac{\mathrm{sinh}^2\frac{\theta }{2}}{1+\epsilon \mathrm{sin}B\pi }\right)\text{if}& \epsilon ^{}=\epsilon \hfill \end{array},`$ (2.43)
and
$`r=\left({\displaystyle \frac{2(\epsilon +\mathrm{sin}B\pi )}{\mathrm{sin}B\pi }}\right)^{\frac{1}{2}}.`$ (2.44)
The matrix $`𝖱^{(\epsilon )}(\theta ;\phi )`$ is a solution of the boundary Yang-Baxter equation
$`R_{12}(\theta _1\theta _2)𝖱_1^{(\epsilon )}(\theta _1;\phi )R_{12}(\theta _1+\theta _2)𝖱_2^{(\epsilon )}(\theta _2;\phi )`$ (2.45)
$`=𝖱_2^{(\epsilon )}(\theta _2;\phi )R_{12}(\theta _1+\theta _2)𝖱_1^{(\epsilon )}(\theta _1;\phi )R_{12}(\theta _1\theta _2).`$
We remark that for $`\phi \pm \mathrm{}`$, the matrix $`𝖱^{(\epsilon )}(\theta ;\phi )`$ becomes diagonal and commutes with linear combinations $`Q\pm \overline{Q}`$ of the supersymmetry charges ,.
In order to determine the scalar factor $`𝖸^{(\epsilon )}(\theta ;\phi )`$, we recall that the full boundary $`S`$ matrix must satisfy boundary unitarity $`𝖲(\theta )𝖲(\theta )=𝕀`$ and boundary cross-unitarity , which can be written in matrix form as
$`tr_0𝖲_0({\displaystyle \frac{i\pi }{2}}+\theta )^{t_0}𝒫_{01}S_{01}(2\theta )^{t_1}=𝖲_1({\displaystyle \frac{i\pi }{2}}\theta ).`$ (2.46)
We observe that the matrix $`𝖱^{(\epsilon )}(\theta ;\phi )`$ satisfies
$`𝖱^{(\epsilon )}(\theta ;\phi )𝖱^{(\epsilon )}(\theta ;\phi )=h(\theta )𝕀,`$ (2.47)
where
$`h(\theta )=\left(c_0+c_1\mathrm{sinh}^2{\displaystyle \frac{\theta }{2}}+c_2\mathrm{sinh}^4{\displaystyle \frac{\theta }{2}}\right)\mathrm{cosh}\theta ,`$ (2.48)
and
$`c_0=\{\begin{array}{cc}r^2\mathrm{cosh}^2\phi & \text{ if }\epsilon =+1\\ r^2\mathrm{sinh}^2\phi & \text{ if }\epsilon =1\end{array},c_1={\displaystyle \frac{r^2e^{\epsilon 2\phi }}{1+\epsilon \mathrm{sin}B\pi }}+2\epsilon ,c_2={\displaystyle \frac{r^2e^{\epsilon 2\phi }}{(1+\epsilon \mathrm{sin}B\pi )^2}}.`$ (2.51)
Also,
$`tr_0𝖱_0^{(\epsilon )}({\displaystyle \frac{i\pi }{2}}+\theta ;\phi )^{t_0}𝒫_{01}R_{01}(2\theta )^{t_1}=g(\theta )𝖱_1^{(\epsilon )}({\displaystyle \frac{i\pi }{2}}\theta ;\phi ),`$ (2.52)
where
$`g(\theta )={\displaystyle \frac{\epsilon \mathrm{sinh}\theta i\mathrm{sin}\pi B}{\mathrm{sinh}\theta }}.`$ (2.53)
Setting
$`𝖸^{(\epsilon )}(\theta ;\phi )=𝖸_0^{(\epsilon )}(\theta )𝖸_1^{(\epsilon )}(\theta ;\phi ),`$ (2.54)
it follows that $`𝖸_0^{(\epsilon )}(\theta )`$ and $`𝖸_1^{(\epsilon )}(\theta ;\phi )`$ must satisfy
$`𝖸_0^{(\epsilon )}(\theta )𝖸_0^{(\epsilon )}(\theta )\mathrm{cosh}\theta =1,𝖸_0^{(\epsilon )}({\displaystyle \frac{i\pi }{2}}+\theta )Y(2\theta )g(\theta )=𝖸_0^{(\epsilon )}({\displaystyle \frac{i\pi }{2}}\theta ),`$ (2.55)
and
$`𝖸_1^{(\epsilon )}(\theta ;\phi )𝖸_1^{(\epsilon )}(\theta ;\phi )\left(c_0+c_1\mathrm{sinh}^2{\displaystyle \frac{\theta }{2}}+c_2\mathrm{sinh}^4{\displaystyle \frac{\theta }{2}}\right)=1,𝖸_1^{(\epsilon )}({\displaystyle \frac{i\pi }{2}}+\theta ;\phi )=𝖸_1^{(\epsilon )}({\displaystyle \frac{i\pi }{2}}\theta ;\phi ),`$
(2.56)
respectively.
For simplicity, we shall henceforth restrict our attention to the case $`\epsilon =+1`$, and so we shall drop the superscript $`(\epsilon )`$. We propose the following integral representations for $`𝖸_0(\theta )`$ and $`𝖸_1(\theta ;\phi )`$:
$`𝖸_0(\theta )`$ $`=`$ $`{\displaystyle \frac{i}{\sqrt{2}\mathrm{sinh}(\frac{\theta }{2}+\frac{i\pi }{4})}}\mathrm{exp}\left({\displaystyle \frac{1}{2}}{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{dt}{t}}{\displaystyle \frac{\mathrm{sinh}(2it\theta /\pi )\mathrm{sinh}(t(1+B))\mathrm{sinh}(tB)}{\mathrm{cosh}^2t\mathrm{cosh}^2\frac{t}{2}}}\right),`$
$`𝖸_1(\theta ;\phi )`$ $`=`$ $`{\displaystyle \frac{1}{r\mathrm{cosh}\phi }}\mathrm{exp}\left(4{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{dt}{t}}{\displaystyle \frac{\mathrm{cosh}(t\zeta /\pi )\mathrm{cosh}(\frac{t}{2}(12B))\mathrm{sinh}(\frac{t}{2}(1+\frac{i\theta }{\pi }))\mathrm{sinh}(\frac{it\theta }{2\pi })}{\mathrm{sinh}t\mathrm{cosh}\frac{t}{2}}}\right),`$
where $`\zeta `$ is a function of the boundary parameter $`\phi `$ defined by
$`\zeta =\mathrm{cos}^1\left(1+e^{2\phi }(1+\mathrm{sin}B\pi )\right).`$ (2.58)
In order to streamline the notation, let us denote the set of boundary parameters $`\{\eta ,\vartheta ,\phi \}`$ by $`\xi `$, and denote the total scalar factor by $`𝖹(\theta ;\xi )`$
$`𝖹(\theta ;\xi )`$ $`=`$ $`𝖲_{ShG}(\theta ;\eta ,\vartheta )𝖸(\theta ;\phi )`$ (2.59)
$`=`$ $`𝖹_0(\theta )𝖹_1(\theta ;\xi ).`$
The proposed SShG boundary $`S`$ matrix is then given by
$`𝖲(\theta ;\xi )=𝖹(\theta ;\xi )𝖱(\theta ;\phi ).`$ (2.60)
We now observe that the boundary $`S`$ matrix is also invariant under the strong-weak duality transformation (2.29). Indeed, it is evident that the matrix $`𝖱(\theta ;\phi )`$ (2.40) has this invariance, if we assume that the parameter $`\phi `$ remains invariant under this transformation. Let us now consider the scalar factor. The part of the scalar factor that does not depend on boundary parameters can be written in the form
$`𝖹_0(\theta )`$ $`=`$ $`𝖷_0(\theta )𝖸_0(\theta )`$
$`=`$ $`{\displaystyle \frac{i}{\sqrt{2}\mathrm{sinh}(\frac{\theta }{2}\frac{i\pi }{4})}}`$
$`\times `$ $`\mathrm{exp}\left({\displaystyle \frac{1}{4}}{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{dt}{t}}{\displaystyle \frac{\mathrm{sinh}(2it\theta /\pi )}{\mathrm{cosh}^2t\mathrm{cosh}^2\frac{t}{2}}}\left[\mathrm{cosh}(t(12B))(1+2\mathrm{cosh}t)+\mathrm{cosh}t\right]\right),`$ (2.62)
in which the duality invariance is manifest. (The factors $`𝖷_0(\theta )`$ and $`𝖸_0(\theta )`$ are not separately invariant.) Finally, the part of the scalar factor which does depend on boundary parameters,
$`𝖹_1(\theta ;\xi )=𝖷_1(\theta ;{\displaystyle \frac{4\eta B}{\pi }})𝖷_1(\theta ;{\displaystyle \frac{4i\vartheta B}{\pi }})𝖸_1(\theta ;\phi ),`$ (2.63)
is also invariant under duality, since each of its factors are separately invariant (provided the boundary parameters $`\eta B`$ and $`\vartheta B`$ are assumed to remain invariant).
## 3 Yang equations
Having specified the bulk and boundary $`S`$ matrices, we are ready to start the TBA program. The first step is to formulate the Yang matrix and relate it to a commuting transfer matrix. Since this is not obvious for the case of boundaries, we begin by reviewing the more familiar case of periodic boundary conditions.
### 3.1 Closed
We consider $`N`$ particles of mass $`m`$ with real rapidities $`\theta _1,\mathrm{},\theta _N`$ and two-particle $`S`$ matrix $`S(\theta )`$ (2.28) in a periodic box of length $`L>>\frac{1}{m}`$. The Yang equation , for particle 1 (which has momentum $`P_1=m\mathrm{sinh}\theta _1`$) is given by
$`\left(e^{iLm\mathrm{sinh}\theta _1}Y_{(1)}𝕀\right)|\theta _1,\mathrm{}\theta _N=0,`$ (3.1)
where $`Y_{(1)}`$ is the “Yang matrix”<sup>7</sup><sup>7</sup>7We remind the reader that we are using the convention explained in Footnote 3.
$`Y_{(1)}=S_{1N}(\theta _1\theta _N)S_{1,N1}(\theta _1\theta _{N1})\mathrm{}S_{12}(\theta _1\theta _2),`$ (3.2)
which acts on $`V^N`$. There are similar equations, and corresponding matrices $`Y_{(i)}`$, for the other particles $`i=2,3,\mathrm{},N`$.
The objective is to diagonalize $`Y_{(i)}`$. The key to this problem is to relate $`Y_{(i)}`$ to an inhomogeneous closed-chain transfer matrix, for which there are well-developed diagonalization techniques. (For reviews, see e.g. \- .) Indeed, consider the transfer matrix (see Figure 1)
$`\tau _{closed}(\theta |\theta _1,\mathrm{},\theta _N)=tr_0\{S_{0N}(\theta \theta _N)\mathrm{}S_{02}(\theta \theta _2)S_{01}(\theta \theta _1)\},`$ (3.3)
with inhomogeneities $`\theta _1,\mathrm{},\theta _N`$. Notice that we have introduced an additional (“auxiliary”) 2-dimensional vector space denoted by $`0`$. The product of $`S`$ matrices inside the trace (the so-called monodromy matrix) acts on $`V^{(N+1)}`$; but after performing the trace over the auxiliary space, one is left with an operator which acts on the (“quantum”) space $`V^N`$. Because $`S(\theta )`$ satisfies the Yang-Baxter equation, the transfer matrix commutes for different values of $`\theta `$
$`[\tau _{closed}(\theta |\theta _1,\mathrm{},\theta _N),\tau _{closed}(\theta ^{}|\theta _1,\mathrm{},\theta _N)]=0.`$ (3.4)
Let us now evaluate this transfer matrix at $`\theta =\theta _1`$. Using the fact that $`S(0)=𝒫`$ (the permutation matrix (2.22)) and $`𝒫^2=𝕀`$, we see that
$`\tau _{closed}(\theta _1|\theta _1,\mathrm{},\theta _N)=tr_0\{(𝒫_{01}𝒫_{01})S_{0N}(\theta _1\theta _N)\mathrm{}(𝒫_{01}𝒫_{01})S_{02}(\theta _1\theta _2)𝒫_{01}\}.`$ (3.5)
Finally, using $`𝒫_{01}S_{0i}𝒫_{01}=S_{1i}`$ and $`tr_0𝒫_{01}=𝕀_1`$, we conclude that $`\tau _{closed}(\theta _1|\theta _1,\mathrm{},\theta _N)=Y_{(1)}`$. In general, we have
$`Y_{(i)}=\tau _{closed}(\theta _i|\theta _1,\mathrm{},\theta _N),i=1,\mathrm{},N.`$ (3.6)
This is the sought-after relation. In order to diagonalize the Yang matrices $`Y_{(i)}`$, it suffices to diagonalize the commuting closed-chain transfer matrix $`\tau _{closed}(\theta |\theta _1,\mathrm{},\theta _N)`$. That calculation, as well as the corresponding bulk TBA analysis, is described in .
### 3.2 Open
We now turn to the case with boundaries, which is our primary interest in this paper. We therefore consider $`N`$ particles of mass $`m`$ with real rapidities $`\theta _1,\mathrm{},\theta _N`$ in an interval of length $`L>>\frac{1}{m}`$, with bulk $`S`$ matrix $`S(\theta )`$ (2.28) and boundary $`S`$ matrix $`𝖲(\theta ;\xi )`$ (2.60). The Yang equation for particle 1 is given by ,
$`\left(e^{2iLm\mathrm{sinh}\theta _1}Y_{(1)}𝕀\right)|\theta _1,\mathrm{}\theta _N=0,`$ (3.7)
where the Yang matrix $`Y_{(1)}`$ is now given by
$`Y_{(1)}=𝖲_1(\theta _1;\xi _{})S_{21}(\theta _1+\theta _2)\mathrm{}S_{N1}(\theta _1+\theta _N)𝖲_1(\theta _1;\xi _+)S_{1N}(\theta _1\theta _N)\mathrm{}S_{12}(\theta _1\theta _2),`$ (3.8)
where the subscripts $`\pm `$ denote the left and right boundaries. (There are similar matrices $`Y_{(i)}`$ for the other particles.) In analogy with the case of periodic boundary conditions, the key to diagonalizing the Yang matrix is to relate it to an inhomogeneous open-chain transfer matrix (see Figure 2)
$`\tau (\theta |\theta _1,\mathrm{},\theta _N)`$ $`=`$ $`tr_0\{𝖲_0(\theta +i\pi ;\xi _+)^{t_0}S_{0N}(\theta \theta _N)\mathrm{}S_{01}(\theta \theta _1)`$ (3.9)
$`\times `$ $`𝖲_0(\theta ;\xi _{})S_{01}(\theta +\theta _1)\mathrm{}S_{0N}(\theta +\theta _N)\},`$
which commutes for different values of $`\theta `$
$`[\tau (\theta |\theta _1,\mathrm{},\theta _N),\tau (\theta ^{}|\theta _1,\mathrm{},\theta _N)]=0.`$ (3.10)
Using the boundary cross-unitarity relation (2.46) as well as the Yang-Baxter equation (2.16), (2.45), one can show that
$`Y_{(i)}=\tau (\theta _i|\theta _1,\mathrm{},\theta _N),i=1,\mathrm{},N.`$ (3.11)
A proof for the case $`N=2`$ is presented in Appendix A. <sup>8</sup><sup>8</sup>8We therefore fill a gap left open in , where it was first observed that the open-chain Yang matrix is related to the Sklyanin transfer matrix; but neither the precise form of the relation nor its proof was given. Hence, in order to diagonalize the Yang matrices $`Y_{(i)}`$, it suffices to diagonalize the commuting open-chain transfer matrix $`\tau (\theta |\theta _1,\mathrm{},\theta _N)`$. It is to this task that we devote the following section.
## 4 Inversion identity and transfer-matrix eigenvalues
In this Section, we consider the problem of determining the eigenvalues of the inhomogeneous open-chain transfer matrix (3.9). Our approach will be to first derive an exact so-called inversion identity. This approach has been used in the past to diagonalize simple (e.g., Ising) closed-chain transfer matrices ,,.
### 4.1 Inversion identity
Instead of working with the “dressed” transfer matrix (3.9), it is more convenient (see Footnote 9 below) to strip away the scalar factors from the bulk and boundary $`S`$ matrices, and to work instead with the “bare” transfer matrix
$`𝔱(\theta |\theta _1,\mathrm{},\theta _N)`$ $`=`$ $`tr_0\{𝖱_0(\theta +i\pi ;\phi _+)^{t_0}R_{0N}(\theta \theta _N)\mathrm{}R_{01}(\theta \theta _1)`$ (4.1)
$`\times `$ $`𝖱_0(\theta ;\phi _{})R_{01}(\theta +\theta _1)\mathrm{}R_{0N}(\theta +\theta _N)\},`$
where $`R(\theta )`$ is given by (2.13) and $`𝖱(\theta ;\phi )`$ is given by (2.40) with $`\epsilon =+1`$.
There are two key points involved in obtaining the inversion identity. The first key point is to observe that the bulk $`S`$ matrix degenerates into a one-dimensional projector for a certain value of $`\theta (=i\pi )`$:
$`S(i\pi ){\displaystyle \frac{1}{2}}\left(\begin{array}{cccc}\hfill 1& \hfill 0& \hfill 0& \hfill 1\\ \hfill 0& \hfill 0& \hfill 0& \hfill 0\\ \hfill 0& \hfill 0& \hfill 0& \hfill 0\\ \hfill 1& \hfill 0& \hfill 0& \hfill 1\end{array}\right).`$ (4.6)
Hence, it is possible to “fuse” ,, in the auxiliary space, and thereby obtain a fusion formula of the form
$`𝔱(\theta |\theta _1,\mathrm{},\theta _N)𝔱(\theta +i\pi |\theta _1,\mathrm{},\theta _N)\stackrel{~}{}𝔱(\theta |\theta _1,\mathrm{},\theta _N)+\mathrm{\Delta },`$ (4.7)
where $`\stackrel{~}{}𝔱(\theta |\theta _1,\mathrm{},\theta _N)`$ is a “fused” open-chain transfer matrix (see Figure 3), and here $`\mathrm{\Delta }`$ represents a product of certain quantum determinants . The fused transfer matrix is constructed from the fused bulk $`S`$ matrix $`\stackrel{~}{R}(\theta )`$ and the fused boundary $`S`$ matrix $`\stackrel{~}{}𝖱(\theta ;\phi )`$ , using the “fused” 3-dimensional (instead of 2-dimensional) auxiliary space.
The second key point is that both $`\stackrel{~}{R}(\theta )`$ and $`\stackrel{~}{}𝖱(\theta ;\phi )`$ can be brought to upper-triangular form by a $`\theta `$-independent similarity transformation. This remarkable fact is presumably due to the fact that $`R(\theta )`$ satisfies the free Fermion condition (2.15) (cf, ). As a result, the fused transfer matrix is proportional to the identity matrix
$`\stackrel{~}{}𝔱(\theta |\theta _1,\mathrm{},\theta _N)𝕀.`$ (4.8)
It follows from the fusion formula that the transfer matrix obeys an exact inversion identity
$`𝔱(\theta |\theta _1,\mathrm{},\theta _N)𝔱(\theta +i\pi |\theta _1,\mathrm{},\theta _N)=f(\theta )𝕀,`$ (4.9)
where $`f(\theta )`$ is a calculable scalar function. We find (see Appendix B for more details)
$`f(\theta )={\displaystyle \frac{16\mathrm{sinh}^2\theta }{\mathrm{sinh}(\theta iB\pi )\mathrm{sinh}(\theta +iB\pi )(1+\mathrm{sin}B\pi )^2\mathrm{sin}^2B\pi }}\{`$ (4.10)
$`(𝔞{\displaystyle \underset{j=1}{\overset{N}{}}}{\displaystyle \frac{\mathrm{cosh}(\frac{1}{2}(\theta \theta _j)iB\pi )}{\mathrm{cosh}(\frac{1}{2}(\theta \theta _j))}}{\displaystyle \frac{\mathrm{cosh}(\frac{1}{2}(\theta +\theta _j)iB\pi )}{\mathrm{cosh}(\frac{1}{2}(\theta +\theta _j))}}`$
$`𝔟{\displaystyle \underset{j=1}{\overset{N}{}}}{\displaystyle \frac{\mathrm{sinh}(\frac{1}{2}(\theta \theta _j)iB\pi )}{\mathrm{sinh}(\frac{1}{2}(\theta \theta _j))}}{\displaystyle \frac{\mathrm{sinh}(\frac{1}{2}(\theta +\theta _j)iB\pi )}{\mathrm{sinh}(\frac{1}{2}(\theta +\theta _j))}})`$
$`\times (𝔠{\displaystyle \underset{j=1}{\overset{N}{}}}{\displaystyle \frac{\mathrm{cosh}(\frac{1}{2}(\theta \theta _j)+iB\pi )}{\mathrm{cosh}(\frac{1}{2}(\theta \theta _j))}}{\displaystyle \frac{\mathrm{cosh}(\frac{1}{2}(\theta +\theta _j)+iB\pi )}{\mathrm{cosh}(\frac{1}{2}(\theta +\theta _j))}}`$
$`𝔡{\displaystyle \underset{j=1}{\overset{N}{}}}{\displaystyle \frac{\mathrm{sinh}(\frac{1}{2}(\theta \theta _j)+iB\pi )}{\mathrm{sinh}(\frac{1}{2}(\theta \theta _j))}}{\displaystyle \frac{\mathrm{sinh}(\frac{1}{2}(\theta +\theta _j)+iB\pi )}{\mathrm{sinh}(\frac{1}{2}(\theta +\theta _j))}})\},`$
where
$`𝔞`$ $`=`$ $`\mathrm{sinh}^2({\displaystyle \frac{1}{2}}(i\pi B+{\displaystyle \frac{i\pi }{2}}\theta ))\left[e^\phi _{}\mathrm{sinh}^2({\displaystyle \frac{1}{2}}(i\pi B+{\displaystyle \frac{i\pi }{2}}))+e^\phi _{}\mathrm{sinh}^2({\displaystyle \frac{1}{2}}(i\pi B+{\displaystyle \frac{i\pi }{2}}+\theta ))\right]\times \left[\phi _{}\phi _+\right]`$
$`𝔟`$ $`=`$ $`\mathrm{sinh}^2({\displaystyle \frac{1}{2}}(i\pi B+{\displaystyle \frac{i\pi }{2}}+\theta ))\left[e^\phi _{}\mathrm{sinh}^2({\displaystyle \frac{1}{2}}(i\pi B+{\displaystyle \frac{i\pi }{2}}))+e^\phi _{}\mathrm{sinh}^2({\displaystyle \frac{1}{2}}(i\pi B+{\displaystyle \frac{i\pi }{2}}\theta ))\right]\times \left[\phi _{}\phi _+\right]`$
$`𝔠`$ $`=`$ $`\mathrm{sinh}^2({\displaystyle \frac{1}{2}}(i\pi B+{\displaystyle \frac{i\pi }{2}}+\theta ))\left[e^\phi _{}\mathrm{sinh}^2({\displaystyle \frac{1}{2}}(i\pi B+{\displaystyle \frac{i\pi }{2}}))+e^\phi _{}\mathrm{sinh}^2({\displaystyle \frac{1}{2}}(i\pi B{\displaystyle \frac{i\pi }{2}}+\theta ))\right]\times \left[\phi _{}\phi _+\right]`$
$`𝔡`$ $`=`$ $`\mathrm{sinh}^2({\displaystyle \frac{1}{2}}(i\pi B{\displaystyle \frac{i\pi }{2}}+\theta ))[e^\phi _{}\mathrm{sinh}^2({\displaystyle \frac{1}{2}}(i\pi B+{\displaystyle \frac{i\pi }{2}})+e^\phi _{}\mathrm{sinh}^2({\displaystyle \frac{1}{2}}(i\pi B+{\displaystyle \frac{i\pi }{2}}+\theta ))]\times [\phi _{}\phi _+].`$
Notice that the function $`f(\theta )`$ is invariant under the duality transformation $`B1B`$. This inversion identity is one of the main results of this paper. We have checked it numerically up to $`N=3`$.
### 4.2 Eigenvalues
We now proceed to determine the eigenvalues of the transfer matrix. First, observe that by virtue of the commutativity property (3.10), the bare transfer matrix $`𝔱(\theta |\theta _1,\mathrm{},\theta _N)`$ has eigenstates $`|\theta _1,\mathrm{}\theta _N`$ which are independent of $`\theta `$,
$`𝔱(\theta |\theta _1,\mathrm{},\theta _N)|\theta _1,\mathrm{}\theta _N=𝔏(\theta |\theta _1,\mathrm{},\theta _N)|\theta _1,\mathrm{}\theta _N,`$ (4.12)
where $`𝔏(\theta |\theta _1,\mathrm{},\theta _N)`$ are the corresponding eigenvalues. Acting on $`|\theta _1,\mathrm{}\theta _N`$ with the inversion identity, we obtain the corresponding identity for the eigenvalues
$`𝔏(\theta |\theta _1,\mathrm{},\theta _N)𝔏(\theta +i\pi |\theta _1,\mathrm{},\theta _N)=f(\theta ).`$ (4.13)
Moreover, one can show that the bare transfer matrix $`𝔱(\theta |\theta _1,\mathrm{},\theta _N)`$ is a periodic function of $`\theta `$ with period $`2\pi i`$ <sup>9</sup><sup>9</sup>9 This is not the case for the dressed transfer matrix $`\tau (\theta |\theta _1,\mathrm{},\theta _N)`$, due to the presence of the scalar factors.
$`𝔱(\theta +2\pi i|\theta _1,\mathrm{},\theta _N)=𝔱(\theta |\theta _1,\mathrm{},\theta _N),`$ (4.14)
whose asymptotic behavior for large $`\theta `$ is given by
$`𝔱(\theta |\theta _1,\mathrm{},\theta _N){\displaystyle \frac{c}{32}}e^{3\theta }𝕀\text{ for }\theta \mathrm{},`$ (4.15)
where
$`c={\displaystyle \frac{4ie^{\phi _{}+\phi _+}}{(1+\mathrm{sin}B\pi )\mathrm{sin}B\pi }}.`$ (4.16)
Correspondingly, the eigenvalues obey
$`𝔏(\theta +2\pi i|\theta _1,\mathrm{},\theta _N)`$ $`=`$ $`𝔏(\theta |\theta _1,\mathrm{},\theta _N),`$
$`𝔏(\theta |\theta _1,\mathrm{},\theta _N)`$ $``$ $`{\displaystyle \frac{c}{32}}e^{3\theta }\text{ for }\theta \mathrm{}.`$ (4.17)
The eigenvalues $`𝔏(\theta |\theta _1,\mathrm{},\theta _N)`$ are uniquely determined by the zeros and poles of $`f(\theta )`$, together with periodicity and asymptotic behavior. Indeed, observe that $`f(\theta )`$ is a product of two factors. Let $`\theta =z_k^+,z_k^{}`$ be zeros of the first, second factors, respectively. Then $`z_k^+`$ obeys
$`{\displaystyle \underset{j=1}{\overset{N}{}}}{\displaystyle \frac{\mathrm{tanh}(\frac{1}{2}(z_k^+\theta _j)iB\pi )}{\mathrm{tanh}(\frac{1}{2}(z_k^+\theta _j))}}{\displaystyle \frac{\mathrm{tanh}(\frac{1}{2}(z_k^++\theta _j)iB\pi )}{\mathrm{tanh}(\frac{1}{2}(z_k^++\theta _j))}}={\displaystyle \frac{\mathrm{sinh}^2(\frac{1}{2}(i\pi B+\frac{i\pi }{2}z_k^+))}{\mathrm{sinh}^2(\frac{1}{2}(i\pi B+\frac{i\pi }{2}+z_k^+))}}`$ (4.18)
$`\times `$ $`\left[{\displaystyle \frac{e^\phi _{}\mathrm{sinh}^2(\frac{1}{2}(i\pi B+\frac{i\pi }{2}))+e^\phi _{}\mathrm{sinh}^2(\frac{1}{2}(i\pi B+\frac{i\pi }{2}+z_k^+))}{e^\phi _{}\mathrm{sinh}^2(\frac{1}{2}(i\pi B+\frac{i\pi }{2}))+e^\phi _{}\mathrm{sinh}^2(\frac{1}{2}(i\pi B+\frac{i\pi }{2}z_k^+))}}\right]\times \left[\phi _{}\phi _+\right],`$
and $`z_k^{}`$ obeys
$`{\displaystyle \underset{j=1}{\overset{N}{}}}{\displaystyle \frac{\mathrm{tanh}(\frac{1}{2}(z_k^{}\theta _j)+iB\pi )}{\mathrm{tanh}(\frac{1}{2}(z_k^{}\theta _j))}}{\displaystyle \frac{\mathrm{tanh}(\frac{1}{2}(z_k^{}+\theta _j)+iB\pi )}{\mathrm{tanh}(\frac{1}{2}(z_k^{}+\theta _j))}}={\displaystyle \frac{\mathrm{sinh}^2(\frac{1}{2}(i\pi B+\frac{i\pi }{2}+z_k^{}))}{\mathrm{sinh}^2(\frac{1}{2}(i\pi B\frac{i\pi }{2}+z_k^{}))}}`$ (4.19)
$`\times `$ $`\left[{\displaystyle \frac{e^\phi _{}\mathrm{sinh}^2(\frac{1}{2}(i\pi B+\frac{i\pi }{2}))+e^\phi _{}\mathrm{sinh}^2(\frac{1}{2}(i\pi B\frac{i\pi }{2}+z_k^{}))}{e^\phi _{}\mathrm{sinh}^2(\frac{1}{2}(i\pi B+\frac{i\pi }{2}))+e^\phi _{}\mathrm{sinh}^2(\frac{1}{2}(i\pi B+\frac{i\pi }{2}+z_k^{}))}}\right]\times \left[\phi _{}\phi _+\right].`$
These are our “magnonic” Bethe Ansatz equations.
It follows <sup>10</sup><sup>10</sup>10Indeed, let us denote the right-hand-side of Eq. (4.20) by $`F(\theta )`$. We observe that both $`f(\theta )`$ and $`F(\theta )`$ have the same periodicity (namely, $`i\pi `$, which is half the period of $`𝔏(\theta |\theta _1,\mathrm{},\theta _N)`$), the same zeros and poles in the strip $`\frac{i\pi }{2}<\theta <\frac{i\pi }{2}`$, and the same asymptotic behavior. (The apparent poles of $`f(\theta )`$ at $`\theta =\pm iB\pi `$ are canceled by corresponding zeros.) Hence, the function $`g(\theta )=F(\theta )/f(\theta )`$ is regular everywhere in the complex $`\theta `$ plane, and thus must be constant by Liouville’s theorem. By considering the limit $`\theta \mathrm{}`$, we see that this constant must be $`1`$. that $`f(\theta )`$ can be represented as
$`f(\theta )={\displaystyle \frac{c^2}{16}}\mathrm{sinh}^2\theta {\displaystyle \frac{\underset{k=0}{\overset{N}{}}\mathrm{sinh}(\theta z_k^+)\mathrm{sinh}(\theta +z_k^+)\mathrm{sinh}(\theta z_k^{})\mathrm{sinh}(\theta +z_k^{})}{_{k=1}^N\mathrm{sinh}^2(\theta \theta _k)\mathrm{sinh}^2(\theta +\theta _k)}}.`$ (4.20)
It now follows by similar arguments that
$`𝔏(\theta |\theta _1,\mathrm{},\theta _N)=c\mathrm{sinh}\theta {\displaystyle \frac{\underset{k=0}{\overset{N}{}}\mathrm{sinh}(\frac{1}{2}(\theta z_k^+))\mathrm{sinh}(\frac{1}{2}(\theta +z_k^+))\mathrm{sinh}(\frac{1}{2}(\theta z_k^{}))\mathrm{sinh}(\frac{1}{2}(\theta +z_k^{}))}{_{k=1}^N\mathrm{sinh}(\frac{1}{2}(\theta \theta _k))\mathrm{sinh}(\frac{1}{2}(\theta +\theta _k))\mathrm{cosh}(\frac{1}{2}(\theta \theta _k))\mathrm{cosh}(\frac{1}{2}(\theta +\theta _k))}}.`$
(4.21)
is the unique solution to the inversion identity (4.13) with the properties (4.17). Note that there are $`N+1`$ pairs of roots $`z_k^\pm `$, whereas in the case of periodic boundary conditions there are only $`N`$. The appearance of the additional pair of roots $`z_0^\pm `$ is due to the fact that the boundary $`S`$ matrix $`𝖱(\theta ;\phi )`$ is not diagonal. The existence of these roots is essential for obtaining the correct asymptotic behavior; and it can be easily checked for the case $`N=0`$.
In summary, the eigenvalues of the bare transfer matrix (4.1) are given by (4.21), where $`z_k^\pm `$ satisfy Eqs. (4.18), (4.19).
### 4.3 Structure of Bethe Ansatz roots
Before performing the thermodynamic ($`N\mathrm{}`$) limit (which is the subject of the next section), it is necessary to first understand the structure of the Bethe Ansatz roots. Following , we observe that the Bethe Ansatz Eqs. (4.18), (4.19) have roots of the form
$`z_k^+=\{\begin{array}{c}x_k+iB\pi \\ x_k+iB\pi +i\pi \end{array},z_k^{}=\{\begin{array}{c}x_kiB\pi \\ x_kiB\pi i\pi \end{array},`$ (4.26)
where $`x_k`$ are real and satisfy
$`{\displaystyle \underset{j=1}{\overset{N}{}}}\left[{\displaystyle \frac{\mathrm{tanh}(\frac{1}{2}(x_k\theta _jiB\pi ))}{\mathrm{tanh}(\frac{1}{2}(x_k\theta _j+iB\pi ))}}{\displaystyle \frac{\mathrm{tanh}(\frac{1}{2}(x_k+\theta _jiB\pi ))}{\mathrm{tanh}(\frac{1}{2}(x_k+\theta _j+iB\pi ))}}\right]{\displaystyle \frac{\mathrm{sinh}^2(\frac{1}{2}(\frac{i\pi }{2}+x_k))}{\mathrm{sinh}^2(\frac{1}{2}(\frac{i\pi }{2}x_k))}}`$ (4.27)
$`\times `$ $`\left[{\displaystyle \frac{e^\phi _{}\mathrm{sinh}^2(\frac{1}{2}(i\pi B+\frac{i\pi }{2}))+e^\phi _{}\mathrm{sinh}^2(\frac{1}{2}(\frac{i\pi }{2}x_k))}{e^\phi _{}\mathrm{sinh}^2(\frac{1}{2}(i\pi B+\frac{i\pi }{2}))+e^\phi _{}\mathrm{sinh}^2(\frac{1}{2}(\frac{i\pi }{2}+x_k))}}\right]\times \left[\phi _{}\phi _+\right]=1,`$
$`k=0,1,\mathrm{},N.`$
Evidently, for each $`x_k`$, there are 4 possible combinations of roots $`(z_k^+,z_k^{})`$. However, by considering the limit $`B0`$, one can argue that only 2 of these combinations are allowed, which we denote by $`ϵ_k=+1`$ and $`ϵ_k=1`$, respectively:
$`ϵ_k`$ $`=`$ $`+1:(z_k^+=x_k+iB\pi ,z_k^{}=x_kiB\pi i\pi ),`$
$`ϵ_k`$ $`=`$ $`1:(z_k^+=x_k+iB\pi +i\pi ,z_k^{}=x_kiB\pi ).`$ (4.28)
Hence, the eigenvalues are specified by $`\{x_k,ϵ_k\}`$, $`k=0,\mathrm{},N`$:
$`𝔏(\theta |\theta _1,\mathrm{},\theta _N)_{ϵ_0\mathrm{}ϵ_N}=c\mathrm{sinh}\theta {\displaystyle \frac{\underset{k=0}{\overset{N}{}}\lambda _{ϵ_k}(\theta x_k)\lambda _{ϵ_k}(\theta +x_k)}{_{k=1}^N\frac{1}{4}\mathrm{sinh}(\theta \theta _k)\mathrm{sinh}(\theta +\theta _k)}},`$ (4.29)
where
$`\lambda _ϵ(\theta )=\mathrm{sinh}({\displaystyle \frac{1}{2}}(\theta ϵiB\pi ))\mathrm{cosh}({\displaystyle \frac{1}{2}}(\theta +ϵiB\pi )),`$ (4.30)
$`ϵ_k=\pm 1`$, and $`x_k`$ satisfy (4.27).
To close this section, we observe that the “dressed” transfer matrix (3.9) is simply related to the “bare” transfer matrix (4.1) by
$`\tau (\theta |\theta _1,\mathrm{},\theta _N)=𝖹(\theta +i\pi ;\xi _+)𝖹(\theta ;\xi _{}){\displaystyle \underset{k=1}{\overset{N}{}}}\left[Z(\theta \theta _k)Z(\theta +\theta _k)\right]𝔱(\theta |\theta _1,\mathrm{},\theta _N),`$ (4.31)
where the scalar factors $`Z(\theta )`$ and $`𝖹(\theta ;\xi )`$ are introduced in Eqs. (2.25), (2.59). Hence, the eigenvalues $`\mathrm{\Lambda }(\theta |\theta _1,\mathrm{},\theta _N)`$ of $`\tau (\theta |\theta _1,\mathrm{},\theta _N)`$ are given by
$`\mathrm{\Lambda }(\theta |\theta _1,\mathrm{},\theta _N)_{ϵ_0\mathrm{}ϵ_N}`$ $`=`$ $`{\displaystyle \frac{𝖹(\theta ;\xi _+)𝖹(\theta ;\xi _{})Z(2\theta )}{g(\frac{i\pi }{2}\theta )}}{\displaystyle \underset{k=1}{\overset{N}{}}}Z(\theta \theta _k)Z(\theta +\theta _k)`$ (4.32)
$`\times `$ $`𝔏(\theta |\theta _1,\mathrm{},\theta _N)_{ϵ_0\mathrm{}ϵ_N},`$
where $`𝔏(\theta |\theta _1,\mathrm{},\theta _N)_{ϵ_0\mathrm{}ϵ_N}`$ is given by (4.29). Here we have used the fact
$`𝖹(\theta +i\pi ;\xi )={\displaystyle \frac{𝖹(\theta ;\xi )Z(2\theta )}{g(\frac{i\pi }{2}\theta )}},`$ (4.33)
which follows from the bulk unitarity (2.7), (2.24) and boundary cross-unitarity (2.36), (2.55), (2.56) relations.
## 5 Thermodynamic Bethe Ansatz analysis
Having obtained the eigenvalues of the transfer matrix and the Bethe Ansatz equations, we can proceed to the derivation of the TBA equations and boundary entropy. We begin by briefly reviewing the general framework. Following , we consider the partition function $`Z_+`$ of the system on a cylinder of length $`L`$ and circumference $`R`$ with left/right boundary conditions denoted by $`\pm `$ (see Figure 4)
$`Z_+`$ $`=`$ $`tre^{RH_+}=e^{RF}`$ (5.1)
$`=`$ $`B_+|e^{LH_P}|B_{}`$
$``$ $`B_+|00|B_{}e^{LE_0}\text{ for }L\mathrm{}.`$
In the first line, Euclidean time evolves along the circumference of the cylinder, and $`H_+`$ is the Hamiltonian for the system with spatial boundary conditions $`\pm `$. In passing to the second line, we rotate the picture, so that time evolves parallel to the axis of the cylinder; $`H_P`$ is the Hamiltonian for the system with periodic boundary conditions, and $`|B_\pm `$ are boundary states which encode initial/final (temporal) conditions. In the third line, we consider the limit $`L\mathrm{}`$; the state $`|0`$ is the ground state of $`H_P`$, and $`E_0`$ is the corresponding eigenvalue. The quantity $`\mathrm{ln}B_+|00|B_{}`$ is the sought-after boundary entropy ,. <sup>11</sup><sup>11</sup>11More precisely, we shall compute the dependence of the boundary entropy on the boundary parameters. The term in the boundary entropy which is “constant” (independent of boundary parameters) seems to be difficult to compute even for simpler models ,. Taking the logarithm of the above expressions for the partition function, one obtains
$`RFLE_0+\mathrm{ln}B_+|00|B_{}.`$ (5.2)
Whereas the free energy $`F`$ has a leading contribution which is of order $`L`$, here we seek the subleading correction which is of order $`1`$.
### 5.1 Thermodynamic limit
We proceed to compute $`F`$ using the TBA approach , -. To this end, we introduce the densities $`P_\pm (\theta )`$ of “magnons”, i.e., of real Bethe Ansatz roots $`\{x_k\}`$ with $`ϵ_k=\pm 1`$, respectively; and also the densities $`\rho _1(\theta )`$ and $`\stackrel{~}{\rho }(\theta )`$ of particles $`\{\theta _k\}`$ and holes, respectively. Computing the logarithmic derivative of the “magnonic” Bethe Ansatz equations (4.27), we obtain <sup>12</sup><sup>12</sup>12The term $`\frac{1}{2\pi L}\mathrm{\Phi }(\theta )`$ originates from the exclusion , of the Bethe Ansatz root $`x_k=0`$.
$`P_+(\theta )+P_{}(\theta )`$ $`=`$ $`{\displaystyle \frac{1}{2\pi }}{\displaystyle _0^{\mathrm{}}}𝑑\theta ^{}\rho _1(\theta ^{})\left[\mathrm{\Phi }(\theta \theta ^{})+\mathrm{\Phi }(\theta +\theta ^{})\right]`$ (5.3)
$`+`$ $`{\displaystyle \frac{1}{2\pi L}}\left[\mathrm{\Phi }(\theta )+2\mathrm{\Psi }(\theta )+\mathrm{\Psi }_{\phi _+}(\theta )+\mathrm{\Psi }_\phi _{}(\theta )\right],`$
where
$`\mathrm{\Phi }(\theta )`$ $`=`$ $`{\displaystyle \frac{1}{i}}{\displaystyle \frac{}{\theta }}\mathrm{ln}\left({\displaystyle \frac{\mathrm{tanh}(\frac{1}{2}(\theta iB\pi ))}{\mathrm{tanh}(\frac{1}{2}(\theta +iB\pi ))}}\right)={\displaystyle \frac{4\mathrm{cosh}\theta \mathrm{sin}B\pi }{\mathrm{cosh}2\theta \mathrm{cos}2B\pi }},`$
$`\mathrm{\Psi }(\theta )`$ $`=`$ $`{\displaystyle \frac{1}{i}}{\displaystyle \frac{}{\theta }}\mathrm{ln}\left({\displaystyle \frac{\mathrm{sinh}(\frac{1}{2}(\frac{i\pi }{2}+\theta ))}{\mathrm{sinh}(\frac{1}{2}(\frac{i\pi }{2}\theta ))}}\right)={\displaystyle \frac{1}{\mathrm{cosh}\theta }},`$
$`\mathrm{\Psi }_\phi (\theta )`$ $`=`$ $`{\displaystyle \frac{1}{i}}{\displaystyle \frac{}{\theta }}\mathrm{ln}\left({\displaystyle \frac{e^\phi \mathrm{sinh}^2(\frac{1}{2}(i\pi B+\frac{i\pi }{2}))+e^\phi \mathrm{sinh}^2(\frac{1}{2}(\frac{i\pi }{2}\theta ))}{e^\phi \mathrm{sinh}^2(\frac{1}{2}(i\pi B+\frac{i\pi }{2}))+e^\phi \mathrm{sinh}^2(\frac{1}{2}(\frac{i\pi }{2}+\theta ))}}\right)`$ (5.4)
$`=`$ $`{\displaystyle \frac{4\mathrm{cosh}\theta \mathrm{cos}\zeta }{\mathrm{cosh}2\theta +\mathrm{cos}2\zeta }},`$
where $`\zeta `$ is defined in (2.58). Defining $`\rho _1(\theta )`$ for negative values of $`\theta `$ to be equal to $`\rho _1(|\theta |)`$, we obtain the final form
$`P_+(\theta )+P_{}(\theta )={\displaystyle \frac{1}{2\pi }}\left(\rho _1\mathrm{\Phi }\right)(\theta )+{\displaystyle \frac{1}{2\pi L}}\left[\mathrm{\Phi }(\theta )+2\mathrm{\Psi }(\theta )+\mathrm{\Psi }_{\phi _+}(\theta )+\mathrm{\Psi }_\phi _{}(\theta )\right],`$ (5.5)
where $``$ denotes convolution
$`\left(fg\right)(\theta )={\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\theta ^{}f(\theta \theta ^{})g(\theta ^{}).`$ (5.6)
We next consider the Yang equations, which imply (see Eqs. (3.7),(3.11))
$`e^{2iLm\mathrm{sinh}\theta _k}\mathrm{\Lambda }(\theta _k|\theta _1,\mathrm{},\theta _N)=1,k=1,\mathrm{},N,`$ (5.7)
where $`\mathrm{\Lambda }(\theta |\theta _1,\mathrm{},\theta _N)`$ is the eigenvalue of the dressed transfer matrix $`\tau (\theta |\theta _1,\mathrm{},\theta _N)`$, which is given by (4.32). Computing the logarithmic derivative, we obtain
$`\rho _1(\theta )+\stackrel{~}{\rho }(\theta )`$ $`=`$ $`{\displaystyle \frac{1}{2\pi }}\{2m\mathrm{cosh}\theta +{\displaystyle _0^{\mathrm{}}}d\theta ^{}\rho _1(\theta ^{})[\mathrm{\Phi }_Z(\theta \theta ^{})+\mathrm{\Phi }_Z(\theta +\theta ^{})]`$ (5.8)
$`+`$ $`{\displaystyle _0^{\mathrm{}}}d\theta ^{}[P_+(\theta ^{})\mathrm{\Phi }_+(\theta \theta ^{})+P_{}(\theta ^{})\mathrm{\Phi }_{}(\theta \theta ^{})`$
$`+`$ $`P_{}(\theta ^{})\mathrm{\Phi }_+(\theta +\theta ^{})+P_+(\theta ^{})\mathrm{\Phi }_{}(\theta +\theta ^{})]`$
$`+`$ $`{\displaystyle \frac{1}{L}}[\mathrm{\Phi }_Z(\theta )2\mathrm{\Phi }_Z(2\theta )+{\displaystyle \frac{}{\theta }}Im\mathrm{ln}𝖹(\theta ;\xi _+)+{\displaystyle \frac{}{\theta }}Im\mathrm{ln}𝖹(\theta ;\xi _{})]\},`$
where
$`\mathrm{\Phi }_Z(\theta )={\displaystyle \frac{}{\theta }}Im\mathrm{ln}Z(\theta ),\mathrm{\Phi }_\pm (\theta )={\displaystyle \frac{}{\theta }}Im\mathrm{ln}\lambda _\pm (\theta ).`$ (5.9)
Using the fact $`\mathrm{\Phi }_\pm (\theta )=\pm \frac{1}{2}\mathrm{\Phi }(\theta )`$, and defining $`P_\pm (\theta )`$ for negative values of $`\theta `$ to be equal to $`P_\pm (|\theta |)`$, we obtain
$`\rho _1(\theta )+\stackrel{~}{\rho }(\theta )`$ $`=`$ $`{\displaystyle \frac{m}{\pi }}\mathrm{cosh}\theta +{\displaystyle \frac{1}{2\pi }}\left(\rho _1\mathrm{\Phi }_Z\right)(\theta )+{\displaystyle \frac{1}{4\pi }}\left((P_+P_{})\mathrm{\Phi }\right)(\theta )`$ (5.10)
$`+`$ $`{\displaystyle \frac{1}{2\pi L}}\left[\mathrm{\Phi }_Z(\theta )2\mathrm{\Phi }_Z(2\theta )+{\displaystyle \frac{}{\theta }}Im\mathrm{ln}𝖹(\theta ;\xi _+)+{\displaystyle \frac{}{\theta }}Im\mathrm{ln}𝖹(\theta ;\xi _{})\right].`$
We now use (5.5) to eliminate $`P_{}`$, and use the expressions (2.59), (2.63) to separate the various factors in $`𝖹(\theta ;\xi )`$ to obtain
$`\rho _1(\theta )+\stackrel{~}{\rho }(\theta )`$ $`=`$ $`{\displaystyle \frac{m}{\pi }}\mathrm{cosh}\theta +{\displaystyle \frac{1}{2\pi }}P_+\mathrm{\Phi }+{\displaystyle \frac{1}{2\pi }}\rho _1\left(\mathrm{\Phi }_Z{\displaystyle \frac{1}{4\pi }}\mathrm{\Phi }\mathrm{\Phi }\right)`$
$`+`$ $`{\displaystyle \frac{1}{2\pi L}}[(\mathrm{\Phi }_Z{\displaystyle \frac{1}{4\pi }}\mathrm{\Phi }\mathrm{\Phi })+2({\displaystyle \frac{}{\theta }}Im\mathrm{ln}𝖹_0(\theta )\mathrm{\Phi }_Z(2\theta ){\displaystyle \frac{1}{4\pi }}\mathrm{\Psi }\mathrm{\Phi })`$
$`+`$ $`\left({\displaystyle \frac{}{\theta }}Im\mathrm{ln}𝖸_1(\theta ;\phi _+){\displaystyle \frac{1}{4\pi }}\mathrm{\Psi }_{\phi _+}\mathrm{\Phi }\right)+\left({\displaystyle \frac{}{\theta }}Im\mathrm{ln}𝖸_1(\theta ;\phi _{}){\displaystyle \frac{1}{4\pi }}\mathrm{\Psi }_\phi _{}\mathrm{\Phi }\right)`$
$`+`$ $`{\displaystyle \frac{1}{i}}{\displaystyle \frac{}{\theta }}\mathrm{ln}𝖷_1(\theta ;{\displaystyle \frac{4\eta _+B}{\pi }})𝖷_1(\theta ;{\displaystyle \frac{4i\vartheta _+B}{\pi }})+{\displaystyle \frac{1}{i}}{\displaystyle \frac{}{\theta }}\mathrm{ln}𝖷_1(\theta ;{\displaystyle \frac{4\eta _{}B}{\pi }})𝖷_1(\theta ;{\displaystyle \frac{4i\vartheta _{}B}{\pi }})].`$
Noting the “bulk” identity
$`\mathrm{\Phi }_Z(\theta ){\displaystyle \frac{1}{4\pi }}\left(\mathrm{\Phi }\mathrm{\Phi }\right)(\theta )=0,`$ (5.12)
and its boundary counterparts
$`{\displaystyle \frac{}{\theta }}Im\mathrm{ln}𝖸_1(\theta ;\phi ){\displaystyle \frac{1}{4\pi }}\left(\mathrm{\Psi }_\phi \mathrm{\Phi }\right)(\theta )`$ $`=`$ $`0,`$
$`{\displaystyle \frac{}{\theta }}Im\mathrm{ln}𝖹_0(\theta )\mathrm{\Phi }_Z(2\theta ){\displaystyle \frac{1}{4\pi }}\left(\mathrm{\Psi }\mathrm{\Phi }\right)(\theta )`$ $`=`$ $`{\displaystyle \frac{1}{4}}\mathrm{\Phi }(\theta )+\mathrm{\Psi }(\theta ),`$ (5.13)
we remain with the rather simple result
$`\rho _1(\theta )+\stackrel{~}{\rho }(\theta )`$ $`=`$ $`{\displaystyle \frac{m}{\pi }}\mathrm{cosh}\theta +{\displaystyle \frac{1}{2\pi }}(P_+\mathrm{\Phi })(\theta )+{\displaystyle \frac{1}{2\pi L}}[{\displaystyle \frac{1}{2}}\mathrm{\Phi }(\theta )+2\mathrm{\Psi }(\theta )`$ (5.14)
$`+`$ $`\kappa (\theta ;{\displaystyle \frac{4\eta _+B}{\pi }})+\kappa (\theta ;{\displaystyle \frac{4i\vartheta _+B}{\pi }})+\kappa (\theta ;{\displaystyle \frac{4\eta _{}B}{\pi }})+\kappa (\theta ;{\displaystyle \frac{4i\vartheta _{}B}{\pi }})],`$
where
$`\kappa (\theta ;F)={\displaystyle \frac{1}{i}}{\displaystyle \frac{}{\theta }}\mathrm{ln}𝖷_1(\theta ;F)={\displaystyle \frac{4\mathrm{cosh}\theta \mathrm{cos}(\pi F/2)}{\mathrm{cosh}2\theta +\mathrm{cos}\pi F}}.`$ (5.15)
The thermodynamic limit of the magnonic Bethe Ansatz equations and the Yang equations, given by (5.5) and (5.14), respectively, are the main results of this subsection. Notice that the former depends on the boundary parameters $`\phi _\pm `$, while the latter depends on the (boundary sinh-Gordon) boundary parameters $`\eta _\pm ,\vartheta _\pm `$.
### 5.2 TBA equations and boundary entropy
The free energy $`F`$ is given by
$`F=ETS,`$ (5.16)
where the temperature is $`T=\frac{1}{R}`$, the energy $`E`$ is
$`E={\displaystyle \underset{k=1}{\overset{N}{}}}m\mathrm{cosh}\theta _k={\displaystyle \frac{L}{2}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\theta \rho _1(\theta )m\mathrm{cosh}\theta ,`$ (5.17)
and the entropy $`S`$ is ,
$`S`$ $`=`$ $`{\displaystyle \frac{L}{2}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}d\theta \{(\rho _1+\stackrel{~}{\rho })\mathrm{ln}(\rho _1+\stackrel{~}{\rho })\rho _1\mathrm{ln}\rho _1\stackrel{~}{\rho }\mathrm{ln}\stackrel{~}{\rho }`$ (5.18)
$`+`$ $`(P_++P_{})\mathrm{ln}(P_++P_{})P_+\mathrm{ln}P_+P_{}\mathrm{ln}P_{}\}.`$
Extremizing the free energy $`(\delta F=0)`$ subject to the constraints
$`\delta P_{}`$ $`=`$ $`\delta P_++{\displaystyle \frac{1}{2\pi }}\delta \rho _1\mathrm{\Phi },`$
$`\delta \stackrel{~}{\rho }`$ $`=`$ $`\delta \rho _1+{\displaystyle \frac{1}{2\pi }}\delta P_+\mathrm{\Phi },`$ (5.19)
(which follow from Eqs. (5.5), (5.14), respectively) we obtain a set of TBA equations which is the same as for the case of periodic boundary conditions ,
$`r\mathrm{cosh}\theta `$ $`=`$ $`ϵ_1(\theta )+{\displaystyle \frac{1}{2\pi }}\left(\mathrm{\Phi }L_2\right)(\theta ),`$
$`0`$ $`=`$ $`ϵ_2(\theta )+{\displaystyle \frac{1}{2\pi }}\left(\mathrm{\Phi }L_1\right)(\theta ),`$ (5.20)
where
$`L_i(\theta )`$ $`=`$ $`\mathrm{ln}\left(1+e^{ϵ_i(\theta )}\right),r=mR,`$
$`ϵ_1`$ $`=`$ $`\mathrm{ln}\left({\displaystyle \frac{\stackrel{~}{\rho }}{\rho _1}}\right),ϵ_2=\mathrm{ln}\left({\displaystyle \frac{P_{}}{P_+}}\right).`$ (5.21)
We next evaluate $`F`$ using also the constraints (5.5), (5.14) and the TBA equations. From the boundary (order $`1`$) contribution, we obtain (see Eq. (5.2)) the boundary entropy
$`\mathrm{ln}B_+|00|B_{}={\displaystyle \frac{1}{4\pi }}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\theta `$ $`\{`$ $`[{\displaystyle \frac{1}{2}}\mathrm{\Phi }(\theta )+2\mathrm{\Psi }(\theta )+\kappa (\theta ;{\displaystyle \frac{4\eta _+B}{\pi }})+\kappa (\theta ;{\displaystyle \frac{4i\vartheta _+B}{\pi }})`$ (5.22)
$`+`$ $`\kappa (\theta ;{\displaystyle \frac{4\eta _{}B}{\pi }})+\kappa (\theta ;{\displaystyle \frac{4i\vartheta _{}B}{\pi }})]L_1(\theta )`$
$`+`$ $`\left[\mathrm{\Phi }(\theta )+2\mathrm{\Psi }(\theta )+\mathrm{\Psi }_{\phi _+}(\theta )+\mathrm{\Psi }_\phi _{}(\theta )\right]L_2(\theta ).`$
In particular, the dependence of the boundary entropy of a single boundary on the boundary parameters is given by <sup>13</sup><sup>13</sup>13For the case $`\epsilon =1`$, we obtain a similar result, except the parameter $`\zeta `$ appearing in the kernel $`\mathrm{\Psi }_\phi (\theta )`$ is now given by $`\zeta =\mathrm{cos}^1\left(1+e^{2\phi }(1\mathrm{sin}B\pi )\right)`$ instead of by Eq. (2.58).
$`s_B(\eta ,\vartheta ,\phi )={\displaystyle \frac{1}{4\pi }}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\theta \left\{\left[\kappa (\theta ;{\displaystyle \frac{4\eta B}{\pi }})+\kappa (\theta ;{\displaystyle \frac{4i\vartheta B}{\pi }})\right]L_1(\theta )+\mathrm{\Psi }_\phi (\theta )L_2(\theta )\right\},`$ (5.23)
where the kernels $`\kappa (\theta ;F)`$ and $`\mathrm{\Psi }_\phi (\theta )`$ are defined in Eqs. (5.15) and (5.4), respectively. The term involving $`L_1`$, which had previously been conjectured , depends on the boundary sinh-Gordon parameters $`\eta ,\vartheta `$. The term involving $`L_2`$, which had not been anticipated, depends on the boundary parameter $`\phi `$ (which appears in $`𝖱(\theta ;\phi )`$, i.e., the non-diagonal part of the boundary $`S`$ matrix). This expression for the boundary entropy is another of the main results of this paper.
## 6 Boundary roaming trajectories
One application of our result (5.23) for the boundary entropy is to obtain boundary roaming trajectories corresponding to $`c<3/2`$ superconformal models. In order to best explain this result, it is helpful to first recall earlier work on bulk and boundary roaming.
Al. Zamolodchikov first considered the TBA equations for the bulk ShG (non-supersymmetric) model with the coupling constant $`\gamma `$ analytically continued to complex values,
$`\gamma ={\displaystyle \frac{\pi }{2}}\pm i\theta _0,\theta _0>>1.`$ (6.1)
The corresponding effective central charge $`c_{eff}(r)`$ interpolates (“roams”) between the values
$`c_p=1{\displaystyle \frac{6}{p(p+1)}},p=3,4,5,\mathrm{}`$ (6.2)
corresponding to the unitary $`c<1`$ minimal models . Indeed, a plot of $`c_{eff}(r)`$ vs. $`\mathrm{log}(r/2)`$ reveals a “staircase” with plateaus at values of $`c_{eff}(r)`$ equal to $`c_p`$.
This result was later generalized to the boundary ShG model: choosing the value of $`r`$ so that $`c_{eff}(r)`$ lies on some plateau, the boundary entropy $`s_B(F)`$ (where $`F`$ is a boundary parameter) interpolates between values corresponding to various conformal boundary conditions .
The original work was also generalized to the bulk SShG (supersymmetric) model. The TBA equations with a similar analytic continuation of the coupling constant
$`\pi B={\displaystyle \frac{\pi }{2}}\pm i\theta _0,\theta _0>>1`$ (6.3)
cause the effective central charge $`c_{eff}(r)`$ to interpolate between the values
$`c_p={\displaystyle \frac{3}{2}}\left(1{\displaystyle \frac{8}{p(p+2)}}\right),p=4,6,8,\mathrm{}`$ (6.4)
corresponding to the even unitary $`c<3/2`$ minimal models , . Precisely this set of TBA equations had been conjectured earlier in , and then further generalized in .
Finally, let us consider the model of primary interest here, namely, boundary SShG. For simplicity, we fix $`\varphi _0=0`$ in the boundary Lagrangian (2.30), which corresponds to $`\vartheta =0`$.<sup>14</sup><sup>14</sup>14 Consider the boundary SSG model first. When $`\varphi _0=0`$ the total Lagrangian respects $`𝖢`$ symmetry due to the $`Z_2`$ symmetry $`\varphi \varphi `$. Therefore, the boundary $`S`$ matrix should respect $`𝖢`$ symmetry, namely the soliton and antisoliton should scatter equally on the boundary. Since the topological sector of the SSG $`S`$ matrix is encoded in the SG part, the boundary parameter $`\vartheta `$ should vanish as it does in the SG model . This holds also for the boundary SShG $`S`$ matrix because the two models are related by the fusion procedure. Due to the roaming limit (6.3), we should rescale the remaining two parameters $`\eta ,\phi `$ so that the boundary entropy can be a function of well-defined (finite) boundary parameters. For this purpose we set $`\eta 0`$ and $`\phi \mathrm{}`$ while keeping $`\theta _0\eta `$ and $`\theta _02\phi `$ finite. Let us introduce new boundary parameters $`f_1`$ and $`f_2`$ defined by (see Eq. (2.58))
$`{\displaystyle \frac{2\theta _0\eta }{\pi }}f_1,1+{\displaystyle \frac{1}{2}}e^{\theta _02\phi }\mathrm{cosh}f_2.`$ (6.5)
We can reexpress the boundary entropy (5.23) in terms of these parameters as
$`s_B=s_B^{(1)}+s_B^{(2)},\text{ where }s_B^{(i)}={\displaystyle \frac{1}{4\pi }}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\theta \mathrm{\Psi }(\theta ;f_i)L_i(\theta ),i=1,2,`$ (6.6)
with
$`\mathrm{\Psi }(\theta ;f)={\displaystyle \frac{4\mathrm{cosh}\theta \mathrm{cosh}f}{\mathrm{cosh}2\theta +\mathrm{cosh}2f}}.`$ (6.7)
To compute the roaming boundary entropy, we fix a value of $`r`$ where $`c_{eff}(r)`$ lies on a plateau (6.4). Then, as we change the boundary roaming parameters $`f_1`$ and $`f_2`$, we check if the boundary entropy interpolates between the values <sup>15</sup><sup>15</sup>15Note that this expression satisfies $`s_B(1,1)=0`$. The correct expression for the conformal boundary entropies has an additional “constant” term (i.e., independent of both $`r`$ and $`s`$); we neglect this term here, since we are mostly interested in differences $`s_B(r,s)s_B(r^{},s^{})`$, for which the constant term cancels.
$`s_B(r,s)`$ $`=`$ $`\mathrm{ln}\left[\left({\displaystyle \frac{\mathrm{sin}(\frac{\pi r}{p})}{\mathrm{sin}(\frac{\pi }{p})}}\right)\left({\displaystyle \frac{\mathrm{sin}(\frac{\pi s}{p+2})}{\mathrm{sin}(\frac{\pi }{p+2})}}\right)\right]`$ (6.8)
$`=`$ $`s_B(r,1)+s_B(1,s).`$
corresponding to conformal boundary states $`(r,s)`$ (which, in turn, correspond <sup>16</sup><sup>16</sup>16We recall that for each bulk primary field $`\mathrm{\Phi }_{(r,s)}`$, there corresponds a conformal boundary state $`|\stackrel{~}{h}_{(r,s)}`$ (which, for brevity, we denote here by $`(r,s)`$ ) such that the partition function $`Z_{(1,1)(r,s)}`$ for the CFT on a cylinder with conformal boundary states $`(1,1)`$ and $`(r,s)`$ is given by $`Z_{(1,1)(r,s)}=\chi _{(r,s)}(q),`$ i.e., the character of $`\mathrm{\Phi }_{(r,s)}`$. In particular, $`Z_{(1,1)(1,1)}=\chi _{(1,1)}(q)`$ is the character of the unit operator. to primary fields $`\mathrm{\Phi }_{(r,s)}`$).
Indeed, we can see clearly from Figure 5 that $`s_B^{(1)}`$ interpolates between boundary entropies of the conformal boundary states
$`(1,a)\{\begin{array}{c}(a2,1)\\ (a,1)\end{array},a\text{ odd }.`$ (6.11)
Similarly $`s_B^{(2)}`$ generates the new flow (see Figure 6) <sup>17</sup><sup>17</sup>17For $`p>4`$, we cannot associate any conformal boundary state to the final plateau (i.e., for asymptotically large values of the boundary parameter $`f_2`$), since there is no state $`(0,1)`$.
$`(1,a)\{\begin{array}{c}(a2,1)\\ (a,1)\end{array},a\text{ even }.`$ (6.14)
While these flows are generated by changing one parameter while fixing the other, we can generate more general flows by changing $`f_1`$ and $`f_2`$ simultaneously. In view of the additivity property (6.8), these two sets of flows can be combined to generate additional flows for the total boundary entropy $`s_B`$
$`(r,s)\{\begin{array}{c}(s,r)\\ (s2,r)\\ (s,r+2)\\ (s2,r+2)\end{array},rs=\text{ odd }.`$ (6.19)
Note that $`rs=`$ even/odd corresponds to the Neveu-Schwarz/Ramond sectors, respectively.
## 7 Discussion
We have presented the exact solution of the boundary SShG model – an integrable QFT whose bulk and boundary $`S`$ matrices are not diagonal. In particular, we have derived an exact inversion identity (4.9) - (LABEL:functionf2), as well as the TBA equations and boundary entropy (5.23). Moreover, we have uncovered a rich pattern of boundary roaming trajectories, which remain to be understood in detail.
Although the boundary SShG model has a special feature which allows it to be solved by an inversion identity (namely, the bulk $`S`$ matrix satisfies by free-Fermion condition (2.15)), it is by no means the only such model. Indeed, there are infinite families of integrable QFTs with $`N=1`$ or $`N=2`$ supersymmetry - that have this property. These models have bulk and boundary $`S`$ matrices which are similar to those of SShG, and therefore, we expect similar inversion identities to hold. We hope to report on these models in the near future .
Finally, we recall that one can readily obtain the Hamiltonian of an integrable open quantum spin chain with $`N`$ spins from any homogeneous open-chain transfer matrix $`𝔱(\theta |0)`$ (4.1). Indeed, the Hamiltonian $``$ is given by
$`{\displaystyle \frac{}{\theta }}𝔱(\theta |0)|_{\theta =0},`$ (7.1)
which commutes with $`𝔱(\theta |0)`$. For the $`R`$ matrices which we have considered here (2.13), (2.40), the corresponding Hamiltonian is that of a certain anisotropic XY chain with both bulk and boundary magnetic fields. By determining the eigenvalues (4.21) of the transfer matrix, we have evidently also solved the corresponding open quantum spin chain. It would be interesting to exploit this solution to determine properties of this model in the thermodynamic limit.
## Acknowledgments
We thank O. Alvarez, D. Bernard, E. Corrigan, G. Delius, P. Dorey, P. Fendley, M. Martins and H. Saleur for helpful comments and/or correspondence. One of us (R.N.) is grateful for the hospitality at the APCTP in Seoul (where this work was initiated) and at the CRM in Montreal (where the results were first reported). This work was supported in part by KOSEF 1999-2-112-001-5 (C.A.) and by the National Science Foundation under Grant PHY-9870101 (R.N.).
## Appendix A Relation of Yang matrix to Sklyanin transfer matrix
In Section 3.2, we stated that the Yang matrix (3.8) is related to the Sklyanin open-chain transfer matrix (3.9) in the following way (3.11):
$`Y_{(i)}=\tau (\theta _i|\theta _1,\mathrm{},\theta _N),i=1,\mathrm{},N.`$ (A.1)
We present here a proof for the case $`N=2`$. Evaluating the transfer matrix at $`\theta =\theta _1`$, we have
$`\tau (\theta _1|\theta _1,\theta _2)`$ $`=`$ $`tr_0\{𝖲_0(\theta _1+i\pi ;\xi _+)^{t_0}S_{02}(\theta _1\theta _2)S_{01}(0)𝖲_0(\theta _1;\xi _{})S_{01}(2\theta _1)S_{02}(\theta _1+\theta _2)\}`$ (A.2)
$`=`$ $`tr_0\{S_{02}(\theta _1\theta _2)𝒫_{01}𝖲_0(\theta _1;\xi _{})(𝒫_{01}𝒫_{01})S_{01}(2\theta _1)(𝒫_{01}𝒫_{01})`$
$`\times `$ $`S_{02}(\theta _1+\theta _2)(𝒫_{01}𝒫_{01})𝖲_0(\theta _1+i\pi ;\xi _+)^{t_0}\}=\mathrm{}`$
In passing to the second line, we have used the cyclic property of the trace, as well as $`S(0)=𝒫`$ and $`𝒫^2=𝕀`$, where $`𝒫`$ is the permutation matrix (2.22).
$`\mathrm{}=𝖲_1(\theta _1;\xi _{})tr_0\{S_{02}(\theta _1\theta _2)S_{01}(2\theta _1)S_{12}(\theta _1+\theta _2)𝒫_{01}𝖲_0(\theta _1+i\pi ;\xi _+)^{t_0}\}=\mathrm{}`$ (A.3)
Here we have used $`𝒫_{01}X_0𝒫_{01}=X_1`$, and the $`𝖯`$ symmetry of the $`R`$ matrix (2.17).
$`\mathrm{}=𝖲_1(\theta _1;\xi _{})tr_0\{S_{12}(\theta _1+\theta _2)S_{01}(2\theta _1)S_{02}(\theta _1\theta _2)𝒫_{01}𝖲_0(\theta _1+i\pi ;\xi _+)^{t_0}\}=\mathrm{}`$ (A.4)
Here we have used the Yang-Baxter equation (2.16).
$`\mathrm{}`$ $`=`$ $`𝖲_1(\theta _1;\xi _{})S_{12}(\theta _1+\theta _2)tr_0\{S_{01}(2\theta _1)(𝒫_{01}𝒫_{01})S_{02}(\theta _1\theta _2)𝒫_{01}𝖲_0(\theta _1+i\pi ;\xi _+)^{t_0}\}`$ (A.5)
$`=`$ $`𝖲_1(\theta _1;\xi _{})S_{12}(\theta _1+\theta _2)tr_0\{S_{01}(2\theta _1)𝒫_{01}𝖲_0(\theta _1+i\pi ;\xi _+)^{t_0}\}S_{12}(\theta _1\theta _2)`$
$`=`$ $`𝖲_1(\theta _1;\xi _{})S_{12}(\theta _1+\theta _2)𝖲_1(\theta _1;\xi _+)S_{12}(\theta _1\theta _2).`$
In passing to the last line, we have used the boundary cross-unitarity relation (2.46) with $`\theta =\frac{i\pi }{2}\theta _1`$, and the crossing relation $`S_{01}(i\pi \theta )^{t_1}=S_{01}(\theta )`$. Comparing the last line to the expression (3.8) for the Yang matrix, we conclude that
$`\tau (\theta _1|\theta _1,\theta _2)=Y_{(1)}.`$ (A.6)
For higher values of $`N`$, the proof is similar.
## Appendix B Derivation of inversion identity
In Section 4.1, we give the important inversion identity (4.9) - (LABEL:functionf2). Here we explain in more detail how we derived it. As already mentioned in text, the main idea is to formulate the fusion formula, following Ref. , to which we shall refer as I. <sup>18</sup><sup>18</sup>18In order to facilitate comparison with , we use here similar notations.
Although the “dressed” bulk $`S`$ matrix $`S(\theta )`$ (2.28) is regular at $`\theta =0`$, the “bare” bulk $`S`$ matrix $`R(\theta )`$ (2.13) has a pole there. In order to avoid complications from this spurious pole, in this Appendix we rescale $`R(\theta )`$ by the factor $`\mathrm{sinh}\theta `$; i.e., we take $`R(\theta )`$ to be given still by (2.13), but now with matrix elements
$`a_\pm (\theta )`$ $`=`$ $`\pm \mathrm{sinh}\theta 2i\mathrm{sin}B\pi ,b(\theta )=\mathrm{sinh}\theta ,`$
$`c(\theta )`$ $`=`$ $`2i\mathrm{sin}B\pi \mathrm{cosh}{\displaystyle \frac{\theta }{2}},d(\theta )=2\mathrm{sin}B\pi \mathrm{sinh}{\displaystyle \frac{\theta }{2}}.`$ (B.1)
Keeping in mind the symmetries (2.17) of the $`R`$ matrix, the unitarity relation (I 2.3) is
$`R_{12}(\theta )R_{12}(\theta )=\zeta (\theta )𝕀,\zeta (\theta )=4\mathrm{cosh}^2{\displaystyle \frac{\theta }{2}}(\mathrm{sinh}^2{\displaystyle \frac{\theta }{2}}+\mathrm{sin}^2B\pi ),`$ (B.2)
and the crossing relation (I 2.4) is
$`R_{12}(\theta )=V_1R_{12}(\theta \rho )^{t_2}V_1,`$ (B.3)
with <sup>19</sup><sup>19</sup>19Alternatively, choosing $`\rho =i\pi `$, one has $`V=𝕀`$.
$`\rho =i\pi ,V=\left(\begin{array}{cc}\hfill 1& \hfill 0\\ \hfill 0& \hfill 1\end{array}\right).`$ (B.6)
The matrix $`R_{12}(\theta )`$ at $`\theta =\rho `$ is proportional to the one-dimensional projector $`\stackrel{~}{P}_{12}^{}`$
$`\stackrel{~}{P}_{12}^{}={\displaystyle \frac{1}{2}}\left(\begin{array}{cccc}\hfill 1& \hfill 0& \hfill 0& \hfill 1\\ \hfill 0& \hfill 0& \hfill 0& \hfill 0\\ \hfill 0& \hfill 0& \hfill 0& \hfill 0\\ \hfill 1& \hfill 0& \hfill 0& \hfill 1\end{array}\right),(\stackrel{~}{P}_{12}^{})^2=\stackrel{~}{P}_{12}^{}.`$ (B.11)
As explained in I, from the corresponding degeneration of the (boundary) Yang-Baxter equation, one can derive identities which allow one to prove that fused (boundary) $`S`$ matrices satisfy generalized (boundary) Yang-Baxter equations.
The fused $`R`$ matrix is given by (I 2.13)
$`R_{<12>3}(\theta )=\stackrel{~}{P}_{12}^+R_{13}(\theta )R_{23}(\theta +\rho )\stackrel{~}{P}_{12}^+,`$ (B.12)
where $`\stackrel{~}{P}_{12}^+=𝕀\stackrel{~}{P}_{12}^{}`$. An important observation (which one can verify by direct calculation) is that the fused $`R`$ matrix can be brought to upper triangular form by a similarity transformation <sup>20</sup><sup>20</sup>20This observation is similar to, but not the same as, the one made by Felderhof . Indeed, in our language, he shows that $`R_{13}(\theta )R_{23}(\theta +\rho )`$ (i.e., the expression for the fused transfer matrix without the projectors $`\stackrel{~}{P}_{12}^+`$) can be brought to triangular form by a (somewhat more complicated) $`\theta `$-independent similarity transformation. Although for the case of periodic boundary conditions both approaches lead to the inversion identity, this appears to be no longer true for the case of boundaries.
$`X_{12}R_{<12>3}(\theta )X_{12}=\text{ upper triangular },`$ (B.13)
where the $`4\times 4`$ matrix $`X`$ is independent of $`\theta `$, and is given by
$`X=\left(\begin{array}{cccc}\frac{1}{\sqrt{2}}& 0& 0& \frac{1}{\sqrt{2}}\\ 0& \mathrm{sin}\frac{B\pi }{2}& \mathrm{cos}\frac{B\pi }{2}& 0\\ 0& \mathrm{cos}\frac{B\pi }{2}& \mathrm{sin}\frac{B\pi }{2}& 0\\ \frac{1}{\sqrt{2}}& 0& 0& \frac{1}{\sqrt{2}}\end{array}\right),X^2=𝕀.`$ (B.18)
It follows that the fused monodromy matrices <sup>21</sup><sup>21</sup>21For simplicity, we consider here the homogeneous case ($`\theta _i=0,i=1,\mathrm{},N`$). (I 4.7), (I 5.4), (I 5.5)
$`T_{<12>}(\theta )`$ $`=`$ $`R_{<12>N}(\theta )\mathrm{}R_{<12>1}(\theta ),`$
$`\widehat{T}_{<12>}(\theta +\rho )`$ $`=`$ $`R_{<12>1}(\theta )\mathrm{}R_{<12>N}(\theta ),`$ (B.19)
also become triangular by the same transformation.
Denoting (as in I) our “bare” boundary $`S`$ matrices $`𝖱(\theta ;\phi _{})`$, $`𝖱(\theta +i\pi ;\phi _+)`$ by $`K^{}(\theta )`$, $`K^+(\theta )`$, respectively, the corresponding fused matrices are given by (I 3.5), (I 3.9)
$`K_{<12>}^{}(\theta )`$ $`=`$ $`\stackrel{~}{P}_{12}^+K_1^{}(\theta )R_{12}(2\theta +\rho )K_2^{}(\theta +\rho )\stackrel{~}{P}_{12}^+,`$
$`K_{<12>}^+(\theta )`$ $`=`$ $`\{\stackrel{~}{P}_{12}^+K_1^+(\theta )^{t_1}R_{12}(2\theta 3\rho )K_2^+(\theta +\rho )^{t_2}\stackrel{~}{P}_{12}^+\}^{t_{12}},`$ (B.20)
since $`M=V^tV=𝕀`$.
Remarkably, the fused $`K`$ matrices are also brought to upper triangular form by the same similarity transformation
$`X_{12}K_{<12>}^{}(\theta )X_{12}=\text{ upper triangular }.`$ (B.21)
It follows that the fused transfer matrix $`\stackrel{~}{}𝔱(\theta )`$, which is given by (I 4.5), (I 4.6)
$`\stackrel{~}{}𝔱(\theta )=tr_{12}K_{<12>}^+(\theta )T_{<12>}(\theta )K_{<12>}^{}(\theta )\widehat{T}_{<12>}(\theta +\rho ),`$ (B.22)
is proportional to the identity matrix,
$`\stackrel{~}{}𝔱(\theta )𝕀,`$ (B.23)
where the proportionality factor is determined from the diagonal elements of the various triangular matrices.
The fusion formula is given by (I 4.17), (I 5.1)
$`𝔱(\theta )𝔱(\theta +\rho )={\displaystyle \frac{1}{\zeta (2\theta +2\rho )}}\left[\stackrel{~}{}𝔱(\theta )+\mathrm{\Delta }\left\{K^+(\theta )\right\}\mathrm{\Delta }\left\{K^{}(\theta )\right\}\delta \left\{T(\theta )\right\}\delta \left\{\widehat{T}(\theta )\right\}\right],`$ (B.24)
where the transfer matrix $`𝔱(\theta )`$ is given by (4.1) (see also (I 4.1), (I 4.2)), and the quantum determinants are given by (I 4.15), (I 5.3), (I 5.7)
$`\delta \left\{T(\theta )\right\}`$ $`=`$ $`\delta \left\{\widehat{T}(\theta )\right\}=\zeta (\theta +\rho )^N,`$
$`\mathrm{\Delta }\left\{K^{}(\theta )\right\}`$ $`=`$ $`tr_{12}\left\{\stackrel{~}{P}_{12}^{}K_1^{}(\theta )R_{12}(2\theta +\rho )K_2^{}(\theta +\rho )V_1V_2\right\},`$
$`\mathrm{\Delta }\left\{K^+(\theta )\right\}`$ $`=`$ $`tr_{12}\left\{\stackrel{~}{P}_{12}^{}V_1V_2K_2^+(\theta +\rho )R_{12}(2\theta 3\rho )K_1^+(\theta )\right\}.`$ (B.25)
Reverting to the original normalization of the $`R`$ matrix by rescaling each of the transfer matrices $`𝔱(\theta )`$ in (B.24) by $`(\mathrm{sinh}\theta )^{2N}`$, introducing the inhomogeneities $`\theta _i`$ in the obvious way, and factoring the result into a product of two factors, we arrive at the results (4.9) - (LABEL:functionf2). <sup>22</sup><sup>22</sup>22We have refrained from giving explicit results for the intermediate steps, which are rather unwieldy and not very illuminating. We have done these computations with the help of Mathematica.
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# Multiphase Smoothed-Particle Hydrodynamics
## 1 Introduction
Since its introduction more than two decades ago by Lucy (1977) and Gingold & Monaghan (1977), Smoothed-Particle Hydrodynamics (SPH) has become one of the standard techniques for modelling astrophysical fluid flow (e.g. Evrard 1988; Hernquist & Katz 1989, hereafter HK89; Thomas & Couchman 1992; Steinmetz & Müller 1993; Couchman, Thomas & Pearce 1995, hereafter CTP95; Shapiro et al. 1996). SPH is fully Lagrangian, with the particles themselves being the framework on which the fluid equations are solved, and so there is no grid to constrain the dynamic range or geometry of the system being modelled. This is of particular importance for phenomena involving the growth of gravitational fluctuations, such as cosmological structure formation (Bertschinger 1998) and star formation (Bhattal et al. 1998), and in an adaptive form (Wood 1981; Nelson & Papaloizou 1994) the SPH algorithm will follow the wide range of densities encountered without difficulty. In contrast, computational demands severely restrict the dynamic range of Eulerian finite-difference methods (e.g. the Piecewise-Parabolic Method of Collela & Woodward 1984), and to date the best three-dimensional Eulerian simulations of galaxy formation have a gas resolution of $`300500`$ kpc (Blanton et al. 1999), although adaptive mesh refinement (AMR) techniques promise to greatly enhance the fixed-grid approach. SPH can be easily integrated with a range of $`N`$-body gravity solvers such as the tree algorithm used by HK89 and the Adaptive Particle-Particle, Particle-Mesh (AP<sup>3</sup>M) algorithm of Couchman (1991).
However, SPH is not without its problems. The need for an artificial viscosity means that SPH resolves shocks poorly in comparison with finite-difference methods. Pairwise artificial viscosities (Monaghan & Gingold 1983) generally give the sharpest resolution of shocks, but also introduce a large shear viscosity unless correction terms are used (Balsara 1995). There is some concern that these terms further degrade the shock capturing ability of SPH (Navarro & Steinmetz 1997; Thacker et al. 1998, hereafter T98). While it is possible to add a physical viscosity to SPH and solve the Navier-Stokes equations directly (e.g. Flebbe et al. 1994), the relatively simple SPH interpolation method is quite sensitive to particle disorder and tends to give large errors in the higher-order dissipative terms. Boundary conditions are also difficult to implement in SPH and do not sit naturally with the method. Generally boundaries are either periodic or situated far from the region of interest, and SPH is unsuitable for simulations in which complex boundary conditions are of critical importance.
Standard implementations of SPH also have a limited ability to resolve steep density gradients, and a number of numerical problems can occur when particles are close to, but not physically part of, a region of higher density. These arise because the usual formulation of SPH assumes that the density gradient across the smoothing kernel of each particle is small. However, this is not true in many situations in which SPH is commonly used. In simulations of galaxy formation, for example, thermal instability causes cold, dense clumps to form within halos of hot gas, and density contrasts of several orders of magnitude can occur within the typical smoothing length of the halo gas. Current implementations of SPH will overestimate the density of the halo gas, leading to the gas cooling excessively and accreting on to the cold clump (T98; Pearce et al. 1999). Tittley, Couchman and Pearce (1999, hereafter TCP99) have also shown that the drag exerted on these protogalaxies by the intracluster medium can be seriously overestimated. These problems are thought to contribute to the overmerging commonly seen in simulations of structure formation including gas (Frenk et al. 1996).
We wish to extend the method to cope with these problems, along with multiphase fluids, such as the intracluster medium and cooling flows, which consist of an emulsion of discrete phases in which there is no correlation between the density of neighbouring particles. We should emphasise however that we are not seeking to model homogeneously mixed materials with different equations of state, such as dust and aerosols (Monaghan 1997). We make two changes to the standard implementation of SPH. Firstly, particles no longer have to try and maintain a fixed number of neighbours, as is required by the algorithm of HK89, and instead the smoothing length is adjusted to keep a density-weighted quantity constant (see Section 2.2). This prevents the smoothing length of low density particles from decreasing as they approach a high density region. Secondly we make the assumption that pressure, not density, is constant across the smoothing kernel. We then summate the local pressure at each particle and calculate the local density from the equation of state. This is a reasonable approach to take when the cooling time of the gas is greater than the local dynamic time, where a local pressure equilibrium can be expected even when the local density gradient is steep. It will not be true in the presence of shocks (although neither is the assumption that density is smoothly varying). As we will see in sections 3.1 and 3.2, however, there is no evidence that the shock capturing ability of SPH is degraded.
The layout of the paper is as follows. In section 2 we detail our new implementation of SPH. In Section 3 we present the results of a number of tests of the new method, and compare the performance with results produced using the SPH implementations of CTP95 and T98. This is followed in Section 4 by a discussion of our results.
## 2 Methodology
In SPH the fluid is represented by particles of known mass, $`m_i`$, and specific energy, $`ϵ_i`$. Other properties must be inferred by averaging over a smoothing sphere that typically extends to enclose a fixed number, $`N_{\mathrm{SPH}}`$, of particles. For a continuous distribution, the average value of some quantity $`A`$ at the location of particle $`i`$ would be
$$A_i=A(𝐫)W(𝐫𝐫_𝐢)dV$$
(1)
where $`W`$ is the smoothing kernel. In SPH the integral is replaced by a sum:
$$A_i=\underset{j}{}\frac{m_j}{\rho _j}A_jW_{ij}$$
(2)
that extends over all particles, $`j`$, within the smoothing sphere. $`\rho _j`$ is the density of particle $`j`$ and so $`m_j/\rho _j`$ is its volume. Here
$$W_{ij}=\frac{1}{h_{ij}^3}W(r_{ij}/h_{ij}),$$
(3)
where $`r_{ij}`$ is the separation of particles $`i`$ and $`j`$, $`r_{ij}=|𝐫_i𝐫_j|`$, and $`h_{ij}`$ is the smoothing length. There are several possible forms for $`W`$. Throughout this paper, we use the standard form described in Thomas & Couchman (1992). The choice of $`h_{ij}`$ determines the size of the region over which the density is to be averaged. Many authors take $`h_{ij}`$ to be a symmetric function (for example the harmonic average) of $`h_i`$ and $`h_j`$ (see T99 and references therein). We prefer to set $`h_{ij}=h_i`$. This has the advantages that we know in advance exactly how far we have to search and always have the same number of neighbours within the smoothing region. The integral of the smoothing kernel over all space is unity, which translates to the condition
$$\underset{j}{}\frac{m_j}{\rho _j}W_{ij}=1,$$
(4)
although this will only be true in a statistical sense.
Equation 2 seems to require the value of $`\rho _j`$, which is not known in advance. However, it is possible to circumvent this by formulating SPH such that $`A`$ is always a multiple of density and a known quantity. Taking $`A=\rho `$ gives the standard SPH estimate for the density in the neighbourhood of a particle
$$\overline{\rho }_i=\underset{j}{}m_jW_{ij}.$$
(5)
Where expressions arise involving derivatives, we use the divergence theorem to transfer the derivative onto the smoothing kernel as follows:
$$A_i=AW\mathrm{d}V=AW\mathrm{d}V$$
(6)
$$.𝐀_i=.𝐀W\mathrm{d}V=𝐀.W\mathrm{d}V$$
(7)
### 2.1 Evaluation of the density
In multiphase SPH, it is important to distinguish between the density of a particle, $`\rho _i`$, and the mean density in the neighbourhood of a particle, $`\overline{\rho }_i`$, defined by Equation 5. The pressure of the gas is expected to be a smooth quantity because (excluding shocks) the sound-crossing time is shorter than the flow time across the smoothing sphere. The density of individual particles, $`\rho _i`$, can therefore be determined from an estimate of the local pressure and the equation of state,
$$P=2/3\rho ϵ$$
(8)
The specific energy $`ϵ`$ is related to the gas temperature $`T`$ by $`ϵ=3kT/2\mu m_H`$, where $`k`$ is Boltzmann’s constant, $`m_H`$ is the mass of the hydrogen atom and $`\mu =0.6`$ is the relative molecular mass. The SPH estimate of the pressure is
$$P_i=2/3\underset{j}{}m_jϵ_jW_{ij}$$
(9)
and the density of a particle can therefore be written
$$\rho _i=\frac{3P_i}{2ϵ_i}=\frac{\underset{j}{}m_jϵ_jW_{ij}}{ϵ_i}.$$
(10)
Variations in $`ϵ_i`$ can cause large variations in density between neighbouring particles, whereas Equation 5, being an average, varies only slightly between neighbours. Note that a cold, high-density clump of particles will contribute a lot of terms to Equation 10, but each of these will be given a low weight due to the presence of the $`ϵ_j`$ term.
The consistency condition, Equation 4, becomes
$$\underset{j}{}\frac{2m_jϵ_jW_{ij}}{3P_j}=1$$
(11)
which will be true if the pressure is slowly-varying, as we have assumed.
### 2.2 The smoothing length
The SPH smoothing length is usually defined in terms of the radius of a sphere, centred on the particle in question, that encloses a fixed number of neighbours. One can either search for such a radius on each time-step, a time-consuming process, or allow the number of neighbours to vary and accept an estimate based on the number of neighbours found on the previous step:
$$h_ih_i\left[\alpha +(1\alpha )\left(\frac{N_{\mathrm{SPH}}}{N_i}\right)^{\frac{1}{3}}\right],$$
(12)
where $`N_i`$ is the actual number of neighbours and $`N_{\mathrm{SPH}}`$ is the desired number (HK89). $`\alpha `$ is a convergence parameter: for a uniform distribution, choosing $`\alpha =1`$ returns the correct value of $`h_i`$ on the next step. Where large density contrasts are present, however, $`\alpha `$ must be reduced to avoid overshooting and convergence is much slower. We use a convergence parameter $`\alpha =0.4`$ and $`N_{\mathrm{SPH}}=32`$.
For the extreme density contrasts we envisage, T98 have shown that $`N_i`$ can oscillate between values which are alternately too low and too high. They solve this problem both by giving lower weight to particles at the extreme edge of the smoothing sphere, and by the introduction of a more sophisticated convergence algorithm. In this paper, we suggest a simpler approach that uses a neighbour count weighted as follows:
$$N_i=\underset{j}{}w_{ij}=\underset{j}{}\frac{2\overline{\rho }_i}{\overline{\rho }_i+\rho _j},$$
(13)
where $`\overline{\rho }_i`$ is the mean density at particle $`i`$ and the sum extends over all neighbours within the smoothing sphere. Note that:
1. For a uniform distribution in which $`\overline{\rho }_i=\rho _j`$, then $`w_{ij}=1`$ and $`N_i`$ is simply equal to the number of neighbours.
2. Particles do not notice neighbours with a much greater than average density ($`w_{ij}0`$ for $`\overline{\rho }_i\rho _j`$). This prevents the common instability whereby changing the radius of the smoothing sphere so as to include or exclude a high-density clump can dramatically change the number count.
3. Neighbours with a lower than average density count at most double ($`w_{ij}2`$ for $`\overline{\rho }_i\rho _j`$). While this could lead to an isolated cold particle in a hot medium having fewer neighbours than is desirable, such a situation is unlikely to arise in practice and $`w_{ij}`$ could be limited to a maximum value of unity if desired.
### 2.3 The equation of motion
In the absence of artificial viscosity, heating and radiative cooling, the equation of motion for a parcel of gas is simply
$$\frac{\mathrm{d}𝐯}{\mathrm{d}t}=\frac{P}{\rho }.$$
(14)
In the same spirit as above, we require an estimate of $`P`$ that does not depend upon the number density of particles. This is
$$P_i=\frac{2}{3}\underset{j}{}m_jϵ_j_iW_{ij}.$$
(15)
It is common practice to use the following identity
$$\frac{P}{\rho }=\left(\frac{P}{\rho }\right)+\frac{P}{\rho ^2}\rho $$
(16)
to symmetrise the pressure force. This expression is not suitable here, both because we are envisaging abrupt changes in density (so that $`\rho `$ is undefined) and because the estimator for $`P/\rho ^2`$ depends upon the density of the particles. Equation 16 is obtained using the general identity
$$\frac{P}{\rho }=\frac{P}{\rho ^\sigma }\left(\frac{1}{\rho ^{1\sigma }}\right)+\frac{1}{\rho ^{2\sigma }}\left(\frac{P}{\rho ^{\sigma 1}}\right)$$
(17)
(Monaghan, 1992) and using $`\sigma =1`$ gives an alternative identity that at first sight seems a little bizarre:
$$\frac{P}{\rho }=\frac{P}{\rho }+\frac{P}{\rho }1.$$
(18)
This leads to the estimator for the force on particle $`i`$,
$$𝐟_i=\underset{j}{}\frac{2}{3}m_im_j\left[\frac{ϵ_j_iW_{ij}}{\rho _i}+\frac{ϵ_i_iW_{ji}}{\rho _j}\right],$$
(19)
where we have chosen to use $`W_{ji}`$ in place of $`W_{ij}`$ in the second term on the right hand side in order to make the force symmetric. This gives
$$𝐟_i=\underset{j}{}\left(𝐟_{ij}𝐟_{ji}\right),$$
(20)
where
$$𝐟_{ij}=\frac{2}{3}m_im_j\frac{ϵ_j_iW_{ij}}{\rho _i}.$$
(21)
The first set of terms in Equation 20 is evaluated when calculating the SPH properties for particle $`i`$; the second set is accumulated in stages when calculating the properties of the neighbours. Note that the only change from the usual SPH formalism is that $`ϵ_i`$ and $`ϵ_j`$ have exchanged places in Equation 19. However, this change means that the force on particle $`i`$ is dependent only upon the local pressure gradient and not upon the density of its neighbours.
### 2.4 Conservation of energy
Conservation of energy is ensured by equating the heating to the $`P\mathrm{d}V`$ work done:
$$\frac{\mathrm{d}ϵ}{\mathrm{d}t}=\frac{P}{\rho }.𝐯.$$
(22)
By making use of the identity
$$.𝐯=.\mathrm{\Delta }𝐯,$$
(23)
where $`\mathrm{\Delta }𝐯=𝐯𝐯_i`$, this can be put into a variety of forms. For example,
$$\frac{\mathrm{d}ϵ}{\mathrm{d}t}=\frac{.(P\mathrm{\Delta }𝐯)}{\rho }+\frac{\mathrm{\Delta }𝐯.P}{\rho }$$
(24)
gives
$$\frac{\mathrm{d}ϵ_i}{\mathrm{d}t}=\frac{2}{3}\underset{j}{}\frac{m_jϵ_j𝐯_{ij}._iW_{ij}}{\rho _i}=\frac{1}{m_i}\underset{j}{}𝐟_{ij}.𝐯_{ij},$$
(25)
where $`𝐯_{ij}=𝐯_j𝐯_i`$ –the second term has vanished because it contains $`𝐯_{ii}`$ which is identically zero. Using
$$\frac{\mathrm{d}ϵ}{\mathrm{d}t}=\frac{P.\mathrm{\Delta }𝐯}{\rho },$$
(26)
we obtain an alternative expression
$$\frac{\mathrm{d}ϵ_i}{\mathrm{d}t}=\frac{2}{3}\underset{j}{}\frac{m_jϵ_i𝐯_{ij}._iW_{ij}}{\rho _j}=\frac{1}{m_i}\underset{j}{}𝐟_{ji}.𝐯_{ji}.$$
(27)
Combining Equations 25 and 27, we find that
$$m_i\frac{\mathrm{d}ϵ_i}{\mathrm{d}t}=\frac{1}{2}\underset{j}{}(𝐯_{ij}.𝐟_{ij}+𝐯_{ji}.𝐟_{ji}),$$
(28)
where $`𝐟_{ij}`$ is the pairwise force given in Equation 21.
Any of Equations 25, 27, or 28 (or any similarly derived equation) may be used as an estimator of $`\mathrm{d}ϵ/\mathrm{d}t`$ (and hence, by rearrangement of Equation 22, for $`.𝐯`$). Indeed, if density is eliminated by substituting Equation 8 it is clear that the equations are identical, provided that our assumption that $`P_i`$ and $`P_j`$ are approximately equal is valid. Overall energy conservation is ensured as all three expressions give
$$\frac{\mathrm{d}E}{\mathrm{d}t}=\underset{i}{}(m_i\frac{\mathrm{d}ϵ_i}{\mathrm{d}t}+𝐯_i.𝐟_i)=0.$$
(29)
Under most conditions the performance of the three equations is indistinguishable, and standard tests of SPH cannot tell them apart. However, when large temperature variations are present the double-sided form can lead to a large scatter in particle entropy. We use Equation 25 here.
### 2.5 Artificial viscosity
In the presence of shocks, we require a mechanism to convert relative motion into heat. In SPH, this is achieved via an artificial pressure, $`Q`$, that is added to the usual one in regions of convergent flow. This has the effect of replacing $`P`$ by $`P+Q`$ in Equations 14 and 22. In the current method, we replace $`𝐟_{ij}`$ in Equations 20 and 27 with $`𝐟_{ij}+𝐠_{ij}`$, where
$$𝐠_{ij}=𝐟_{ij}_{ij}(\beta _{ij}+\alpha )$$
(30)
(c.f. Monaghan & Gingold 1983). Here the pairwise Mach number
$$_{ij}=\{\begin{array}{cc}0,\hfill & 𝐫_{ij}.𝐯_{ij}>0\hfill \\ \frac{h_i|𝐫_{ij}.𝐯_{ij}|}{c_{ij}(r_{ij}^2+0.01h_i^2)},\hfill & 𝐫_{ij}.𝐯_{ij}<0\hfill \end{array}$$
(31)
the particle separation, $`r_{ij}=|𝐫_{ij}|=|𝐫_j𝐫_i|`$, and the sound speed
$$c_{ij}=\frac{c_i+c_j}{2}$$
(32)
where
$$c=\sqrt{(10ϵ/9)}.$$
(33)
We adopt the values $`\alpha =1`$, $`\beta =2`$. The use of $`c_{ij}`$ in Equation 31 rather than $`c_i`$ helps to limit the degree to which cold particles shock against dense clumps.
In addition, a shear-correcting term (Balsara 1995) can be applied to limit the damping of shear flows. Following Steinmetz (1996) we use
$$_{ij}\stackrel{~}{}_{ij}=_{ij}k_i$$
(34)
where
$$k_i=\frac{|.𝐯_i|}{|.𝐯_i|+|\times 𝐯_i|+0.0001c_i/h_i}.$$
(35)
For compressional flows $`k`$ = 1 and $`_i`$ is unchanged, while in shearing flows $`k0`$ and the artificial viscosity vanishes. Only the cosmological galaxy formation test described in Section 3.7 uses this correction term.
### 2.6 Cooling
In the presence of artificial viscosity and cooling, the Equation of Energy Conservation becomes
$$\frac{\mathrm{d}ϵ}{\mathrm{d}t}=\frac{P+Q}{\rho }.𝐯\xi ,$$
(36)
where $`\xi `$ is the emissivity (the emission rate per unit volume). The cooling function is interpolated from Sutherland & Dopita (1993). To minimise problems of cooling during shock-heating (see Hutchings & Thomas 2000), we allow the gas to cool only at the end of a time-step after the artificial viscosity term has been applied. The cooling is assumed to occur at constant density (the time-step ensures that this condition is approximately satisfied), as described in Thomas & Couchman (1992).
## 3 TESTS
The tests discussed in the following sections were carried out using Hydra<sup>1</sup><sup>1</sup>1This code is in the public domain and can be downloaded from http://hydra.mcmaster.ca/. (Couchman, Thomas and Pearce, 1995), an AP<sup>3</sup>M+SPH code modified to use the SPH method detailed in this paper. The simulations were performed on single-processor Sun and Intel workstations.
### 3.1 Shock tube
A standard test of gas dynamical codes is the Sod shock (Sod 1978), which has been applied to SPH by many authors (e.g. Monaghan & Gingold 1983; Rasio & Shapiro 1991; T98). Analytic results are given by Hawley, Smarr & Wilson (1984) and Rasio & Shapiro (1991). This test is often carried out in one dimension, but this does not properly test particle interpenetration; we perform a three dimensional test, with the normal cubical simulation volume altered to match the geometry of a shock tube. Dimensions are $`6\times 6\times 120`$ in code units and all boundaries are periodic.
In a Sod shock two regions of gas with different densities are brought into contact, resulting in a shock wave propagating into the low-density gas and a rarefaction wave propagating into the high density gas. Between these two regions is a contact discontinuity, where the pressure is constant but the density jumps. Following T98 we use the initial conditions
$$\begin{array}{cccc}\rho _l=4\hfill & P_l=1\hfill & v_l=0.\hfill & \hfill \mathrm{for}x<0\\ \rho _r=1\hfill & P_r=0.1795\hfill & v_r=0.\hfill & \hfill \mathrm{for}x0\end{array}$$
(37)
giving a shock Mach number $`1.4`$. Both regions are allowed to evolve at constant temperature before being brought into contact, to allow the gas to relax to a physically realistic initial state. A total of 7343 equal mass particles are used.
Thirty timesteps are required to reach the state shown in Figure 1, by which time the shock front has moved around 13 code units to the right. Both the shock and the contact discontinuity are broadened over a range $`\mathrm{\Delta }x3h`$. The results are in good agreement with the analytic solutions, although in common with other implementations of SPH the gradient of the rarefaction wave is too shallow. Flow is reasonably smooth, and there is no post-shock ringing. Figure 2 compares the results from our code with those of T98 and CTP95. There is very little difference between the three implementations, although the CTP95 implementation produces a broader shock than the other two codes. This is due to the choice of artificial viscosity – CTP95 uses an artificial viscosity based on the divergence of the local velocity field, which is shown in T98 to be worse at shock capturing than the pairwise artificial viscosity used in T98 and our code. All three implementations model the rarefaction wave similarly, as viscous terms do not apply in expanding regions.
With $`1.4`$, this test represents a fairly weak shock. A strong adiabatic shock presents a more demanding test, as the pressure jump across the shock, and hence the Mach number, are infinite. In this case the jump conditions are
$$\begin{array}{cccc}\rho _l=1\hfill & P_l=0\hfill & v_l=1.333\hfill & \hfill \mathrm{for}x<0\\ \rho _r=4\hfill & P_r=1.333\hfill & v_r=0.333\hfill & \hfill \mathrm{for}x0\end{array}$$
(38)
although in our code we require particles to have a small minimum temperature to prevent numerical divergences in Equation 10, leading to $`P_l`$ being slightly greater than zero and the Mach number remaining large but finite.
The density profile across the shock is shown in Figure 3. The new code is clearly capable of handling such shocks, and although some post-shock oscillation is visible it is not noticeably different from the results obtained from the other codes. As the performance of our code is so close to that of T98 there is no evidence from shock-tube experiments that our assumption that pressure is smoothly varying has degraded the shock-capturing ability of the code.
### 3.2 Adiabatic collapse
One of the primary requirements for any hydrodynamics code that includes gravity is the ability to correctly follow the shock-heating of cold gas during gravitational collapse. A common test problem for SPH codes is the collapse of an initially isothermal sphere of gas (Evrard 1988; HK89; Steinmetz & Müller 1993; T98), with an initial density profile
$$\rho (r)=\frac{M(R)}{2\pi R^2}\frac{1}{r}.$$
(39)
To create this profile we translate particles radially from a uniform grid, which gives a lower sampling error than a random distribution of particles. Initial particle temperatures are set to the code minimum. A total of 8184 particles are used, with the gravitational softening set as 0.02R. Following Evrard (1988) results are presented in normalised units, with density, internal energy, velocity, pressure and time normalised by $`3M/4\pi R^3`$, $`GM/R`$, $`(GM/R)^{1/2}`$, $`\rho u`$ and $`(R^3/GM)^{1/2}`$ respectively.
The gas in the sphere initially has negligible thermal energy, and collapses due to the lack of pressure support. During the collapse a central bounce occurs, causing a shock wave to propagate outwards through the gas, and the sphere should eventually reach virial equilibrium, with the ratio of thermal and gravitational energy approaching a value of $`0.5`$. Figure 4 compares the evolution of the thermal, kinetic, total and gravitational potential energies during the collapse for the multiphase and standard methods. The performance of the two codes is very similar throughout the bounce and subsequent expansion, with only relatively minor differences apparent after time $`t1`$.
Profiles of the system at time $`t=0.8`$ are shown in Figure 5, at which time the shock is located at $`r/r^{}0.2`$. The temperature jump across the shock is well modelled by both codes, and the post-shock conditions are very similar. However, the density profile across the shock produced by the multiphase code is noticeably shallower. This is a result of the overestimation of the local pressure due to the steep pressure gradient in the shock, and is analogous to the smoothing of the density profile in the presence of steep density gradients that affects the standard implementation of SPH. This error is unimportant in the force calculations, which are dominated by the artificial viscosity during shocks, but may be significant when radiative cooling is implemented, as it leads to to excess cooling in the shock.
In such circumstances it may be preferable to use the standard SPH estimate of the density $`\overline{\rho }_i`$ for calculating the emissivity in regions where the local pressure gradient is steep. This switch can be implemented easily using either a pairwise condition based on the pressure of particles $`i`$ and $`j`$ or an SPH estimation of the local pressure gradient. Both methods effectively restrict the use of $`\overline{\rho }_i`$ to the vicinity of the shock. When this approach is used, the density profile across the shock in Figure 5 closely resembles that of the standard code, as would be expected, with the other profiles remaining unchanged.
### 3.3 Density estimation in a two-phase medium
In the standard formulation of SPH arbitrarily steep density gradients are smoothed over a distance representative of the local value of $`h`$. This is a cause for concern in simulations of cosmological structure formation, where cold dense clumps of gas – galaxies – form within diffuse halos of hotter gas and the density of halo gas can be overestimated by an order of magnitude or more in the steep density gradients around the largest galaxies.
The smoothing of the density is a result of two effects. Firstly, when Equation 5 is used, $`\overline{\rho }h^3`$ for a fixed number of neighbours, and the estimated density becomes closely tied to the choice of smoothing length. Secondly, most adaptive forms of SPH update the smoothing length each timestep to try and keep the number of neighbouring particles approximately constant (e.g. the method of HK89). Particles act as tracers of the underlying mass distribution, and a high-density region will therefore contain more particles than a region of lower density. When the HK89 algorithm is used, a particle close to a dense clump of gas will need to search little further than the edge of the dense region to find its required number of neighbours. If Equation 5 and the HK89 algorithm are used together then the estimate of the density of the particle becomes strongly dependent on its separation from the clump.
Simply estimating the density using Equation 10 is not enough to cure this problem. While the high density particles will be at low temperature in a region of pressure equilibrium, and will therefore contribute little to the estimate of the pressure, the size of the smoothing sphere, and hence weight given to neighbouring particles by the smoothing kernel $`W_{ij}`$, still depends on the distance to the clump. It is therefore necessary to weight the neighbour count using Equation 13, so that a low-density particle close to a dense clump will assign a very low weight to particles in the clump and will search for neighbours as if the clump was not present.
The benefits of our approach can be seen in Figure 6. Here we plot the change in smoothing length and the SPH estimate of the density as a hot gas particle moves past a cold, dense clump, grazing the surface of the clump at closest approach. The clump contains 420 particles, and has a central density 200 times that of the surrounding gas, with temperatures set so that the two phases are in pressure equilibrium. For the purposes of this test all inter-particle forces have been turned off, so that the particles move at constant velocity. Initially, the density of the hot particle is $`1.15`$ and the smoothing length $`h0.75`$, both in code units.
Panel (a) shows the results from the standard implementation of SPH, in which the density is estimated using Equation 5 and the smoothing length is updated using the algorithm of HK89. Once the smoothing kernel of the hot particle overlaps the cold clump it finds many more than the desired number of neighbours and the HK89 algorithm starts to decrease the smoothing length in response, with $`h`$ reaching a minimum of $`0.19`$ shortly after closest approach (this delay is due to the convergence parameter $`\alpha `$ in Equation 12 limiting the rate at which $`h`$ can change). There is a corresponding increase in the estimate of the density, which peaks slightly earlier at $`\rho 107`$, an overestimate of some two orders of magnitude. These values of $`h`$ and $`\rho `$ are typical of cold particles on the surface of the clump, indicating that despite the large difference in temperature between the hot and cold particles the standard implementation of SPH treats them identically.
Panel (b) shows the result of changing to the $`h`$-update method described in Section 3.3, with the density still estimated using Equation 5. The peak density has only dropped slightly to $`\rho 80`$, while the smoothing length remains fairly constant until shortly before the closest approach, when the density of the hot gas particle has increased to a point where particles in the clump become significantly weighted, leading to a decrease smoothing length.
Panel (c) shows the combination of the new estimate of the density, given by Equation 10, and the HK89 $`h`$-update algorithm. The peak density is reduced to $`\rho 10`$, which occurs when the smoothing length reaches it’s minimum and the cold particles in the clump are weighted most strongly.
Finally, panel (d) shows the results when the new estimate of the density is combined with the new smoothing length update algorithm. The smoothing length remains virtually constant throughout the transit, with $`h`$ decreasing to $`0.72`$ shortly after closest approach and $`\rho `$ only increasing from $`1.15`$ to $`1.61`$.
The new $`h`$-update algorithm does imply some additional computational expense when strong density gradients are present. In the example presented here the hot particle has in excess of 450 neighbours when the smoothing sphere overlaps the whole of the cold clump, far more than the 32 required by the HK89 algorithm. The overall computational expense will not be nearly as severe as this might suggest, however, as this idealised test follows a single particle next to a dense clump, where the weighting of neighbouring particles is the most extreme. For a distribution of particles, such as the cosmological structure formation discussed in Section 3.7, the overhead is typically on the order of $`10\%`$ per timestep.
### 3.4 Force estimation in a two-phase medium
Figure 7 shows the SPH estimate of the density profile across the boundary between two regions of different density, and demonstrates the sharpness with which density contrasts can be resolved by our method. To create the high density region we first evolve a cubical volume of gas at constant temperature, to ensure a relaxed particle distribution. We then extract a spherical region and compress it radially to achieve the desired density. This is then inserted back into the cubical simulation volume and the particles are allowed to move at constant radius from the centre of the box for a few time-steps to ensure a fully stable initial state. Here the sphere has a radius $`r=0.108`$ in code units, and is 100 times more dense than the surrounding gas, with the temperature again set so that the two regions are in pressure equilibrium. Away from the boundary both methods correctly estimate the density, as we would expect. However, as discussed in the previous section, the standard implementation of SPH clearly overestimates the density close to the dense region, whereas the new algorithm gives a density contrast that is much sharper.
In addition to the problems in estimating the density, unphysical forces can also occur in the presence of steep density gradients. In implementations of SPH which try to keep the number of neighbours constant, a particle near a density gradient will find many, if not all, of its neighbours in the region of higher density. The pressure gradient at the particle will therefore be highly asymmetric, leading, from Equation 14, to a force that acts to push the particle away from the high density region. The strength of this force will depend on the magnitude of the density contrast, saturating once the particle finds all its neighbours in the high density region, and a dense clump of gas will therefore rapidly empty the surrounding region of particles.
We can see the effect of this asymmetric pressure gradient in Figure 8, where we plot the radial velocities after 10 time-steps of particles in or near the dense clump.
The scatter in velocity due to thermal motion of the hot gas is around $`\pm 0.05`$ in both cases. The standard SPH code produces a large outward velocity in the hot gas, which evacuates a space between the two phases within a few tens of time-steps. This effect can be clearly seen in Figure 9, where the number density of particles drops to zero between $`0.108r0.13`$. While our method also produces excess velocities at the boundary, implying that the spurious pressure gradient has still not been completely eliminated, the magnitude of the outward velocity has been greatly reduced and the effect is far less systematic.
### 3.5 Drag
Drag resulting from gas dynamical forces is an important factor in simulations of cosmological structure formation, as incorrectly estimating the drag can bias both the distribution of matter in clusters and the size and number of objects formed. Frenk et al. (1996) have suggested that excessive drag can worsen the overmerging problem seen in N-body simulations, and TCP99 have shown that standard implementations of SPH can significantly overestimate the drag on a cold clump of gas moving through hot gas representative of the intracluster medium. This excess drag is caused by accretion of gas from the ICM on to a shell around the clump, where it is held by forces arising from the miscalculation of the pressure gradient around the clump (a discussion of this effect can be found in TCP99). This accretion leads to an increase in the effective radius of the clump, and hence the drag. This drag is most severe when the clump is moving subsonically, and at higher velocities the accretion of gas decreases, largely vanishing when the Mach number $`2`$.
In order to examine the drag introduced by our implementation of SPH, we consider the case of a clump of cold gas initially at rest in a stream of fast moving, diffuse gas (this is identical to the method of TCP99, except we work in the rest-frame of the cold clump). Assuming collisions are inelastic, the clump will be accelerated to the flow velocity $`v_0`$ at a rate
$$v=v_0\left[\frac{e^{kt}1}{e^{kt}+1}\right]$$
(40)
where $`k`$ is a constant
$$k=\frac{2\pi r^2\rho _0v_0}{M}$$
(41)
in which $`\rho _0`$ is the density of the diffuse gas, $`v_0`$ the flow velocity, r is the radius of the clump and M is the mass of the clump (all in code units). For our tests $`M=50`$, $`r=0.5`$, the velocity and density of the flow are set to unity and the Mach number of the flow is determined by setting the temperature of the gas as required. We use a Mach number $`0.5`$ here, which is well inside the regime in which TCP99 show that drag becomes excessive. Figure 10 compares the results from our new code with those obtained the codes of CTP95 and T98 (the code used by TCP99 is essentially the same as that used in T98).
The results of our code are clearly an improvement, being close to the prediction of Equation 40. This is due to the elimination of the accretion that occurs in the other codes, as the asymmetric pressure gradients causing accretion are greatly reduced by our method. The difference between CTP95 and T98 is mainly a result of the kernel smoothing used in T98, which is shown in TCP99 to further increase the effective cross section of a dense clump. At higher Mach numbers the differences between the codes are less pronounced as the accretion decreases once $`>1`$, although T98 continues to give the largest drag of the three codes.
### 3.6 Multiphase cooling flows
Clusters of galaxies contain large quantities of hot, X-ray emitting gas, often concentrated around the most massive galaxy in the cluster. In $`70\%90\%`$ of cases, gas in the centre of the cluster has a radiative cooling time less than the Hubble time, $`H_0^1`$ (Edge, Stewart & Fabian 1992; White, Jones & Forman 1997), and as this gas cools surrounding gas will move inwards to maintain pressure support, initiating a large-scale motion known as a cooling flow (see Fabian 1994 for a review). The mass deposition rate can be as high as $`10^3M_{}\mathrm{yr}^1`$ (Fabian et al. 1985), and is often observed to vary with radius roughly as $`\dot{M}(<r)r`$ within the radius $`r_{\mathrm{cool}}`$ in which the cooling time is less than a Hubble time (e.g. Thomas, Fabian and Nulsen 1987, hereafter TFN87). This is generally taken as evidence for a multiphase flow, in which thermal instability is causing the denser gas to cool out of the flow at larger radii.
Modelling a multiphase flow is impossible with the usual implementation of SPH as density is a locally averaged quantity, making the flow inherently single phase. In contrast, our method allows a wide range of densities, provided that a local pressure equilibrium exists. This is a reasonable assumption for particles with temperatures above $`10^6`$K, which will have cooling times much greater than the local sound crossing time (Nulsen, 1988). Particles below $`10^6`$K will cool to $`10^4`$K within a few time-steps, at which point they are removed from the flow.
To test the ability of our code to model a fully multiphase environment, we examine a simple constant-pressure, spherically-symmetric cooling flow. The distribution of phases in the flow is described by the fractional volume distribution $`f(\rho ,r)`$ introduced by Nulsen (1986, hereafter N86), where $`f\mathrm{d}\rho `$ is the fractional volume occupied by phases in the density range $`\rho `$ to $`\rho +\mathrm{d}\rho `$. N86 considered many analytic forms for $`f`$. Here, we take
$$f(\rho ,r)=\frac{(3\alpha )}{\rho _0}\left(\frac{\rho }{\rho _0}\right)^{(4\alpha )}\rho >\rho _0$$
(42)
with $`\rho _0`$ being a minimum density and $`\alpha `$ the temperature dependence of the cooling function $`\mathrm{\Lambda }T^\alpha `$. For the purposes of this test we replace the Sutherland & Dopita (1993) cooling function with a pure power law in which $`\alpha =0.5`$. Then Equation 42 can be integrated (N86) and gives a mass deposition profile
$$\dot{M}r^\eta $$
(43)
where
$$\eta =\frac{3(2\alpha )}{(52\alpha )}=\frac{9}{8}.$$
(44)
Particles are placed randomly within a cubical simulation of volume $`(200\mathrm{k}\mathrm{p}\mathrm{c})^3`$, and are then allowed to evolve at constant temperature until spurious fluctuations arising from the initial particle distribution have died away. Particles are then ordered in terms of their distance from the centre of the box, and translated radially so as to match the mass profile given by
$$M\rho _0(r)r^3$$
(45)
where the density profile
$$\rho _0(r)r^{\frac{3}{(52\alpha )}}r^{3/4}$$
(46)
is derived in TFN87. Particles translated outside the bounds of the box are discarded. Particle densities are drawn at random from the volume fraction distribution given by Equation 42, with particle temperatures set so as to maintain constant pressure. Here, the outer temperature is $`T5\times 10^7K`$, the inner temperature $`T10^6K`$, the central density is $`0.1\mathrm{cm}^3`$ and the average outer density is $`5\times 10^3\mathrm{cm}^3`$. A total of 20252 particles are used.
Figure 11 shows the mass deposition profile $`\dot{M}(<r)`$, produced by the multiphase method after 500 timesteps. Roughly 11% of the gas has cooled from the flow. The line marks a least-squares fit to the data, with slope $`\eta =1.26\pm 0.01`$. The points within $`15`$kpc of the centre of the flow have been excluded from this fit, as there are too few particles here for the mass deposition profile to be well sampled, as have particles with $`r>100\mathrm{k}\mathrm{p}\mathrm{c}`$ which lie in the corners of the cubical simulation volume. The slope is slightly steeper than the theoretical index of $`\eta =9/8`$. Figure 12 shows the mass distribution function $`\rho f\mathrm{d}\rho `$ for particles at radii of between 50 and 60 kpc. In theory we would expect the distribution of phases given by Equation 42 to remain unchanged with time, giving $`\rho f\mathrm{d}\rho \rho ^{2.5}`$. The dashed line represents a least-squares fit to the data, with gradient $`2.75\pm 0.06`$. This is again steeper than we would expect. This is probably due to the phases not comoving, an assumption made in N86. There is no condition in our code to enforce this, and Equation 14 implies that high-density particles will receive a smaller pressure force than low-density ones. The artificial viscosity will limit the degree to which particles can interpenetrate, ensuring that the flow is largely comoving, but it may not be sufficient to ensure that no slippage occurs. Applying a large bulk viscosity term in Equation 30 forces the particles to comove, and flattens both slopes towards their theoretical values. However, this is not suitable for general application as it degrades the shock capturing ability of the method.
The mass deposition profile produced by the standard implementation of SPH is shown in Figure 13. 500 timesteps have passed and this time about 5% of the gas has cooled. In contrast with the multiphase results, this implementation produces a mass deposition profile that is centrally concentrated, as would be expected from a single-phase flow. This is supported by the mass distribution function, which is plotted in Figure 14. The gas has clearly evolved back to a single-phase state, with all particles having a density close to the mean, and no trace of the original density distribution remains. In addition, the gradient of the density profile has steepened from an initial value of $`\rho r^{0.75}`$ (as given by Equation 46) to $`\rho r^{1.15\pm 0.03}`$, close to the theoretical single-phase gradient $`\rho r^{1.2}`$ (Thomas, 1988).
When compared to the standard implementation of SPH, it is clear that our method performs well. Both the mass deposition profile and the distribution of densities in the flow are close to the theoretical values, whereas the standard implementation fails to reproduce the properties of the flow correctly.
### 3.7 Cosmological galaxy formation
Our final test involves a simulation of the formation of a cluster of galaxies. We use the initial conditions from the fiducial simulation of Kay et al. (2000), who used it to test in detail the effect of varying a number of numerical and physical parameters. The simulation uses the standard cold dark matter (SCDM) cosmology, with $`\mathrm{\Omega }=1`$, $`\mathrm{\Lambda }=0`$ and $`h=0.5`$<sup>2</sup><sup>2</sup>2$`H_0=100h\mathrm{kms}^1\mathrm{Mpc}^1`$. The baryon fraction was set from primordial nucleosynthesis constraints, $`\mathrm{\Omega }_bh^2=0.015`$ (Copi, Schramm & Turner 1995), and an unevolving gas metallicity of $`0.5Z_{}`$ was used. The initial fluctuation amplitude was set so that the model produces the same number density of rich clusters as is observed today, with $`\sigma _8`$, the present-day linear rms fluctuation on a scale of 8$`h^1\mathrm{Mpc}`$, set to 0.6 (Eke, Cole & Frenk 1996; Vianna & Liddle 1996). Until $`z1`$ we adopt a comoving $`\beta `$-spline gravitational softening equivalent to a Plummer softening of $`20`$$`h^1\mathrm{kpc}`$, after which it is switched to a fixed physical softening of $`10`$ $`h^1\mathrm{kpc}`$. The minimum SPH resolution is set to match the constraints imposed by the gravitational softening. $`32^3`$ particles of both dark matter and gas were used, giving a mass per particle of $`8\times 10^9h^1M_{}`$ for the dark matter and $`5\times 10^8h^1M_{}`$ for the gas. The simulation was started at $`z24`$ and was evolved to the present day.
For the purposes of this simulation, Equation 36 incorporates a source term $``$ to reflect the heating of gas due to a photoionising background. We assume that the background has a standard power-law form
$$J(\nu )=J_{21}(z)\times 10^{21}\left(\frac{\nu }{\nu _{H_I}}\right)\mathrm{ergs}\mathrm{s}^1\mathrm{cm}^2\mathrm{sr}^1\mathrm{Hz}^1,$$
(47)
where
$$J_{21}(z)=\frac{J_{21}^0}{1+(5/(1+z)^4)}$$
(48)
is the flux at the H<sub>I</sub> Lyman limit (Navarro & Steinmetz, 1997); we take $`J_{21}^0=1`$ and $`\alpha =1`$ here. Photoheating is implemented following Theuns et al. (1997). The ultraviolet background has the effect of imposing a minimum temperature on the gas, rapidly heating it to $`10^4K`$ and ensuring that pressure gradients in the gas remain shallow. If the effects of photoionisation are not included, then problems can occur when particles which have been in free expansion since the start of the simulation, cooling to very low temperatures with $`T(1+z)^2`$, encounter the accretion shock at the outer edges of halos. This shock is generally poorly-resolved, and neighbouring particles that have passed through the shock can make an overwhelming contribution to the estimate of the pressure of the cold particle, which can lead to the density being overestimated, potentially by several orders of magnitude. Incorporating photoionisation serves as a simple way to limit this effect; an alternative is to use the standard SPH estimate of the density for calculating the emissivity in regions where the pressure gradient is steep, as discussed in Section 3.2.
Figures 15 and 16 illustrate the temperature-density distribution of baryonic matter at $`z=0`$ produced by the multiphase and standard codes. We have divided the gas into two phases; a galaxy phase containing gas which is overdense by at least a factor of $`500`$ and has a temperature of $`10^3`$K$`T10^5`$K and a hot phase which includes all gas with a temperature of $`T10^5`$K and the gas with $`T10^3`$K which is not sufficiently overdense to be considered part of the galaxy phase. Almost all the particles in the galaxy phase lie along a line with $`T12\times 10^4`$K, which marks the point at which the radiative cooling and photoheating rates balance. Most of these particles are in the form of dense clumps with a size similar to the gravitational softening length – we refer to these clumps as galaxies. The properties of these galaxies are calculating by first extracting all particles within the galaxy phase from the simulation volume and then running a ’friends-of-friends’ group finder (Davis et al. 1985). The most significant difference between the two simulations is in the mass of the largest galaxy, which is nearly $`50`$% more massive in the single-phase simulation. Kay et al. (2000) find that this is a result of excessive cooling, as decoupling the hot and cold phases (Pearce et al. 1999) reduces the final mass of the galaxy to a more reasonable value.
The hot phase consists of particles that have been shock heated during gravitational collapse, with the bulk kinetic energy of the gas being converted to heat. Figure 15 shows a significant quantity of hot ($`T>10^7`$K), high-density ($`\rho /\rho >10^4`$) gas which is not present in Figure 16 but is seen in the single-phase simulations when radiative cooling is not allowed. This gas is located near the centre of large dark matter halos, close to the central galaxy, and overestimation of the already high gas density by the standard method leads to the gas rapidly cooling and being accreted by the galaxy. In contrast, the multiphase method correctly estimates the density of the gas, resulting in a slower cooling rate and the gas remaining at high temperatures throughout the halo. This effect can be seen in more detail in Figures 17 and 18, which show the radial temperature profile of the gas around the most massive galaxy found in the simulation volume. The size of the galaxy is on the order of the softening length, $`10`$$`h^1\mathrm{kpc}`$, and is not shown here. The results from the standard implementation clearly show two of the problems inherent in the method. Firstly, particle temperatures drop sharply within $`r<100`$kpc - this is clear evidence for overcooling. No such effect is visible in Figure 17 and the lack of overcooling in our method is probably the principle reason for the difference in mass of the largest objects in the two simulations. Secondly, there is a complete absence of particles at radii $`r<50`$kpc, whereas particles are found all the way in to $`r10`$kpc in the multiphase simulations. This is an example of the effect examined in Section 3.4, with particles close to the central galaxy being forced away by an artificially asymmetric pressure gradient in the single-phase method.
## 4 DISCUSSION AND CONCLUSIONS
We have presented a multiphase implementation of Smoothed-Particle Hydrodynamics (SPH), along with a number of tests to compare the performance of our method with standard implementations of SPH. The usual SPH formalism assumes that density is a smooth quantity, varying negligibly on distances on the order of a typical smoothing length. This is clearly not true in many situations in which SPH is applied, such as simulations of galaxy formation, in which large density contrasts are present. However, the pressure of the gas is expected to be a much smoother quantity because in almost all situations the sound-crossing time is shorter than the flow time across the smoothing sphere. We therefore summate the local pressure at each particle, and calculate the density from the equation of state.
One situation in which our assumption will definitely not be true is in the presence of shocks, in which density variations are generally smaller than the variations in pressure. Sections 3.1 and 3.2 demonstrate that our method handles shocks acceptably, although the estimate of the density in the shock can be inaccurate when strong shocks are resolved poorly. This does not alter the force calculation, which is dominated by the artificial viscosity under such conditions, but is potentially significant when radiative cooling is implemented, and it may be may be preferable to use the standard estimate of the density for calculating the emissivity in regions where the local pressure gradient is steep. In Section 3.3 we examine the degree to which steep density gradients can be resolved by the two methods. Standard implementations of SPH are shown to severely overestimate the density of particles close to a region of high density, while these particles have their density correctly estimated by our method. In addition, we show in Section 3.4 that unphysical forces can occur in the presence of steep density gradients, and while such forces are not completely eliminated by our method, they are greatly reduced. Section 3.5 examines the drag introduced by our implementation of SPH. We find that at low Mach numbers drag is greatly reduced compared to the SPH implementations of CTP95 and T98, while at higher velocities the results are in broad agreement with the findings of TCP99, depending mainly on the choice of smoothing kernel symmetrization.
In Section 3.6 we examine a simple spherically-symmetric, constant pressure cooling flow. The fractional volume distribution $`f\mathrm{d}\rho `$ of phases within the flow is taken to be a pure power-law, which is shown by N86 to remain unchanged with time and deposit mass as $`\dot{M}(<r)r^{9/8}`$. The standard implementation of SPH is shown to be unable to reproduce the expected behaviour of this cooling flow, with conditions rapidly returning to a single-phase state. Our code performs well, although the gradient of both the mass distribution $`\rho f\mathrm{d}\rho `$ and the mass deposition profile $`\dot{M}`$ are too steep. This is probably due to phases not comoving in the flow, which is an assumption made in N86. Application of a large bulk viscosity to force particles to comove appears to reduce the problem, although this is not a suitable for application generally as it results in the shock capturing ability of the code being significantly degraded.
In Section 3.7 we examine the formation of a cluster of galaxies. The idealised problems examined in Sections 3.3 and 3.4 are shown to have real analogues in the simulation using the standard SPH implementation, with overcooling resulting from overestimating the density of halo gas being clearly visible, and halo gas being forced away from the central galaxy. No such effects are visible in the simulation using our method. The lack of overcooling is also apparent in the masses of the galaxies formed, with the largest galaxy being nearly 50% more massive in the single-phase simulation, in agreement with the findings of Pearce et al. (1999), who found that decoupling the galaxy from the hot halo gas produced a similar effect.
Our method is an alternative to the standard formulation of SPH. In simulations without large density contrasts, the two give very similar results. However it represents a significant improvement over the standard implementation of SPH when the gas component cannot be assumed to be a single-phase fluid, such as galaxy formation and cluster formation. In addition, fully multiphase fluid flow can be modelled, allowing SPH to be applied to simulations of cooling flows and the intracluster medium.
## Acknowledgments
BWR acknowledges the support of a PPARC postgraduate studentship. PAT is a PPARC Lecturer Fellow. The authors would like to thank Rob Thacker and Scott Kay for supplying initial conditions used in testing, and Andy Fabian for useful suggestions. We would also like to thank the referee, James Wadsley, for helpful suggestions that have greatly improved this paper. Parts of this work was conducted on the SGI Origin platform using COSMOS Consortium facilities, funded by HEFCE, PPARC and SGI.
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# 1 Introduction
## 1 Introduction
An invariant regularization scheme is necessary for the treatment of the ultraviolet divergence in quantum gauge theories. The dimensional regularization is known as the most powerful and popular method, but it is not available for the theory preserving the symmetry which depends on the space-time dimension, like chiral gauge theory or topological field theory. In such a case, the hybrid regularization based on the higher covariant derivative (HCD) method is expected to be useful.
The HCD method is a partial regularization itself because the higher derivative terms are introduced in a covariant way; HCD terms render the propagators less divergent but the vertices more divergent. It is easily shown by the calculation of the superficial degree of divergence that some diagrams at one-loop level are left unregularized. We have to introduce an additional regularization scheme to regularize the divergence from these diagrams.
The Pauli-Villars (PV) regularization is suitable for the additional regularization from the standpoint of the invariant method, because the PV regulators are constructed in a chiral invariantly or parity invariantly by an infinite number of them. Applying such PV fields to the hybrid regularization, the method is expected to be an invariant regularization scheme that is available for the theory containing higher symmetries like the supersymmetric gauge theories. It has not been verified, however, whether such special PV fields actually regularize the intrinsically divergent theory in the framework of the hybrid regularization: at least the Yang-Mills (YM) theory must be regularized properly.
In recent years, on the other hand, it was pointed out that the original Slavnov’s hybrid regularization scheme does not give the correct value of the coefficient of the renormalization group (RG) $`\beta `$-function when the YM theory is regularized by this scheme . Nevertheless, this problem was overcome by the minor modification of the scheme , and it was confirmed that the modification is proper at the one-loop level by an explicit calculation .
Since only the logarithmic divergence plays an important role in the calculation of the RG functions, the other divergence, the quadratic divergence in the case of four dimensions, has not been seriously considered. If the regularization works properly, the quadratic divergence ought to be canceled out. But in this scheme, the cancellation is not trivial because the HCD term contributes to the quadratic divergence and then increases the complexity of the divergence . So it is worthy to show the cancellation of the quadratic divergence for the consistency of the hybrid regularization scheme.
In this note, we confirm that the quadratic divergence is actually canceled out in the YM theory with the hybrid regularization of the HCD and the PV method. By an explicit calculation of the vacuum polarization tensor, it is shown that the higher derivative terms for the ghost fields are necessary for the complete regularization of this method. Using an infinitely many PV fields as the additional regulator, we also check the consistency of such PV fields in the quadratic divergence when they are used in the hybrid regularization.
## 2 The regularization method
We consider the SU($`N`$) Yang-Mills theory in four dimensional Euclidean space-time. The action is given by
$$S=S_{\mathrm{YM}}+S_{\mathrm{GF}},$$
(1)
where
$`S_{\mathrm{YM}}`$ $`={\displaystyle \frac{1}{4}}{\displaystyle }\mathrm{d}^4xF_{\mu \nu }^aF^{\mu \nu }{}_{}{}^{a},`$ (2)
$`S_{\mathrm{GF}}`$ $`={\displaystyle \mathrm{d}^4x\left[\frac{\xi _0}{2}b^ab^ab^a(^\mu A_\mu ^a)+\overline{c}^a(_\mu D^\mu c)^a\right]},`$ (3)
with the field strength $`F_{\mu \nu }^a=_\mu A_\nu ^a_\nu A_\mu ^a+gf^{abc}A_\mu ^bA_\nu ^c`$ and the covariant derivative $`D_\mu ^{ac}=\delta ^{ac}_\mu +gf^{abc}A_\mu ^b.`$ Here $`A_\mu ^a`$, $`c^a`$, $`\overline{c}^a`$ and $`b^a`$ denote the gauge field, ghost, anti-ghost and auxiliary field respectively, $`\xi _0`$ is the gauge-fixing parameter and $`f^{abc}`$ is the structure constant of the gauge group.
The hybrid regularization consists of the following two steps: first we introduce HCD terms and next PV fields. The HCD terms improve the behavior of propagators at large momentum, rendering the theory less divergent at the cost of the emergence of new vertices, and the theory is reduced to superrenormalizable, i.e. there are just a finite number of divergent loops. As see later, all the diagrams except one-, two-, three- and four-point functions at one-loop level are convergent with a suitable choice of the HCD action. We deal with the remaining divergence by a PV type of regularization.
### 2.1 Introduction of HCD terms
We first regularize the action by an addition of the HCD term;
$$S_\mathrm{\Lambda }=S_{\mathrm{YM}}+S_{\mathrm{HCD}}+S_{\mathrm{GF}}^H.$$
(4)
$`S_{\mathrm{HCD}}`$ is the HCD action and $`S_{\mathrm{GF}}^H`$ is the modified gauge-fixing action when we use the HCD method. The explicit forms are
$`S_{\mathrm{HCD}}`$ $`={\displaystyle \frac{1}{4\mathrm{\Lambda }^4}}{\displaystyle \mathrm{d}^4x(D^2F_{\mu \nu })^a(D^2F^{\mu \nu })^a},`$ (5)
$`S_{\mathrm{GF}}^H`$ $`={\displaystyle \mathrm{d}^4x\left[\frac{\xi _0}{2}b^ab^ab^aH(^2/\mathrm{\Lambda }^2)(^\mu A_\mu ^a)+\overline{c}^aH(^2/\mathrm{\Lambda }^2)(_\mu D^\mu c)^a\right]},`$ (6)
where $`\mathrm{\Lambda }`$ is a cutoff parameter which has mass dimension of one, and $`H(^2/\mathrm{\Lambda }^2)`$ is a dimensionless function and must be a polynomial of $`^2/\mathrm{\Lambda }^2`$ to ensure the locality of the gauge field. Its explicit form is determined by the behavior of the gauge propagator which is obtained from (4) as follows:
$$\frac{\mathrm{\Lambda }^4}{p^4(p^4+\mathrm{\Lambda }^4)}(p^2\delta _{\mu \nu }p_\mu p_\nu )+\frac{\xi _0}{p^4H^2(p^2/\mathrm{\Lambda }^2)}p_\mu p_\nu .$$
(7)
The first term has the momentum degree of $`6`$, so the second term must be the same degree or less to ensure the convergence of the diagrams; $`H^2`$ behaves $`p^4`$ or higher at large $`p`$. While the familiar form of the propagator must be recovered in the limit of $`\mathrm{\Lambda }\mathrm{}`$; $`H^2`$ converges to unity at large $`\mathrm{\Lambda }`$. On these conditions, the simplest choice of the $`H^2`$ in momentum space is
$$H^2\left(\frac{p^2}{\mathrm{\Lambda }^2}\right)=1+\frac{p^4}{\mathrm{\Lambda }^4}.$$
(8)
Since the first and the second term of (7) has the same denominator in this choice of $`H^2`$, when we work in the Feynman gauge ($`\xi _0=1`$), the gauge propagator is reduced to
$$\frac{\mathrm{\Lambda }^4}{p^2(p^4+\mathrm{\Lambda }^4)}\delta _{\mu \nu }.$$
(9)
We take the Feynman gauge in the explicit calculation of diagrams in the following section.
The HCD action is generally introduced in the form of $`\frac{1}{4\mathrm{\Lambda }^{2n}}\mathrm{d}^4x(D^nF_{\mu \nu })^2`$. In such a case, the superficial degree of divergence is written
$$\omega =42n(L1)E_A$$
(10)
where $`L`$ and $`E_A`$ are the number of loops and the number of the external lines of the gauge field. For all the diagrams higher than two-loops ($`L2`$), $`n2`$ always gives negative $`\omega `$. This means that we may remove the higher loops by a suitable choice of $`n`$. For one-loop ($`L=1`$), $`\omega `$ is not always negative for any $`n`$. As we take the higher $`n`$, though the propagator becomes more convergent at large momentum, the vertices are more divergent and increase their complexities. Then the calculation of the diagram is more complex even though at one-loop level. The most economical choice is $`n=2`$ which leads the HCD action (5).
### 2.2 Introduction of Pauli-Villars fields
So far, all the diagrams except one-, two-, three- and four-point functions at one-loop level are convergent by the HCD action. We deal with the remaining divergence by a PV type of regularization.
Consider the following generating functional:
$$\begin{array}{c}Z[J,\chi ,\eta ,\overline{\eta }]=𝒟A_\mu 𝒟b𝒟\overline{c}𝒟c\mathrm{exp}[S_\mathrm{\Lambda }S_J]\hfill \\ \hfill \underset{j=1}{\overset{\mathrm{}}{}}\left[det{}_{}{}^{\frac{\alpha _j}{2}}𝐀_{j}^{}\right]\left[det{}_{}{}^{\frac{\alpha _j}{2}}𝐀_{j}^{}\right]\underset{i=1}{\overset{\mathrm{}}{}}\left[det{}_{}{}^{\gamma _i}𝐂_{i}^{}\right]\left[det{}_{}{}^{\gamma _i}𝐂_{i}^{}\right],\end{array}$$
(11)
where $`S_J`$ is a source term consisting from $`J`$, $`\chi `$, $`\eta `$ and $`\overline{\eta }`$ which is the source of $`A_\mu `$, $`b`$, $`\overline{c}`$ and $`c`$ respectively. $`\left[det{}_{}{}^{\frac{\alpha _j}{2}}𝐀_{j}^{}\right]`$ and $`\left[det{}_{}{}^{\gamma _i}𝐂_{i}^{}\right]`$ are PV determinants for the gauge and ghost field respectively, defined by
$`\left[det{}_{}{}^{\frac{\alpha _j}{2}}𝐀_{j}^{}\right]`$ $`={\displaystyle 𝒟A_j{}_{\mu }{}^{}𝒟b_j\mathrm{exp}[S_{M_j}S_{b_j}]},`$ (12)
$`\left[det{}_{}{}^{\gamma _i}𝐂_{i}^{}\right]`$ $`={\displaystyle 𝒟\overline{c}_i𝒟c_i\mathrm{exp}[S_{m_i}]},`$ (13)
where $`\{\alpha _j\}`$ and $`\{\gamma _i\}`$ are real parameters to be fixed for each indices and the explicit forms of the functions in the exponents are
$`S_{M_j}`$ $`={\displaystyle \frac{1}{2}}{\displaystyle \mathrm{d}^4x\mathrm{d}^4y}`$
$`A_j{}_{\mu }{}^{a}(x)[{\displaystyle \frac{\delta ^2S_\mathrm{\Lambda }}{\delta A_\mu ^a(x)\delta A_\nu ^b(y)}}M_j^2\delta ^{ab}g^{\mu \nu }\delta (xy)]A_j{}_{\nu }{}^{b}(y),`$ (14)
$`S_{b_j}`$ $`={\displaystyle }\mathrm{d}^4x[{\displaystyle \frac{\xi _j}{2}}b_j^ab_j^ab_j^a\stackrel{~}{H}(D^\mu A_j{}_{\mu }{}^{a})],`$ (15)
$`S_{m_i}`$ $`={\displaystyle \mathrm{d}^4x\left[\overline{c}_i^a\stackrel{~}{H}(D_\mu D^\mu c_i)^am_i^2\overline{c}_i^ac_i^a\right]}.`$ (16)
The field $`A_j_\mu ^a`$ is a PV field for the gauge of mass $`M_j`$, $`b_j^a`$ an auxiliary field for $`A_j_\mu ^a`$, $`\overline{c}_i^a`$ and $`c_i^a`$ PV fields for the ghost and anti-ghost of mass $`m_i`$. We introduce a gauge-fixing parameter $`\xi _j`$ for the correct regularization of the theory . $`\stackrel{~}{H}`$ is the HCD term for the ‘gauge-fixing function’ for the PV fields, which has the form
$$\stackrel{~}{H}=\left(1+\frac{D^4}{\mathrm{\Lambda }^4}\right)^{\frac{1}{2}}.$$
(17)
These PV fields have the same form as the ones that we used in the Chern-Simons gauge theory . The idea is to regularize the theory with the pairs of the two types of the PV fields $`A_j`$ and $`A_j`$. In the Chern-Simons gauge theory to construct the parity invariant regulator, the two fields are related by the parity transformation and represented by the slightly different actions, but in this case they have the same action except the sign of the index.
The reason why we have to introduce an infinite number of the PV fields comes from the idea ‘to regularize with the pair’. Imagine when the gauge field is regularized by the PV pairs which we prepared above. Since the number of the gauge field is one, introducing one pair corresponds to subtracting double the divergence. Then to remedy the over subtraction we introduce another pair of opposite statistics which means adding double the divergence. To remedy the over addition we have to introduce the third pair. Such steps correspond to introducing fermionic PV fields $`(\alpha _j=1)`$ and bosonic PV fields $`(\alpha _j=+1)`$ alternately. We repeat such steps alternately until the divergence is removed. Namely, we cannot regulate the theory by a finite number of PV pairs, but we need an infinite number. Then we take the following as PV conditions:
$`M_j`$ $`=M|j|,`$ $`\alpha _j`$ $`=(1)^j.`$ (18)
In the same way, we take the PV conditions for ghost and anti-ghost such as $`m_i=m|i|`$ and $`\gamma _i=(1)^i`$.
The generating functional (11) is invariant under the BRST transformations in the reference , and this regularization manifestly preserves BRST invariance.
### 2.3 Feynman rules
The regularized action $`S_\mathrm{\Lambda }`$ is decomposed into the kinetic part $`K`$ and the vertex part $`V`$ as follows :
$$\mathrm{d}^4x\mathrm{\Psi }(x)(K+V+M^2)\mathrm{\Phi }(x),$$
(19)
where $`\mathrm{\Psi }(x)`$ and $`\mathrm{\Phi }(x)`$ denotes an arbitrary field and $`M`$ its mass term. Since $`K`$ and $`V`$ consist of the $`\mathrm{\Lambda }`$-free part from the YM term (we denote with suffix ‘0’) and the $`\mathrm{\Lambda }`$-dependent part from HCD term (with suffix ‘$`\mathrm{\Lambda }`$’) we formally decompose
$`K`$ $`=K_0+{\displaystyle \frac{1}{\mathrm{\Lambda }^4}}K_\mathrm{\Lambda },`$ $`V`$ $`=V_0+{\displaystyle \frac{1}{\mathrm{\Lambda }^4}}V_\mathrm{\Lambda }.`$ (20)
Then the propagators are written in the form
$$\frac{1}{K+M^2}=\frac{1}{K_0+M^2}\left(1\frac{K_\mathrm{\Lambda }}{K_0+M^2}\mathrm{\Lambda }^4+O(\mathrm{\Lambda }^8)\right).$$
(21)
Using this decomposition, the Feynman rules are written by the order of $`\mathrm{\Lambda }`$ and we calculate the quantum corrections order by order of $`\mathrm{\Lambda }`$.
## 3 One-Loop Contributions
Now we calculate the one-loop vacuum polarization tensor order by order in $`\mathrm{\Lambda }`$ up to $`\mathrm{\Lambda }^4`$. There are eleven diagrams in $`\mathrm{\Lambda }^0`$ order and twelve diagrams in $`\mathrm{\Lambda }^4`$ order. Each diagram has the quadratic divergence which must cancel in totally. We consider this divergence calculating each diagram under the Feynman gauge $`\xi _j=1`$. The calculation is carried out under the same rules that we take in references whose summaries are the following:
1. Take the same assignment for the internal momentum among the graphically same diagrams.
2. Take the infinite sum of the graphically same diagrams under the PV condition (18).
3. If there is no massless term (whose index $`j`$ or $`i`$ is zero) for the infinite sum, add the lacking terms and subtract the same ones to balance.
All the diagrams are classified into three groups by the kind of the internal line in the diagram and calculation is carried out in each group.
First we consider the diagrams that contains only $`A_j_\mu `$ fields in the internal lines. The index $`j`$ runs from $`\mathrm{}`$ to $`\mathrm{}`$ except zero for these diagrams, if we take the infinite sum from $`\mathrm{}`$ to $`\mathrm{}`$ we have to add the diagrams which contains ‘$`A_0_\mu `$’ field in the internal lines. Fortunately, the diagrams in which the gauge field runs have the same structure of them and we take the infinite sum from $`\mathrm{}`$ to $`\mathrm{}`$ without any extra terms. Then the total of the quadratic divergence from these diagrams is calculated
$$\frac{g^2c_v\delta ^{ab}}{8\pi ^2}\left(\frac{M^2}{40}C_2\frac{9M^6}{154\mathrm{\Lambda }^4}C_4\right)\delta _{\mu \nu }.$$
(22)
Where $`C_2`$ and $`C_4`$ are the constants arising from the infinite sum of the index $`j`$ .
For the diagrams that contains the $`b_j`$ fields, there is no diagrams from the field $`b`$ to compensate the diagrams consisting of $`b_0`$ because any vertex of the form $`A_\mu A_\nu b`$ does not exist in the theory. So we have to add the $`b_0`$ diagrams to take the infinite sum and subtract the same after the summations to balance out. The total contributions are calculated
$$\begin{array}{c}\frac{g^2c_v\delta ^{ab}}{8\pi ^2}\left(\frac{M^2}{40}C_2\frac{9M^6}{154\mathrm{\Lambda }^4}C_4\right)\delta _{\mu \nu }\hfill \\ \hfill +g^2c_v\delta ^{ab}\frac{\mathrm{d}^4k}{(2\pi )^4}\frac{1}{k^2(kp)^2}\left[2k^2\delta _{\mu \nu }3k_\mu k_\nu +\frac{1}{\mathrm{\Lambda }^4}\left(2k^6\delta _{\mu \nu }4k^4k_\mu k_\nu \right)\right].\end{array}$$
(23)
The first line comes from the infinite sum with index $`j`$ and the second line is the counter terms which we introduced to take the sum.
The situation does not alter in the diagrams containing $`\overline{c}_i`$ and $`c_i`$ in the internal lines. Since there are some differences in the vertices between ghost fields and PV fields for ghosts, $`\overline{c}`$ and $`c`$ do not play the role of $`\overline{c}_0`$ and $`c_0`$. Then we have to add some massless terms for the infinite sums and subtract the same ones in the same way as the diagrams with $`b_j`$. We calculate the contributions from these diagrams as follows:
$$g^2c_v\delta ^{ab}\frac{\mathrm{d}^4k}{(2\pi )^4}\frac{1}{k^2(kp)^2}\left[2k^2\delta _{\mu \nu }3k_\mu k_\nu +\frac{1}{\mathrm{\Lambda }^4}\left(2k^6\delta _{\mu \nu }4k^4k_\mu k_\nu \right)\right].$$
(24)
The each contribution has the quadratic divergence proportional to the constant $`C_2`$ and $`C_4`$ after the infinite sum with indices, but these divergence cancel out in total within this group. Then only the quadratic divergence from the counter terms remain as in (24).
It is easy to see from (22), (23) and (24) that the quadratic divergence of the vacuum polarization tensor disappears from the theory.
Here we notice that the function $`f(^2/\mathrm{\Lambda }^2)`$ in the reference does not cancels the quadratic divergence completely. In that reference, since the modified gauge-fixing action is inserted through the form $`\frac{\xi _0}{2f}b^2b^\mu A_\mu `$ instead of (6), so the naively extended gauge-fixing action for the PV field is written by $`\frac{\xi _j}{2f}b_j^2b_j^\mu A_j_\mu `$. This action, however, breaks the BRST invariance because $`f`$ contains usual derivative $`_\mu `$.This symmetry breaking effects to the cancellation of the quadratic divergence. In such a case, all the Feynman rules related to the auxiliary, ghost and their PV fields are modified in $`\mathrm{\Lambda }^4`$ order, and then the quadratic divergence corresponds to (23) and (24) are calculated as follows:
$$\begin{array}{c}\frac{g^2c_v\delta ^{ab}}{8\pi ^2}\left(\frac{M^2}{40}C_2\frac{6M^6}{154\mathrm{\Lambda }^4}C_4\right)\delta _{\mu \nu }\hfill \\ \hfill +g^2c_v\delta ^{ab}\frac{\mathrm{d}^4k}{(2\pi )^4}\frac{1}{k^2(kp)^2}\left[2k^2\delta _{\mu \nu }3k_\mu k_\nu \right],\end{array}$$
(24)
$$g^2c_v\delta ^{ab}\frac{\mathrm{d}^4k}{(2\pi )^4}\frac{1}{k^2(kp)^2}\left[2k^2\delta _{\mu \nu }3k_\mu k_\nu \right].$$
(24)
This result shows that the complete cancellation at $`\mathrm{\Lambda }^4`$ does not occur though the contribution at $`\mathrm{\Lambda }^0`$ cancels. We may choose the higher covariant derivative function $`\stackrel{~}{f}`$ instead of $`f`$ to avoid this difficulty as we extended $`H`$ to $`\stackrel{~}{H}`$. If we substitute $`\stackrel{~}{f}=1+\frac{D^4}{\mathrm{\Lambda }^4}`$ in place of $`f=1+\frac{^4}{\mathrm{\Lambda }^4}`$, some new diagrams arise from the new vertices among $`b_j`$ under the naively treatment of $`\stackrel{~}{f}`$. These diagrams, however, give no quadratic divergence to (24) in total, and then the divergence is not removed. This failure may be caused by the treatment of the non-local contribution in the function $`\stackrel{~}{f}`$. We may have to consider the exact treatment of such terms when we start from the action described by $`f`$.
## 4 Conclusion
In this note we consider the cancellation of the quadratic divergence of the YM theory regularized by the hybrid regularization consisting of the HCD method and the infinitely many PV fields. By an explicit calculation of the vacuum polarization tensor up to $`\mathrm{\Lambda }^4`$ order, all the quadratic divergence exactly cancels in all orders. This result shows that the quadratic divergence is regularized by the hybrid regularization as well as the logarithmic divergence.
The divergence cancels order by order in $`\mathrm{\Lambda }^4`$ and the cancellation mechanism is the same in all orders: the combination of (22) and (24) cancels with (23). We expect that this mechanism works in all the higher orders than $`\mathrm{\Lambda }^4`$ e.g. in the order of $`\mathrm{\Lambda }^8`$ and the quadratic divergence completely cancels out.
In our calculation, the higher derivative term for the ghost field, $`H`$, plays an important role in the cancellation of the quadratic divergence. Since any contribution of $`H`$ does not appear in the superficial degree of divergence $`\omega `$ in (10), the necessity of $`H`$ is unclear. So the simplest choice of the higher derivative term to improve the longitudinal part of the gauge propagator is the function $`f`$ in the reference . In such a case, the treatment of the non-local contribution is so problematic that the cancellation of the quadratic divergence is not shown by the calculation. $`H`$, however, is introduced instead of $`f`$ at the beginning, such a difficulty does not arise and we can demonstrated the cancellation clearly.
We are also interested in the coefficient of the $`\beta `$-function with this regularization scheme. Since $`H`$ gives the same effect with $`f`$ and does not give any contribution to the logarithmic divergence in the order of $`\mathrm{\Lambda }^0`$ we get the familiar value of the coefficient. We will discuss in detail the logarithmic divergence of this theory elsewhere .
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# I INTRODUCTION
## I INTRODUCTION
Within the factorization framework, perturbative QCD has been applied to various processes involving large momentum transfers, both in the spacelike $`q^2=Q^2<0`$ (for reviews, we refer to ) and the timelike $`q^2>0`$ regions (see, for example, ). Note that the running coupling constant $`\alpha _s(\mu ^2)`$ is usually defined with reference to some Euclidean (spacelike) configuration of momenta of scale $`\mu `$. For large spacelike $`q`$, this produces no special complications. One simply uses the renormalization group to resum the logarithmic corrections $`(\alpha _s(\mu ^2)\mathrm{ln}(Q^2/\mu ^2))^N`$ that appear in higher orders of perturbation theory, arriving at an expansion in the effective coupling constant $`\alpha _s(Q^2)`$ which, in the 1-loop approximation, is given by
$$\alpha _s(Q^2)=\frac{4\pi }{(112N_f/3)\mathrm{ln}(Q^2/\mathrm{\Lambda }^2)},$$
(1)
with $`N_f`$ being the number of active flavors and $`\mathrm{\Lambda }`$ denoting $`\mathrm{\Lambda }^{\text{QCD}}`$. In general, the $`\mathrm{\Lambda }`$-parameterization of $`\alpha _s(Q^2)`$ is a series expansion in $`1/L`$ (where $`L=\mathrm{ln}(Q^2/\mathrm{\Lambda }^2)`$), and the definition of $`\mathrm{\Lambda }`$ is fixed only if the $`O(1/L^2)`$-term is added to Eq.(1) . Continuing the logarithms into the region of timelike $`q`$, one should deal with the $`i\pi `$ terms: $`\mathrm{ln}(Q^2/\mu ^2)\mathrm{ln}(Q^2/\mu ^2)\pm i\pi `$, which may produce large higher-order corrections. In the case of the $`R`$-ratio for $`e^+e^{}hadrons`$ process, this problem was discussed in refs. . It was shown there that, by using the $`\mathrm{\Lambda }`$-parameterization for $`\alpha _s(Q^2)`$ in the spacelike region, it is possible to construct for $`R(q^2)`$ an expansion in the timelike region in which all the $`(\pi ^2/L^2)^N`$-terms are resummed explicitly, and, what is most important, the transformation into the timelike region reduces the magnitude of each particular term of the $`1/L`$ expansion. Another well-studied example related to the analytic continuation into the timelike region is the cross section of the Drell-Yan (DY) process $`AB\gamma ^{}X`$. In this case, the $`i\pi `$ factors associated with the continuation of the Sudakov double logarithms $`(\alpha _s\mathrm{ln}^2(Q^2/m^2))^N`$ result in a $`\pi ^2`$-enhanced correction which gives rise to the $`K`$-factor increasing in turn the result of the perturbative QCD calculation by the factor of 3 to bring the DY cross section in agreement with experiment.
For elastic form factors, existing experimental data show a considerable enhancement of the timelike form factors over their spacelike counterparts. In the present paper, we study possible sources of such an enhancement. To disentangle different aspects of the analytic continuation into the timelike region, we proceed step by step, beginning with the simplest cases and then going on to more complicated ones. We start with a discussion of the analytic continuation into the timelike region of the UV logarithms $`\mathrm{ln}(Q^2/\mu _R^2)`$ inducing the $`Q^2`$-dependence of the running coupling constant $`\alpha _s(Q^2)`$. We take the cleanest case of $`R(e^+e^{}hadrons)`$, in which no other types of logarithms appear and review in Section II the continuation procedure for $`R(e^+e^{}hadrons)`$ as given in refs.. In section III, we consider another fundamental process: $`\gamma ^{}\gamma \pi ^0`$. At the leading logarithm level, only the collinear logarithms $`\mathrm{ln}(Q^2/\mu _F^2)`$ are important while $`\alpha _s`$ can be treated as a constant. So, this is another “clean situation” which gives an opportunity to concentrate on the study of the analytic continuation of the collinear logarithms which induce the $`Q^2`$-dependence of the pion distribution amplitude $`\phi _\pi (x,Q^2)`$. In Section IV, we briefly discuss the effects due to the analytic continuation of the Sudakov double logarithms. We consider first the cross section of the Drell-Yan process $`AB\mu ^+\mu ^{}X`$. In this case, the double logs $`\mathrm{ln}^2(Q^2/\mu ^2)`$ appear on a diagram by diagram basis but cancel after resumming over all diagrams of a given order. However, the $`\pi ^2`$ terms generated by the analytic continuation survive and, as already mentioned, produce an enhancement due to the $`K`$-factor. We contrast this outcome with the case of the hard contribution to the pion electromagnetic form factor, in which the induced $`\pi ^2`$ terms cancel together with the double logs. For this reason, the modification of the hard term of the pion form factor in the timelike region is only affected by the analytic continuation of the UV and collinear logarithms. These effects are discussed in Section V. In Section VI, we study the analytic continuation of the hard pQCD contribution to the nucleon form factor. Both in the pion and the nucleon case, we find that the effects due to the continuation into the timelike region are very small. Experimentally, however, the timelike nucleon form factor is essentially larger than its spacelike counterpart. This discrepancy may be regarded as an indication that the hard contribution does not dominate the form factors at accessible momentum transfers. An alternative scenario discussed in many papers is that in the few GeV<sup>2</sup> region the form factors are dominated by the soft mechanism. In Section VII, we study the analytic continuation effects for the soft contribution to the pion electromagnetic form factor within the local quark-hadron duality model motivated by the QCD sum rule analysis of refs.. We show that at the one loop level, there are explicit non-canceled double logarithms $`\mathrm{ln}^2(Q^2/\mu ^2)`$ which produce the $`\pi ^2`$ terms in the timelike region, giving rise to a $`K`$-factor-type enhancement.
## II Continuation of $`\alpha _s`$ into the timelike region: $`𝐑(𝐞^+𝐞^{}\mathrm{𝐡𝐚𝐝𝐫𝐨𝐧𝐬},𝐬)`$
The ratio $`R(s)=\sigma (e^+e^{}\text{hadrons})/\sigma (e^+e^{}\mu ^+\mu ^{})`$, characterizing the total cross section of $`e^+e^{}`$ annihilation into hadrons, provides the simplest example of the analytic continuation of the effective QCD coupling constant $`\alpha _s`$ into the timelike region. The standard procedure (see, e.g., and references cited therein) is to calculate the Adler function $`D(Q^2)`$ by taking the derivative $`D(Q^2)=Q^2d\mathrm{\Pi }/dQ^2`$ of the vacuum polarization $`\mathrm{\Pi }(Q^2)`$ related to $`R(s)`$ by
$$R(s)=\frac{1}{2\pi i}\left(\mathrm{\Pi }(s+iϵ)\mathrm{\Pi }(siϵ)\right).$$
(2)
In perturbative QCD, $`D(Q^2)`$ is given by the $`\alpha _s(Q^2)`$-expansion:
$$D^{\text{QCD}}(Q^2)=\underset{q}{}e_q^2\left\{1+\frac{\alpha _s(Q^2)}{\pi }+d_2\left(\frac{\alpha _s(Q^2)}{\pi }\right)^2+d_3\left(\frac{\alpha _s(Q^2)}{\pi }\right)^3+\mathrm{}\right\}.$$
(3)
In the $`\overline{\mathrm{MS}}`$ scheme, the coefficients $`d_i`$ are known up to $`i=3`$ . Using Eq. (2) and the definition of $`D`$, one can relate $`R^{\text{QCD}}(s)`$, the perturbative QCD version of $`R(s)`$, directly to $`D^{\text{QCD}}(Q^2)`$
$$R^{\text{QCD}}(s)=\frac{1}{2\pi i}_{siϵ}^{s+iϵ}D^{\text{QCD}}(\sigma )\frac{d\sigma }{\sigma }.$$
(4)
The integration contour in Eq.(4) goes below the real axis from $`siϵ`$ to some point $`Q^2`$ in the deep spacelike region and then above the real axis to $`s+iϵ`$, i.e. in the region where the function $`D(s)`$ is analytic.
In a shorthand notation, $`DR\mathrm{\Phi }[D]`$. The actual calculation is very simple if one represents $`\alpha _s(Q^2)`$ through an expansion in $`1/\mathrm{ln}(Q^2/\mathrm{\Lambda }^2)`$, i.e., $`via`$ the $`\mathrm{\Lambda }`$-parameterization. The latter results from the QCD Gell-Mann-Low equation
$$L\mathrm{ln}(Q^2/\mathrm{\Lambda }^2)=\frac{4\pi }{b_0\alpha _s}+\frac{b_1}{b_0^2}\mathrm{ln}\left(\frac{\alpha _s}{4\pi }\right)+\mathrm{\Delta }+\frac{b_2b_0b_1^2}{b_0^3}\frac{\alpha _s}{4\pi }+O(\alpha _s^2),$$
(5)
where $`b_k`$ are $`\beta `$-function coefficients:
$`b_0=112N_f/3`$ , $`b_1=10238N_f/3`$ , $`b_2^{MS}=2857/25033N_f/18+325N_f^2/54`$ .
Inverting (5) by iterations and reexpanding the result in $`1/L`$ we get the $`\mathrm{\Lambda }`$-parameterization for the running coupling constant
$$\alpha _s(Q^2)=\frac{4\pi }{b_0L}\left\{1\frac{L_1}{L}+\frac{1}{L^2}\left[L_1^2\frac{b_1}{b_0^2}L_1+\frac{b_2b_0b_1^2}{b_0^4}\right]+O(1/L^3)\right\},$$
(6)
where $`L_1=(b_1/b_0^2)\mathrm{ln}(b_0L)\mathrm{\Delta }`$ . To fix the functional dependence of $`\alpha _s(Q^2)`$ on $`Q^2`$, one should specify the integration constant $`\mathrm{\Delta }`$. The standard (or “popular”) choice is
$$\mathrm{\Delta }^{pop}=\frac{b_1}{b_0^2}\mathrm{ln}b_0$$
(7)
which gives the shortest expression $`(b_1/b_0^2)\mathrm{ln}(L)`$ for $`L_1`$. A clear disadvantage of this choice is that it guarantees a rather large $`1/L^2`$ correction to $`\alpha _s`$, which results in a large difference between $`\mathrm{\Lambda }^{LO}`$ and $`\mathrm{\Lambda }^{NLO}`$. As argued in Ref. , a more appropriate (optimal) choice is
$$\mathrm{\Delta }^{opt}=\frac{b_1}{b_0^2}\mathrm{ln}b_0\overline{L},$$
(8)
where $`\overline{L}`$ is the average value of the logarithm $`L`$ within the region under study, e.g., $`\overline{L}=4`$ corresponding to $`\alpha _s/\pi 0.1`$. For this choice, the ratio $`L_1/L`$ is smaller than $`7\%`$ and Eq.( 6) has $`1\%`$ accuracy in the whole region $`L>3`$, with the total correction to the simplest formula (1) being less than $`10\%`$.
The $`\mathrm{\Lambda }`$-parameters corresponding to different $`\mathrm{\Delta }`$’s are related by
$$\mathrm{\Lambda }_2=\mathrm{\Lambda }_1e^{(\mathrm{\Delta }_1\mathrm{\Delta }_2)/2}.$$
(9)
In particular,
$$\mathrm{\Lambda }^{opt}=\mathrm{\Lambda }^{pop}/\overline{L}^{b_1/2b_0^2}.$$
(10)
Taking $`\overline{L}=4`$ we get $`\mathrm{\Lambda }^{opt}|_{\overline{L}=4}\mathrm{\Lambda }^{pop}/1.73`$. In connection with the discussion above, we want to stress here that preparing to analytically continue an approximate expression it makes sense to take care of the convergence quality of the original expansion in the spacelike region. If there are corrections which are under our full control and we can make them small, then we should use this opportunity and make them small.
Now one can substitute $`\alpha _s(Q^2)`$ in Eq.(3) by its $`\mathrm{\Lambda }`$-parameterization to get an $`1/L`$ expansion for the Adler function $`D(Q^2)`$. For each term of this expansion, the integral (4) can be calculated explicitly (see also )
$`11,`$ (11)
$`{\displaystyle \frac{1}{L_\sigma }}{\displaystyle \frac{1}{\pi }}(\pi /2\mathrm{arctan}(L_s/\pi ))|_{s>\mathrm{\Lambda }^2}={\displaystyle \frac{1}{\pi }}\mathrm{arctan}(\pi /L_s)={\displaystyle \frac{1}{L_s}}\left\{1{\displaystyle \frac{1}{3}}{\displaystyle \frac{\pi ^2}{L_s^2}}+\mathrm{}\right\},`$ (12)
$`{\displaystyle \frac{\mathrm{ln}(L_\sigma /L_0)}{L_\sigma ^2}}{\displaystyle \frac{\mathrm{ln}(\sqrt{L_s^2+\pi ^2}/L_0)(L_s/\pi )(\pi /2\mathrm{arctan}(L_s/\pi ))+1}{L_s^2+\pi ^2}}|_{s>\mathrm{\Lambda }^2}`$ (13)
$`={\displaystyle \frac{\mathrm{ln}(\sqrt{L_s^2+\pi ^2}/L_0)(L_s/\pi )\mathrm{arctan}(\pi /L_s)+1}{L_s^2+\pi ^2}}={\displaystyle \frac{L_s/L_0}{L_s^2}}\left\{1{\displaystyle \frac{\pi ^2}{L_s^2}}+\mathrm{}\right\}+{\displaystyle \frac{5}{6}}{\displaystyle \frac{\pi ^2}{L_s^4}}+\mathrm{},`$ (14)
$`{\displaystyle \frac{1}{L_\sigma ^2}}{\displaystyle \frac{1}{L_s^2+\pi ^2}}={\displaystyle \frac{1}{L_s^2}}\left\{1{\displaystyle \frac{\pi ^2}{L_s^2}}+\mathrm{}\right\},`$ (15)
$`{\displaystyle \frac{1}{L_\sigma ^n}}(1)^n{\displaystyle \frac{1}{(n1)!}}\left({\displaystyle \frac{d}{dL_s}}\right)^{n2}{\displaystyle \frac{1}{L_s^2+\pi ^2}}={\displaystyle \frac{1}{L_s^n}}\left\{1{\displaystyle \frac{\pi ^2}{L_s^2}}{\displaystyle \frac{n(n+1)}{6}}+\mathrm{}\right\},`$ (16)
where $`L_s=\mathrm{ln}(s/\mathrm{\Lambda }^2)`$, $`L_\sigma =\mathrm{ln}(\sigma /\mathrm{\Lambda }^2)`$, and we assume that $`s>0`$. Furthermore, $`L_0=e^{\mathrm{\Delta }b_0^2/b_1}/b_0`$ is a constant depending on the $`\mathrm{\Delta }`$-choice in the $`\mathrm{\Lambda }`$-parameterization.
Using (6) and incorporating Eqs.(11)-(16) (as well as their generalizations for $`\mathrm{ln}^2L/L^3,\mathrm{ln}L/L^3`$, etc.) one obtains the expansion for $`R^{\text{QCD}}(s)`$
$$R^{\text{QCD}}(s)=\underset{q}{}e_q^2\left\{1+\underset{k=1}{}d_k\mathrm{\Phi }[(\alpha _s/\pi )^k]\right\}$$
(17)
in which all the $`(\pi ^2/L^2)^N`$-terms are resummed.
As noted in Ref., the application of the $`\mathrm{\Phi }`$-operation normally violates nonlinear relations: $`\mathrm{\Phi }[1/L^2](\mathrm{\Phi }[1/L])^2`$, etc. However, it respects linear relations $`\mathrm{\Phi }[A+B]=\mathrm{\Phi }[A]+\mathrm{\Phi }[B]`$, $`\mathrm{\Phi }[\lambda A]=\lambda \mathrm{\Phi }[A]`$ and
$$\mathrm{\Phi }\left[\frac{dD}{dL_\sigma }\right]=\frac{d}{dL_s}\mathrm{\Phi }[D].$$
(18)
In particular, this relation was used to explicitly obtain $`\mathrm{\Phi }[1/L^n]`$ in Eq.(16). As a result, expansion (17) is not an expansion in powers of some particular parameter. A priori, there is no reason to believe that a power series expansion is better than any other. In fact, expansion (17) converges better than the generating expansion (4) for $`D(\sigma )`$ because, as it follows from Eqs. (12)-(16), $`\mathrm{\Phi }[\alpha _s^N]`$ is always smaller than $`\alpha _s^N`$. Moreover, $`(\mathrm{\Phi }[\alpha _s^{N+1}])^{1/(N+1)}<(\mathrm{\Phi }[\alpha _s^N])^{1/N}`$, i.e., the effective expansion parameter decreases in higher orders. Thus, if one succeeded in obtaining a good $`\alpha _s^N`$ expansion for $`D(\sigma )`$ (with all $`d_N`$ being small numbers), then the resulting $`\mathrm{\Phi }[\alpha _s^N]`$ expansion for $`R^{\text{QCD}}(s)`$ is even better, and the best thing to do is to leave it as it is.
The timelike analogue of the simplest $`\mathrm{\Lambda }`$-parameterization for $`\alpha _s(Q^2)`$ (Eq.(1)) is then
$$\stackrel{~}{\alpha }_s(q^2)=\frac{4}{b_0}\left[\frac{\pi }{2}\mathrm{arctan}\left(\frac{\mathrm{ln}(q^2/\mathrm{\Lambda }^2)}{\pi }\right)\right]|_{s>\mathrm{\Lambda }^2}=\frac{4}{b_0}\mathrm{arctan}\left(\frac{\pi }{\mathrm{ln}(q^2/\mathrm{\Lambda }^2)}\right).$$
(19)
This function has a finite value both at $`q^2=\mathrm{\Lambda }^2`$ and $`q^2=0`$.
The well-known deficiency of the perturbative expansion for $`D^{\text{QCD}}(Q^2)`$ in powers of $`\alpha _s(Q^2)`$ is the presence of the unphysical singularity at $`Q^2=\mathrm{\Lambda }^2`$ induced by the Landau pole of $`1/\mathrm{ln}(Q^2/\mathrm{\Lambda }^2)`$. As a consequence, $`R^{\text{QCD}}(s)`$ as calculated from Eq. (4), also has unphysical features: namely, it does not vanish on the negative real axis. In particular, substituting $`1/L_\sigma `$ into the integral (4) and taking negative $`s`$ we get
$$\frac{1}{L_\sigma }|_{s<0}\theta (\mathrm{\Lambda }^2s0),$$
(20)
which results in an unphysical cut of $`\mathrm{\Pi }^{\text{QCD}}(s)`$ in the region $`\mathrm{\Lambda }^2s0`$. Furthermore, applying Eq. (4) to the pole term $`D^{\text{pole}}(Q^2)=\mathrm{\Lambda }^2/(Q^2\mathrm{\Lambda }^2)`$ one obtains the result coinciding with the rhs of Eq.(20). Hence, if one now postulates that $`D^{\text{QCD}}(Q^2)`$ is given by integrating $`R^{\text{QCD}}(s)`$ over the physical region $`s>0`$ only, i.e., if one takes
$$\stackrel{~}{D}^{\text{QCD}}(Q^2)=Q^2_0^{\mathrm{}}\frac{R^{\text{QCD}}(s)}{(s+Q^2)^2}𝑑s$$
(21)
(this transformation will be denoted as $`R\stackrel{~}{D}`$), then $`\stackrel{~}{D}^{\text{QCD}}(Q^2)`$ is free from the unphysical singularities at $`Q^2=\mathrm{\Lambda }^2`$. For instance, combining the two transformations $`(DR\stackrel{~}{D})(D\stackrel{~}{D})`$ one would get
$$\frac{4\pi }{b_0\mathrm{ln}(Q^2/\mathrm{\Lambda }^2)}\frac{4\pi }{b_0}\left(\frac{1}{\mathrm{ln}(Q^2/\mathrm{\Lambda }^2)}\frac{\mathrm{\Lambda }^2}{Q^2\mathrm{\Lambda }^2}\right)\overline{\alpha }_s(Q^2),$$
(22)
which coincides with the pole-free expression for the running coupling constant proposed by Shirkov and Solovtsov . However, since the $`DR`$ operation does not respect nonlinear relations, the $`D\stackrel{~}{D}`$ transformation acting on $`1/L_\sigma ^n`$ would not produce the $`n`$th power of the rhs of Eq.(22). Hence, $`\overline{\alpha }_s`$ cannot serve as an expansion parameter of a power series. Noting that both $`DR`$ and $`R\stackrel{~}{D}`$ convert derivatives with respect to the logarithm of the initial variable into derivatives with respect to the logarithm of the resulting variable we obtain
$$\frac{1}{L_{Q^2}^n}=(1)^n\frac{1}{(n1)!}\frac{d^{n1}}{L_{Q^2}^{n1}}\frac{1}{L_{Q^2}}(1)^n\frac{1}{(n1)!}\frac{d^{n1}}{dL_{Q^2}^{n1}}\left(\frac{1}{L_{Q^2}}\frac{\mathrm{\Lambda }^2}{Q^2\mathrm{\Lambda }^2}\right).$$
(23)
This relation was given in a recent paper by Shirkov , see also Ref. for a related discussion of perturbation theory expansions in the timelike and spacelike regions.
For moderate values of $`Q^2`$, the modification due to the continuation into the timelike region is numerically rather significant: for $`\alpha _s0.3`$ the $`\pi ^2/L^2`$-terms change $`\alpha _s`$ by more than 20%, i.e., they are more important (for an optimal choice of the $`\mathrm{\Delta }`$-parameter) than the 2-loop corrections in the $`\mathrm{\Lambda }`$-parameterization (6). On the other hand, the difference between $`\stackrel{~}{\alpha }_s(Q^2)`$ and the modified spacelike coupling $`\overline{\alpha }_s(Q^2)`$ (taken at mirror momenta) is rather small (less than 10%) for all values of $`Q^2`$.
Thus, using the $`\mathrm{\Lambda }`$-parameterization for the effective QCD coupling constant in the spacelike region, we obtained an explicit expansion for the timelike quantity $`R^{\text{QCD}}(s)`$. One may question, though, the reliability of the above formulas in the region of small momenta $`|q|\mathrm{\Lambda }`$. In particular, a rapid change of $`\stackrel{~}{\alpha }_s`$ in the small-$`q^2`$ region (compare $`\stackrel{~}{\alpha }_s(\mathrm{\Lambda }^2)=2\pi /b_0`$ and $`\stackrel{~}{\alpha }_s(0)=4\pi /b_0`$) is as suspicious as the Landau pole of $`\alpha _s(Q^2)`$. Evidently, they both are artifacts of the analytic continuation procedure applied outside the applicability region. It is well known that the physical $`R(s)`$ vanishes below the two-pion threshold and approaches the perturbative value only for values of $`s`$ marginally larger than $`\mathrm{\Lambda }^2`$. So, one may argue that a more realistic procedure is to integrate $`R^{QCD}(s)`$ in the dispersion relation (21) from some effective threshold $`s_0`$ rather than from zero. Taking, e.g., $`s_0=\mathrm{\Lambda }^2`$, one would get another effective spacelike coupling, call it $`\widehat{\alpha }_s(Q^2)`$. It vanishes at $`Q^2=0`$, but is essentially constant $`\widehat{\alpha }_s(Q^2)/\pi 0.1`$ in a wide range $`\mathrm{\Lambda }^2Q^230\mathrm{\Lambda }^2`$ of momenta. Hence, $`\widehat{\alpha }_s(Q^2)`$ effectively “freezes” at small momenta (see also ).
## III Collinear logarithms and distribution amplitudes in the timelike region
The logarithmic dependence on the large momentum scale $`Q^2`$ may also appear through mass logarithms $`\mathrm{ln}(Q^2/m^2)`$, where $`m`$ is some mass or an infrared regularization parameter. Note that the standard pQCD factorization
$$T(Q^2/m^2)=t(Q^2/\mu ^2)\phi (\mu ^2)$$
(24)
works only in a single-logarithm situation, when there may appear just one $`\mathrm{ln}(Q^2/m^2)`$ factor per each loop. These collinear logarithms can be absorbed into the renormalization of the long-distance function (distribution amplitude) $`\phi (\mu ^2)`$. In particular, taking $`\mu ^2=Q^2`$, one arrives at the description in terms of $`Q^2`$-dependent functions $`\phi (Q^2)`$. Again, if the large momentum is timelike, the collinear logarithms $`\mathrm{ln}(Q^2/m^2)`$ acquire the imaginary part $`\pm i\pi `$, and we may ask how one should define the $`Q^2`$-dependent distribution amplitudes $`\phi (Q^2)`$ in the timelike region.
To approach this problem, let us consider the simplest example of a hard exclusive process: $`\pi ^0`$ production in $`\gamma ^{}\gamma `$ collisions. Its pQCD expansion starts at zero order in $`\alpha _s`$
$$t_0(x,Q^2)=\frac{1}{xQ^2},$$
(25)
and the leading pQCD result for the large-$`Q^2`$ behavior of the form factor is
$$F_{\gamma ^{}\gamma \pi }(Q^2)=\frac{4\pi }{3}_0^1\frac{\phi _\pi (x)}{xQ^2}𝑑x\frac{4\pi f_\pi }{3Q^2}I_0.$$
(26)
The nonperturbative information here is accumulated in the same integral
$$I_0=\frac{1}{f_\pi }_0^1\frac{\phi _\pi (x)}{x}𝑑x$$
(27)
that appears in the one-gluon-exchange diagram for the pion electromagnetic form factor . The value of $`I`$ depends on the shape of the pion distribution amplitude $`\phi _\pi (x)`$. In particular, using the asymptotic form
$$\phi _\pi ^{\text{as}}(x)=6f_\pi x(1x)$$
(28)
gives $`I_0^{\text{as}}=3`$. If one takes instead the Chernyak–Zhitnitsky model
$$\phi _\pi ^{CZ}(x)=30f_\pi x(1x)(12x)^2,$$
(29)
the integral $`I_0`$ increases by a sizable factor of 5/3: $`I_0^{CZ}=5`$. This difference can be used for an experimental discrimination between the two competing models for the pion distribution amplitude.
At one loop, the $`\overline{\mathrm{MS}}`$ coefficient function for the $`\gamma ^{}\gamma \pi ^0`$ form factor was calculated in refs. and was found to be
$$t(x,Q^2;\mu ^2)=\frac{1}{xQ^2}\left\{1+C_F\frac{\alpha _s}{2\pi }\left[\left(\frac{3}{2}+\mathrm{ln}x\right)\mathrm{ln}(Q^2/\mu ^2)+\frac{1}{2}\mathrm{ln}^2x\frac{x\mathrm{ln}x}{2(1x)}\frac{9}{2}\right]\right\}.$$
(30)
In full compliance with the factorization theorems (see also ), the one-loop contribution contains no Sudakov double logarithms $`\mathrm{ln}^2Q^2`$ of the large momentum transfer $`Q`$. Physically, this result is due to the color neutrality of the pion. In the axial gauge, the Sudakov double logarithms appear in the box diagram but they are canceled by similar terms from the quark self-energy corrections. In Feynman gauge, the double logarithms $`\mathrm{ln}^2Q^2`$ do not appear in any one-loop diagram. It is easy to check that the term containing the logarithm $`\mathrm{ln}(Q^2/\mu ^2)`$ has the form of a convolution
$$\frac{1}{xQ^2}C_F\frac{\alpha _s}{2\pi }\left(\frac{3}{2}+\mathrm{ln}x\right)=\underset{0}{\overset{1}{}}\frac{1}{\xi Q^2}V(\xi ,x)𝑑\xi $$
(31)
of the lowest-order (“Born”) term $`t_0(\xi ,Q^2)=1/\xi Q^2`$ and the kernel
$$V(\xi ,x)=\frac{\alpha _s}{2\pi }C_F\left[\frac{\xi }{x}\theta (\xi <x)\left(1+\frac{1}{x\xi }\right)+\frac{\overline{\xi }}{\overline{x}}\theta (\xi >x)\left(1+\frac{1}{\xi x}\right)\right]_+$$
(32)
governing the evolution of the pion distribution amplitude. The “+”-operation is defined here, as usual , by
$$[F(\xi ,x)]_+=F(\xi ,x)\delta (\xi x)\underset{0}{\overset{1}{}}F(\zeta ,x)𝑑\zeta .$$
(33)
Since the asymptotic distribution amplitude is the eigenfunction of the evolution kernel $`V(\xi ,x)`$ corresponding to zero eigenvalue
$$_0^1V(\xi ,x)\phi ^{\text{as}}(x)𝑑x=0,$$
(34)
the coefficient $`(\frac{3}{2}+\mathrm{ln}x)`$ of the $`\mathrm{ln}(Q^2/\mu ^2)`$ term vanishes after the $`x`$-integration with $`\phi ^{\text{as}}(x)`$. Hence, the size of the one-loop correction for the asymptotic distribution amplitude is $`\mu `$-independent and is determined only by the remaining terms (for a detailed discussion of their structure, see Ref.).
In this section, we want to concentrate on the $`Q^2`$-dependence induced by collinear logarithms, which in this process start to appear at the one-loop level. The UV logarithms shifting the argument of $`\alpha _s`$ appear only at two-loop order. Hence, analyzing the leading collinear logarithms $`(\alpha _s\mathrm{ln}(Q^2/\mu ^2))^N`$ we will treat $`\alpha _s`$ as a constant. The factorization theorem means essentially that the leading logarithms $`(\alpha _s\mathrm{ln}(Q^2/\mu ^2))^N`$ exponentiate in higher orders producing a factor which can be absorbed into the renormalization of the pion distribution amplitude
$$\phi (\mu ^2)\mathrm{exp}[\mathrm{ln}(Q^2/\mu ^2)V]\phi (\mu ^2).$$
(35)
Now, taking a timelike momentum $`Q^2=q^2`$, we would get an extra $`\pm i\pi `$ term: $`\mathrm{ln}(Q^2/\mu ^2)\mathrm{ln}(q^2/\mu ^2)\pm i\pi `$ and
$$\mathrm{exp}[\mathrm{ln}(Q^2/\mu ^2)V]\mathrm{exp}[\mathrm{ln}(q^2/\mu ^2)V]\mathrm{exp}[\pm i\pi V].$$
(36)
The first exponential corresponds to the standard evolution of the pion distribution amplitude from the scale $`\mu ^2`$ to the scale $`q^2`$. The second exponential is specific for the timelike kinematics. In our approximation, it is $`q^2`$-independent and can be treated as a conversion factor for the transition from a “spacelike” distribution amplitude $`\phi `$ to its timelike counterparts $`\stackrel{~}{\phi }_\pm `$
$$\stackrel{~}{\phi }_\pm =\mathrm{exp}[\pm i\pi V]\phi .$$
(37)
In general, “timelike” distribution amplitudes have both real and imaginary parts. However, since $`V\phi ^{\text{as}}=0`$, the spacelike asymptotic distribution amplitude does not differ from its timelike counterpart.
To estimate the effect of phases, let us consider the case when the spacelike distribution amplitude is given by the Chernyak-Zhitnitsky (CZ) model , which can be represented as
$$\phi ^{\text{CZ}}=\phi ^{\text{as}}+\phi _2,$$
(38)
where $`\phi ^{\text{as}}=6f_\pi x(1x)`$ and $`\phi _2=24f_\pi x(1x)(15x(1x))`$ is the next eigenfunction of the $`V`$ kernel corresponding to the eigenvalue $`\gamma _2=\frac{25}{18}\alpha _s/\pi `$. The timelike distribution amplitude is then
$$\stackrel{~}{\phi }_\pm ^{\text{CZ}}=\phi ^{\text{as}}+e^{\pm i(25/18)\alpha _s}\phi _2$$
(39)
and the $`I`$ integral for this function is
$$\stackrel{~}{I}_\pm ^{\text{CZ}}=3+2e^{\pm i(25/18)\alpha _s}.$$
(40)
Its absolute magnitude
$$\left|\stackrel{~}{I}^{\text{CZ}}\right|=5\sqrt{1\frac{24}{25}\mathrm{sin}^2\left(\frac{25}{36}\alpha _s\right)}$$
(41)
is slightly smaller (by 2% if $`\alpha _s=0.3`$) than the spacelike value $`I^{\text{CZ}}=5`$.
## IV Sudakov logarithms and $`K`$-factor
Small radiative corrections in the timelike version of the $`\gamma ^{}\gamma \pi ^0`$ process are in strong contrast with the large $`K`$-factor value found for the Drell–Yan process $`AB\gamma ^{}X`$. These corrections originate from the Sudakov double logarithms $`(\alpha _s\mathrm{ln}^2(Q^2/m^2))^N`$. In the spacelike region, the double logarithms due to the virtual gluon exchanges exponentiate into the Sudakov form factor
$$S(Q^2/m^2)=e^{\alpha _s\mathrm{ln}^2(Q^2/m^2)/3\pi }$$
(42)
(again, we treat $`\alpha _s`$ as a constant). In the DY process, the photon momentum is timelike, and the logarithm $`\mathrm{ln}(Q^2/m^2)`$ acquires the $`\pm i\pi `$ additional term, so that one has
$$L^2L^2\pm 2i\pi L+\pi ^2.$$
(43)
The imaginary parts of the two conjugate diagrams shown in Fig.1a,b cancel, the double log $`L^2`$ from Fig.1a (b) is also canceled by the real gluon emission diagram Fig.1c (d), while the $`\pi ^2`$-term survives and leads, after exponentiation, to a large $`K`$ factor $`\mathrm{exp}[2\pi \alpha _s/3]2`$. The crucial technical observation here is that the real emission diagrams give $`L^2`$ without $`\pi ^2`$-terms. This can be easily understood looking at the reduced diagrams for the virtual vertex correction and real gluon emission. Take for definiteness, the Feynman gauge. Then the virtual vertex correction diagram Fig.1a contains the $`\mathrm{ln}^2(s/m^2)`$ term, where $`s=(xp_A+yp_B)^2=xyS`$ is timelike, and the resulting contribution contains a $`\pi ^2`$ term. The real emission diagram Fig.1c, in turn, contains the $`\mathrm{ln}^2(u/m^2)`$ term, where $`u=(xp_Ayp_B)^2=xyS`$ is now spacelike, and there is no $`\pi ^2`$ term in this contribution.
For the hard pQCD contribution to the pion electromagnetic form factor (which is considered in more detail in the next section), the situation is completely different. In this case, the initial state of the hard subprocess is represented by a $`q\overline{q}`$ pair with momenta $`xp`$ and $`(1x)p`$. After the hard scattering subprocess, one deals with a $`q\overline{q}`$ pair with the final momentum $`p^{}`$ shared in fractions $`yp^{}`$ and $`(1y)p^{}`$. In Feynman gauge, the double logarithms $`\mathrm{ln}^2(Q^2/\mu ^2)`$, where $`Q^2=(pp^{})^2`$, appear when the reduced diagrams have the structure of those shown in Fig.2. One can easily check that the relevant momentum transfers in all four cases have the structure $`t_{ij}=(x_ipy_jp^{})^2=x_iy_jQ^2`$, resulting in the double logs $`\mathrm{ln}^2(t_{ij}/\mu ^2)`$. When the momentum transfer $`q=p^{}p`$ is spacelike, all $`t_{ij}`$’s are spacelike, whereas for a timelike $`q`$, all $`t_{ij}`$’s are timelike as well. In the latter case, one has $`\pi ^2`$ terms for each particular diagram. The double logarithms in diagrams 2a and 2b (2c and 2d) differ in sign because the soft gluon interacts in the final state with quarks of opposite color charge. Hence, due to the color neutrality of the pion, the double logs $`\mathrm{ln}^2(Q^2/\mu ^2)`$ cancel for the sum of the diagrams of a given order. For timelike $`q`$, they cancel together with the accompanying $`\pi ^2`$ terms.
Thus, even for a timelike momentum transfer $`q`$, there is no $`K`$-factor for the pQCD hard contribution to the pion electromagnetic form factor. After cancellation of the Sudakov double logs, only the evolution-related collinear logarithms remain, and the situation is rather similar to the simplest case of the $`\gamma ^{}\gamma \pi ^0`$ form factor.
## V Pion form factor in the perturbative QCD approach
The general pQCD factorization formula for the pion electromagnetic form factor at large momentum transfer reads
$`F_\pi ^{\text{HARD}}(Q^2)={\displaystyle _0^1}𝑑x{\displaystyle _0^1}\phi (x,\mu _F^2,\mu _R^2)T(x,y,Q^2,\mu _F^2,\mu _R^2)\phi (y,\mu _F^2,\mu _R^2)𝑑y,`$ (44)
where $`\mu _F`$ is the factorization scale for the collinear logarithms and $`\mu _R`$ is the renormalization scale for the UV logarithms. The hard scattering amplitude is given by an expansion in $`\alpha _s`$
$`T(x,y,Q^2)={\displaystyle \frac{2\pi C_F\alpha _s(\mu _R^2)}{xyQ^2N_c}}\left[1+{\displaystyle \frac{\alpha _s(\mu _R^2)}{2\pi }}T^{(1)}(x,y,Q^2,\mu _F^2,\mu _R^2)+O\left(\alpha _s^2\right)\right].`$ (45)
The one-loop correction $`T^{(1)}(x,y,Q^2,\mu _F^2,\mu _R^2)`$ was calculated using the dimensional regularization in several papers which differ from each other by a particular choice of renormalization and factorization prescriptions. These differences (and also typos and mistakes) were discussed in Refs. . In the $`\overline{MS}`$ subtraction scheme, supplemented by the requirement that both the $`\alpha _s`$ and $`\phi _\pi (x)`$ are process-independent functions, the one-loop correction has the form
$`T^{(1)}`$ $`=`$ $`C_FT^F(x,y,Q^2,\mu _F^2)+{\displaystyle \frac{b_0}{2}}T^\beta (x,y,Q^2,\mu _R^2)+\left(C_FN_c/2\right)T^A(x,y),`$ (46)
$`T^F`$ $`=`$ $`\left[3+\mathrm{ln}(xy)\right]\mathrm{ln}\left({\displaystyle \frac{Q^2}{\mu _F^2}}\right)+{\displaystyle \frac{1}{2}}\mathrm{ln}^2(xy)+{\displaystyle \frac{5}{2}}\mathrm{ln}(xy){\displaystyle \frac{x\mathrm{ln}x}{2(1x)}}{\displaystyle \frac{y\mathrm{ln}y}{2(1y)}}{\displaystyle \frac{14}{3}},`$ (47)
$`T^\beta `$ $`=`$ $`\mathrm{ln}\left({\displaystyle \frac{Q^2}{\mu _R^2}}\right)\mathrm{ln}(xy)+{\displaystyle \frac{5}{3}},`$ (48)
$`T^A`$ $`=`$ $`\text{Li}\text{2}(1x)\text{Li}\text{2}(x)+\mathrm{ln}(1x)\mathrm{ln}\left({\displaystyle \frac{y}{1y}}\right){\displaystyle \frac{5}{3}}`$ (51)
$`+{\displaystyle \frac{1}{(xy)^2}}((x+y2xy)\mathrm{ln}(1x)+2xy\mathrm{ln}(x)+{\displaystyle \frac{(1x)x^2+(1y)y^2}{xy}}`$
$`\times [\mathrm{ln}(1x)\mathrm{ln}(y)\text{Li}\text{2}(1x)+\text{Li}\text{2}(x)])+\{xy\}`$
(we use here notations similar to those of Ref. ). As usual, Li<sub>2</sub> is the dilogarithm (Spence) function.
We already discussed in the previous section that all the Sudakov double logarithms $`\mathrm{ln}^2(Q^2/\mu ^2)`$ cancel and that only the collinear single-logarithms $`\mathrm{ln}(Q^2/\mu _F^2)`$ remain. Comparing Eq. (51) with the one loop correction to the $`\gamma ^{}\gamma \pi ^0`$ hard scattering amplitude, Eq. (30), one can easily notice many similarities in the structure of the one-loop corrections in these two cases. In particular, the coefficient $`[3+\mathrm{ln}(xy)]`$ in front of the evolution logarithm $`\mathrm{ln}(Q^2/\mu _F^2)`$ is the sum $`[3/2+\mathrm{ln}x]+[3/2+\mathrm{ln}y]`$ of terms corresponding to the convolution of the tree level term $`1/\xi \eta `$ with the kernels $`V(x,\xi )\delta (\eta y)`$ and $`\delta (x\xi )V(\eta ,y)`$ (see Eq. (31) ) inducing the evolution of the pion distribution amplitudes $`\phi (x,\mu _F^2)`$ and $`\phi (y,\mu _F^2)`$. In a sense, the collinear logarithms indicate that the pion structure is probed at a scale proportional to $`Q`$. However, one should remember that since the asymptotic wave function does not evolve, the coefficient accompanying the evolution logarithm $`\mathrm{ln}(Q^2/\mu _F^2)`$ vanishes if the pion wave function has the asymptotic shape. As a result, the choice of $`\mu _F`$ in that case does not affect the size of the one-loop correction. The latter comes from several sources which can be identified in a way similar to the detailed analysis of the one loop correction for the $`\gamma ^{}\gamma \pi ^0`$ form factor given in Ref. .
In addition to the evolution term proportional to $`\mathrm{ln}(Q^2/\mu _F^2)`$, there is a rather large positive correction due to the $`\frac{1}{2}\mathrm{ln}^2(xy)`$ term and even larger negative contributions corresponding to the constant term $`14/3`$ and the logarithmic term $`\frac{5}{2}\mathrm{ln}(xy)`$. As explained in Ref. , in the $`\gamma ^{}\gamma \pi ^0`$ case, the $`\frac{1}{2}\mathrm{ln}^2x`$ term is a result of a positive $`\mathrm{ln}^2x`$ evolution-related contribution and a negative $`\frac{1}{2}\mathrm{ln}^2x`$ Sudakov-related term. As we emphasized earlier, the Sudakov $`\mathrm{ln}^2Q^2`$ double logs should cancel, otherwise there is no pQCD factorization. However, when several scales are involved, like $`Q^2`$ and $`xyQ^2`$ in our case, there may be a remnant like $`\mathrm{ln}^2(xy)`$. In the pion form factor case, there is another scale $`xQ^2`$, the quark virtuality, whence the single logarithms $`\mathrm{ln}x+\mathrm{ln}y`$. The latter give a rather large negative contribution. There are also large negative constants ($`9/2`$ in the $`\gamma ^{}\gamma \pi ^0`$ case and $`14/3`$ in the pion form factor case), which are another (and numerically very important) manifestation of the Sudakov effects in the impact parameter space. In full analogy with the results of Ref. , these \[and the $`\mathrm{ln}^2(xy),\mathrm{ln}(xy)`$ terms\] result from convoluting the $`b`$-space version $`K_0(\sqrt{xyQ^2b^2})`$ of the one-gluon exchange propagator and the $`b`$-space Sudakov form factors $`S(x,bQ)`$, $`S(y,bQ)`$ (exact one-loop expressions are given in ). In the practically important case of the asymptotic wave function, the total correction due to the $`T^F`$ term is negative and equal to $`(71/18)\alpha _s/\pi `$; as one could expect, it is approximately twice larger than that in the $`\gamma ^{}\gamma \pi ^0`$ case.
The situation is reversed in the case of the UV related $`T^\beta `$ term: it is dominated by large positive contributions. In full accordance with the renormalization group, the UV logarithm $`\mathrm{ln}(Q^2/\mu _R^2)`$ is accompanied by the $`\beta `$-function coefficient $`b_0`$. It generates the running of the effective QCD coupling $`\alpha _s`$, “suggesting” that we should use some scale proportional to $`Q^2`$ as its argument. According to Brodsky, Lepage and Mackenzie , one should choose the argument of the effective coupling constant in such a way as to absorb all the terms proportional to $`b_0`$ from the next loop correction. Taken literally, the BLM prescription in our case corresponds to using the (rescaled) gluon virtuality
$$\mu _R^2=xyQ^2e^{5/3}xy(Q/2.3)^2$$
(52)
as the argument of $`\alpha _s`$. The rescaling factor $`e^{5/6}2.3`$ reflects the fact that the $`\overline{MS}`$ scheme measures the momenta in “wrong” units. To cure this effect, one may introduce a version of the minimal subtraction scheme which measures momenta in more “physical” units $`\mathrm{\Lambda }_{\text{PHYS}}=e^{5/6}\mathrm{\Lambda }_{\overline{MS}}`$. This choice is similar to using the $`\alpha _V`$ coupling of Brodsky et al. (note that their relation $`\alpha _V(Q)=\alpha _s^{\overline{MS}}(e^{5/6}Q)[1\frac{2}{3}N_c\alpha _s^{\overline{MS}}/\pi ]`$ includes a NLO correction). One should remember, however, that the actual expansion parameter for switching from the leading to the next-to-leading level is $`1/\mathrm{ln}(Q^2/\mathrm{\Lambda }^2)`$ rather than $`\alpha _s`$ as a whole. As a result, the “non-physical” nature of the $`\overline{MS}`$ scheme is almost totally compensated by the non-optimal “popular” choice for the analytic form of $`\alpha _s(Q^2)`$. As discussed in Section II, $`\mathrm{\Lambda }^{opt}\mathrm{\Lambda }^{pop}/1.74`$. As a result, $`\mathrm{\Lambda }_{\text{PHYS}}^{opt}1.3\mathrm{\Lambda }_{\overline{MS}}^{pop}`$. Due to the compensation of two opposite corrections, the standard $`\mathrm{\Lambda }_{\overline{MS}}^{pop}`$ parameter is rather close to the genuine “$`\mathrm{\Lambda }_{\text{QCD}}`$” parameter of the $`\mathrm{PHYS}^{opt}`$ scheme in which the coupling $`\alpha _s(k^2)`$ corresponding to the vertex with the gluon momentum $`k`$ is given by $`4\pi /b_0\mathrm{ln}(k^2/\mathrm{\Lambda }^2)`$ without sizable next-to-leading order corrections. In other words, using the NLO expression for $`\alpha _s`$ in the popular form is equivalent to adding a negative term $`(b_1/b_0^2)\mathrm{ln}L`$ to $`T^\beta `$, partially compensating the “5/3” constant. For $`L4`$, this reduces 5/3 by a factor of 3. Choosing “PHYS” vs. $`\overline{MS}`$ and “opt” vs. “pop” one reduces both types of corrections which iterate in higher orders. As stated earlier, if the size of some corrections is under our control, it is preferable to keep them small rather than rely on cancellation of large terms. The closeness of $`\mathrm{\Lambda }_{\text{PHYS}}^{opt}`$ to $`\mathrm{\Lambda }_{\overline{MS}}^{pop}`$ means that discussing the pQCD applicability region one should compare the $`\mathrm{\Lambda }^2`$ parameter of the $`\overline{MS}^{pop}`$ scheme with the actual (unrescaled) gluon virtuality $`xyQ^2`$.
However, taking the argument of the effective coupling constant proportional to $`xyQ^2`$ one faces the following problem: since the integration is over all the momentum fractions in the range $`0x,y1`$, the “short-distance” amplitude in this case always gets contributions from the infrared region of arbitrarily small virtualities. In this sense, such an “inside the integral” BLM prescription contradicts the spirit of the pQCD factorization ideology which aims at a perfect separation of the short-distance and long-distance effects (at least in perturbation theory). The consistent pQCD approach is to apply the BLM prescription to the form factor as a whole, i.e., “outside the integral”. In this case, one should choose $`\mu _R`$ from the requirement that one should get zero for
$$\mathrm{ln}(Q^2/\mu _R^2)\mathrm{ln}(xy)+\frac{5}{3},$$
where the “averaging” procedure $`\mathrm{}`$ stands for integration with $`\phi _\pi (x)\phi _\pi (y)/xy`$. This gives a universal $`x,y`$-independent scale $`\mu _R=a_RQ`$, which depends now on the shape of the distribution amplitude. For the asymptotic wave function, as we have seen, the average value of $`\mathrm{ln}x`$ is $`3/2`$, hence the “outside the integral” BLM scale is (see also )
$$\mu _R^2|_{\phi =\phi ^{\text{as}}}=Q^2e^{5/33}(e^{5/3}Q^2)/20.$$
(53)
As argued above, in the “pop” treatment, the factor $`e^{5/3}`$ is largely compensated by the NLO corrections to $`\alpha _s(Q^2)`$, and, hence the essential virtuality of the “hard gluon” exchanged between the quarks is “only” by a factor of 20 smaller than $`Q^2`$, the nominal momentum transfer to the pion. Nevertheless, despite this sizable rescaling factor, the pQCD factorization approach is fully consistent in the asymptotic sense: for a sufficiently large $`Q^2`$ one can calculate the short-distance amplitude perturbatively in terms of an arbitrarily small expansion parameter $`\alpha _s(Q^2/20)`$. For comparison, in the case of Chernyak-Zhitnitsky wave function
$$\mathrm{ln}x|_{\phi ^{CZ}}=\frac{7}{3}$$
and $`a_R^{CZ}a_R^{\text{as}}/2.3`$: the essential gluon virtualities are 100 times smaller than $`Q^2`$. In this case, one should not expect early applicability of pQCD.
We would like to emphasize that the reason for such a drastic shift of the BLM scale to very low $`\mu _R`$ values is the positive large value of the $`T^\beta `$ correction: for $`\mu _R=Q`$ and $`\phi _\pi =\phi _\pi ^{\text{as}}`$, the $`T^\beta `$ term contributes the NLO correction $`10\alpha _s/\pi `$. One may be tempted to combine the large positive $`T^\beta `$ term and a sizable negative $`T^F`$ term to end up with a smaller total correction $`6\alpha _s/\pi `$. Physically, though, these corrections have a completely different nature: as argued above, the $`T^F`$ term comes primarily from the Sudakov effects. Since the latter exponentiate, one deals here with the $`e^{k_F\alpha _s}`$ type series in which the sign of the corrections alternates. On the other hand, the UV corrections form a geometric series summed into $`1/(1k_\beta \alpha _s)`$. Hence, there is no doubt that a partial cancellation of the $`\alpha _s`$ terms will be followed by an amplified total correction at the $`\alpha _s^2`$ level. Leaving the physically unrelated Sudakov and UV corrections separate and addressing the region of experimentally accessible values of $`Q^210`$GeV<sup>2</sup>, one should take $`\alpha _s`$ at an infrared scale $`\mathrm{\Lambda }`$ (where it freezes at a value close to 0.3) and supplement the result by the exponential $`e^{4\alpha _s/\pi }0.7`$ of the negative one-loop corrections induced by the $`b`$-space Sudakov effects. Turning to the timelike momenta, we cannot find any sources of enhancement: for the $`|k|\mathrm{\Lambda }`$ region we see no other choice rather than to take the frozen value $`\alpha _s0.3`$ both in the spacelike and timelike regions, while for large momenta, the continuation of $`\alpha _s(a_F^2Q^2)`$ converts $`1/L`$ into $`1/(L+i\pi )`$ and the ratio $`|F_{\mathrm{hard}}^{\mathrm{timelike}}|/F_{\mathrm{hard}}^{\mathrm{spacelike}}`$ is $`1/\sqrt{1+\pi ^2/L^2}`$, i.e., the timelike term is suppressed compared to the spacelike one.
Since the structure of the evolution corrections for the pion form factors is essentially identical to that of the $`\gamma ^{}\gamma \pi ^0`$ form factor, to continue the evolution logarithms, we can use the approach outlined in Section III. The first step is to write the solution of the evolution equation as an expansion over Gegenbauer polynomials $`C_n^{3/2}(2x1)`$, the eigenfunctions of the LO kernel $`V^{(0)}(x,y;\alpha _s)`$:
$$\varphi (x,Q^2)=\varphi ^{\text{as}}(x)\left[1+\underset{k=1}{\overset{\mathrm{}}{}}a_{2k}^\pi (Q^2)C_{2k}^{3/2}(2x1)\right],$$
(54)
where $`\varphi ^{\text{as}}(x)=6x(1x)`$ is the asymptotic distribution amplitude of the pion. The Gegenbauer moments $`a_{2k}^\pi (Q^2)`$ have a simple evolution pattern
$$a_{2k}^\pi (Q^2)=a_{2k}^\pi (\mu _0^2)\mathrm{exp}\left[\gamma _{2k}\mathrm{ln}\left(Q^2/\mu _0^2\right)\right]$$
(55)
(we treat $`\alpha _s`$ as a constant here), with $`\gamma _{2k}`$ being the corresponding anomalous dimensions and $`a_{2k}(\mu _0^2)`$ the Gegenbauer moments of the initial distribution amplitude
$$\varphi _0(x)=\varphi ^{\text{as}}(x)\left[1+\underset{k=1}{\overset{\mathrm{}}{}}a_{2k}^\pi (\mu _0^2)C_{2k}^{\frac{3}{2}}(2x1)\right].$$
(56)
This representation is very convenient to perform the analytic continuation to the timelike region of $`Q^2`$. Indeed, changing $`Q^2q^2`$, one obtains the natural shift $`\mathrm{ln}(Q^2)\mathrm{ln}(q^2)+i\pi `$, so that
$$a_{2k}^\pi (Q^2)a_{2k}^\pi (q^2)=a_{2k}^\pi (|q^2|)e^{i\delta _{2k}},$$
(57)
where
$$\delta _{2k}\pi \gamma _{2k}.$$
(58)
From (57) it is obvious that the only interesting and potentially enhancing effect is due to the phases $`\delta _{2k}`$, since they can destroy some fine tuning of the coefficients $`a_{2k}^\pi (\mu _0^2)`$ and produce a positive interference. But in order to realize this possibility, one should start with a situation when there are negative coefficients, say, $`a_2^\pi (\mu _0^2)<0`$, while the corresponding phase is close to $`\pi `$, e.g., $`\delta _2\pi `$. Such a situation is hard to imagine. Even the (unrealistic) CZ distribution amplitude has $`a_2^\pi (\mu _0^2=(0.5\text{GeV})^2)=2/3`$, while other models are closer to the asymptotic distribution amplitude, though all models provide $`a_2>0`$. Furthermore, the phase $`\delta _2`$ is $`25/18\alpha _s`$, so one needs a prohibitively large value $`\alpha _s2.5`$ for the coupling constant. In Fig. 3 we plot the ratio TL/SL for the pion form factor in the CZ model, taking the frozen value $`\alpha _s(Q^2)=0.3`$. As one can see, the absolute value of $`F_\pi ^{\text{HARD}}(q^2)`$ in the timelike region is reduced.
To conclude, the perturbative contribution to the pion form factor, $`F_\pi ^{\text{HARD}}(Q^2)`$, with a realistic distribution amplitude (which is close to the asymptotic one, see ) produces no sizable effects in analytically continuing to timelike $`q^2`$ values. The only potential effect is due to the substitution $`\overline{\alpha }_s(Q^2)\stackrel{~}{\alpha }_s(q^2)`$which results in a 10%-reduction of the form factor.
## VI Evolution phases of the nucleon distribution amplitudes and the nucleon form factor
The nucleon form factor in the leading $`\alpha _s`$ order can be cast in the form
$$Q^4G_\mathrm{M}^\mathrm{N}(Q^2)=\frac{1}{54}\left[4\pi \overline{\alpha }_s(Q^2)\right]^2|f_\mathrm{N}|^2_0^1[dx]_0^1[dy]\left[2\underset{i=1}{\overset{7}{}}e_iT_i(x_j,y_j)+\underset{i=8}{\overset{14}{}}e_iT_i(x_j,y_j)\right],$$
(59)
where the amplitudes $`T_i(x_j,y_j)=\varphi ^{(N)}(\{x\},Q^2)T_\mathrm{H}^i(\{x\},\{y\})\varphi ^{(N)}(\{y\},Q^2)`$ represent convolutions of $`T_\mathrm{H}^i`$ with the appropriate distribution amplitudes $`\varphi ^{(N)}(\{x\},Q^2)`$ evaluated term by term for each contributing diagram (marked by the index “$`i`$”). The nucleon distribution amplitude can be represented as an expansion over symmetrized combinations $`\stackrel{~}{\mathrm{\Phi }}(\{x\})`$ of Appell polynomials (for more details, we refer to )
$$\varphi ^{(N)}(\{x\},Q^2)=\varphi ^{\text{as}}(\{x\})\underset{n=0}{\overset{\mathrm{}}{}}B_n(Q^2)\stackrel{~}{\mathrm{\Phi }}_n(\{x\}),$$
(60)
with
$$B_n(Q^2)=B_n(\mu _0^2)\mathrm{exp}\left[\stackrel{~}{\gamma }_n\mathrm{ln}\left(Q^2/\mu _0^2\right)\right],$$
(61)
and eigenfunctions
$$\stackrel{~}{\mathrm{\Phi }}_k(x_i)=\underset{m,n=0}{\overset{m+n=M}{}}c_{mn}^k_{mn}(5,2,2;x_1,x_3),$$
(62)
where $`_{mn}(5,2,2;x_1,x_3)`$ are the Appell polynomials One can also expand $`\stackrel{~}{\mathrm{\Phi }}_n(\{x\})`$ over the polynomials proposed in Ref.; the particular choice of the basis is not essential for our purposes.. Here
$$\stackrel{~}{\gamma }_n(M)=\frac{\alpha _s}{4\pi }\left(\frac{3}{2}C_\mathrm{F}+2\eta _n(M)C_\mathrm{B}\right),$$
(63)
are the associated anomalous dimensions of trilinear quark operators with the quantum numbers of the nucleon containing external derivatives . In Eq. (60) $`\varphi ^{\text{as}}(\{x\})=120x_1x_2x_3`$ denotes the asymptotic distribution amplitude of the nucleon and $`B_n(\mu _0^2)`$ are expansion coefficients for some initial distribution amplitude
$$\varphi _0^{(N)}(\{x\})=\varphi ^{\text{as}}(\{x\})\underset{n=0}{\overset{\mathrm{}}{}}B_n(\mu _0^2)\stackrel{~}{\mathrm{\Phi }}_n(\{x\}).$$
(64)
Again, the representation given by Eq. (64) is very convenient to analyze the analytic continuation of the (hard part of) the nucleon form factor into the timelike region of $`Q^2`$. Continuing $`Q^2q^2`$ one obtains in Eq. (61) the same shift as in Eq. (57) with $`\delta _n\alpha _s(\beta _0/4)\gamma _n`$. Specifying the particular values of the coefficients $`B_n`$ we can calculate the ratio of timelike to the spacelike form factors for several models known in the literature.
In Fig. 4 we display the ratio of the timelike to the spacelike form factors of the nucleon for three different nucleon distribution amplitudes: Chernyak–Ogloblin–Zhitnitsky, heterotic and Gari–Stefanis . From this figure, it is clear that there is no enhancement due to the analytic continuation except from a marginal factor of order 1.03 for the GS model.
## VII Soft terms for the pion form factor in the local duality approach
So far, we have discussed only the hard pQCD contributions to the hadronic form factors. But as argued in Refs. , the dominant contribution at intermediate values of the momentum transfers $`Q^210\text{GeV}^2`$ is generated by the soft contribution which involves no hard gluon exchanges. As a model for the soft contribution, we assume the local duality (LD) approximation in which it is assumed that the pion form factor is dual to the free quark spectral density
$$f_\pi ^2F_\pi (Q^2)=\frac{1}{\pi ^2}_0^{s_0}_0^{s_0}\rho _0(s,s^{},Q^2)𝑑s𝑑s^{}.$$
(65)
The latter is given by
$$\rho _0(s,s^{},t)=\frac{3}{4}\left[t^2\left(\frac{d}{dt}\right)^2+\frac{1}{3}t^3\left(\frac{d}{dt}\right)^3\right]\frac{1}{\sqrt{\left(s+s^{}+t\right)^24ss^{}}}.$$
(66)
Here the duality interval $`s_0`$ corresponds to the effective threshold for the higher excited states and the “continuum” in the channels with the axial current quantum numbers.
In principle, the value of $`s_0`$ is fixed by the ratio of the nonperturbative power corrections to the (leading) perturbative term in the OPE for the correlator. In what follows, we use the value $`s_00.7`$GeV<sup>2</sup> which has been extracted in the pioneering paper from the QCD sum rule analysis of the correlator of two axial currents. The LD prescription for this correlator just implies the relation
$$s_0=4\pi ^2f_\pi ^2,$$
(67)
between $`s_0`$ and the pion decay constant $`f_\pi `$. This relation ensures that the Ward identity for the pion form factor
$$F_\pi (0)=1$$
(68)
is satisfied within the LD approach.
Performing the integral on the rhs of Eq.(65) we get the explicit analytic expression for the pion form factor
$$F_\pi ^{\text{LD}}(Q^2)=1\frac{1+6s_0/Q^2}{\left(1+4s_0/Q^2\right)^{3/2}},$$
(69)
originally obtained in . Note that for $`t0.6\text{GeV}^2`$, expression (69) is in good agreement with existing data (see Fig. 5).
A simplified version of the LD model is based on using the “duality triangle” instead of the “duality square”. In this approach , one uses the variables $`S=s_1+s_2`$ and $`s_1s_2`$, introducing the reduced spectral density
$$\overline{\rho }_0(S,Q^2)_0^S\rho _0(Ss^{},s^{},Q^2)𝑑s^{}.$$
(70)
The LD relation, eq. (65), is then substituted by its “triangle” version (TrLD)
$$F_\pi (Q^2)F_\pi ^{\text{TrLD}}(Q^2)=\frac{1}{\pi ^2f_\pi ^2}_0^{S_0}\overline{\rho }_0(S,Q^2)𝑑S$$
(71)
with $`S_0=\sqrt{2}s_0`$. The latter condition means that the areas of the integration regions over $`s`$ and $`s^{}`$ in the two approaches are the same (see and for more details).
Using (66) and (70) we can easily calculate the spectral density $`\overline{\rho }_0(S,Q^2)`$
$$\overline{\rho }_0(S,Q^2)=\frac{S^2\left(2S+3Q^2\right)}{2\left(2S+Q^2\right)^3}$$
(72)
producing
$$F_\pi ^{\text{TrLD}}(Q^2)=\frac{1}{\sqrt{2}\left(1+Q^2/2S_0\right)^2}.$$
(73)
As one can see from the left part of Fig. 5 the difference between the two models in the region of interest ($`Q^21\text{GeV}^2`$) is very small.
### A Sudakov effects due to the electromagnetic vertex in $`F_\pi ^{\text{TrLD}}(Q^2)`$
The crucial feature of the soft contribution is that it is accompanied by the Sudakov form factor. In other words, the double logarithms $`\alpha _s\mathrm{ln}^2(Q^2)`$ do not cancel in this case.
The one-loop radiative corrections to the spectral density $`\overline{\rho }_0(S,t=Q^2)`$ has been calculated by one of us (A.P.B.) To analyze the Sudakov effects in the Feynman gauge, we need only the result for the gluon correction to the electromagnetic vertex (accompanied by the appropriate $`1/2`$-insertions of self-energies into the quark lines, which gives an UV-finite result):
$`{\displaystyle \frac{\mathrm{\Delta }^{\text{EM-vertex}}\overline{\rho }(S,t)}{\left(\alpha _sC_F/2\pi \right)\overline{\rho }_0(S,t)}}`$ $`=`$ $`2\left(\text{Li}\text{2}\left({\displaystyle \frac{S}{2S+t}}\right)\mathrm{ln}\left(1+{\displaystyle \frac{t}{S}}\right)\mathrm{ln}\left(2+{\displaystyle \frac{t}{S}}\right)\right)\left(1+{\displaystyle \frac{t^2\left(6S+5t\right)}{S^2\left(2S+3t\right)}}\right)`$ (74)
$`+`$ $`\left(\mathrm{ln}^2\left(2+{\displaystyle \frac{t}{S}}\right)\text{Li}\text{2}(1)\right)\left(1+{\displaystyle \frac{2t^2\left(6S+5t\right)}{S^2\left(2S+3t\right)}}\right)\mathrm{ln}\left(2+{\displaystyle \frac{t}{S}}\right)\left({\displaystyle \frac{1}{2}}+{\displaystyle \frac{2t\left(2S+5t\right)}{S\left(2S+3t\right)}}\right)`$ (75)
$`+`$ $`2\left(\text{Li}\text{2}\left({\displaystyle \frac{t}{2S+t}}\right)\mathrm{ln}\left(1+{\displaystyle \frac{2S}{t}}\right)\mathrm{ln}(2)\right){\displaystyle \frac{t^2\left(6S+5t\right)}{S^2\left(2S+3t\right)}}\mathrm{ln}\left(1+{\displaystyle \frac{2S}{t}}\right){\displaystyle \frac{t^2\left(42S+55t\right)}{8S^2\left(2S+3t\right)}}{\displaystyle \frac{5t}{4S}}.`$ (76)
The leading asymptotics of this expression in the large $`t`$ regime is<sup>§</sup><sup>§</sup>§A similar correction was obtained in the light-cone QCD sum rule approach .
$$\frac{\mathrm{\Delta }^{\text{EM-vertex}}\overline{\rho }(S,t)}{\overline{\rho }_0(S,t)}\underset{t\mathrm{}}{}\frac{\alpha _sC_F}{2\pi }\left[\mathrm{ln}\left(1+\frac{t}{2S}\right)\right]^2+O\left(\mathrm{ln}\frac{t}{S}\right).$$
(77)
So, we model the Sudakov corrections in the following way
$$\overline{\rho }_0^{\text{Sudakov}}(S,Q^2)\overline{\rho }_0(S,Q^2)\mathrm{exp}\left[\frac{\alpha _sC_F}{2\pi }\mathrm{ln}^2\left(1+\frac{Q^2}{2S}\right)\right].$$
(78)
The modified spectral density is then used to model the soft term corrected by the Sudakov effects
$$F_\pi ^{\text{TrLD}\text{Sudakov}}(Q^2)\frac{1}{\pi ^2f_\pi ^2}_0^{S_0}\overline{\rho }_0^{\text{Sudakov}}(S,Q^2)𝑑S.$$
(79)
On the right part of Fig. 5 we show for comparison predictions for $`Q^2F_\pi ^{\text{TrLD}}(Q^2)`$ and for $`Q^2F_\pi ^{\text{TrLD-Sudakov}}(Q^2)`$. One can see from this figure that the Sudakov effects in the electromagnetic vertex reduce (as expected) the soft contribution in the spacelike region by 6–20%.
### B Model dependence of the soft term in timelike region
As we have seen in the previous subsection for spacelike values of the momentum transfer ($`Q^2>0`$) both LD models give rather close results for the pion form factor at $`Q^21\text{GeV}^2`$. But if we analytically continue these two models into the timelike region ($`q^2=Q^2>0`$), we obtain absolutely different results for both Re$`[F_\pi (q^2)]`$ and Im$`[F_\pi (q^2)]`$:
Re$`[F_\pi ^{\text{LD}}(q^2)]`$ $`=`$ $`1\theta (q^24s_0){\displaystyle \frac{16s_0/q^2}{\left(14s_0/q^2\right)^{3/2}}},`$ (80)
Im$`[F_\pi ^{\text{LD}}(q^2)]`$ $`=`$ $`\theta (4s_0q^2){\displaystyle \frac{16s_0/q^2}{\left(4s_0/q^21\right)^{3/2}}},`$ (81)
Re$`[F_\pi ^{\text{TrLD}}(q^2)]`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}\left(1q^2/2S_0\right)^2}},`$ (82)
Im$`[F_\pi ^{\text{TrLD}}(q^2)]`$ $`=`$ $`0.`$ (83)
We see that in the resonance region ($`q^2<4\text{GeV}^2`$) the differences between these two models are rather large, and we can actually say nothing about the true behavior of $`F_\pi (q^2)`$ in this region. On the other hand, in the region $`q^26\text{GeV}^2`$ the differences between the two models are less than the experimental uncertainties, and hence we can use them, at least as a first approximation, to model $`F_\pi (q^2)`$. Furthermore, in the case of the ‘triangle LD’ we have an explicit analytic expression for the Sudakov effects which we can now continue into the timelike region
Re$`[F_\pi ^{\text{TrLD}\text{Sudakov}}(q^2)]`$ $`=`$ $`{\displaystyle \frac{1}{\pi ^2f_\pi ^2}}{\displaystyle _0^{S_0}}\overline{\rho }_{0,\text{Re}}^{\text{Sudakov}}(S,q^2)𝑑S;`$ (84)
Im$`[F_\pi ^{\text{TrLD}\text{Sudakov}}(q^2)]`$ $`=`$ $`{\displaystyle \frac{1}{\pi ^2f_\pi ^2}}{\displaystyle _0^{S_0}}\overline{\rho }_{0,\text{Im}}^{\text{Sudakov}}(S,q^2)𝑑S,`$ (85)
with
$`\overline{\rho }_{0,\text{Re}}^{\text{Sudakov}}(S,q^2)`$ $`=`$ $`\stackrel{~}{\rho }_0(S,q^2)\mathrm{cos}\left(\stackrel{~}{\alpha }_s(q^2)C_F\mathrm{ln}\left[{\displaystyle \frac{q^2}{2S}}1\right]\right);`$ (86)
$`\overline{\rho }_{0,\text{Im}}^{\text{Sudakov}}(S,q^2)`$ $`=`$ $`\stackrel{~}{\rho }_0(S,q^2)\mathrm{sin}\left(\stackrel{~}{\alpha }_s(q^2)C_F\mathrm{ln}\left[{\displaystyle \frac{q^2}{2S}}1\right]\right);`$ (87)
$`\stackrel{~}{\rho }_0(S,q^2)`$ $``$ $`\overline{\rho }_0(S,q^2)\mathrm{exp}\left[{\displaystyle \frac{\stackrel{~}{\alpha }_s(q^2)C_F}{2\pi }}\left(\mathrm{ln}^2\left[{\displaystyle \frac{q^2}{2S}}1\right]\pi ^2\right)\right].`$ (88)
(Here $`\stackrel{~}{\alpha }_s(q^2)`$ is the lowest-order model expression (19) for $`\alpha _s`$ in the timelike regime.)
Using this model we obtain results, depicted on the lhs of Fig. 7. After adding the analytically continued expression for the hard scattering (perturbative) part, including also transverse momentum effects (Sudakov+intrinsic effects) , we arrive at the result, shown on the rhs of Fig. 7.
## VIII Conclusions
In this paper, we investigated various aspects of the analytic continuation procedure from the spacelike to the timelike region of momentum transfers for several processes in QCD. We concentrated on studying several types of logarithmic contributions $`\mathrm{ln}(Q^2)`$ capable of producing $`\pm i\pi `$ in the timelike region. In the case of the ultraviolet logarithms, we reviewed the construction of the effective QCD coupling constant for the timelike region. The major result here is that the transition from a spacelike to the mirror timelike momentum only decreases the effective coupling constant. Studying the collinear logarithms, we established that in this case each eigenfunction $`\varphi _n(x)`$ of the evolution equation acquires a phase factor $`e^{i\delta _n}`$. The phase vanishes for the asymptotic wave function, and there are no changes in this most realistic situation. But even in the case of the Chernyak-Zhitnitsky wave function, the interference effects are very small and, again, they decrease rather than increase the timelike contribution compared to the spacelike one. In the case of the pion electromagnetic form factor, we emphasized that the $`\pi ^2`$ terms which may appear in the timelike region on the diagram by diagram level cancel in the total sum together with the double logarithms which generated them. Thus, we found no sources for the $`K`$-factor-type enhancements in the hard gluon exchange perturbative QCD contributions to the hadronic form factors. However, the situation completely changes if one considers the soft contribution. We investigated the simplest case of the pion electromagnetic form factor. To this end, we incorporated the local duality model suggested by the QCD sum rule studies performed earlier in the lowest (zero) order in $`\alpha _s`$. We included the $`\alpha _s`$ correction which, as expected, contains the Sudakov double logarithms. In the timelike region the latter produce non-canceling positive $`\pi ^2`$ terms which result in a $`K`$-factor-type enhancement. Our results for the soft contribution are in good agreement with existing experimental data on the pion electromagnetic form factor both in the spacelike and the timelike regions. We regard this agreement as another indication that soft contributions dominate the form factors at currently accessible momentum transfers.
## ACKNOWLEDGMENTS
This work was supported in part by the US Department of Energy under contract DE-AC05-84ER40150; Russian Foundation for Fundamental Research (grant N 00-02-1669), Heisenberg–Landau Program and by the COSY Forschungsproject Jülich/Goeke. Two of us (A.B. and A.R.) are highly indebted to Prof. Klaus Goeke for the warm hospitality in Institut für Theoretische Physik II, Ruhr-Universität Bochum, where the major part of this work was done.
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# Preheating with extra dimensions
## I Introduction
Physics in higher-dimensions has received much attention in an attempt to unify the interactions in Nature originating from Kaluza and Klein. For example, the ten-dimensional $`E_8\times E_8`$ heterotic superstring theory would be a strong candidate to describe the real world. The strong coupling limit of this theory was found to be the eleven-dimensional supergravity, which is equivalent to the low-energy effective M-theory. Since the spacetime we recognize is four-dimensional, we conventionally utilize some mechanisms of dimensional reduction assuming that extra dimensions are compactified on some manifolds (Kaluza-Klein reductions). Recently, Randall and Sundrum proposed an alternative way of compactifications based on brane models, which was originally introduced as a solution to the hierarchy problem between the weak and Planck scale. In the brane scenario, gravity works in the five-dimensional bulk of spacetime while matter fields are confined in four-dimensions. In this paper, however, we adopt conventional Kaluza-Klein reduction from a higher-dimensional action whose dimension is larger than five.
One of the most important topics which plagues such higher-dimensional theories is the stability of extra dimensions. The internal dimensions need to be held static against small fluctuations in order to settle in the present universe which would be a direct product of the four-dimensional Minkowski spacetime $`M^4`$ and a small internal space $`K^d`$. In this respect, the basic argument is to introduce a cosmological constant in the higher-dimensional action and keep extra dimensions static with the existence of some fields. In such models including the Candelas-Weinberg (CW) model (sphere compactification with a cosmological constant and the one-loop quantum correction), it has been recognized that the present vacuum is static against linear perturbations and even non-linear large perturbations. It was also found that stability of the internal space is preserved against a quantum tunneling without the topology change.
From a cosmological point of view, it is natural to ask whether the internal space is held static during an inflationary epoch. Amendola et al. considered stability of compactification in the CW model with the existence of an inflaton field $`\varphi `$. In old, new, and extended inflation models, the system exhibits semiclassical properties in which the stability is preserved as long as the transition probability for the scale of the internal space to tunnel through its potential barrier is smaller than that of inflaton. In the chaotic inflation model, the expansion of the internal space can be classically avoided if we choose the initial value of inflaton and parameters of the model appropriately. Amendola et al. also obtained the upper bound for the present scale of the internal space as $`b_{}\stackrel{<}{}\sqrt{d}\times 10^5l_{\mathrm{pl}}`$, where $`l_{\mathrm{pl}}`$ is the Planck length, by the requirement of successful chaotic inflation and stability of compactifications. As for the inflationary scenario in generalized Einstein theories including Brans-Dicke and induced gravity theories, Berkin and Maeda analyzed the dynamics of inflation in new and chaotic inflationary models in the presence of a dilaton field $`\sigma `$ with a potential $`U(\sigma )=0`$. In the context of large internal dimensions, several authors recently investigated inflation with the higher-dimensional fundamental Planck scale in the TeV region.
After the inflationary period ends, the system enters a reheating stage. It has been recognized that reheating will turn on by an explosive particle production which is called preheating. As compared with other inflationary models, the chaotic inflationary scenario typically leads to the strong amplification of a scalar field $`\chi `$ coupled to inflaton with the interaction $`\frac{1}{2}g^2\varphi ^2\chi ^2`$ due to parametric resonance during the oscillating stage of inflaton. In this scenario, a lot of works have already been done using analytic approaches based on the Mathieu or Lamé equations and numerical computations by mean field approximations or full lattice simulations. The existence of the preheating stage is important in the sense that it would affect the baryogenesis in grand unified scale, topological defect formation, nonthermal phase transition, and gravitational waves.
Recently, Mazumdar and Mendes considered preheating in generalized Einstein theories including higher-dimensional theories in the massive chaotic inflationary model $`V(\varphi )=\frac{1}{2}m^2\varphi ^2`$ with a scalar field $`\chi `$ coupled to inflaton. After dimensional reductions, we have a dilaton field $`\sigma `$ which corresponds to the radius of extra dimensions. They investigated the multi-field system of scalar fields $`\varphi `$, $`\chi `$, $`\sigma `$ in the case where dilaton does not have its own potential. It was pointed out that long-wave modes ($`k0`$) of the fluctuation of dilaton can be enhanced due to the growth of metric perturbations. Although their scheme is based on the torus compactifications which have zero curvature, compactifications on the sphere give rise to a potential term due to the curvature of the internal space. It is worth investigating whether the stability of compactification is preserved or not during preheating, when dilaton has its own potential. In this paper, we adopt the dilaton potential based on the CW model in the presence of a massive inflaton field, and analyze the evolution of scalar field fluctuations including the backreaction effect of created particles. We include metric perturbations explicitly for the evolution equations, and also investigate whether this model predicts the density perturbation observed by the Cosmic Background Explorer (COBE) satellite.
This paper is organized as follows. In the next section, we describe our model and consider the dynamics of inflation in the presence of the dilaton field. In Sec. III, we investigate the parametric amplification of field fluctuations during preheating including metric perturbations. It is particular of interest to study the evolutions of super-Hubble dilaton and metric perturbations. We present our conclusion and discussion in the final section.
## II Inflation with extra dimensions
We investigate a model in $`D=d+4`$ dimensions with a cosmological constant $`\overline{\mathrm{\Lambda }}`$ and a single scalar field $`\overline{\varphi }`$,
$`S={\displaystyle d^Dx\sqrt{\overline{g}}\left[\frac{1}{2\overline{\kappa }^2}\overline{R}2\overline{\mathrm{\Lambda }}+\overline{}(\overline{\varphi })\right]},`$ (1)
where $`\overline{\kappa }^2/8\pi \overline{G}`$ and $`\overline{R}`$ are the $`D`$-dimensional gravitational constant and a scalar curvature with respect to the $`D`$-dimensional metric $`\overline{g}_{MN}`$, respectively. The Lagrangian density for the minimally coupled $`\overline{\varphi }`$ field is written as
$`\overline{}(\overline{\varphi })={\displaystyle \frac{1}{2}}\overline{g}^{MN}_M\overline{\varphi }_N\overline{\varphi }\overline{V}(\overline{\varphi }),`$ (2)
where $`\overline{V}(\overline{\varphi })`$ is a potential of the $`\overline{\varphi }`$ field.
We compactify extra dimensions to a $`d`$-dimensional sphere $`S^d`$. Then the metric $`\overline{g}_{MN}`$ is expressed as
$`ds_D^2=\overline{g}_{MN}dx^Mdx^N=\widehat{g}_{\mu \nu }dx^\mu dx^\nu +b^2ds_d^2,`$ (3)
where $`\widehat{g}_{\mu \nu }`$ is a four-dimensional metric, $`b`$ is a scale of the $`d`$-dimensional sphere whose present value is $`b_{}`$, and $`ds_d^2`$ is a line element of the $`d`$-unit sphere. After dimensional reduction, the action $`(\text{1})`$ yields
$`S={\displaystyle d^4x\sqrt{\widehat{g}}\left(\frac{b}{b_{}}\right)^d\left[\frac{1}{2\kappa ^2}\left\{\widehat{R}+d(d1)\frac{_\mu b_\nu b}{b^2}\widehat{g}^{\mu \nu }+\frac{d(d1)}{b^2}\right\}+V_d^0\left\{\widehat{}(\widehat{\varphi })2\overline{\mathrm{\Lambda }}\right\}\right]},`$ (4)
where $`\kappa ^2/8\pi `$ is Newton’s gravitational constant which is expressed as $`\kappa ^2/8\pi =\overline{\kappa }^2/(8\pi V_d^0)`$ with the present volume of the internal space $`V_d^0`$, and $`\widehat{R}`$ is a scalar curvature with respect to $`\widehat{g}_{\mu \nu }`$. The action $`(\text{4})`$ is different from the form of the Einstein-Hilbert action due to the time-dependent term $`(b/b_{})^d`$. In order to obtain the usual form, we perform the following conformal transformation,
$`\widehat{g}_{\mu \nu }=\mathrm{exp}\left(d{\displaystyle \frac{\sigma }{\sigma _{}}}\right)g_{\mu \nu },`$ (5)
where $`\sigma `$ is the so-called dilaton field which is defined by
$`\sigma `$ $`=`$ $`\sigma _{}\mathrm{ln}\left({\displaystyle \frac{b}{b_{}}}\right),`$ (6)
$`\sigma _{}`$ $`=`$ $`\left[{\displaystyle \frac{d(d+2)}{2\kappa ^2}}\right]^{1/2}.`$ (7)
Then the four-dimensional action in the Einstein frame can be described as
$`S={\displaystyle d^4x\sqrt{g}\left[\frac{1}{2\kappa ^2}R\frac{1}{2}g^{\mu \nu }_\mu \sigma _\nu \sigma U(\sigma )\frac{1}{2}g^{\mu \nu }_\mu \varphi _\nu \varphi \mathrm{exp}\left(d\frac{\sigma }{\sigma _{}}\right)V(\varphi )\right]},`$ (8)
where $`R`$ is a scalar curvature related with $`g_{\mu \nu }`$, and a scalar field $`\varphi `$ is defined by $`\varphi =\sqrt{V_d^0}\widehat{\varphi }`$. The potential $`U(\sigma )`$ for the $`\sigma `$ field is expressed as
$`U(\sigma )={\displaystyle \frac{\overline{\mathrm{\Lambda }}}{\kappa ^2}}e^{d\sigma /\sigma _{}}{\displaystyle \frac{d(d1)}{2\kappa ^2b_{}^2}}e^{(d+2)\sigma /\sigma _{}}.`$ (9)
The second term in Eq. $`(\text{9})`$ appears due to the curvature of the internal space by compactifications on the sphere $`S^d`$. However, since the potential $`(\text{9})`$ lacks a local minimum to stabilize the dilaton field, one needs to introduce quantum correction effects which are so-called Casimir effects.
Adding a one-loop effective action which is proportional to $`e^{2(d+2)\sigma /\sigma _{}}`$ to the potential $`U(\sigma )`$ and imposing the conditions that the $`\sigma `$ field has a local minimum at $`\sigma =0`$ and its extremum is zero, $`U(\sigma )`$ can be written in the following form :
$`U(\sigma )=\alpha \left[{\displaystyle \frac{2}{d+2}}e^{2(d+2)\sigma /\sigma _{}}+e^{d\sigma /\sigma _{}}{\displaystyle \frac{d+4}{d+2}}e^{(d+2)\sigma /\sigma _{}}\right],`$ (10)
with
$`\alpha ={\displaystyle \frac{d(d1)(d+2)}{2\kappa ^2(d+4)b_{}^2}}.`$ (11)
The first, second, and third terms in $`(\text{10})`$ are due to the Casimir energy, the cosmological constant, and the curvature of the internal space, respectively. The potential $`U(\sigma )`$ in the action $`(\text{8})`$ has a local minimum at $`\sigma =0`$ and a local maximum at $`\sigma _c(>0)`$ which depends on the extra dimension $`d`$. In order to reach the final state $`\sigma =0`$ which corresponds to the present scale of the internal space $`b=b_{}`$, the initial value of $`\sigma `$ is required to be $`0<\sigma _I<\sigma _c`$ (we assume $`\sigma _I>0`$), where the subscript $`I`$ denotes the initial value. Then $`\sigma `$ evolves toward the minimum of its potential, and begins to oscillate around $`\sigma =0`$. Since extra dimensions are compactified on the sphere, this gives rise to the four-dimensional Kaluza-Klein field $`\psi _{lm}`$ whose mass is given by $`M_l^2=l(l+d1)e^{(d+2)\sigma /\sigma _{}}/b_{}^2`$. It was suggested in Ref. that Kaluza-Klein modes can be excited by the oscillation of the $`\sigma `$ field in the flat Friedmann-Robertson-Walker (FRW) background. Later, we found that catastrophic enhancement of Kaluza-Klein modes does not occur relevantly for any values of $`\sigma _I`$ and the quantum number $`l1`$. Hence we only consider the case of $`l=0`$ in this paper. For a complete study, however, we have to take into account the existence of Kaluza-Klein modes with $`l1`$.
In the presence of the inflaton field $`\varphi `$, the effective potential for the dilaton field is described by the action $`(\text{8})`$ as follows
$`U_1(\sigma ,\varphi )=\alpha \left[{\displaystyle \frac{2}{d+2}}e^{2(d+2)\sigma /\sigma _{}}+e^{d\sigma /\sigma _{}}\left\{1+{\displaystyle \frac{V(\varphi )}{\alpha }}\right\}{\displaystyle \frac{d+4}{d+2}}e^{(d+2)\sigma /\sigma _{}}\right].`$ (12)
The stability of compactification during inflation with a potential $`(\text{12})`$ in several models of inflation was analyzed in Ref. . Since we are interested in the model where strong parametric amplification of scalar fields can be expected during preheating, we adopt the quadratic potential of chaotic inflation,<sup>*</sup><sup>*</sup>*In the model of the self-coupling potential $`V(\varphi )=\frac{1}{4}\lambda \varphi ^4`$, we will give some discussions in the final section.
$`V(\varphi )={\displaystyle \frac{1}{2}}m^2\varphi ^2.`$ (13)
In Fig. 1, we depict the effective potential $`(\text{12})`$ with $`(\text{13})`$. As was pointed out in Ref. , $`U_1(\sigma ,\varphi )`$ has either two local extrema or no local extrema for a fixed value of $`\varphi `$. When $`\varphi `$ is smaller than some critical value $`\varphi _c`$, the potential barrier which prevents the $`\sigma `$ field from going toward infinity exists, and the scalar field evolves toward the potential minimum at $`\varphi =\sigma =0`$. However, when $`\varphi >\varphi _c`$, this barrier disappears and the internal space grows without limit. The critical value $`\varphi _c`$ can be obtained by solving the equation $`U_1/\sigma =^2U_1/\sigma ^2=0`$ as
$`\varphi _c^2={\displaystyle \frac{2\alpha }{m^2}}\left[\left(1+{\displaystyle \frac{2}{d}}\right)\left({\displaystyle \frac{1}{2}}\right)^{2/(d+2)}1\right].`$ (14)
In order to result in the present vacuum $`\sigma =0`$, the inflaton $`\varphi `$ is constrained as
$`\varphi ^2<\varphi _c^2={\displaystyle \frac{d(d1)(d+2)}{8\pi (d+4)}}\left[\left(1+{\displaystyle \frac{2}{d}}\right)\left({\displaystyle \frac{1}{2}}\right)^{2/(d+2)}1\right]\left({\displaystyle \frac{m_{\mathrm{pl}}}{m}}\right)^2{\displaystyle \frac{1}{b_{}^2}},`$ (15)
where we have used Eq. $`(\text{11})`$. In the chaotic inflationary scenario, the initial value of inflaton is required to be $`\varphi _I\stackrel{>}{}3m_{\mathrm{pl}}`$ in order to obtain the number of e-foldings greater than 60. Further, the density perturbation observed by the COBE satellite constrains the coupling of inflaton as $`m10^6m_{\mathrm{pl}}`$. Then the condition $`(\text{15})`$ leads to the following bound for the present value of the internal space:
$`b_{}^2\stackrel{<}{}{\displaystyle \frac{d(d1)(d+2)}{72\pi (d+4)}}\left[\left(1+{\displaystyle \frac{2}{d}}\right)\left({\displaystyle \frac{1}{2}}\right)^{2/(d+2)}1\right]{\displaystyle \frac{10^{12}}{m_{\mathrm{pl}}^2}}.`$ (16)
For example, when $`d=2`$ and $`d=6`$, $`b_{}\stackrel{<}{}5\times 10^4/m_{\mathrm{pl}}`$ and $`b_{}\stackrel{<}{}1.1\times 10^5/m_{\mathrm{pl}}`$, respectively. For large values of $`d`$, Eq. $`(\text{16})`$ reads
$`b_{}\stackrel{<}{}\sqrt{d}\times 10^5/m_{\mathrm{pl}}.`$ (17)
As was suggested in Ref. , this value is by about ten orders of magnitude smaller than the experimental bound $`b_{}\stackrel{<}{}10^{17}/m_{\mathrm{pl}}`$. It is worth mentioning that such theoretical bound is derived by the analysis of stability of compactification during inflation. In what follows, we use the values of $`b_{}`$ which satisfy the relation $`(\text{16})`$.
Let us consider the dynamics of inflation in the presence of the $`\sigma `$ field. In the flat FRW background: $`ds^2=dt^2+a^2(t)d𝐱^\mathrm{𝟐}`$, we find that the homogeneous parts of scalar fields and the scale factor satisfy the following equations of motion by the action $`(\text{8})`$:
$`\ddot{\varphi }+3H\dot{\varphi }+e^{d\sigma /\sigma _{}}V^{}(\varphi )=0,`$ (18)
$`\ddot{\sigma }+3H\dot{\sigma }+U^{}(\sigma ){\displaystyle \frac{d}{\sigma _{}}}e^{d\sigma /\sigma _{}}V(\varphi )=0,`$ (19)
$`H^2\left({\displaystyle \frac{\dot{a}}{a}}\right)^2={\displaystyle \frac{\kappa ^2}{3}}\left[{\displaystyle \frac{1}{2}}\dot{\varphi }^2+e^{d\sigma /\sigma _{}}V(\varphi )+{\displaystyle \frac{1}{2}}\dot{\sigma }^2+U(\sigma )\right],`$ (20)
where $`H`$ is the Hubble expansion rate.
When the initial value of inflaton is larger than $`\varphi _c`$, the last term in the l.h.s. of Eq. $`(\text{19})`$ dominates over the third term and the $`\sigma `$ field evolves toward $`\sigma =\mathrm{}`$ (see Fig. 2). This is the situation we want to avoid. For the values of $`\varphi <\varphi _c`$, there exists a local minimum at $`\sigma =\sigma _1`$ and a local maximum at $`\sigma =\sigma _2`$ with $`\sigma _2>\sigma _1>0`$. As long as the initial value of $`\sigma `$ exists in the range of $`\sigma <\sigma _2`$, the $`\sigma `$ field evolves toward the potential minimum at $`\sigma =\sigma _1`$. The value of $`\sigma _1`$ decreases to zero as the $`\varphi `$ field moves toward the potential minimum at $`\varphi =0`$. For the initial values of $`\sigma `$ and $`\varphi `$ which are finally trapped in the potential minimum at $`\sigma =\varphi =0`$, one may consider that the dynamics of inflation is altered in the presence of the $`\sigma `$ field. In this case, however, we can numerically confirm that the third term in Eq. $`(\text{19})`$ rapidly makes the $`\sigma `$ field shift toward the local minimum at $`\sigma =\sigma _1`$ for the values of $`b_{}`$ which satisfy the condition of Eq. $`(\text{16})`$. Then $`\sigma `$ begins to roll down along the valley of $`\sigma =\sigma _1`$, and decreases as inflaton approaches its potential minimum. This behavior is found in Fig. 2. The dynamics of inflation is hardly affected by the presence of the $`\sigma `$ field, and the system can be effectively described by one scalar field $`\varphi `$.
In Fig. 3, we plot the evolution of both the inflaton field and the number of e-foldings $`N\mathrm{ln}(a/a_I)`$ during inflation. After the rapid decrease of $`\sigma `$ at the initial stage, inflaton slowly rolls down along its potential in the usual manner, which results in sufficient inflation to solve several cosmological puzzles. We also find that the number of e-foldings exceeds $`N60`$ for the initial value of $`\varphi _I=3m_{\mathrm{pl}}`$. The inflationary period ends when inflaton decreases to $`\varphi 0.2m_{\mathrm{pl}}`$, which corresponds to the time $`mt20`$ in Fig. 3. In the next section, we investigate the dynamics of field and metric perturbations during preheating.
## III Preheating with extra dimensions
After inflation, the system enters the preheating stage during which fluctuations of scalar fields will grow by parametric resonance. We introduce another massless scalar field $`\chi `$ coupled to inflaton, and adopt the following modified potential instead of $`(\text{13})`$:
$`V(\varphi ,\chi )={\displaystyle \frac{1}{2}}m^2\varphi ^2+{\displaystyle \frac{1}{2}}g^2\varphi ^2\chi ^2+\stackrel{~}{g}^2\varphi ^3\chi .`$ (21)
Note that we include the interaction term $`\stackrel{~}{g}^2\varphi ^3\chi `$ which often appears in supergravity models in addition to the standard term $`\frac{1}{2}g^2\varphi ^2\chi ^2`$. This provides a way to escape from an inflationary suppression of the $`\chi `$ field as we will show later.
When we consider fluctuations of scalar fields, metric perturbations should be also taken into account for a consistent study of preheating. In fact, inclusion of metric perturbations can change the evolution of field fluctuations significantly in broad classes of models. In this paper, we adopt the perturbed metric in the longitudinal gauge in the flat FRW background:
$`ds^2=(1+2\mathrm{\Phi })dt^2+a^2(t)(12\mathrm{\Psi })\delta _{ij}dx^idx^j,`$ (22)
where $`\mathrm{\Phi }`$ and $`\mathrm{\Psi }`$ are gauge-invariant potentials.
Decomposing scalar fields into homogeneous and gauge-invariant fluctuational parts as $`\varphi (t,𝐱)\varphi (t)+\delta \varphi (t,𝐱)`$, $`\sigma (t,𝐱)\sigma (t)+\delta \sigma (t,𝐱)`$, $`\chi (t,𝐱)\chi (t)+\delta \chi (t,𝐱)`$, and expanding scalar field fluctuations and metric fluctuations by Fourier modes, we obtain the following perturbed equations (see e.g.,):
$`\mathrm{\Phi }_k=\mathrm{\Psi }_k,`$ (23)
$`\dot{\mathrm{\Phi }}_k+H\mathrm{\Phi }_k={\displaystyle \frac{\kappa ^2}{2}}(\dot{\varphi }\delta \varphi _k+\dot{\sigma }\delta \sigma _k+\dot{\chi }\delta \chi _k),`$ (24)
$`3H\dot{\mathrm{\Phi }}_k+\left[{\displaystyle \frac{k^2}{a^2}}+3H^2{\displaystyle \frac{\kappa ^2}{2}}(\dot{\varphi }^2+\dot{\sigma }^2+\dot{\chi }^2)\right]\mathrm{\Phi }_k`$ (25)
$`=`$ $`{\displaystyle \frac{\kappa ^2}{2}}\left(\dot{\varphi }\delta \dot{\varphi }_k+U_{1,\varphi }\delta \varphi _k+\dot{\sigma }\delta \dot{\sigma }_k+U_{1,\sigma }\delta \sigma _k+\dot{\chi }\delta \dot{\chi }_k+U_{1,\chi }\delta \chi _k\right),`$ (26)
$`\delta \ddot{\varphi }_k+3H\delta \dot{\varphi }_k+\left[{\displaystyle \frac{k^2}{a^2}}+e^{d\sigma /\sigma _{}}(m^2+g^2\chi ^2+6\stackrel{~}{g}^2\varphi \chi )\right]\delta \varphi _k=4\dot{\varphi }\dot{\mathrm{\Phi }}_k+2(\ddot{\varphi }+3H\dot{\varphi })\mathrm{\Phi }_kU_{1,\sigma \varphi }\delta \sigma _kU_{1,\chi \varphi }\delta \chi _k,`$ (27)
$`\delta \ddot{\sigma }_k+3H\delta \dot{\sigma }_k+\left[{\displaystyle \frac{k^2}{a^2}}+U^{\prime \prime }(\sigma )+{\displaystyle \frac{d^2}{\sigma _{}^2}}e^{d\sigma /\sigma _{}}\left({\displaystyle \frac{1}{2}}m^2\varphi ^2+{\displaystyle \frac{1}{2}}g^2\varphi ^2\chi ^2+\stackrel{~}{g}^2\varphi ^3\chi \right)\right]\delta \sigma _k`$ (28)
$`=`$ $`4\dot{\sigma }\dot{\mathrm{\Phi }}_k+2(\ddot{\sigma }+3H\dot{\sigma })\mathrm{\Phi }_kU_{1,\varphi \sigma }\delta \varphi _kU_{1,\chi \sigma }\delta \chi _k,`$ (29)
$`\delta \ddot{\chi }_k`$ $`+`$ $`3H\delta \dot{\chi }_k+\left({\displaystyle \frac{k^2}{a^2}}+g^2\varphi ^2e^{d\sigma /\sigma _{}}\right)\delta \chi _k=4\dot{\chi }\dot{\mathrm{\Phi }}_k+2(\ddot{\chi }+3H\dot{\chi })\mathrm{\Phi }_kU_{1,\sigma \chi }\delta \sigma _kU_{1,\varphi \chi }\delta \varphi _k,`$ (30)
where $`U_{1,\varphi \sigma }`$, $`U_{1,\chi \sigma }`$, and $`U_{1,\chi \varphi }`$ are expressed as $`U_{1,\varphi \sigma }=de^{d\sigma /\sigma _{}}(m^2+g^2\chi ^2+3\stackrel{~}{g}^2\varphi \chi )\varphi /\sigma _{}`$, $`U_{1,\chi \sigma }=de^{d\sigma /\sigma _{}}(g^2\chi +\stackrel{~}{g}^2\varphi )\varphi ^2/\sigma _{}`$, and $`U_{1,\chi \varphi }=e^{d\sigma /\sigma _{}}(2g^2\varphi \chi +3\stackrel{~}{g}^2\varphi ^2)`$, respectively.
The relation $`(\text{23})`$ indicates that the anisotropic stress vanishes at linear order. Eliminating the $`\dot{\mathrm{\Phi }}_k`$ term in Eqs. $`(\text{24})`$ and $`(\text{26})`$, we find
$`\left({\displaystyle \frac{k^2}{a^2}}{\displaystyle \frac{\kappa ^2}{2}}{\displaystyle \underset{J}{}}\dot{\phi }_J^2\right)\mathrm{\Phi }_k={\displaystyle \frac{\kappa ^2}{2}}{\displaystyle \underset{J}{}}\left(\dot{\phi }_J\delta \dot{\phi }_{Jk}+3H\dot{\phi }_J\delta \phi _{Jk}+U_{1,\phi _J}\delta \phi _{Jk}\right),`$ (31)
where $`\phi _J(J=1,2,3)`$ correspond to the scalar fields $`\varphi `$, $`\sigma `$, $`\chi `$, respectively. Eq. $`(\text{31})`$ shows that metric perturbations are known when evolutions of scalar fields are determined. When field fluctuations are amplified, it is expected that this stimulates the growth of metric perturbations by Eq. $`(\text{31})`$. The enhancement of metric perturbations will also assist the excitation of field perturbations as is found by Eqs. $`(\text{27})`$-$`(\text{30})`$.
Parametric amplification of field fluctuations affects evolutions of the background quantities. Since the $`\chi `$ fluctuation generally grows faster than other field fluctuations, we include this contribution in the background equations as
$`\ddot{\varphi }+3H\dot{\varphi }+e^{d\sigma /\sigma _{}}(m^2+g^2\chi ^2+3\stackrel{~}{g}^2\varphi \chi )\varphi =0,`$ (32)
$`\ddot{\sigma }+3H\dot{\sigma }+U^{}(\sigma ){\displaystyle \frac{d}{\sigma _{}}}e^{d\sigma /\sigma _{}}\left({\displaystyle \frac{1}{2}}m^2\varphi ^2+{\displaystyle \frac{1}{2}}g^2\varphi ^2\chi ^2+\stackrel{~}{g}^2\varphi ^3\chi \right)=0,`$ (33)
$`\ddot{\chi }+3H\dot{\chi }+e^{d\sigma /\sigma _{}}\left(g^2\varphi ^2\chi +\stackrel{~}{g}^2\varphi ^3\right)=0,`$ (34)
$`H^2={\displaystyle \frac{\kappa ^2}{3}}\left[{\displaystyle \frac{1}{2}}\dot{\sigma }^2+U(\sigma )+{\displaystyle \frac{1}{2}}\dot{\varphi }^2+{\displaystyle \frac{1}{2}}\dot{\chi }^2+e^{d\sigma /\sigma _{}}\left({\displaystyle \frac{1}{2}}m^2\varphi ^2+{\displaystyle \frac{1}{2}}g^2\varphi ^2\chi ^2+\stackrel{~}{g}^2\varphi ^3\chi \right)\right],`$ (35)
where the spatial average of the $`\chi `$ fluctuation is defined by
$`\chi ^2={\displaystyle \frac{1}{2\pi ^2}}{\displaystyle k^2|\delta \chi _k|^2𝑑k}.`$ (36)
Let us examine evolutions of the background fields and the scale factor. As is found in the previous section, the $`\sigma `$ field rapidly decreases during inflation compared with the $`\varphi `$ field, and the condition $`\sigma \sigma _{}=\sqrt{d(d+2)/16\pi }m_{\mathrm{pl}}`$ holds at the beginning of preheating. Then, in the stage where the $`\chi `$ fluctuation is not significantly enhanced, Eqs. $`(\text{32})`$-$`(\text{35})`$ can be approximately written as
$`\ddot{\varphi }+3H\dot{\varphi }+m^2\varphi =0,`$ (37)
$`\ddot{\sigma }+3H\dot{\sigma }+{\displaystyle \frac{2(d1)}{b_{}^2}}\sigma =0,`$ (38)
$`\ddot{\chi }+3H\dot{\chi }+g^2\varphi ^2\chi +\stackrel{~}{g}^2\varphi ^3=0,`$ (39)
$`H^2={\displaystyle \frac{\kappa ^2}{3}}\left({\displaystyle \frac{1}{2}}\dot{\varphi }^2+{\displaystyle \frac{1}{2}}m^2\varphi ^2\right).`$ (40)
Making use of the time-averaged relation $`\frac{1}{2}\dot{\varphi }^2_T=\frac{1}{2}m^2\varphi ^2_T`$ during the oscillating stage of inflaton, the evolution of inflaton is analytically expressed by Eqs. $`(\text{37})`$ and $`(\text{40})`$ as
$`\varphi =\mathrm{\Phi }(t)\mathrm{sin}mt,\mathrm{with}\mathrm{\Phi }(t)={\displaystyle \frac{m_{\mathrm{pl}}}{\sqrt{3\pi }mt}}.`$ (41)
The coherent oscillation of inflaton begins when $`\mathrm{\Phi }(t_i)0.2m_{\mathrm{pl}}`$, and we set the initial time as $`mt=\pi /2`$ as in Ref. . The scale factor evolves as $`at^{2/3}`$ since the system is dominated by the oscillation of the massive inflaton field. Although the $`\sigma `$ field oscillates with a frequency $`\sqrt{2(d1)}/b_{}`$, its amplitude is very small relative to that of the $`\varphi `$ field. For example, in the simulation of Fig. 2, the amplitude of $`\sigma `$ at the start of preheating is found to be about $`10^5m_{\mathrm{pl}}`$.
If we neglect metric perturbations, the evolution of the $`\delta \chi _k`$ fluctuation can be studied analytically at the linear stage of preheating. Ignoring the r.h.s. of Eq. $`(\text{30})`$ and introducing a new scalar field $`\delta X_k=a^{3/2}\delta \chi _k`$, Eq. $`(\text{30})`$ reads
$`{\displaystyle \frac{d^2}{dt^2}}\delta X_k+\left[{\displaystyle \frac{k^2}{a^2}}+g^2\varphi ^2{\displaystyle \frac{3}{4}}\left({\displaystyle \frac{2\ddot{a}}{a}}+{\displaystyle \frac{\dot{a}^2}{a^2}}\right)\right]\delta X_k=0.`$ (42)
The last term in Eq. $`(\text{42})`$ which corresponds to the pressure term can be neglected during the oscillating stage of inflaton. Then Eq. $`(\text{42})`$ is reduced to the well-known Mathieu equation,
$`{\displaystyle \frac{d^2}{dz^2}}\delta X_k+\left(A_k2q\mathrm{cos}2z\right)\delta X_k=0,`$ (43)
where $`z=mt`$ and
$`A_k=2q+{\displaystyle \frac{k^2}{(ma)^2}},q={\displaystyle \frac{g^2\mathrm{\Phi }^2(t)}{4m^2}}.`$ (44)
Then the value of $`q`$ at the beginning of preheating is estimated as
$`q_i10^{10}\times g^2,`$ (45)
where we used $`\mathrm{\Phi }(t_i)0.2m_{\mathrm{pl}}`$ and $`m10^6m_{\mathrm{pl}}`$. The coupling $`g\stackrel{<}{}10^5`$ yields $`q_i\stackrel{<}{}1`$, which is generally called the narrow resonance. In this case, parametric resonance is weak in an expanding universe. However, when $`q_i1`$, it was pointed out in Ref. that the $`\chi `$ particle production can be efficient in spite of the decrease of $`q`$ due to cosmic expansion, which was later confirmed by numerical calculations in Ref. . In this case, the $`\delta \chi _k`$ field initially lies in the broad resonance regime as long as $`k`$ is not so large relative to $`ma`$, and it jumps over many stability and instability bands with the decrease of $`q`$. This was termed stochastic resonance in Ref. , in which the $`\delta \chi _k`$ fluctuation increases stochastically overcoming the diluting effect by the expansion of the universe. For $`g\stackrel{>}{}3\times 10^4`$, the backreaction effect of created $`\chi `$ particles becomes important, which results in the termination of parametric resonance. In this case, since the coherent oscillation of the inflaton field is broken by the growth of the $`\delta \chi _k`$ fluctuation, the analytical method based on the Mathieu equation is no longer applied. In this respect, several numerical works have been done by making use of mean field approximations or fully nonlinear calculations. In the Hartree approximation the final variance of the $`\chi `$ field is estimated as $`\chi ^2_fq^{1/2}`$, while in the fully nonlinear calculations it was found to be $`\chi ^2_fq^1`$ for the case of $`q1`$.
As for the $`\delta \sigma _k`$ field, there exist resonance terms in the l.h.s. of Eq. $`(\text{29})`$ which may lead to the enhancement of the fluctuation of dilaton. If we neglect the effect of metric perturbations in the r.h.s. of Eq. $`(\text{29})`$ and making use of the relation $`|\sigma |\sigma _{}`$ during preheating, the equation of the $`\delta \sigma _k`$ field can be approximately written as
$`\delta \ddot{\sigma }_k+3H\delta \dot{\sigma }_k+\left[{\displaystyle \frac{k^2}{a^2}}+{\displaystyle \frac{2(d1)}{b_{}^2}}+{\displaystyle \frac{8\pi d}{d+2}}\left({\displaystyle \frac{m}{m_{\mathrm{pl}}}}\right)^2\varphi ^2\right]\delta \sigma _k=0.`$ (46)
Defining a new scalar field $`\delta \mathrm{\Sigma }_k=a^{3/2}\delta \sigma _k`$ and ignoring the contribution from the pressure term, we obtain
$`{\displaystyle \frac{d^2}{dz^2}}\delta \mathrm{\Sigma }_k+\left(A_k2q\mathrm{cos}2z\right)\delta \mathrm{\Sigma }_k=0,`$ (47)
where
$`A_k=2q+{\displaystyle \frac{2(d1)}{b_{}^2m^2}}+{\displaystyle \frac{k^2}{(ma)^2}},q={\displaystyle \frac{2d}{3(d+2)z^2}}.`$ (48)
For $`d1`$, $`q_i\stackrel{<}{}1`$ at the beginning of preheating. Moreover, since $`q`$ decreases as $`qt^2`$, parametric resonance is very weak. Namely, in the unperturbed metric, analytic estimates indicate that the dilaton fluctuation does not grow during preheating. We also find from Eq. $`(\text{27})`$ that the enhancement of the inflaton fluctuation can not be expected in the absence of metric perturbations.
Let us proceed to the case where metric perturbations are taken into account. From Eq. $`(\text{24})`$, we can expect that the growth of the $`\delta \chi _k`$ field will enhance metric perturbations. On the other hand, it was pointed out in Refs. that the amplitude of super-Hubble fluctuations in the $`\delta \chi _k`$ field is severely damped during inflation in the case where $`g\varphi `$ is much larger than the Hubble expansion rate $`H`$ with a model of $`V(\varphi ,\chi )=\frac{1}{2}m^2\varphi ^2+\frac{1}{2}g^2\varphi ^2\chi ^2`$. Later, Bassett et al. showed that inclusion of the interaction $`\stackrel{~}{g}^2\varphi ^3\chi `$ protects super-Hubble $`\delta \chi _k`$ fluctuations from being suppressed. In what follows, we will consider both cases of $`\stackrel{~}{g}=0`$ and $`\stackrel{~}{g}0`$ separately.
### A Case of $`\stackrel{~}{g}=0`$
Let us first estimate the amplitude of super-Hubble $`\delta \chi _k`$ modes at the beginning of preheating. When $`\stackrel{~}{g}=0`$, the adiabatic solution for $`\delta \chi _k`$ during inflation is expressed as
$`\delta \chi _k={\displaystyle \frac{a^{3/2}}{\sqrt{2\omega _k}}}e^{i\omega _kt},`$ (49)
where $`\omega _k^2=k^2/a^2+g^2\varphi ^2e^{d\sigma /\sigma _{}}k^2/a^2+g^2\varphi ^2`$. In order to lead to efficient $`\chi `$ particle production, the resonance parameter is required to be $`q=g^2\varphi ^2/4m^21`$. In this case, the effective mass of the $`\delta \chi _k`$ field is much larger than the Hubble expansion rate $`Hm`$ during inflation. Then the amplitude of the super-Hubble $`\delta \chi _k`$ field for modes relevant for structure formation is estimated as $`|\delta \chi _k|a^{3/2}/\sqrt{g\varphi }`$, which exponentially decreases during inflation. On the other hand, since the effective mass of the $`\delta \varphi _k`$ field in the l.h.s. of Eq. $`(\text{27})`$ is comparable to the Hubble rate $`H`$, the super-Hubble inflaton fluctuation is not affected by the suppression during inflation.
As for the $`\delta \sigma _k`$ field, its effective mass for super-Hubble modes is given by Eq $`(\text{46})`$ as
$`m_{\mathrm{eff}}^2=\left[{\displaystyle \frac{2(d1)}{b_{}^2m^2}}+{\displaystyle \frac{8\pi d}{d+2}}\left({\displaystyle \frac{\varphi }{m_{\mathrm{pl}}}}\right)^2\right]m^2.`$ (50)
For the typical scale $`b_{}`$ which is determined by Eq. $`(\text{16})`$ and the initial value of inflaton $`\varphi _I\stackrel{>}{}3m_{\mathrm{pl}}`$, $`m_{\mathrm{eff}}`$ is estimated as $`m_{\mathrm{eff}}^2\stackrel{>}{}100m^210H^2`$. Hence the super-Hubble $`\delta \sigma _k`$ fluctuation will be also affected by the inflationary suppression, which is relevant for small $`b_{}`$ and large initial values of $`\varphi `$. We can roughly estimate the amplitude of super-Hubble $`\delta \sigma _k`$ modes by Eq. $`(\text{29})`$ during inflation. Neglecting the contributions of the $`\chi `$ field and the time derivative terms of $`\sigma `$ and $`\delta \sigma _k`$, we obtain the amplitude of $`\delta \sigma _k`$ at the start of preheating as
$`|\delta \sigma _k(t_i)|{\displaystyle \frac{2}{(d1)}}\sqrt{{\displaystyle \frac{\pi d}{d+2}}}{\displaystyle \frac{\varphi }{m_{\mathrm{pl}}}}(b_{}m)^2|\delta \varphi _k(t_i)|.`$ (51)
For example, when $`d=6`$, since $`b_{}`$ is constrained as $`b_{}\stackrel{<}{}1.1\times 10^5/m_{\mathrm{pl}}`$, Eq. $`(\text{51})`$ yields
$`|\delta \sigma _k(t_i)|\stackrel{<}{}10^3|\delta \varphi _k(t_i)|.`$ (52)
This means that the suppression effect of super-Hubble modes is weak compared with the $`\delta \chi _k`$ field as long as $`b_{}`$ is not much smaller than its upper bound.
Let us estimate the impact on metric perturbations by the growth of field fluctuations. First, we introduce the power spectrum of $`\mathrm{\Phi }_k`$:
$`𝒫(k)={\displaystyle \frac{k^3}{2\pi ^2}}|\mathrm{\Phi }_k|^2={\displaystyle \frac{|\stackrel{~}{\mathrm{\Phi }}_k|^2}{2\pi ^2}},`$ (53)
where $`\stackrel{~}{\mathrm{\Phi }}_kk^{3/2}\mathrm{\Phi }_k`$. Defining new scalar fields $`\stackrel{~}{\phi }_J\phi _J/m_{\mathrm{pl}}`$ and $`\delta \stackrel{~}{\phi }_{Jk}k^{3/2}\delta \phi _{Jk}/m_{\mathrm{pl}}`$ ($`J=1,2,3`$), we obtain the following relation from Eq. $`(\text{24})`$:
$`\dot{\stackrel{~}{\mathrm{\Phi }}}_k+H\stackrel{~}{\mathrm{\Phi }}_k=4\pi (\dot{\stackrel{~}{\varphi }}\delta \stackrel{~}{\varphi }_k+\dot{\stackrel{~}{\sigma }}\delta \stackrel{~}{\sigma }_k+\dot{\stackrel{~}{\chi }}\delta \stackrel{~}{\chi }_k).`$ (54)
For super-Hubble modes $`kaH`$, the amplitude of $`\delta \stackrel{~}{\chi }_k`$ is written as $`|\delta \stackrel{~}{\chi }_k|\overline{k}^{3/2}(m/m_{\mathrm{pl}})\sqrt{m/(2g\varphi )}`$ where $`\overline{k}k/(ma_i)`$ with $`a_i`$ the scale factor at the onset of preheating. Since the cosmological modes correspond to $`\overline{k}e^{60}10^{26}`$, $`|\delta \stackrel{~}{\chi }_k|`$ is estimated as $`|\delta \stackrel{~}{\chi }_k|\stackrel{<}{}10^{45}`$ for the broad resonance case $`q1`$. The homogeneous part of the $`\chi `$ field is also affected by this strong suppression \[see Eq. $`(\text{34})`$ with $`\stackrel{~}{g}=0`$\]. At the beginning of preheating, the $`\dot{\chi }\delta \chi _k`$ term in the r.h.s. of Eq. $`(\text{24})`$ is very small relative to the $`\dot{\varphi }\delta \varphi _k`$ term. Although super-Hubble $`\delta \chi _k`$ fluctuations exhibit parametric amplification during preheating, it increases only by the factors $`10^410^5`$ for the coupling of $`g=3\times 10^410^2`$. Hence we can expect that the excitement of the $`\delta \chi _k`$ fluctuation hardly affects the evolution of super-Hubble metric perturbations by analytic estimates.
We are also concerned with whether super-Hubble $`\delta \sigma _k`$ fluctuations are enhanced or not during preheating. Although the inflationary suppression for $`\delta \sigma _k`$ is not so significant as compared with the $`\delta \chi _k`$ case, we have to keep in mind that $`\delta \sigma _k`$ can not be enhanced unless metric perturbations are taken into account. In order to amplify super-Hubble metric perturbations, we generally require some scalar fields such as $`\chi `$ which exhibit parametric amplification even in the absence of metric perturbations. However, when $`\stackrel{~}{g}=0`$, the $`\delta \chi _k`$ fluctuation in large scales is strongly suppressed. It is expected that the super-Hubble fluctuation of dilaton will be held static during preheating, because the $`\sigma `$ field will not play a dominant role to stimulate the enhancement of metric perturbations.
In order to verify the above estimates, we numerically solved perturbed equations $`(\text{23})`$-$`(\text{30})`$ along with background equations $`(\text{32})`$-$`(\text{35})`$. In Fig. 4, we plot the evolutions of field perturbations $`\delta \stackrel{~}{\varphi }_k`$, $`\delta \stackrel{~}{\sigma }_k`$, $`\delta \stackrel{~}{\chi }_k`$, and the metric perturbation $`\stackrel{~}{\mathrm{\Phi }}_k`$ during inflation and preheating for a super-Hubble mode $`\overline{k}=10^{26}`$ with $`g=1.0\times 10^3`$, $`d=6`$, $`b_{}=1.0\times 10^5/m_{\mathrm{pl}}`$, and $`m=10^6m_{\mathrm{pl}}`$. The initial values of homogeneous scalar fields are chosen as $`\varphi _I=3.0m_{\mathrm{pl}}`$, $`\sigma _I=0.1m_{\mathrm{pl}}`$, and $`\chi _I=1.0\times 10^3m_{\mathrm{pl}}`$. As for the initial field perturbations, we take
$`|\delta \phi _{Jk}|={\displaystyle \frac{1}{\sqrt{2\omega _{Jk}}}},|\delta \dot{\phi }_{Jk}|=\omega _{Jk}|\delta \phi _{Jk}|,`$ (55)
where $`\omega _{Jk}^2k^2/a^2+m_{\phi _J}^2`$ ($`J=1,2,3)`$ with $`m_{\phi _J}`$ is the effective mass of the each scalar field in the l.h.s. of Eqs. $`(\text{27})`$-$`(\text{30})`$.
In Fig. 4, the $`\delta \chi _k`$ fluctuation is exponentially suppressed during inflation ($`0<mt\stackrel{<}{}20`$), yielding $`\delta \stackrel{~}{\chi }_k10^{45}`$ at the beginning of preheating. Although $`\delta \chi _k`$ is enhanced by parametric resonance for $`mt\stackrel{>}{}20`$, the final value is very small as $`\delta \stackrel{~}{\chi }_k10^{40}`$. The amplitude of $`\delta \sigma _k`$ is by three orders of magnitude smaller than that of $`\delta \varphi _k`$ as is analytically estimated by Eq. $`(\text{51})`$ at the end of inflation. We also find in Fig. 4 that both of super-Hubble $`\delta \sigma _k`$ and $`\delta \varphi _k`$ fluctuations do not grow during preheating, which means that field fluctuations are not assisted by the presence of metric perturbations. As is found in the inset of Fig. 4, super-Hubble metric perturbations remain almost constant during preheating. In the case of $`\stackrel{~}{g}=0`$, the usual prediction of the inflationary spectrum in large scales is not likely to be modified, and the analysis neglecting metric perturbations gives almost the same results as compared with the perturbed metric case. As a result, the fluctuation of dilaton on super-Hubble scales can not be amplified during preheating.
Let us next consider smaller scales which are within the Hubble radius at the beginning of preheating. Since this corresponds to the modes $`k\stackrel{>}{}a_iH_i`$, the condition $`k^2/a^2>g^2\varphi ^2`$ holds in most stage of inflation. Hence the amplitude of these modes during inflation is approximately expressed by Eq. $`(\text{49})`$ as
$`|\delta \chi _k|{\displaystyle \frac{1}{a\sqrt{2k}}}.`$ (56)
The r.h.s. of Eq. $`(\text{56})`$ decreases slower compared with the super-Hubble modes. In fact, for the modes of $`k^2/a_i^2\stackrel{>}{}g^2\varphi _i^2`$ at the start of preheating, we obtain
$`|\delta \stackrel{~}{\chi }_k(t_i)|{\displaystyle \frac{1}{\sqrt{2}}}{\displaystyle \frac{m}{m_{\mathrm{pl}}}}\overline{k}.`$ (57)
For sub-Hubble modes $`\overline{k}=k/(ma_i)\stackrel{>}{}1`$, we find $`|\delta \stackrel{~}{\chi }_k(t_i)|\stackrel{>}{}10^6`$, which is much larger than in the super-Hubble case. However, since the homogeneous part of the $`\chi `$ field is severely damped, the $`\dot{\chi }\delta \chi _k`$ term in the r.h.s. of Eq. $`(\text{24})`$ is still much smaller than the $`\dot{\varphi }\delta \varphi _k`$ term at the beginning of preheating. This indicates that parametric amplification of the $`\delta \chi _k`$ fluctuation will not lead to the growth of sub-Hubble metric perturbations. We have numerically confirmed that metric perturbations and the $`\delta \sigma _k`$ fluctuation for sub-Hubble modes $`1\stackrel{<}{}k/(ma_i)\stackrel{<}{}100`$ are not relevantly enhanced for the coupling of $`3\times 10^4\stackrel{<}{}g\stackrel{<}{}10^2`$ (see Fig. 5). However, we have to caution that including the second order metric backreaction effect will lift the homogeneous $`\chi `$ field, which may assist the enhancement of metric and field fluctuations in sub-Hubble modes as was mentioned in Ref. . The full backreaction issues are left for the future work.
In the case of $`\stackrel{~}{g}=0`$, we argue that the dilaton fluctuation as well as metric perturbations in both super- and sub-Hublle scales can not be strongly amplified during preheating.
### B Case of $`\stackrel{~}{g}0`$
If the coupling $`\stackrel{~}{g}^2\varphi ^3\chi `$ is taken into account, the suppression of the super-Hubble $`\delta \chi _k`$ fluctuation can be avoided. For super-Hubble modes, neglecting derivative terms of $`\delta \chi _k`$ and $`\chi `$ fields as well as the $`\delta \sigma _k`$ term in Eq. $`(\text{30})`$, we find the following relation at the end of inflation:
$`\delta \chi _k3\left({\displaystyle \frac{\stackrel{~}{g}}{g}}\right)^2\delta \varphi _k.`$ (58)
As for the homogeneous part of the $`\chi `$ field, Eq. $`(\text{39})`$ implies
$`\chi \left({\displaystyle \frac{\stackrel{~}{g}}{g}}\right)^2\varphi .`$ (59)
We find from Eqs. $`(\text{58})`$ and $`(\text{59})`$ that both of super-Hubble $`\delta \chi _k`$ fluctuations and the homogeneous $`\chi `$ field are not severely suppressed compared with the $`\stackrel{~}{g}=0`$ case. Then we can expect that the growth of the $`\dot{\chi }\delta \chi _k`$ term in the r.h.s. of Eq. $`(\text{24})`$ during preheating may lead to the enhancement of super-Hubble metric perturbations.
Although larger values of $`\stackrel{~}{g}`$ will surely escape the inflationary suppression, we have to take care that this may prevent the successful inflationary scenario. During inflation, the frequency $`\mathrm{\Omega }_\varphi `$ of the inflaton condensate is estimated by making use of Eq. $`(\text{59})`$ as
$`\mathrm{\Omega }_\varphi ^2m^2+g^2\chi ^2+3e^{d\sigma /\sigma _{}}\stackrel{~}{g}^2\varphi \chi m^2\left[1{\displaystyle \frac{2\stackrel{~}{g}^2\varphi ^2}{m^2}}\left({\displaystyle \frac{\stackrel{~}{g}}{g}}\right)^2\right],`$ (60)
where we used the relation $`e^{d\sigma /\sigma _{}}1`$. Eq. $`(\text{60})`$ indicates that larger values of $`\stackrel{~}{g}`$ make $`\mathrm{\Omega }_\varphi ^2`$ negative, and will lead to the unphysical result that $`\varphi `$ increases by negative instability. In order to avoid this, $`\mathrm{\Omega }_\varphi ^2`$ should be positive at the start of inflation ($`\varphi _I3m_{\mathrm{pl}}`$). Then the ratio $`\stackrel{~}{g}/g`$ is constrained as
$`{\displaystyle \frac{\stackrel{~}{g}}{g}}\stackrel{<}{}{\displaystyle \frac{5\times 10^4}{\sqrt{g}}},`$ (61)
where we used $`m10^6m_{\mathrm{pl}}`$. For an efficient $`\chi `$ particle production, the coupling $`g`$ is required to be $`g\stackrel{>}{}3.0\times 10^4`$ , which leads to the constraint: $`\stackrel{~}{g}/g\stackrel{<}{}2.9\times 10^2`$. The upper bound of $`\stackrel{~}{g}/g`$ decreases with the increase of $`g`$ as is found by Eq. $`(\text{61})`$. For example, $`\stackrel{~}{g}/g\stackrel{<}{}1.6\times 10^2`$ for $`g=1.0\times 10^3`$, and $`\stackrel{~}{g}/g\stackrel{<}{}5.0\times 10^3`$ for $`g=1.0\times 10^2`$. Although larger values of $`g`$ are favorable for the rapid growth of the $`\delta \chi _k`$ fluctuation, this simultaneously results in the stronger suppression for $`\delta \chi _k`$ and $`\chi `$ in Eqs. $`(\text{58})`$ and $`(\text{59})`$.
We have numerically examined the dynamics of preheating in the coupling regimes $`g=3\times 10^410^2`$ with $`\stackrel{~}{g}`$ constrained by Eq. $`(\text{61})`$, and also checked that the inflationary period proceeds in the usual manner. In Fig. 6, the evolutions of field perturbations $`\delta \varphi _k`$, $`\delta \chi _k`$, $`\delta \sigma _k`$ and the metric perturbation $`\mathrm{\Phi }_k`$ are depicted for a super-Hubble mode $`\overline{k}=10^{26}`$ with $`g=5.0\times 10^4`$, $`\stackrel{~}{g}/g=2.0\times 10^2`$. We start integrating about 60 e-foldings before the beginning of preheating, and choose initial conditions $`\varphi _I=3.0m_{\mathrm{pl}}`$, $`\sigma _I=0.1m_{\mathrm{pl}}`$, and $`\chi _I=(\stackrel{~}{g}/g)^2\varphi _I`$ for the homogeneous part, and use Eq. $`(\text{55})`$ for the fluctuational parts. In this case, the analytic estimation $`(\text{58})`$ implies the relation $`|\delta \chi _k(t_i)|10^3|\delta \varphi _k(t_i)|`$ at the onset of preheating ($`mt20`$), which can be easily confirmed in Fig. 6(a). The super-Hubble $`\delta \chi _k`$ fluctuation starts to grow from $`mt30`$ by parametric resonance, and catches up the $`\delta \varphi _k`$ fluctuation at $`mt120`$. At this stage, the backreaction effect of the produced $`\chi `$ particle begins to destroy the coherent oscillation of the $`\varphi `$ field \[see the inset of Fig. 6(a)\]. In spite of this, the amplification of the $`\delta \chi _k`$ fluctuation still takes place before the oscillation of $`\varphi `$ is completely broken at $`mt220`$. For $`140\stackrel{<}{}mt\stackrel{<}{}170`$, the super-Hubble $`\delta \varphi _k`$ fluctuation is enhanced by about two orders of magnitude. This occurs in the perturbed metric case where the r.h.s. of Eq. $`(\text{27})`$ stimulates the excitement of the $`\delta \varphi _k`$ fluctuation. However, since the increase of $`\delta \chi _k`$ is weakened by the backreaction effect of created particles, the period during which the $`\delta \varphi _k`$ fluctuation is enhanced does not continue long. In Fig. 6(b), we find that the super-Hubble metric perturbation $`\mathrm{\Phi }_k`$ begins to oscillate for $`mt\stackrel{>}{}180`$, which is due to the enhancement of field fluctuations. However, $`\mathrm{\Phi }_k`$ does not increase even by one order of magnitude from the beginning of preheating. This is mainly because the backreaction effect restricts the rapid increase of $`\delta \chi _k`$ soon after the super-Hubble $`\delta \chi _k`$ fluctuation catches up $`\delta \varphi _k`$. Although one may think that larger values of $`\stackrel{~}{g}`$ will lead to the strong amplification of $`\mathrm{\Phi }_k`$, Eq. $`(\text{61})`$ constrains the coupling as $`\stackrel{~}{g}/g\stackrel{<}{}2.2\times 10^2`$ with $`g=5.0\times 10^4`$, in which case the super-Hubble metric perturbation can not be strongly excited.
We have also investigated other values of the coupling $`g`$, and numerical results exhibit the similar behavior. The large scale cosmic microwave background (CMB) anisotropies will not be significantly modified with the existence of the preheating phase, even taking into account the coupling $`\stackrel{~}{g}^2\varphi ^3\chi `$. However, the additional enhancement of the super-Hubble metric perturbation found in Fig. 6(b) for $`mt\stackrel{>}{}180`$ may give some small imprints in the CMB spectrum.
The $`\delta \sigma _k`$ fluctuation in super-Hubble scales can be amplified a little in a short stage as in the case of $`\delta \varphi _k`$. In Fig. 6(b), we find that $`\delta \sigma _k`$ increases by about two orders of magnitude during $`140\stackrel{<}{}mt\stackrel{<}{}220`$. However, for $`g\stackrel{>}{}3.0\times 10^4`$ and values of $`\stackrel{~}{g}`$ which satisfy the condition of $`(\text{61})`$, numerical calculations imply that the growth of super-Hubble $`\delta \sigma _k`$ modes is hardly expected except in the case where $`\stackrel{~}{g}/g`$ is close to its upper bound. Even when $`\stackrel{~}{g}/g`$ is close to its upper bound as in the case of $`g=5.0\times 10^4`$, $`\stackrel{~}{g}/g=2.0\times 10^2`$, the enhancement is found to be weak. Moreover, final fluctuations $`\delta \stackrel{~}{\sigma }_k`$ in super-Hubble modes are typically smaller than $`\delta \stackrel{~}{\varphi }_k`$ and $`\delta \stackrel{~}{\chi }_k`$ as is found in Fig. 6.
As for sub-Hubble scales, amplifications of metric perturbations and the $`\delta \sigma _k`$ fluctuation are relevant compared with the super-Hubble case. In Fig. 7, we show the evolutions of $`\mathrm{\Phi }_k`$ and $`\delta \sigma _k`$ during preheating for a sub-Hubble mode $`\overline{k}=5`$ with $`g=5.0\times 10^4`$, $`\stackrel{~}{g}/g=2.0\times 10^2`$. In this case, the sub-Hubble fluctuation of $`\delta \stackrel{~}{\chi }_k`$ is larger than in the super-Hubble case at the start of preheating. We find in Fig. 7 that $`\mathrm{\Phi }_k`$ increases by more than one order of magnitude, which indicates that metric preheating can be vital in small scales in the presence of the $`\stackrel{~}{g}^2\varphi ^3\chi `$ term. The dilaton fluctuation in sub-Hubble modes is also enhanced with the growth of metric perturbations. However, the final $`\delta \sigma _k`$ fluctuation does not exceed its fluctuation at the onset of preheating. If we choose smaller values of $`\stackrel{~}{g}`$, the enhancement of sub-Hubble $`\mathrm{\Phi }_k`$ and $`\delta \sigma _k`$ modes becomes weaker. For the couplings which range $`3\times 10^4\stackrel{<}{}g\stackrel{<}{}10^2`$, we numerically find that the sub-Hubble dilaton fluctuation is not relevantly amplified for the value of $`\stackrel{~}{g}`$ which is smaller by one order of magnitude than its upper bound given by Eq. $`(\text{61})`$. When the standard coupling $`g^2\varphi ^2\chi ^2`$ dominates over the coupling $`\stackrel{~}{g}^3\varphi ^3\chi `$ (namely $`g\stackrel{~}{g}`$), we conclude that the fluctuation of dilaton can be held static both in the sub- and super-Hubble scales.
## IV Concluding remarks and discussions
In this paper, we have studied preheating after inflation with a quadratic inflaton potential $`V(\varphi )=\frac{1}{2}m^2\varphi ^2`$ in the presence of a dilaton field $`\sigma `$ which represents the scale of compactifications in a higher-dimensional generalized Kaluza-Klein theory. We consider the Candelas-Weinberg model where extra dimensions are compactified on the sphere with a cosmological constant and a one-loop quantum correction (Casimir effects). In the chaotic inflation model, a potential barrier which prevents the growth of the internal space disappears for large values of inflaton. However, the fine-tuned initial conditions and parameters of the model naturally lead to successful inflation. We find that the existence of dilaton during inflation hardly affects the evolution of inflaton, and the chaotic inflationary scenario proceeds in the usual manner as long as initial conditions are chosen so that dilaton does not go beyond the potential barrier.
At the stage of preheating after inflation, another scalar field $`\chi `$ coupled to inflaton can be amplified by parametric resonance due to the oscillation of inflaton. In addition to the standard coupling $`\frac{1}{2}g^2\varphi ^2\chi ^2`$, we have also included the coupling $`\stackrel{~}{g}^2\varphi ^3\chi `$ by which the exponential suppression of the super-Hubble $`\chi `$ fluctuation can be avoided during inflation. We include metric perturbations explicitly for scalar field equations, and investigate how the fluctuation of dilaton will be amplified both in super- and sub-Hubble scales. Neglecting metric perturbations, the equation for the dilaton fluctuation is reduced to the form of Mathieu equation as in the case of the $`\chi `$ fluctuation at the linear stage of preheating. Since the resonance parameter $`q`$ is smaller than unity during the whole stage of preheating, the dilaton fluctuation does not exhibit parametric amplification in the rigid spacetime case.
In the perturbed metric case, it is generally expected that field resonances will stimulate the enhancement of metric perturbations $`\mathrm{\Phi }_k`$. In the case of $`\stackrel{~}{g}=0`$, however, since low momentum modes of the $`\chi `$ field fluctuation are severely suppressed during inflation, super-Hubble metric perturbations are hardly affected by parametric amplification of the field perturbation. We have numerically verified that super-Hubble metric perturbations remain almost constant during preheating, and also found that the dilaton fluctuation in super-Hubble modes can not be enhanced. As for sub-Hubble modes with $`\stackrel{~}{g}=0`$, the $`\chi `$ field fluctuation is not suppressed relative to super-Hubble modes. However, since the source term in the equation of $`\mathrm{\Phi }_k`$ contains the time derivative of the homogeneous $`\chi `$ field, we can not expect the strong amplification of dilaton fluctuation as well as metric perturbations.
If the coupling $`\stackrel{~}{g}^2\varphi ^3\chi `$ is taken into account, both the homogeneous and super-Hubble fluctuational parts of the $`\chi `$ field can escape the inflationary suppression. These are about $`(\stackrel{~}{g}/g)^2`$ times those of inflaton at the onset of preheating. In this case, the $`\chi `$ fluctuation can typically increase to the order of the inflaton fluctuation, after which the backreaction effect of created $`\chi `$ particles onto the $`\varphi `$ field begins to be relevant. This restricts the further excitement of the $`\chi `$ fluctuation, and amplifications of the super-Hubble metric perturbation and the dilaton fluctuation are found to be weak even for large values of $`\stackrel{~}{g}`$ which are close to its upper bound. As for sub-Hubble scales, the dilaton fluctuation can be modestly amplified by the growth of metric perturbations. However, the dilaton fluctuation does not grow for the value of $`\stackrel{~}{g}`$ which is smaller by one order of magnitude than its upper bound given by Eq. $`(\text{61})`$. We argue that the stability of compactifications can be preserved during preheating in the quadratic chaotic inflationary scenario as long as $`\stackrel{~}{g}g`$.
We have found that the enhancement of the dilaton perturbation is intimately related with the generation of metric perturbations. Although we only considered the backreaction due to field fluctuations, we should also include second order metric perturbations to the background equations for a consistent study. In fact, the effective momentum tensor formalism in Ref. gives rise to the coupling of the metric $`\mathrm{\Phi }`$ and the fluctuation $`\delta \chi _k`$ in the equation of the homogeneous $`\varphi `$ field, which leads to the growth of the super-Hubble $`\delta \chi _k`$ fluctuation as well as the homogeneous $`\chi `$ field. Although it is expected that metric perturbations in super-Hubble scales remain well within linear regimes in the case of $`\stackrel{~}{g}=0`$ as in Ref. , inclusion of second order metric backreaction effects may lead to the distortion of the large-scale CMB spectrum as well as the enhancement of sub-Hubble metric and field fluctuations in the presence of the interaction $`\stackrel{~}{g}^2\varphi ^3\chi `$. In addition to this, rescattering effects of field and metric fluctuations (i.e. mode-mode coupling) will be important at the final stage of preheating. Although we do not consider these complicated issues in this paper, we should take into account the full backreaction and the rescattering effects for a complete study of preheating.
In this paper, we have restricted ourselves to the massive inflaton model as a first step toward understanding the dynamics of preheating with extra dimensions. We found that amplification of the super-Hubble dilaton fluctuation is weak in this model. However, in other models of inflation, there may be a possibility that extra dimensions will be unstable during preheating. Indeed, it was recently suggested that super-Hubble metric perturbations can be strongly enhanced in broad classes of models. One of such models is the massless chaotic inflation model $`V(\varphi )=\frac{1}{4}\lambda \varphi ^4`$ with the interaction $`\frac{1}{2}g^2\varphi ^2\chi ^2`$. In this model, even when the effective mass of the $`\chi `$ field is comparable to the Hubble expansion rate $`H`$, the super-Hubble $`\chi `$ field fluctuation can be excited by parametric resonance. This will lead to the enhancement of the dilaton fluctuation in long wave modes through the amplification of super-Hubble metric perturbations. It was also pointed out in Ref. that the negative coupling $`\frac{1}{2}g^2\varphi ^2\chi ^2`$ with $`g^2<0`$ or the negative nonminimal coupling $`\frac{1}{2}\xi R\chi ^2`$ with $`\xi <0`$ will escape the inflationary suppression of the super-Hubble $`\chi `$ fluctuation. In addition to this, although we did not consider the interaction between $`\sigma `$ and $`\chi `$ fields, including this coupling may strengthen parametric resonance of scalar fields. It is quite interesting to investigate the evolution of metric and field fluctuations in broad classes of models with several interactions in the sense that we can constrain the inflaton potential in terms of distortions from the CMB spectrum.
Although we have investigated compactifications on the sphere, there exist several ways of compactifications on other manifolds. For example, the dilaton does not have its own potential in the torus compactification, which would lead to the growth of extra dimensions during inflation as was analyzed in Ref. . In this case, since the resonance term $`g^2\varphi ^2e^{d\sigma /\sigma _{}}`$ in the l.h.s. of Eq. $`(\text{30})`$ will be suppressed during inflation, the ordinary picture of preheating may be modified in the presence of extra dimensions. It is also of interest to investigate the dynamics of inflation and preheating in more realistic models of compactifications such as Calabi-Yau manifolds, because it is possible to judge whether such compactifications are appropriate or not from the cosmological point of view. These issues are under consideration.
## ACKNOWLEDGMENTS
The author would like to thank Bruce A. Bassett, Kei-ichi Maeda, Shinji Mukohyama, Takashi Torii, Kunihito Uzawa, Fermin Viniegra, and Hiroki Yajima for useful discussions and comments. This work was supported partially by a Grant-in-Aid for Scientific Research Fund of the Ministry of Education, Science and Culture (No. 09410217), and by the Waseda University Grant for Special Research Projects.
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# Point-Form Analysis of Elastic Deuteron Form Factors
## I Introduction
Electron scattering is considered to be an ideal tool to study the electromagnetic structure of hadronic systems. Relativity cannot be ignored for momentum transfers that provide information about the structure of the hadrons at the scale of a few tenths of a fermi. In order to understand hadronic systems at this scale, consistent relativistic models of both the hadronic dynamics and the hadronic electromagnetic current operator are required. If the dynamics and the current operator satisfy cluster properties , then the information learned about the structure of the simplest two- and three-body systems provides the essential components of models needed to treat complex targets.
Elastic electron-deuteron scattering is the simplest reaction that must be accurately modeled in order to constrain the dynamical generators and current operators that are needed to model complex systems. Because the deuteron is an isoscalar target, it might be expected that it can be accurately described by a pure impulse approximation. Unfortunately, pure impulse approximations are not consistent with current covariance or current conservation. An important goal is to find a physically motivated extension of the impulse approximation that is consistent with current conservation and covariance and is also qualitatively consistent with experiment.
In the one-photon-exchange approximation the experimental observables can be expressed in terms of matrix elements of the hadronic electromagnetic current operator between the initial and final eigenstates of the hadronic Hamiltonian. The general form of these matrix elements is
$$p^{}j^{}\mu _j^{}|\widehat{J}^\mu (0)|pj\mu _j,$$
(1)
where $`|p^{}j^{}\mu _j^{}`$ and $`|pj\mu _j`$ are eigenstates of the four-momentum, spin, and three-component of the spin, in reference frames related by a boost with momentum transfer $`Q=p^{}p`$. $`\widehat{J}^\mu (0)`$ is the hadronic current density at $`x=0`$.
The dynamical constraints on the current are current conservation,
$$[\widehat{P}_\mu ,\widehat{J}^\mu (0)]=0,$$
(2)
and current covariance,
$$U(\mathrm{\Lambda },a)\widehat{J}^\mu (x)U^{}(\mathrm{\Lambda },a)=(\mathrm{\Lambda }^1)^\mu {}_{\nu }{}^{}\widehat{J}_{}^{\nu }(\mathrm{\Lambda }x+a).$$
(3)
Here $`U(\mathrm{\Lambda },a)`$ is the unitary representation of the inhomogeneous Lorentz group, whose existence is required by relativistic invariance . $`\widehat{P}^\mu `$ is the four-momentum operator, with $`U(I,a)=e^{i\widehat{P}a}`$.
In applications there are essentially two approaches used to compute the hadronic current matrix elements. These are the covariant and direct interaction approaches. Each approach has its own advantages and disadvantages.
Covariant approaches assume that the underlying theory is a local quantum field theory. For the case of the deuteron the input is a covariant current vertex of the form
$$0|T(\mathrm{\Psi }(x_1)\mathrm{\Psi }(x_2)\widehat{J}^\mu (x)\overline{\mathrm{\Psi }}(y_1)\overline{\mathrm{\Psi }}(y_2))|0,$$
(4)
which is the vacuum expectation value of a time-ordered product of nucleon and current fields. Assuming the existence of an underlying quantum theory, Mandelstam showed how to extract the desired current matrix elements from the vertex. The Fourier transform of the vertex has pole terms on the deuteron mass shell. The residue includes a pair of Bethe-Salpeter amplitudes and a current matrix element. The Bethe-Salpeter normalization condition can be used with the solution of the homogeneous Bethe-Salpeter equation to remove the amplitudes from the residue. What remains is the desired current matrix element.
Quasipotential equations are based on these same concepts, but they introduce constraints designed to preserve the physical singularities. Current matrix elements are extracted using the constrained amplitudes and vertex functions.
Covariant methods are appealing because of their formal connection to a quantum field theory; however in most applications it is necessary to model the vertex and Bethe-Salpeter kernel, and to replace the two-point Green’s function with the free two-point function. Quasi-potential methods lead to simpler calculations than the full Bethe-Salpeter approach , but the reductions complicate cluster properties.
Direct interaction approaches attempt to construct consistent models of $`U(\mathrm{\Lambda },a)`$ and $`\widehat{J}^\mu (x)`$ directly. The transformation properties require that both $`U(\mathrm{\Lambda },a)`$ and $`\widehat{J}^\mu (x)`$ have an interaction dependence. Dirac addressed the problem of constructing $`U(\mathrm{\Lambda },a)`$ by including interactions in some of the infinitesimal generators of $`U(\mathrm{\Lambda },a)`$. He introduced the notion of forms of dynamics which minimize the number of interaction-dependent generators. The three main forms are: the instant form, where the interactions are in the Hamiltonian and Lorentz boost generators; the point form, where the interactions are in the four-momentum; and the front form, where interactions appear in the operators that generate transformations transverse to a fixed light front (a three-dimensional hyperplane tangent to the light cone.) But while Dirac identified the different possibilities for putting interactions in selected generators, he did not show how to actually construct the Poincaré generators with interactions.
The first exact construction of Poincaré generators with interactions was due to Bakamjian and Thomas using Dirac’s instant form. There are Bakamjian-Thomas-like constructions in each of the forms of dynamics, and they are scattering equivalent.
While explicit dynamical models of current operators are difficult to construct, consistent current matrix elements can be obtained by prescriptions that evaluate selected independent matrix elements using single nucleon currents. The remaining current matrix elements can then be determined by using covariance, current conservation and discrete symmetries. These generate the needed dynamical contributions to the current matrix elements.
Direct interaction approaches provide an exact treatment of the symmetries associated with special relativity, but are not directly related to an underlying field theory. A number of direct interaction applications to elastic electron-deuteron scattering exist in the literature. To date most applications have used Dirac’s instant- or front-forms of the dynamics. The point form of relativistic quantum mechanics has important simplifying features that are useful in modeling electron scattering. The purpose of this paper is to investigate the hadronic current operator in Dirac’s point form of dynamics.
In section II we discuss some of the features of point-form dynamics and construct a mass operator for the deuteron. Section III deals with current operators, their relation to observables, and the point-form spectator approximation. Then in section IV the numerical results are discussed and compared with other methods. Section V presents our conclusions.
## II Point-Form Relativistic Quantum Mechanics
Unlike nonrelativistic quantum mechanics, where all the interactions can be put in the Hamiltonian operator, for relativistic quantum mechanics it is necessary that at least three generators contain interactions. This can already be seen by examining the commutator of the Lorentz boost generators with the momentum generators. Such a commutator produces the Hamiltonian; if the Hamiltonian contains interactions, then some combination of boost and momentum generators must also contain interactions. In the instant form additional interactions are put in the boost generators, leaving the momentum generators free of interactions, while in the point form the additional interactions are in the momentum generators, with the boost generators free of interactions. The front form puts interactions in a mixture of Lorentz and momentum generators.
Even though all forms of dynamics are scattering equivalent, each has certain advantages that are useful for specific applications. The goal of this paper is to analyze elastic deuteron form factors using the point form. The point form has a number of features that set it aside from the other forms. First, all of the interactions are in the Hamiltonian and momentum generators, that is, the four-momentum operator. Since there are no interactions in the boost or angular momentum generators, the Lorentz generators are all kinematic and the theory is manifestly Lorentz covariant. It is convenient to write the Poincaré commutation relations not in terms of the ten generators, but rather in terms of the four-momentum operators that contain the interactions, and global kinematic Lorentz transformations:
$`[\widehat{P}_\mu ,\widehat{P}_\nu ]`$ $`=`$ $`0;`$ (5)
$`U_\mathrm{\Lambda }\widehat{P}_\mu U_\mathrm{\Lambda }^1`$ $`=`$ $`(\mathrm{\Lambda }^1)_\mu ^\nu \widehat{P}_\nu ;`$ (6)
where $`U_\mathrm{\Lambda }U_\mathrm{\Lambda }(\mathrm{\Lambda },0)`$ is a unitary operator representing the Lorentz transformation $`\mathrm{\Lambda }`$. These rewritten Poincaré relations will be called the point-form equations, and are the fundamental equations that have to be satisfied for the system of interest. The mass operator is given by $`\widehat{M}=\sqrt{\widehat{P}\widehat{P}}`$ and must have a spectrum that is bounded from below.
Since the interactions are all in the four-momentum operators, which are the generators of space-time translations, the nonrelativistic Schrödinger equation can be generalized to a Lorentz covariant relativistic Schrödinger equation, namely
$$i\mathrm{\Psi }_x/x^\mu =\widehat{P}_\mu \mathrm{\Psi }_x,$$
(7)
where $`x=x^\mu `$ is the four-vector space-time point. If the four-momentum operator does not depend explicitly on space-time, this equation becomes the eigenvalue equation
$$\widehat{P}_\mu \mathrm{\Phi }=p_\mu \mathrm{\Phi }.$$
(8)
Finally, as will be shown in the following paragraphs, it is possible to define states with the property that angular momentum can be coupled in exactly the same way as is done in nonrelativistic quantum mechanics.
The simplest example of a system satisfying the point form equations is a one-particle system with mass $`m`$ and spin $`j`$. If $`|p,\sigma `$ is an eigenstate of four-momentum $`p`$ (with $`pp=m^2`$) and spin projection $`\sigma `$, then
$`\widehat{P}_\mu |p,\sigma `$ $`=`$ $`p_\mu |p,\sigma `$ (9)
$`U_\mathrm{\Lambda }|p,\sigma `$ $`=`$ $`{\displaystyle \underset{\sigma ^{}}{}}|\mathrm{\Lambda }p,\sigma ^{}D_{\sigma ^{}\sigma }^j(R_W)\sqrt{{\displaystyle \frac{v_0^{}}{v_0}}},`$ (10)
with $`R_W`$ a Wigner rotation defined by $`R_W=B^1(\mathrm{\Lambda }v)\mathrm{\Lambda }B(v)`$, and $`B(v)`$ a canonical spin (rotationless) boost (see reference ) with argument $`v=p/m`$. $`D_{\sigma ^{}\sigma }^j(R_W)`$ is a Wigner D function, and the eigenstates are normalized to
$$p^{},\sigma ^{}|p,\sigma =\delta ^3(𝐩^{}𝐩)\delta _{\sigma ^{}\sigma },$$
(11)
relativity requiring the $`\sqrt{v_0^{}/v_0}`$ factor.
States of many noninteracting particles are tensor products of one-particle states; however a problem arises when such many-particle states are Lorentz transformed. As can be seen from Eq. (10) each state is Lorentz transformed by its own Wigner rotation, which in general are different. This means that these multiparticle states cannot be directly coupled together as is the case nonrelativistically. Such a problem is resolved by coupling the single-particle states in the overall rest frame and boosting. It is convenient to label the state by the system’s four-velocity $`v`$:
$`|v,𝐤_𝐢,\mu _i`$ $`:=`$ $`U_{B(v)}(|k_1,\mu _1\mathrm{}|k_n,\mu _n)`$ (12)
$`=`$ $`{\displaystyle \left(|p_1,\sigma _1\mathrm{}|p_n,\sigma _n\underset{i}{}\left[D_{\sigma _i,\mu _i}^{j_i}(R_{W_i})\sqrt{\frac{(v_i^{})_0}{(v_i)_0}}\right]\right)},`$ (13)
where $`p_i=B(v)k_i`$, $`𝐤_𝐢=\mathrm{𝟎}`$, and $`R_{W_i}=B^1(p_i/m)B(v)B(k_i/m)`$. Under Lorentz transformations, using the definition, Eq. (13), such velocity states transform as
$$U_\mathrm{\Lambda }|v,𝐤_𝐢,\mu _i=|\mathrm{\Lambda }v,R_W𝐤_𝐢,\mu _i^{}\underset{i}{}\left[D_{\mu _i^{},\mu _i}^{j_i}(R_W)\sqrt{\frac{(v_i^{})_0}{(v_i)_0}}\right];$$
(14)
where the Wigner rotation $`R_W=B^1(\mathrm{\Lambda }v)\mathrm{\Lambda }B(v)`$ is the same in all the arguments of the D functions and all the internal momentum vectors $`𝐤_𝐢`$. That means all the spins as well as the orbital angular momenta can be coupled together exactly as is done nonrelativistically. This property will be used in the following paragraphs for coupling the nucleon spins together with the relative orbital angular momentum to get the spin of the deuteron. From the relation between external and internal momenta, it follows that the velocity states defined in Eq. (13) are eigenstates of the noninteracting mass operator $`\widehat{M}_{\mathrm{free}}`$ and free four-velocity operator $`\widehat{V}_\mu `$:
$`\widehat{M}_{\mathrm{free}}|v,𝐤_𝐢,\mu _i`$ $`=`$ $`{\displaystyle \underset{i}{}}\sqrt{m_i^2+𝐤_{𝐢}^{}{}_{}{}^{2}}|v,𝐤_𝐢,\mu _i;`$ (15)
$`\widehat{V}_\mu |v,𝐤_𝐢,\mu _i`$ $`=`$ $`v_\mu |v,𝐤_𝐢,\mu _i.`$ (16)
The Bakamjian-Thomas procedure is implemented in the point form by writing $`\widehat{P}_\mu =\widehat{M}\widehat{V}_\mu `$, where now $`\widehat{M}`$ is the sum of free and interacting mass operators, $`\widehat{M}=\widehat{M}_{\mathrm{free}}+\widehat{M}_{\mathrm{int}}`$. $`\widehat{M}`$ takes the place of the center of momentum Hamiltonian $`\widehat{h}=\widehat{H}\frac{\widehat{P}^2}{2M}`$ in nonrelativistic quantum mechanics; note however that even though there is only one operator containing interactions, namely the mass operator, that nevertheless there are interactions in all four components of the four-momentum operator.
In order that the four-momentum operator satisfy the point-form equations, Eqs. (5,6), the interacting mass operator must satisfy certain conditions. To satisfy Eq. (5), the mass operator must commute with the four-velocity operator, defined in Eq. (16):
$$[\widehat{M},\widehat{V}^\mu ]=0.$$
(17)
This has the consequence that mass and four-velocity can be simultaneously diagonalized. Eigenstates of the four-momentum operator can thus be written as the mass times the four-velocity. Since the four-velocity is purely kinematic, it can be factored from the wave function leaving the covariant Schrödinger equation, Eq. (8), to become a mass operator eigenvalue equation,
$$\widehat{M}\mathrm{\Phi }=\lambda \mathrm{\Phi }.$$
(18)
Moreover, even though the four-momentum is conserved in reactions, the total four-momentum is not the sum of the four-momenta of the individual particles. Rather what is conserved is the overall four-velocity of the individual particles, and the mass is then “off-shell”, not unlike the situation with Feynman diagrams. This is to be contrasted with the instant form, where the three-momentum of all the individual particles give the total three-momentum of the system, while the energy is “off-shell”.
The mass operator must also satisfy the other point form equation, Eq. (6), implying the mass operator is a Lorentz scalar. On velocity states this means the kernel of the mass operator must be rotationally invariant and independent of $`𝐯^2`$, exactly the condition put on nonrelativistic Hamiltonians in order that they be Galilei invariant.
For a two-body system such as the deuteron, the relevant Hilbert space is the tensor product of proton and neutron Hilbert spaces, $`H=H_pH_n`$. In that case the velocity states can be written as $`|v,𝐤,\mu _p,\mu _n`$, where $`𝐤=𝐤_\mathrm{𝟏}=𝐤_\mathrm{𝟐}`$, and $`\mu _p`$ and $`\mu _n`$ are the eigenvalues of the three-components of the canonical spins of the proton and neutron respectively. Because with velocity states the angular momenta can all be coupled together, these states can also be written as $`|v,|𝐤|,j,\mu _j,l,s`$, as in the nonrelativistic case. The mass of the two particle state, from Eq. (15), is $`2\sqrt{m^2+𝐤^2}`$; $`j`$ is the total angular momentum, while $`l`$ and $`s`$ are the orbital and spin angular momentum respectively.
It is advantageous to express the interacting mass operator in terms of a mass squared operator with matrix elements:
$`v,|𝐤|,j,\mu _j,l,s|\widehat{M}_I^2|v^{},|𝐤^{}|,j^{},\mu _j^{},l^{},s^{}`$ (19)
$`=`$ $`\delta (vv^{})\delta _{\mu _j\mu _j^{}}\delta _{jj^{}}k,l,s(m_I^j)^2k^{},l^{},s^{}.`$ (20)
A mass operator with a kernel of the form Eq. (20) will satisfy Eq. (17) and thus the Poincaré commutation relations, Eqs. (5,6). The kernel of $`\widehat{M}_I^2`$ is taken to be
$$k,l,s(m_I^j)^2k^{},l^{},s^{}:=4mk,l,sv_{nn}^jk^{},l^{},s^{},$$
(21)
where $`v_{nn}^j`$ is a nucleon-nucleon interaction. The mass is then defined by
$$\widehat{M}:=\sqrt{\widehat{M}^2};\widehat{M}^2:=4(𝐤^2+m^2)+\widehat{M}_I^2.$$
(22)
Denoting the eigenvalue of the interacting mass operator by $`\lambda ^2`$, the equation
$$\widehat{M}^2\mathrm{\Phi }=(4m^2+4𝐤^2+4mv_{nn}^j)\mathrm{\Phi }=\lambda ^2\mathrm{\Phi }$$
(23)
can be rewritten in the form of the nonrelativistic Schrödinger equation,
$$\left(\frac{𝐤^2}{m}+v_{nn}^j\right)\mathrm{\Phi }=\left(\frac{\lambda ^2}{4m}m\right)\mathrm{\Phi }.$$
(24)
This defines a relativistic model of the two-nucleon system. It can be shown that this model leads to a small correction to the nonrelativistic binding energy and has scattering observables identical to the corresponding nonrelativistic model. Equation (20) shows that the solution of Eq. (24) leads to simultaneous eigenstates of the mass, velocity, spin, and z-component of spin. The Poincaré transformation properties of the deuteron eigenstates are given by
$$\widehat{P}^\mu |v,m_D,j,\mu _j=\lambda v^\mu |v,m_D,j,\mu _j$$
(25)
and
$$U_\mathrm{\Lambda }|v,m_D,j,\mu _j=\underset{\mu _j^{}}{}|\mathrm{\Lambda }v,m_D,j,\mu _jD_{\mu _j^{}\mu _j}^j(R_W(\mathrm{\Lambda },v))\sqrt{\frac{(\mathrm{\Lambda }v)_0}{v_0}},$$
(26)
where
$$v,|𝐤|,j,\mu _j,l,s|v^{},m_D,j^{},\mu _j^{}=\delta (vv^{})\delta _{\mu _j^{}\mu _j}\delta _{j^{}j}\mathrm{\Psi }_{ls}^j(|𝐤|).$$
(27)
$`\mathrm{\Psi }_{ls}^j(|𝐤|)`$ is the nonrelativistic wavefunction associated with one of the two chosen nonrelativistic potentials. (There are analogous formulas for the scattering states.) This provides the desired point-form dynamics.
## III Current Operators, Form Factors, and Elastic Observables
The second key element in a theoretical description of electron-scattering is a conserved, covariant hadronic current density $`\widehat{J}^\mu (x)`$. In the point form the dynamical Poincaré transformations are the space-time translations. Translational covariance can be realized by using the dynamical translation operators to define $`\widehat{J}^\mu (x)`$ in terms of $`\widehat{J}^\mu (0)`$:
$$\widehat{J}^\mu (x):=e^{i\widehat{P}x}\widehat{J}^\mu (0)e^{i\widehat{P}x}.$$
(28)
The density $`\widehat{J}^\mu (0)`$ is assumed to transform as a four-vector with respect to the free Lorentz transformations.
We now want to define the model current operator in terms of measured one-body current operators. This is done as follows. The deuteron matrix elements of $`\widehat{J}^\mu (0)`$ are defined in terms of their Breit-frame values with $`Q`$ in the $`\widehat{𝐳}`$ direction:
$$𝐐/2,1,\mu _j^{}|\widehat{J}^\mu (0)|𝐐/2,1,\mu _j.$$
(29)
For $`\mu =0,1,2`$ the current matrix elements are defined in terms of the single-nucleon current matrix elements:
$`𝐐/2,1,\mu _j^{}|\widehat{J}^\mu (0)|𝐐/2,1,\mu _j=`$ (30)
$`𝐐/2,1,\mu _j^{}|\left(\widehat{J}_p^\mu (0)\widehat{I}_n+\widehat{I}_p\widehat{J}_n^\mu (0)\right)|𝐐/2,1,\mu _j.`$ (31)
Current conservation requires that
$$\underset{\mu }{}Q_\mu 𝐐/2,1,\mu _j^{}|\widehat{J}^\mu (0)|𝐐/2,1,\mu _j=0,$$
(32)
which generates a dynamical contribution $`\widehat{J}_{pn}^\mu (0)`$ to the $`\widehat{𝐳}`$ component of the current:
$`𝐐/2,1,\mu _j^{}|\widehat{J}_{pn}^3(0)|𝐐/2,1,\mu _j`$ (33)
$`=𝐐/2,1,\mu _j^{}|\left(\widehat{J}_p^3(0)\widehat{I}_n+\widehat{I}_p\widehat{J}_n^3(0)\right)|𝐐/2,1,\mu _j.`$ (34)
These relations define the components of the Breit-frame matrix elements of $`\widehat{J}^\mu (0)`$. The remaining deuteron matrix elements of $`\widehat{J}^\mu (x)`$ are fixed by kinematic Lorentz covariance and dynamical space-time translational covariance. Although the current matrix element is defined in the Breit frame, the expression for the general current matrix element is Lorentz covariant, as can be seen in reference , Eq. 3.31.
The computation of the matrix elements is carried out by inserting single-particle intermediate states in the velocity basis that was used to formulate the dynamical model in the previous section. The deuteron wavefunction in the basis Eq. (27) has the form
$$v,|𝐤|,j,\mu _j,l,s|v^{},m_D,j,\mu _j^{}=\delta (vv^{})\delta _{\mu _j^{}\mu _j}\delta _{j\mathrm{\hspace{0.33em}1}}\delta _{s\mathrm{\hspace{0.33em}1}}[\delta _{l\mathrm{\hspace{0.33em}0}}u_0(k)+\delta _{l\mathrm{\hspace{0.33em}2}}u_2(k)],$$
(35)
where $`u_0(k)`$ and $`u_2(k)`$ are the nonrelativistic S and D state deuteron wavefunctions. Transformation coefficients are used to express this in terms of single-particle basis elements:
$`v_1,\mu _1,v_2,\mu _2|v^{},\mu _j^{},m_D`$ (36)
$`=\delta ^3[𝐯^{}𝐯(𝐯_\mathrm{𝟏},𝐯_\mathrm{𝟐})]{\displaystyle \frac{\delta [kk(𝐯_\mathrm{𝟏},𝐯_\mathrm{𝟐})]}{k^2}}\left|{\displaystyle \frac{(𝐯,𝐤)}{(𝐯_\mathrm{𝟏},𝐯_\mathrm{𝟐})}}\right|^{1/2}`$ (37)
$`\times D_{\mu _1\mu _1^{}}^{1/2}[B^1(v_1)B(v)B(k_1)]D_{\mu _2\mu _2^{}}^{1/2}[B^1(v_2)B(v)B(k_2)]Y_{l\mu _l}[\widehat{𝐤}_1(𝐯_\mathrm{𝟏},𝐯_\mathrm{𝟐})]`$ (38)
$`\times {\displaystyle \frac{1}{2}},\mu _1,{\displaystyle \frac{1}{2}},\mu _2|1,\mu _sl,\mu _l,1,\mu _s|1,\mu _j^{}\left\{\delta _{l\mathrm{\hspace{0.33em}0}}u_0[k(𝐯_\mathrm{𝟏},𝐯_\mathrm{𝟐})]+\delta _{l\mathrm{\hspace{0.33em}2}}u_2[k(𝐯_\mathrm{𝟏},𝐯_\mathrm{𝟐})]\right\}.`$ (39)
These expressions can be used to compute the current matrix element
$`v,m_D,1,\mu _j|\widehat{J}_{SA}^\mu (0)|v^{},m_D,1,\mu _j^{}`$ (40)
$`={\displaystyle v,m_D,1,\mu _j|v_1,\mu _1,v_2,\mu _2v_1^{},\mu _1^{},v_2^{},\mu _2^{}|v^{},m_D,1,\mu _j^{}}`$ (41)
$`\times \left[v_1,\mu _1|\widehat{J}_1^\mu (0)|v_1^{},\mu _1^{}\delta ^3(𝐯_\mathrm{𝟐}^{}𝐯_\mathrm{𝟐})\delta _{\mu _2^{}\mu _2}+v_2,\mu _2|\widehat{J}_2^\mu (0)|v_2^{},\mu _2^{}\delta ^3(𝐯_\mathrm{𝟏}^{}𝐯_\mathrm{𝟏})\delta _{\mu _1^{}\mu _1}\right],`$ (42)
where the nucleon current matrix elements are given by Eq. (31). After integrating out the delta functions, one is left with a final three-dimensional integral:
$`v^{},m_D,1,\mu _j^{}|\widehat{J}^\mu (0)|v,m_d,1,\mu _j`$ $`=`$ $`{\displaystyle \underset{\mu _1^{}\mu _1}{}}{\displaystyle \underset{\mu _2^{}\mu _2}{}}{\displaystyle \underset{\mu _s^{}\mu _s}{}}{\displaystyle \underset{l^{}l}{}}{\displaystyle \underset{\mu _l^{}\mu _l}{}}{\displaystyle \underset{\sigma _1^{}\sigma _1}{}}{\displaystyle d^3k}`$ (43)
$`C_{\mu _s^{}\mu _1^{}\mu _2^{}}^{\mathrm{1\hspace{0.33em}1}/\mathrm{2\hspace{0.33em}1}/2}C_{\mu _s\mu _1\mu _2}^{\mathrm{1\hspace{0.33em}1}/\mathrm{2\hspace{0.33em}1}/2}`$ $`\times `$ $`C_{\mu _j^{}\mu _l^{}\mu _s^{}}^{1l\mathrm{\hspace{0.33em}1}}C_{\mu _k\mu _l\mu _s}^{1l^{}\mathrm{\hspace{0.33em}1}}`$ (44)
$`\times Y_{l^{}\mu _l^{}}^{}(\theta ^{},\varphi ^{})u_l^{}(|𝐤^{}|)`$ $`\times `$ $`Y_{l\mu _l}(\theta ,\varphi )u_l(|𝐤|)`$ (45)
$`\times D_{\mu _1^{}\sigma _1^{}}^{\mathrm{\hspace{0.33em}1}/2}\{R_W^1[k_1^{},B(v^{})]\}`$ $`\times `$ $`D_{\mu _2^{}\sigma _2^{}}^{\mathrm{\hspace{0.33em}1}/2}\{R_W^1[k_2^{},B(v^{})]\}`$ (46)
$`\times D_{\sigma _1\mu _1}^{1/2}\{R_W[k_1,B(v)]\}`$ $`\times `$ $`D_{\sigma _2\mu _2}^{1/2}\{R_W[k_2,B(v)]\}`$ (47)
$`\times \overline{u}(p_1^{}\sigma _1^{})\{\gamma ^\mu F_1[(p_1^{}p_1)^2]`$ $`+`$ $`i{\displaystyle \underset{\nu }{}}\sigma ^{\mu \nu }{\displaystyle \frac{(p_1^{}p_1)_\nu }{2m_N}}F_2[(p_1^{}p_1)^2]\}u(p_1\sigma _1)`$ (48)
$`+\{1`$ $``$ $`2\};`$ (49)
where the C’s are Clebsch-Gordan coefficients and the conventions for the spinors, gamma and sigma matrices are those of Bjorken and Drell .
In this form it can be seen that the momentum of the unstruck particle (the spectator) is unchanged, while the struck particle’s momentum is changed, but the impulse given to the struck particle is not the impulse given to the deuteron. For this reason we call this the point-form spectator approximation (PFSA). It should not be confused with the use of the term spectator approximation in, for example, reference . The practical advantage of the PFSA is that the steps above can be generalized to any hadronic target. Moreover, the current matrix element is generally Lorentz covariant and can be evaluated in any frame.
Because the interactions in the point form are in the four-momentum, in the PFSA the momentum transfer seen by the scattered nucleon is not the same as the momentum transfer seen by the nucleus. In Appendix A we show that the relationship between the momentum $`Q`$ transferred to the deuteron and the momentum transferred to the interacting nucleon is
$`|(p_1^{}p_1)^2|`$ $`=`$ $`Q^2{\displaystyle \frac{4(m_N^2+𝐤_{}^2)}{m_D^2}}(1+{\displaystyle \frac{Q^2}{4m_D^2}})`$ (50)
$`>`$ $`Q^2{\displaystyle \frac{4m_N^2}{m_D^2}}(1+{\displaystyle \frac{Q^2}{4m_D^2}})>Q^2.`$ (51)
That is, the point-form momentum transfer seen by an individual nucleon will be greater in magnitude than the total deuteron momentum transfer $`Q^2`$.
Two important implications follow from the equations (50) and (51) above. First, the PFSA momentum transfer depends on the internal momentum $`𝐤`$, which is a variable of integration. Thus in the PFSA, form factors depending on $`(p_i^{}p_i)^2`$ must remain inside the integral. Second, since $`|(p_i^{}p_i)^2|>Q^2`$, the deuteron form factors will fall off faster in the point-form calculations than in forms where $`|(p_i^{}p_i)^2|=Q^2`$.
The input to the PFSA are single-nucleon current operators. The general structure of these operators follows from covariance, parity, hermiticity, and time-reversal symmetry. For a spin-1/2 target the conditions imply that all matrix elements can be expressed in terms of the Dirac form factors, $`F_1(Q^2)`$ and $`F_2(Q^2)`$. The general expression has the form
$$p,\nu |J^\mu (0)|p^{},\nu ^{}=\overline{u}_\nu ^{}(p^{})\left[F_1(Q^2)\gamma ^\mu +F_2(Q^2)\underset{\alpha }{}\frac{iQ_\alpha \sigma ^{\mu \alpha }}{2m}\right]u_\nu (p),$$
(52)
where $`u_\nu (p)`$ and $`\overline{u}_\nu ^{}(p^{})`$ are Dirac spinors. In this form the one-body matrix elements are easily evaluated in any kinematic frame.
The Sachs electric and magnetic form factors of the nucleons are
$`G_E(Q^2)`$ $`=`$ $`\sqrt{1+\tau }𝐐/2,{\displaystyle \frac{1}{2}},{\displaystyle \frac{1}{2}}|\widehat{J}^0(0)|𝐐/2,{\displaystyle \frac{1}{2}},{\displaystyle \frac{1}{2}};`$ (53)
$`G_M(Q^2)`$ $`=`$ $`\sqrt{{\displaystyle \frac{1+\tau }{\tau }}}𝐐/2,{\displaystyle \frac{1}{2}},{\displaystyle \frac{1}{2}}|\widehat{J}^1(0)|𝐐/2,{\displaystyle \frac{1}{2}},{\displaystyle \frac{1}{2}};`$ (54)
where $`\tau =Q^2/4m^2`$. (Here the standard frame is the Breit frame, where the nucleon enters with momentum $`Q/2`$ and exits with momentum $`Q/2`$, both along the z-axis, which is also the axis along which the spin projection is measured.) These Dirac and Sachs form factors are related by
$`F_1(Q^2)`$ $`=`$ $`{\displaystyle \frac{1}{1+\tau }}[G_E(Q^2)+\tau G_M(Q^2)];`$ (55)
$`F_2(Q^2)`$ $`=`$ $`{\displaystyle \frac{1}{1+\tau }}[G_M(Q^2)G_E(Q^2)].`$ (56)
The input we use to define the model PFSA current are the single-nucleon form factor parameterizations of Gari-Krümpelmann and Mergell-Meissner-Drechsel .
The experimental observables for the deuteron and nucleon are well known. The elastic observables $`A(Q^2)`$ and $`B(Q^2)`$ are extracted from the Rosenbluth formula for the cross-section of unpolarized scattering in the lab frame,
$$\frac{d\sigma }{d\mathrm{\Omega }}=\frac{\alpha ^2\mathrm{cos}^2(\theta /2)}{4E^2\mathrm{sin}^4(\theta /2)}\frac{E^{}}{E}[A(Q^2)+B(Q^2)\mathrm{tan}^2(\theta /2)],$$
(57)
where $`\alpha `$ is the fine-structure constant, $`\theta `$ the scattering angle, and $`E`$ and $`E^{}`$ the initial and final energies. For the nucleons, it can be shown that
$`A(Q^2)`$ $`=`$ $`{\displaystyle \frac{G_E^2(Q^2)+\tau G_M^2(Q^2)}{1+\tau }};`$ (58)
$`B(Q^2)`$ $`=`$ $`2\tau G_M^2(Q^2).`$ (59)
For spin-1/2 particles, measurements of $`A(Q^2)`$ and $`B(Q^2)`$ suffice to determine $`G_E`$ and $`G_M`$. Various models of the nucleon form factors have been constructed.
The deuteron has three independent form factors. A common classification is to denote them as the charge monopole $`G_E`$, magnetic dipole $`G_M`$, and electric quadrupole $`G_Q`$ form factors. As current matrix elements, these are defined in the Breit frame:
$`G_E`$ $`=`$ $`{\displaystyle \frac{1}{3}}𝐐/2,1,0|\widehat{J}^0(0)|𝐐/2,1,0`$ (60)
$`+`$ $`{\displaystyle \frac{2}{3}}𝐐/2,1,1|\widehat{J}^0(0)|𝐐/2,1,1;`$ (61)
$`G_M`$ $`=`$ $`\sqrt{{\displaystyle \frac{2}{\eta }}}𝐐/2,1,1|\widehat{J}^1(0)|𝐐/2,1,0;`$ (62)
$`G_Q`$ $`=`$ $`{\displaystyle \frac{1}{2\eta }}[𝐐/2,1,0|\widehat{J}^0(0)|𝐐/2,1,0`$ (63)
$``$ $`𝐐/2,1,1|\widehat{J}^0(0)|𝐐/2,1,1];`$ (64)
where $`\eta =Q^2/4m_D^2`$. These form factors have the static limits
$`G_E(0)`$ $`=`$ $`e;`$ (65)
$`\underset{Q^20}{lim}G_M(Q^2)`$ $`=`$ $`e{\displaystyle \frac{m_D}{m_N}}\mu _D;`$ (66)
$`\underset{Q^20}{lim}G_Q(Q^2)`$ $`=`$ $`em_D^2Q_D;`$ (67)
where $`e`$ is the charge, $`\mu _D`$ the magnetic dipole moment, and $`Q_D`$ the electric quadrupole moment of the deuteron.
The Rosenbluth formula alone cannot determine all three of the deuteron’s form factors. The other independent observable normally measured is the deuteron tensor polarization $`T_{20}`$, defined as
$$T_{20}:=\sqrt{2}\frac{d\sigma ^1d\sigma ^0}{d\sigma },$$
(68)
where $`d\sigma ^\mu `$ refers to the differential cross-section with helicity $`\mu `$. Conventionally it is displayed at a 70 angle in the lab frame. The deuteron elastic observables are:
$`A(Q^2)`$ $`=`$ $`G_E^2+{\displaystyle \frac{8}{9}}\eta ^2G_Q^2+{\displaystyle \frac{2}{3}}\eta G_M^2;`$ (69)
$`B(Q^2)`$ $`=`$ $`{\displaystyle \frac{4}{3}}\eta (1+\eta )G_M^2;`$ (70)
$`T_{20}(Q^2)`$ $`=`$ $`\sqrt{2}\eta {\displaystyle \frac{\frac{4}{9}\eta G_Q^2+\frac{4}{3}G_QG_E+\frac{1}{3}fG_M^2}{A(Q^2)+B(Q^2)\mathrm{tan}^2(\theta /2)}};`$ (71)
where $`f=1/2+(1+\eta )\mathrm{tan}^2(\theta /2)`$.
## IV Numerical Results and Comparisons
One purpose of this work was to test the point-form spectator approximation on the simplest nucleus, the deuteron, where realistic interactions and nucleon form factors are available. Comparisons with other models are then an indication of the relative size of the required two-body currents.
For the nucleon-nucleon interaction, we have used the Argonne $`v_{18}`$ and Reid ’93 potentials to construct a mass operator $`\widehat{M}`$. The S (l=0) and D (l=2) state wavefunctions are plotted in Figures 1 and 2. The only significant differences between the wavefunctions these potentials produce in configuration space occur below 0.4 fm for the S wave and below 1.0 fm for the D wave. In momentum space the wavefunctions do not exhibit significant differences up to 5 $`\mathrm{fm}^1`$, about 1 GeV, above which they do differ noticeably. The effects on the choice of interaction may therefore be expected to be relevant at higher momentum transfers, but as will be seen, these high-momentum differences in the wavefunction make only slight differences in the calculations.
The PFSA currents are constructed using the Gari-Krümpelmann and Mergell-Meissner-Drechsel parameterizations of the nucleon form factors. At the range of momentum transfer under consideration, the parameterizations give very similar results for the proton form factors and the neutron magnetic form factor. The neutron electric form factor, however, varies significantly between the two.
The deuteron form factor $`G_E`$ has been calculated using both form factor parameterizations and both nucleon-nucleon potentials. The absolute values of the results are displayed in Figure 3. The results are independent of the nucleon-nucleon potential used, except for small variations at high momenta. The primary differences in $`G_E`$ and $`G_Q`$ are due to the nucleon form factor parameterizations. (Figure 4 compares the G-K and MMD neutron form factors. Note that the major difference is in the parameterization of the neutron electric form factor; the neutron magnetic form factor parameterizations, and the proton parameterizations as well, are very similar.) For $`G_E`$, both the Gari-Krümpelmann and the Mergell-Meissner-Drechsel form factors agree at low momentum transfers and have zeroes near 0.8 $`\mathrm{GeV}^2`$. The G-K form factors predict a second zero near 5.5 $`\mathrm{GeV}^2`$ while MMD predicts a second zero between 6 and 7 $`\mathrm{GeV}^2`$. Because the two form factors are almost identical except for the parameterization of the neutron electric form factor, this would suggest that the neutron form factor is the dominant cause of the differences in the calculations of $`G_E`$.
Figure 5 illustrates the dependence of the magnitude of the form factor $`G_M`$ on the potential and on the nucleon form factors used. Both parameterizations predict the same behavior up to the first zero, this time at 1.6 $`\mathrm{GeV}^2`$, and within the range studied, fall off with almost identical behavior. Comparison with Figure 3 would suggest that the neutron electric form factor has little effect on the calculation of $`G_M`$. A further comparison to experimental data can be made by examining the static limit of $`G_M`$. Equation (66) relates this limit to the deuteron magnetic moment, and Table I displays the results. In the static limit, the parameterization of the nucleon form factors does not affect the results, while the choice of nucleon-nucleon interaction does. This is expected, as the form factors must approach precise limits as $`Q^20`$, while the momentum-space wavefunctions have no such constraints.
This procedure is repeated for $`G_Q`$ in Figure 6. As was the case with $`G_E`$, there is little difference due to the potential used, but a noticeable difference between the predictions of the G-K and MMD parameterizations. The G-K form factors show a zero between 4.5 and 5.0 $`\mathrm{GeV}^2`$, while the MMD form factors produce a zero approximately one $`\mathrm{GeV}^2`$ higher. Further, the magnitude of the G-K results is greater than that of MMD almost everywhere throughout. The different neutron electric form factor parameterizations is the primary cause of the differences in the results. The deuteron electric quadrupole moment (Eq. 67), displayed in Table 1, differs from the experimental result, the calculated values approaching about 90% of the experimental value, as opposed to 99% for the magnetic moment calculations. This is consistent with other models .
To summarize, these point-form calculations imply that the deuteron form factors are essentially independent of the potential (Argonne $`v_{18}`$ or Reid ’93) used, but depend more significantly on the parameterizations of the form factors, and in particular on the neutron electric form factor, as this is the only substantial difference between the G-K and MMD parameterizations. The static moments are similar to predictions in other realistic models with the electric quadrupole moment differing with experiment by about 10%.
Figure 7 displays the results for $`A(Q^2)`$ up to 2 $`\mathrm{GeV}^2`$ for both potentials and both form factor parameterizations; Figure 8 extends the calculations to 8 $`\mathrm{GeV}^2`$. The data come from Refs. . Differences among the various calculations begin to appear at intermediate momentum transfers. For $`A(Q^2)`$, with $`Q^2`$ between 0.5 and 3 GeV<sup>2</sup>, the PFSA combined with the Gari-Krümpelmann form factors fit the data fairly closely, while the Mergell-Meissner-Drechsel form factors produce results that fall short by an order of magnitude. This pattern occurs in other impulse and spectator calculations as well. In the front-form calculations of Chung et al, the fit for an earlier G-K parameterization (using the Argonne $`v_{14}`$ potential) is even better, while the Höhler (on which the newer MMD form factors were based) calculations again fall an order of magnitude short. In Lev et al the two curves are closer, though the G-K remains higher and fits the data out to 2.0 GeV<sup>2</sup>. In the nonrelativistic calculations of Carlson and Schiavilla , which cover the range 0–2.4 GeV<sup>2</sup> and use only the Höhler form factors, impulse approximations using various potentials (including the Argonne $`v_{18}`$) all fall nearly an order of magnitude short in the intermediate range. In the work of Van Orden et al, which contains an impulse approximation that fits most of the data for all three form factors quite closely, the variation from the $`A(Q^2)`$ data starts at 2.0 GeV<sup>2</sup> and is an order of magnitude short at high (8 GeV<sup>2</sup>) momentum transfers.
The elastic observable $`B(Q^2)`$, related directly to the magnetic form factor $`G_M`$, is displayed in Figure 9 up to 8 $`\mathrm{GeV}^2`$ using both potentials and both sets of parameterizations. The data come from Refs. . The PFSA calculation of $`B(Q^2)`$ in the intermediate region 0.5–3 GeV<sup>2</sup>, a region which contains all presently available data, fits that data poorly, though the differences between the two form factor parameterizations are less marked. Both parameterizations exhibit a zero at 1.5$`\pm `$0.1 GeV<sup>2</sup>, causing a wide discrepancy with experiment; data suggest a zero nearer 1.9 GeV<sup>2</sup>. Again Chung et al. and Van Orden et al. fit the $`B(Q^2)`$ data quite well. In contrast, Chung et al., using the Paris and the Bonn wavefunctions instead of the Argonne $`v_{14}`$, produce results quite similar to the point form’s. Carlson and Schiavilla obtain a zero at 2.2 GeV<sup>2</sup> only with the Nijmegen potential; their other potentials reproduce the zero at 1.6 GeV<sup>2</sup>. Lev et al. produce zeroes between 1–2 GeV<sup>2</sup> in the 0–4 GeV<sup>2</sup> region using various potentials and parameterizations. The position of the zero in all these forms seems to be the most salient feature of calculations of $`B(Q^2)`$, affecting as it does the deviation from data in the 1–2 GeV<sup>2</sup> range.
The tensor polarization, $`T_{20}(Q^2)`$, is displayed in Figure 10, with data from Refs. . As with $`B(Q^2)`$, measurements of $`T_{20}(Q^2)`$ have only investigated the low (0–0.5 GeV) and intermediate (0.5–1.5 GeV) ranges of momentum transfer. In both ranges, the MMD and G-K parameterizations produce identical results in the PFSA. In the intermediate range, the results fall slightly below the data. The impulse approximations of Carlson and Schiavilla as well as Lev et al. do this too, while Van Orden et al. and Chung et al. produce curves that fit modern data quite closely.
Finally, it is instructive to compare the results obtained in point-form dynamics to the results of nonrelativistic impulse calculations to get some idea of the nature of the relativistic effects; and to compare the results in the point form to the same relativistic calculations done assuming $`|(p_1^{}p_1)^2|=Q^2`$ (that is, pulling the nucleon form factors outside of the integral) to examine how the point-form momentum transfer affects the results.
The results (using the Argonne $`v_{18}`$ potential and the Gari-Krümpelmann form factors) for $`A(Q^2)`$ are displayed in Figure 11. For low momentum transfers, all three agree. In the range of 1.0–5.0 GeV<sup>2</sup>, the nonrelativistic and PFSA calculations decrease similarly, the nonrelativistic curve consistently higher. The constant-$`Q^2`$ calculation in this region gradually rises from the PFSA to the nonrelativistic curve. Above 5.0 GeV<sup>2</sup>, all three curves systematically decrease, nonrelativistic above constant-$`Q^2`$ above PFSA. While a relativistic treatment is needed as a matter of principle at high momentum transfers, it is clear that in these calculations the effects of combining the PFSA with point-form quantum models has a tendency to reduce the structure function $`A(Q^2)`$ at high $`Q^2`$. One clear cause of this is that
$$|(p_i^{}p_i)^2|>Q^2;$$
(72)
the magnitude of the point-form momentum transfer is greater than the magnitude of the nonrelativistic momentum transfer. Because the form factors depend on the magnitude of the momentum transfer, they therefore drop off more quickly in the point form. This reduces the point-form results in comparison to the constant-$`Q^2`$ calculations, as Figure 11 shows: as the momentum transfer increases, the two calculations diverge further from each other.
The point-form and nonrelativistic results for $`B(Q^2)`$ and $`T_{20}(Q^2)`$ do not exhibit as dramatic differences as they did for $`A(Q^2)`$. Again in the graphs of $`B(Q^2)`$ and $`T_{20}(Q^2)`$, Figures 12 and 13, one sees the similar but increasingly divergent results obtained from the constant-$`Q^2`$ and the PFSA methods.
## V Conclusion
This work has used the point form of relativistic dynamics to calculate elastic deuteron form factors. The point form stands somewhat between the covariant approaches and direct-interaction approaches mentioned in the introduction in that it is on the one hand manifestly covariant (because the Lorentz generators are kinematic) and it is the mass that is “off-shell” (rather than the energy as is the case with the instant form.) On the other hand the point form is one of the forms of dynamics listed by Dirac, in which all of the interactions are in the four-momentum generators. Moreover there is a natural way in which one-body currents can be introduced in the point form (called the point-form spectator approximation) that satisfies the correct Poincaré and charge conservation properties.
We have shown that the PFSA produces results consistent with other impulse and spectator approximations. Within the range Schiavilla and Riska examined, for example, their impulse approximation and the PFSA (using G-K form factors) predict nearly identical results for $`A(Q^2)`$ and $`T_{20}(Q^2)`$; and though the zero they predict for $`B(Q^2)`$ falls near 2.0 GeV<sup>2</sup> rather than the PFSA’s 1.6, the fall-off from the data begins near 0.5 GeV<sup>2</sup> in both. The calculations of Kobushkin and Syamtomov (before using their approach of reduced transition amplitudes) and the PFSA (G-K) results for $`A(Q^2)`$ and $`B(Q^2)`$ nearly duplicate each other, as do the results for $`T_{20}(Q^2)`$ except in the high-momentum range, where no data are currently available. And although the calculations of Chung et al. with earlier G-K form factors and the Argonne $`v_{14}`$ potential fall quite close to the data for all three observables, their calculations using Höhler form factors and potentials show the same salient points: the deviation from data beginning in $`A(Q^2)`$ near 1 GeV<sup>2</sup> and in $`B(Q^2)`$ near 0.5 GeV<sup>2</sup>; the location of the first zero in $`B(Q^2)`$ between 1.5–2.0 GeV<sup>2</sup>; and the first minimum in $`T_{20}(Q^2)`$ around 0.8 GeV<sup>2</sup>. The work of Van Orden et al. gives results similar to Chung et al., except that the high-$`Q^2`$ behavior of $`A(Q^2)`$ and $`B(Q^2)`$ is nearly level at $`10^8`$, while Chung et al. and the PFSA show gradual decreases at $`10^{11}`$ for $`A(Q^2)`$ and at $`10^8`$ for $`B(Q^2)`$. In contrast, Lev et al. calculate that the Höhler form factors produce results for $`A(Q^2)`$ that lie close to the data, while the G-K form factors fit $`A(Q^2)`$ up to 2 GeV<sup>2</sup> but produce results increasingly higher than the data thereafter. Their results for $`B(Q^2)`$ and $`T_{20}(Q^2)`$ are similar to those of Van Orden et al. and Chung et al.
This work also addressed the sensitivity of PFSA results to different nucleon-nucleon interactions and different parameterizations of the nucleon form factors. In almost every instance it was found that the two nucleon-nucleon potentials produced only slight, if any, differences in the form factors and elastic observables. This may not be surprising considering that the Argonne $`v_{18}`$ and Reid ’93 nucleon-nucleon interactions produce nearly identical momentum-space wavefunctions on the momentum scale of interest.
Much more pronounced were the differences between the Gari-Krümpelmann and the Mergell-Meissner-Drechsel parameterizations of the nucleon form factors. As the momentum transfers become higher, the two often predict significantly different results. These are most notable in the deuteron form factors $`G_E`$ and $`G_Q`$, which are sensitive to the neutron electric form factor; the G-K parameterization, whose neutron electric form factor falls off markedly more rapidly than the MMD, produces deuteron form factor zeroes in the intermediate range that occur at higher momentum transfers in the MMD results. That this phenomenon is due to the neutron electric form factor is supported by the similarity of results for $`G_M`$, where the nucleon magnetic form factors dominate the calculations.
Aside from differences due to varying potentials and nucleon form factors, the consequences of the point form’s nontrivial momentum transfer have also been examined. We have shown that the point-form momentum transferred to a nucleon is greater than the $`Q^2`$ transferred to the deuteron, and that its deviation increases with increasing $`Q^2`$. This results in a lowering of the deuteron form factors and elastic scattering observables compared to nonrelativistic calculations. The greater magnitude of the point-form momentum transfer causes a quicker fall-off of the nucleon form factors, and some deviation from nonrelativistic calculations was attributed to this.
Additional dynamically consistent two-body currents must be added to the PFSA in order to bring the calculations into agreement with data at intermediate to high momentum transfers. Such currents have not been considered in this preliminary work; but other calculations that include dynamical two-body currents (see references ) suggest that the addition of dynamical currents is capable of reconciling the differences between various impulse or spectator approximations with data.
In the point form it is quite easy to interpret Feynman diagrams for nucleon-nucleon scattering with the production of a photon as a current matrix element satisfying the requirements given in Section III. However, while the addition of such current matrix elements may produce better agreement with data, it does not provide a systematic procedure for constructing two-body currents. What is needed is a procedure for constructing conserved currents from one-body currents and the dynamical mass operator. Models based on such a procedure are being developed.
## A Point-form Momentum Transfer
The momentum transfer $`(p_i^{}p_i)`$ seen by nucleon $`i`$ can be computed following Ref. . Suppose that the momentum transfer is along the z-axis.
$`B(v_{in})`$ $`=`$ $`\left(\begin{array}{cccc}\mathrm{cosh}\mathrm{\Delta }/2& 0& 0& \mathrm{sinh}\mathrm{\Delta }/2\\ 0& 1& 0& 0\\ 0& 0& 1& 0\\ \mathrm{sinh}\mathrm{\Delta }/2& 0& 0& \mathrm{cosh}\mathrm{\Delta }/2\end{array}\right),`$ (A5)
$`B(v_{out})`$ $`=`$ $`\left(\begin{array}{cccc}\mathrm{cosh}\mathrm{\Delta }/2& 0& 0& \mathrm{sinh}\mathrm{\Delta }/2\\ 0& 1& 0& 0\\ 0& 0& 1& 0\\ \mathrm{sinh}\mathrm{\Delta }/2& 0& 0& \mathrm{cosh}\mathrm{\Delta }/2\end{array}\right),`$ (A10)
are the boosts that take the deuteron from the center of momentum frame to the Breit frame (where $`𝐏_{\mathrm{tot}}^{}=𝐏_{\mathrm{tot}}`$.) For elastic scattering,
$$\mathrm{sinh}\mathrm{\Delta }/2=\sqrt{\frac{Q^2}{4m_D^2}}.$$
(A11)
The initial energies and momenta are then:
$`E_1`$ $`=`$ $`\omega \mathrm{cosh}\mathrm{\Delta }/2+k_z\mathrm{sinh}\mathrm{\Delta }/2;`$ (A12)
$`p_{1z}`$ $`=`$ $`k_z\mathrm{cosh}\mathrm{\Delta }/2+\omega \mathrm{sinh}\mathrm{\Delta }/2;`$ (A13)
$`E_2`$ $`=`$ $`\omega \mathrm{cosh}\mathrm{\Delta }/2k_z\mathrm{sinh}\mathrm{\Delta }/2;`$ (A14)
$`p_{2z}`$ $`=`$ $`k_z\mathrm{cosh}\mathrm{\Delta }/2+\omega \mathrm{sinh}\mathrm{\Delta }/2;`$ (A15)
where $`\omega `$ and $`𝐤`$ are center of momentum variables.
In this notation, $`k_z`$ refers to the relative z-axis momentum of particle one. This convention gives rise to the following relations:
$`\omega ^{}`$ $`=`$ $`\omega \mathrm{cosh}\mathrm{\Delta }k_z\mathrm{sinh}\mathrm{\Delta };`$ (A16)
$`k_z^{}`$ $`=`$ $`k_z\mathrm{cosh}\mathrm{\Delta }\omega \mathrm{sinh}\mathrm{\Delta };`$ (A17)
where the minus signs are used when particle one is struck, the plus signs when particle two is struck. Suppose for illustration that particle one is struck. The final energies and momenta will then be
$`E_1^{}`$ $`=`$ $`\omega \mathrm{cosh}3\mathrm{\Delta }/2k_z\mathrm{sinh}3\mathrm{\Delta }/2;`$ (A18)
$`p_{1z}^{}`$ $`=`$ $`k_z\mathrm{cosh}3\mathrm{\Delta }/2\omega \mathrm{sinh}3\mathrm{\Delta }/2;`$ (A20)
$`E_2^{}=E_2;p_{2z}^{}=p_{2z}.`$
Now some hyperbolic trigonometry reveals that
$$(p_1^{}p_1)^2=4(k_z^2\omega ^2)\mathrm{sinh}^2\mathrm{\Delta }.$$
(A21)
Since
$$\mathrm{sinh}\mathrm{\Delta }=2\sqrt{\frac{Q^2}{4m_D^2}}\sqrt{1+\frac{Q^2}{4m_D^2}},$$
(A22)
and
$$k_z^2\omega ^2=k_z^2m_n^2𝐤^2=(m_N^2+𝐤_{}^2),$$
(A23)
the resulting Equation 51 is established.
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# 1 Introduction
## 1 Introduction
Anti-de Sitter black holes play a major role in the AdS/CFT conjecture (see for a review), and they have also received interest in the context of brane-world scenarios based on the setup of Randall and Sundrum . The purpose of this note is to briefly investigate the basic property of black holes, Hawking radiation, in anti-de Sitter spacetimes.
In the 2+1 dimensional special case, the BTZ black holes , Hawking radiation of a massles conformally coupled scalar field was investigated in detail by Hyun et. al. . Here we will instead take a more generic approach (while sacrifying detail), and extend some standard ways to analyze Hawking radiation from Schwarzschild black holes to the anti-de Sitter case in a generic dimension. The standard methods in question are to investigate eternal black holes and the role of different vacua based on different boundary conditions at the horizons, and to investigate radiation in a spherical collapse geometry while the quantum field is in a natural initial vacuum. Customarily, the focus is on finding the leading thermal characteristics of the radiation. For this purpose it is sufficient to invoke a geometric optics approximation which neglects backscattering of outgoing waves from the spacetime curvature and essentially reduces the problem to a two dimensions. The neglected effects would give rise to a gray-body factor in the thermal emission spectrum. The gray-body factors can be found by investigating absorption by AdS black holes.
We will also investigate the back-reaction to the geometry by taking into account the selfinteraction effect analyzed by Kraus, Parikh and Wilczek (see also ). This approach views Hawking radiation as a tunneling process across the horizon. Here we extend the approach of to the anti-de Sitter case, and find similar results.
## 2 Hawking radiation from eternal black holes
One standard method to derive the Hawking radiation is to use a Bogoliubov transformation between two basis of annihilation and creation operators corresponding to mode expansions of the field operator in two preferred coordinate systems used to describe a black hole: the asymptotic coordinates and the Kruskal coordinates. In the asymptotic coordinates, the AdS<sub>d+1</sub> black hole metric is (see e.g. )
$$ds^2=F(r)dt^2+\frac{dr^2}{F(r)}+r^2d\mathrm{\Omega }_{d1}^2$$
(1)
where
$$F(r)=1\frac{\mu }{r^{d2}}+r^2,$$
(2)
and we work in units where the AdS radius $`l=1`$. The parameter $`\mu `$ is proportional to the ADM mass $`M`$ of the black hole,
$$M=\frac{(d1)A_{d1}}{16\pi G_{d+1}}\mu $$
(3)
where $`G_{d+1}`$ is Newton’s constant in $`d+1`$ dimensions and $`A_{d1}=2\pi ^{d/2}/\mathrm{\Gamma }(d/2)`$ is the volume of a unit $`(d1)`$-sphere. The explicit formula for the horizon radius $`r_H`$ can be found by solving the polynomial equation $`F(r_H)=0`$. For example, in $`4+1`$ dimensions,
$$r_H=\sqrt{\frac{1}{2}\left(\sqrt{1+4\mu }1\right)}.$$
(4)
Then, the Hawking temperature $`T_H`$ can be found by looking at the periodicity of the Euclidian section of the metric near $`r_H`$. The generic result is
$$T_H=\frac{1}{4\pi }F^{}(r_H)$$
(5)
where $`{}_{}{}^{}=d/dr`$. E.g. in 4+1 dimensions,
$$T_H=\frac{\sqrt{1+4\mu }}{2\pi r_H}.$$
(6)
To investigate mode solutions of a field equation, it is useful to introduce the tortoise coordinate $`r_{}`$,
$$r_{}=_{r_H}^r\frac{d\widehat{r}}{F(\widehat{r})}.$$
(7)
Again, in $`4+1`$ dimensions the explicit formula is
$$r_{}=\frac{1}{\sqrt{1+4\mu }}\left(r_0\mathrm{arctan}\left(\frac{r}{r_0}\right)+\frac{1}{2}r_H\mathrm{ln}\left(\frac{rr_H}{r+r_H}\right)\right),$$
(8)
where $`r_H`$ is the horizon radius and $`r_0`$ is a shorthand notation denoting a radius
$$r_0=\sqrt{\frac{1}{2}\left(\sqrt{1+4\mu }+1\right)}.$$
(9)
(In 4+1 dimensions, the equation $`F(r)=0`$ has 4 complex zeroes: 2 real zeroes at $`\pm r_H`$ and 2 imaginary zeroes at $`\pm ir_0`$.) It is also convenient to introduce the null coordinates
$`u`$ $`=`$ $`tr_{}`$
$`v`$ $`=`$ $`t+r_{}`$ (10)
where $`r_{}`$ is the tortoise coordinate. Then the metric takes the form
$$ds^2=F(r)dudv+r^2d\mathrm{\Omega }_{d1}^2$$
(11)
and the solutions to the wave equation are infalling and outgoing partial waves.
In the region outside the past and future horizons, the Kruskal coordinates $`U,V`$ are defined as follows:
$`U`$ $`=`$ $`\mathrm{exp}\left(2\pi T_Hu\right)`$
$`V`$ $`=`$ $`\mathrm{exp}\left(2\pi T_Hv\right)`$ (12)
In $`U,V`$ coodinates, the metric takes the form
$$ds^2=\frac{F(r)e^{4\pi T_Hr_{}}}{(2\pi T_H)^2}dUdV+r^2d\mathrm{\Omega }_{d1}^2.$$
(13)
In 4+1 dimensions, the explicit form is
$$ds^2=\frac{r_H^2}{1+4\mu }(1+\frac{r_0^2}{r^2})(r_H+r)^2\mathrm{exp}\left(\frac{2r_0}{r_H}\mathrm{arctan}\left(\frac{r}{r_0}\right)\right)dUdV+r^2d\mathrm{\Omega }_3^2.$$
(14)
In Kruskal coordinates, the metric can be extended over the whole spacetime, except for the origin $`r=0`$, which corresponds to the true curvature singularity of the black hole.
Next, we will focus on s-waves, adopt the geometric optics approximation, and truncate to two dimensions, following the classic paper by Unruh . The discussion also overlaps with <sup>3</sup><sup>3</sup>3We thank D. Klemm for bringing this reference to our attention..
For a quantum field in the black hole background, there are two canonical choices for a natural vacuum state. If one wants to mimic a situation where the black hole is created by collapsing matter, one requires the field to be in a vacuum corresponding absence of positive energy modes in the $`U`$ and $`v`$ coordinates near the past horizon $`V=0`$. This boundary condition refers only to the past of the asymptotic region of spacetime, an is known as the Unruh vacuum. Another vacuum choice, the Hartle-Hawking vacuum, refers to a mixture of boundary conditions in the past and future horizons. Now one requires the absence of positive energy $`U`$ modes near $`V=0`$, and the absence of positive energy $`V`$ modes near $`U=0`$. Physically, this mimics a black hole in thermal equilibrium with an external heat bath. The task is to compare these two vacuua with the Boulware vacuum which is a natural vacuum for a fiducial observer in the asymptotic region. The Boulware vacuum corresponds to absence of positive energy $`u`$ and $`v`$ modes. To complete the discussion of vacua, we will also need to take into account the reflective boundary condition at the boundary of the anti-de Sitter space. We will do at the end of this section.
To compute the Bogoliubov transformations between the different natural modes, we use the standard trick . Let us focus on the outgoing modes first. To begin with, we define the modes
$`\varphi _{+,\omega }`$ $`=`$ $`\mathrm{exp}(i\omega u)=\theta (U)(U)^{\frac{i\omega }{2\pi T_H}},(U<0)`$ (15)
$`\varphi _{,\omega }`$ $`=`$ $`\mathrm{exp}(i\omega u)=\theta (U)U^{\frac{i\omega }{2\pi T_H}},(U>0).`$ (16)
Then, to find a complete basis for positive energy $`U`$ modes we consider the linear combinations which extend over the whole $`V=0`$ line,
$`\varphi _{1,\omega }`$ $`=`$ $`\varphi _{+,\omega }+C_1\varphi _{,\omega }^{}`$ (17)
$`=`$ $`\theta (U)(U)^{\frac{i\omega }{2\pi T_H}}+C_1\theta (U)U^{\frac{i\omega }{2\pi T_H}}`$
$`\varphi _{2,\omega }`$ $`=`$ $`\varphi _{,\omega }+C_2\varphi _{+,\omega }^{}`$ (18)
$`=`$ $`\theta (U)U^{\frac{i\omega }{2\pi T_H}}+C_2\theta (U)(U)^{\frac{i\omega }{2\pi T_H}}`$
We then demand that $`\varphi _{1,\omega }`$ is a positive energy Kruskal mode: that for $`\omega >0`$ it must be analytic in the lower half complex $`U`$-plane. This condition is satisfied if the coefficient $`C_1`$ is
$$C_1=\mathrm{exp}\left(\frac{\omega }{2T_H}\right)$$
(19)
A similar condition for $`\varphi _{2,\omega }`$ fixes the coefficient $`C_2`$:
$$C_2=\mathrm{exp}\left(\frac{\omega }{2T_H}\right)$$
(20)
Hence $`C_1=C_2C`$.
From this we can compute the (unnormalized) $`\alpha `$ and $`\beta `$ Bogoliubov coefficients which denote the overlap of a positive energy $`U`$ mode with positive and negative energy $`u`$ modes:
$`\alpha =(\varphi _{1,\omega },\varphi _{+,\omega })=(\varphi _{2,\omega },\varphi _{,\omega })`$ $``$ $`1`$ (21)
$`\beta =(\varphi _{1,\omega },\varphi _{,\omega }^{})=(\varphi _{2,\omega },\varphi _{+,\omega }^{})`$ $``$ $`C`$ (22)
(for $`i=1,2`$). The ratio of the two coefficients is thus
$$\left|\frac{\beta }{\alpha }\right|^2=\left|C\right|^2=\mathrm{exp}\left(\frac{\omega }{T_H}\right)$$
(23)
Using the normalization condition of the Bogoliubov coefficients,
$$\left|\alpha \right|^2\left|\beta \right|^2=1$$
(24)
we find that the average occupation number for positive energy $`u`$ modes, seen by a fiducial observer when the quantum field is in a vacuum with respect to positive energy $`U`$ modes, simplifies to the expected form of a Bose-Einstein distribution,
$$\overline{n}_\omega =\left|\beta \right|^2=\frac{1}{\mathrm{exp}\left(\frac{\omega }{T_H}\right)1}$$
(25)
where $`T_H`$ is the Hawking temperature of the AdS<sub>d+1</sub> black hole.
A similar relation holds between the $`V`$ and $`v`$ modes also. Now, we take into account the reflection from the boundary of adS space. Ref. considered different possibilities for the boundary condition at infinity: Dirichlet, Neumann, and Robin boundary conditions. What is the preferred boundary condition? Let us leave the geometric optics approximation for the moment and consider mode solutions to the exact wave equation for a free scalar field in the adS-black hole space time. In Minkowski signature, the mode solutions can fall into two categories, nonnormalizable solutions $`\varphi _{\omega ,\stackrel{}{k}}^{()}`$ and normalizable solutions $`\varphi _{\omega ,\stackrel{}{k}}^{(+)}`$, with the asymptotic behavior
$$\varphi _{\omega ,\stackrel{}{k}}^{(\pm )}(t,r,\mathrm{\Omega })r^{2h_\pm }\stackrel{~}{\varphi }^\pm (t,\mathrm{\Omega })(r\mathrm{}),$$
(26)
where $`h_\pm `$ are parameters related to the mass $`\mu =`$ and the dimension of the space $`d`$ by
$$2h_\pm =\frac{1}{2}(d\pm \sqrt{d^2+4\mu ^2}).$$
(27)
The quantized field $`\varphi `$ is expanded as a linear combination of the normalizable modes. Their decay behavior at the boundary corresponds to reflection. The exact mode solutions $`\varphi ^{(+)}`$ are easy to find in $`2+1`$ dimensions in terms of hypergeometric functions . Near the black hole horizon, the normalizable modes reduce to a form
$$\varphi ^{(+)}(e^{i\omega u}+e^{i\omega v+i2\theta _0})e^{in\varphi }$$
(28)
where $`\theta _0`$ is a phase shift factor, its exact form can be found in . Thus, in the geometric optics approximation, the modes which take into account the reflection from the boundary and are appropriate to a fiducial observer are a linear combination of the positive energy $`u`$ and $`v`$ modes,
$$\varphi _\omega =e^{i\omega u}+e^{i\omega v+i2\theta _0}.$$
(29)
For Kruskal modes, one must consider the corresponding linear combination of the positive energy $`U`$ and $`V`$ modes. Thus, with the reflective boundary condition, the appropriate vacuum is the Hartle-Hawking vacuum. Then, a fiducial observer sees a thermal spectrum for both infalling and outgoing modes. In the eternal adS geometry, the Unruh vacuum is not well defined with respect to the boundary condition at infinity. It can be viewed as an artificial construction describing the very onset of radiation, where only outgoing modes are thermally excited and they have not yet reflected back from the We will return to this issue in Section 4.
## 3 Hawking radiation as tunneling
Recently, a method to describe Hawking radiation as a tunneling process, where a particle moves in dynamical geometry, was developed by Kraus and Wilczek and elaborated upon by Parikh and Wilczek . This method also gives a leading correction to the emission rate arising from loss of mass of the black hole correponding to the energy carried by the radiated quantum. This method was also investigated in the context of black holes in string theory , and it was demonstrated that in the string picture of microstates of the black hole, the correction to the emission rate corresponds to a difference between counting of states in the microcanonical and canonical ensembles. However, in all these investigations the black holes have had asymptotically flat spacetime geometry. We now extend the investigation to black holes in AdS spacetime. We will base our treatment on the presentation of .
A convenient trick in the method of is to write the black hole metric in a coordinate system where constant time slices are flat, without a singularity at the horizon. In these coordinates, the Schwarzschild metric takes the form
$$ds^2=(1\frac{2M}{r})dt^2+2\sqrt{\frac{2M}{r}}dtdr+dr^2+r^2d\mathrm{\Omega }^2.$$
(30)
These coordinates were first introduced 80 years ago by Painlevé , but then disappeared from general knowledge, until they were independently rediscovered in and used to investigate black hole quantum mechanics. We will now derive an analogue of the Painlevé coordinates for AdS black holes, which we shall refer to as the AdS-Painlevé coordinates.
By analogue to the asymptotically flat black holes, the AdS-Painlevé coordinates should have the property that constant time slices of the AdS black hole metric (1) will have the same geometry as constant time slices of a global AdS<sub>d+1</sub> metric
$$ds^2=(1+r^2)dt^2+\frac{dr^2}{(1+r^2)}+r^2d\mathrm{\Omega }_{d1}^2.$$
(31)
Thus, we perform a coordinate transformation $`t=\widehat{t}+f(r)`$ so the metric (1) takes the form
$$ds^2=F(r)d\widehat{t}^2+2f^{}(r)F(r)d\widehat{t}dr+\left(\frac{1}{F(r)}F(r)(f^{}(r))^2\right)dr^2+r^2d\mathrm{\Omega }^2$$
(32)
and then demand that on constant $`\widehat{t}`$ slices the metric reduces to
$$ds^2=(1+r^2)^1dr^2+r^2d\mathrm{\Omega }^2.$$
(33)
This implies that
$$f^{}(r)=\frac{1}{r^{\frac{d2}{2}}F(r)}\sqrt{\frac{\mu }{1+r^2}}$$
(34)
so the AdS-Painlevé metric reads as follows:
$$ds^2=F(r)d\widehat{t}^2+\frac{2}{r^{\frac{d2}{2}}}\sqrt{\frac{\mu }{1+r^2}}d\widehat{t}dr+(1+r^2)^1dr^2+r^2d\mathrm{\Omega }^2.$$
(35)
Now we move on to discuss Hawking radiation. The (s-wave) quanta of a massless scalar field follow radial light-like geodesics
$$\dot{r}=\frac{\sqrt{\mu }}{r^{\frac{d2}{2}}}\sqrt{1+r^2}\pm (1+r^2)$$
(36)
where the $`(+)`$ sign corresponds to an outgoing geodesic and the $`()`$ sign corresponds to an ingoing geodesic, respectively. Next we take into account the response of the background geometry to an emitted quantum of frequency $`\omega `$. We keep the total mass $`M`$ of the spacetime fixed, but in order to take into account the energy carried by the quantum, we replace $`\mu `$ in (36) by $`\mu ^{}`$,
$$\mu ^{}\frac{16\pi }{(d1)A_{d1}}\left(M\omega \right).$$
(37)
Note that at the horizon,
$$\dot{r}|_{r_H}=0.$$
(38)
As the particle travels across the horizon from $`r_{in}`$ to $`r_{out}`$, its action<sup>4</sup><sup>4</sup>4Note that in the local point particle description used in this section, the issue of boundary conditions at infinity does not arise. receives an imaginary contribution
$$\mathrm{Im}S=\mathrm{Im}_{r_{in}}^{r_{out}}p_r𝑑r=\mathrm{Im}_{r_{in}}^{r_{out}}_H\frac{dH}{\dot{r}}𝑑r$$
(39)
where on the last line we switched the order of integration and used the Hamilton’s equation $`\dot{r}=\frac{dH}{dp_r}`$. Next, we substitute from (36) the radial velocity along the outgoing geodesic, and use $`dH=d(M\omega )=d\omega `$:
$$\mathrm{Im}S=\mathrm{Im}_0^\omega 𝑑\omega ^{}_{r_{in}}^{r_{out}}𝑑r\frac{1+\sqrt{\frac{\mu ^{}}{r^{d2}}(1+r^2)^1}}{F(r)}.$$
(40)
The only imaginary contribution to the radial integral comes from the pole at $`r_H`$. Then
$$\mathrm{Im}S=\pi _0^\omega 𝑑\omega ^{}\frac{2r_H}{r_H^2d+(d2)}.$$
(41)
On the other hand, after solving for $`r_H`$ as a function of $`\mu `$, we can derive that
$$r_H^{d2}\frac{dr_H}{d\mu }=\frac{r_H}{r_H^2d+(d2)}.$$
(42)
Substituting this into the integral yields
$`\mathrm{Im}S`$ $`=`$ $`\pi {\displaystyle \frac{(d1)A_{d1}}{16\pi }}{\displaystyle _{\mu \frac{16\pi }{(d1)A_{d1}}\omega }^\mu }𝑑\mu ^{}2r_H^{d2}{\displaystyle \frac{dr_H}{d\mu ^{}}}`$ (43)
$`=`$ $`{\displaystyle \frac{1}{8}}A_{d1}\left(r^{d1}(M)r^{d1}(M\omega )\right)={\displaystyle \frac{1}{2}}\mathrm{\Delta }S_{BH}`$
where $`\mathrm{\Delta }S_{BH}=S_{BH}(M)S_{BH}(M\omega )`$ is the difference of the entropies of the black hole before and after the emission. Thus, the tunneling probability for the particle is
$$\mathrm{\Gamma }=\mathrm{exp}(\mathrm{\Delta }S_{BH}).$$
(44)
If we Taylor expand $`\mathrm{\Delta }S_{BH}`$ in $`\omega `$, the leading term gives the thermal Boltzmann factors $`\mathrm{exp}(\omega /T_H)`$ for the emanating radiation. The second term represents corrections from the response of the background geometry to the emission of a quantum. The same result holds for emission from asymptotically flat black holes .
## 4 Particle creation by a collapsing spherical shell in AdS
We will now turn to a third way to analyze Hawking radiation from black holes, and investigate particle creation by a collapsing spherical body which forms a black hole in AdS. We will base our treatment on the discussion in , which in turn follows and . As in , the starting point is that we assume that in the remote past the spherical body is distended so much that it deforms the anti-de Sitter space. Thus, in the beginning we can assume that a quantum field is in a vacuum constructed with respect to global coordinates in AdS space. Now a convenient choice for the global coordinates is given by
$$ds^2=(\mathrm{sec}\rho )^2(dt^2+d\rho ^2)+(\mathrm{tan}\rho )^2d\mathrm{\Omega }_{d1}^2.$$
(45)
The (normalizable) mode solutions can be found e.g. in :
$$\varphi _{n,l}^{(+)}=e^{i\omega t}Y_{l,\{m\}}(\mathrm{\Omega })(\mathrm{cos}\rho )^{2h_+}(\mathrm{sin}\rho )^lP_n^{(l+\frac{d}{2}1,2h_+\frac{d}{2})}(\mathrm{cos}2\rho ),$$
(46)
where $`Y_{l,\{m\}}(\mathrm{\Omega })`$ is a spherical harmonic on $`S^{d1}`$, $`P_n^{(l+\frac{d}{2}1,2h_+\frac{d}{2})}`$ is a Jacobi polynomial, and
$$2h_+=\frac{d}{2}+\frac{1}{2}\sqrt{d^2+4m^2}\frac{d}{2}+\frac{1}{2}\nu $$
(47)
Due to the boundary conditions at the origin and at the boundary of AdS, the spectrum is discrete, with
$$\omega =2h_++2n+2l;n=0,1,2,\mathrm{}$$
(48)
Again, we will focus on $`s`$-waves ($`l=0`$). As in flat space, we expect the quantum to experience a strong redshift as it propagates across the collapsing body. Thus, in the remote past, we are most interested in high frequency modes. In the high frequency limit, the mode solution (46) reduces to a simplified form
$$\varphi _n^{(+)}(\mathrm{cos}\rho )^{(d1)/2}e^{i\omega _nt}\mathrm{cos}(\omega _n\rho \frac{d\pi }{4})$$
(49)
where we suppressed normalization factors. In other words, they take a form of a standing wave, a superposition of an ingoing and outgoing spherical wave, with a discrete spectrum. Note that using the radial coordinate $`r`$, the overall factor $`(\mathrm{cos}\rho )^{(d1)/2}(1/r)^{(d1)/2}`$ as $`r>>1`$, so we recover the expected overall decay factor for the amplitude. Now, we will add into the picture the collapsing body and try to compute the redshift due to the passage of the wave across it. As in , we will use the geometric optics approximation, and truncate the analysis to two dimensions to the $`t,r`$-plane by suppressing the overall decay factor of the waveform. For simplicity, we assume that the collapsing body is a thin shell of radius $`R`$, with $`R`$ monotonically decreasing in time. The truncated metric inside and outside the shell takes the form
$$ds^2=F_\pm dt_\pm ^2+F_\pm ^1(r)dr^2$$
(50)
where $`F_+(r)=1\frac{\mu }{r^{d2}}+r^2`$ and $`F_{}(r)=1+r^2`$. We then define the tortoise coordinates $`r_{}^\pm `$,
$$r_{}^\pm =\frac{dr}{F_\pm (r)},$$
(51)
and the null coordinates
$`u`$ $`=`$ $`t_+r_{}^+,v=t_++r_{}^+`$
$`U`$ $`=`$ $`t_{}r_{}^{},V=t_{}+r_{}^{}`$ (52)
so that the interior and exterior metrics are conformal to a flat metric. The tortoise coordinate in the interior is $`r_{}^{}=\mathrm{arctan}r`$, so the origin $`r=0`$ corresponds to $`r_{}^{}=0`$. In terms of the null coordinates $`U,V`$, the origin is then at $`VU=0`$. Note also that $`r_{}^{}=\rho `$, where $`\rho `$ is the coordinate that appears in the global metric (45). The exterior and interior null coordinates are related by
$`v`$ $`=`$ $`\beta (V)`$
$`U`$ $`=`$ $`\alpha (u)`$ (53)
where $`\alpha (u)`$ and $`\beta (V)`$ are to be determined below. In the $`(t,r)`$ coordinates, the passage of a wave across the shell turns to a reflection condition at the origin:
$$v=\beta (V)=\beta (U)=\beta (\alpha (u)).$$
(54)
Thus, in the asymptotic region (near the boundary), the waves have a phase structure
$$\stackrel{~}{\varphi }^{(+)}e^{i\omega _nv}e^{i\omega _n\beta (\alpha (u))}.$$
(55)
To find the functions $`\alpha ,\beta `$, we match the interior and exterior metrics across the collapsing shell at $`r=R(\tau )`$. Here $`\tau `$ denotes the shell time, which is related to the time coordinates $`t_\pm `$ in the interior and exterior of the shell through
$$ds^2=[F_\pm dt_\pm ^2+F_\pm ^1dr^2]_{|r=R(\tau )}=d\tau ^2.$$
(56)
It is easiest to consider the derivatives
$`\alpha ^{}(u)`$ $`=`$ $`{\displaystyle \frac{dU}{du}}={\displaystyle \frac{\dot{U}}{\dot{u}}}`$
$`\beta ^{}(V)`$ $`=`$ $`{\displaystyle \frac{dv}{dV}}={\displaystyle \frac{\dot{v}}{\dot{V}}},`$ (57)
(at the shell) where $`=d/d\tau `$. Using the definition (56), we obtain
$`{\displaystyle \frac{dU}{du}}`$ $`=`$ $`{\displaystyle \frac{F_+(R)[\sqrt{F_{}(R)+\dot{R}^2}\dot{R}]}{F_{}(R)[\sqrt{F_+(R)+\dot{R}^2}\dot{R}]}}`$
$`{\displaystyle \frac{dv}{dV}}`$ $`=`$ $`{\displaystyle \frac{F_{}(R)[\sqrt{F_+(R)+\dot{R}^2}+\dot{R}]}{F_+(R)[\sqrt{F_{}(R)+\dot{R}^2}+\dot{R}]}}.`$ (58)
As the radius of the shell approaches the horizon, we can approximate
$`F_+(R)`$ $``$ $`4\pi T_H(Rr_H)`$
$`F_{}(R)`$ $``$ $`F_{}(r_H)A.`$ (59)
Then, we can approximate
$$\frac{dU}{du}\frac{2\pi T_H(Rr_H)}{A\dot{R}}[\sqrt{A+\dot{R}^2}\dot{R}]$$
(60)
where we used $`\sqrt{\dot{R}^2}=\dot{R}`$ since $`\dot{R}<0`$ as the shell is collapsing and $`|\dot{R}|0`$ as a function of shell time. In the above, $`\dot{R}=\dot{R}_{|r_H}`$. Next, we relate $`U`$ to $`Rr_H`$ by expanding
$$UU(r_H)+\frac{dU}{dR}_{|R=r_H}(Rr_H)$$
(61)
and evaluating the derivative $`dU/dR`$ at the horizon, using the chain rule and the definitions () and (). We obtain
$$(Rr_H)(UU(r_H))\frac{A\dot{R}}{[\sqrt{A+\dot{R}^2}\dot{R}]}.$$
(62)
We substitute this to (), and obtain
$$\frac{dU}{du}2\pi T_H(UU(r_H))\kappa U+const.$$
(63)
where $`\kappa `$ is the surface gravity of the black hole. Integration then gives
$$\alpha (u)=e^{\kappa u}+const.$$
(64)
A similar calculation for $`dv/dV`$ gives
$$\frac{dv}{dV}\frac{A}{\dot{R}[\sqrt{A+\dot{R}^2}+\dot{R}]}c(=const.).$$
(65)
By integration,
$$\beta (V)=cV+const.$$
(66)
Thus, we find that in the asymptotic region the waves have a phase structure
$$\stackrel{~}{\varphi }^{(+)}e^{i\omega _nv}e^{i\omega _nc(e^{\kappa u}+const.)}.$$
(67)
To obtain modes where the outgoing wave is of standard form, we invert functionally and write
$$\stackrel{~}{\varphi }^{(+)}e^{i\omega _n\kappa ^1\mathrm{ln}[(v_0v)/c]}e^{i\omega _nu},$$
(68)
which is valid only for $`v<v_0`$. Now, we move back to $`(d+1)`$ dimensions and compare with the high-frequency limit of the global modes (49). Note that far in the asymptotic region, $`r\mathrm{}`$, the exterior tortoise coordinate reduces to the same form as $`\rho `$, $`r_{}^+\rho `$, so we can write the global mode as
$$\varphi ^{(+)}(\mathrm{cos}\rho )^{(d2)/2}(e^{i\omega _nvid\pi /2}e^{i\omega _nu}).$$
(69)
We want to compare this with the modes that we found in the collapsing shell geometry,
$$\stackrel{~}{\varphi }^{(+)}(\mathrm{cos}\rho )^{(d2)/2}(e^{i\omega _n\kappa ^1\mathrm{ln}[(v_0v)/c]}e^{i\omega _nu})(v<v_0),$$
(70)
where we added the overall decay factor. The Bogoliubov transformation follows the discussion in , and as a result we find that the outgoing modes are thermally excited, if the field is in a global vacuum. Thus, the global vacuum resembles an Unruh vacuum.
Naturally, the above result only applies to the onset of Hawking radiation as the black hole has formed, and does not address the issue of subsequent evolution. Since the AdS space can be viewed as a finite box, what will happen is that the outgoing radiation cannot escape to infinity but will slowly fill the box. Subsequently, the black hole will come to a thermal equilibrium with the surrounding thermal bath. This situation is descibed by the Hartle-Hawking vacuum in the eternal geometry. Note however that in dimensions $`d>2`$, if the initial size of the black hole is too small, the equilibrium will not be established before the black hole evaporates completely. See for a thorough analysis.
Note added
As we were finalizing this paper, the paper appeared, discussing Hawking radiation in the optical collapse geometry for spherically symmetric black holes.
Acknowledgment
We would like to thank Jorma Louko for useful comments.
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# 1 Introduction.
## 1 Introduction.
One of the most important problems in Quantum Field Theory is the study of the bound state spectra of non-abelian gauge theories. There are several approaches to this problem. For QCD like theories lattice gauge theory is probably the most popular approach since the approximation does not break the most important symmetry, gauge symmetry. Similarly for supersymmetric theories, supersymmetric DLCQ (SDLCQ)is probably the most powerful approach since the approximation does not break the most important symmetry, supersymmetry. In this paper will consider supersymmetric theories and follow this latter approach.
Long ago ’t Hooft showed that two dimensional models can be a powerful laboratory for studying the bound state problem. These models remain popular to this day since they are easy to solve and share many properties with their four dimensional cousins, most notably stable bound states. Supersymmetric two dimensional models are particularly attractive since they are also super-renormalizable. Given that the dynamics of gauge field is responsible for strong interaction and the formation of bound states, it comes as no surprise that a great deal of effort has gone into investing bound states of pure glue in supersymmetric models . While such theories capture the essential properties of the mass spectrum and some of them are relevant for the string theory , the wavefunctions are quite different from the ones for mesons and baryons. Extensive study of the meson spectrum of non–supersymmetric theories has been done (see for a review), but the problem has not been addressed in a context of supersymmetric models. In this paper we will introduce seven new two dimensional supersymmetric models that have not been previously studied that are particularly useful for studying mesons within a two dimensional supersymmetric laboratory. Throughout this paper we use a word “meson” to indicate the group structure of the state. Namely we define a meson as a bound state whose wavefunction can be written as a linear combination of parton chains, each chain starts and ending with a creation operator in fundamental representation. In supersymmetric theories the states, defined this way, can have either bosonic or fermionic statistics.
To simplify the calculation we will consider only the large $`N`$ limit , which has proven to be a powerful approximation for bound state calculations. While baryons can be constructed in this limit , they have an infinite number of partons and thus practical calculations for such states are complicated. In this paper we concentrate our attention on the mesonic spectrum. Note that throughout this paper we completely ignore the zero mode problem , however it is clear that considerable progress on this issue could be made following our earlier work on the zero modes of the two dimensional supersymmetric model with only adjoint fields .
The paper has the following organization. In section 2 we consider the three dimensional SQCD and dimensionally reduce it to $`1+1`$. We perform the light-cone quantization of the resulting theory by applying canonical commutation relations at fixed $`x^+`$ and choosing the light-cone gauge ($`A^+=0`$) for the vector field. After solving the constraint equations we end up with a model containing $`4`$ dynamical fields. We construct the supercharge for the dimensionally reduced theory and observe that it can also be used to define models with less supersymmetry. In particular, in section 3 we study the mesonic spectrum of systems without dynamical quarks. We find that one of these systems (we call it $`A\lambda `$) has many light states in SDLCQ, which probably give rise to a continuous spectrum in the continuum limit. The other system, containing only dynamical gluinos, seems not to have a well-defined bound state problem: all masses are pushed to infinity in the continuum limit. In section 4 we study the systems which do not include the adjoint scalar. They all share the same properties: a well-defined continuum spectrum and the existing of a critical value of the coupling constant at which the lowest mass bound state becomes massless. Finally, in section 5 we study the remaining theories which include the adjoint boson and at least one of the fundamental fields. We find that all this models have a continuous spectrum, which seems to be a general property of two dimensional supersymmetric systems with adjoint scalars .
## 2 Supersymmetric Systems with fundamental matter.
We consider the family of supersymmetric models in two dimensions which can be obtained as the result of dimensional reduction of SQCD<sub>2+1</sub> and possible truncation of some fields in the resulting two dimensional theory. Our starting point is the three dimensional action:
$`S`$ $`=`$ $`{\displaystyle }d^3x\text{tr}({\displaystyle \frac{1}{4}}F_{\mu \nu }F^{\mu \nu }+{\displaystyle \frac{i}{2}}\overline{\mathrm{\Lambda }}\mathrm{\Gamma }^\mu D_\mu \mathrm{\Lambda }+D_\mu \xi ^{}D^\mu \xi +i\overline{\mathrm{\Psi }}D_\mu \mathrm{\Gamma }^\mu \mathrm{\Psi }`$ (1)
$``$ $`g^{}[\overline{\mathrm{\Psi }}\mathrm{\Lambda }\xi +\xi ^{}\overline{\mathrm{\Lambda }}\mathrm{\Psi }]).`$
This action describes the system of a gauge field $`A_\mu `$ and its superpartner $`\mathrm{\Lambda }`$, both taking values in the adjoint representation of $`SU(N)`$, and two complex fields, a scalar $`\xi `$ and a Dirac fermion $`\mathrm{\Psi }`$, transforming according to the fundamental representation of the same group. Thus in matrix notation the covariant derivatives are given by:
$$D_\mu \mathrm{\Lambda }=_\mu \mathrm{\Lambda }+ig[A_\mu ,\mathrm{\Lambda }],D_\mu \xi =_\mu \xi +ig^{}A_\mu \xi ,D_\mu \mathrm{\Psi }=_\mu \mathrm{\Psi }+ig^{}A_\mu \mathrm{\Psi }.$$
(2)
The action (1) is invariant under supersymmetry transformations which are parameterized by a two–component Majorana fermion $`\epsilon `$:
$`\delta A_\mu ={\displaystyle \frac{i}{2}}\overline{\epsilon }\mathrm{\Gamma }_\mu \mathrm{\Lambda },\delta \mathrm{\Lambda }={\displaystyle \frac{1}{4}}F_{\mu \nu }\mathrm{\Gamma }^{\mu \nu }\epsilon ,`$
$`\delta \xi ={\displaystyle \frac{i}{2}}\overline{\epsilon }\mathrm{\Psi },\delta \mathrm{\Psi }={\displaystyle \frac{1}{2}}\mathrm{\Gamma }^\mu \epsilon D_\mu \xi .`$ (3)
We introduced the commutator of the Dirac matrices:
$$\mathrm{\Gamma }^{\mu \nu }=\frac{1}{2}\left[\mathrm{\Gamma }^\mu \mathrm{\Gamma }^\nu \right].$$
Using standard techniques one can evaluate the Noether current corresponding to these transformations:
$`\overline{\epsilon }q^\mu `$ $`=`$ $`{\displaystyle \frac{i}{4}}\overline{\epsilon }\mathrm{\Gamma }^{\alpha \beta }\mathrm{\Gamma }^\mu \text{tr}\left(\mathrm{\Lambda }F_{\alpha \beta }\right)+{\displaystyle \frac{i}{2}}D^\mu \xi ^{}\overline{\epsilon }\mathrm{\Psi }+{\displaystyle \frac{i}{2}}\xi ^{}\overline{\epsilon }\mathrm{\Gamma }^{\mu \nu }D_\nu \mathrm{\Psi }`$ (4)
$``$ $`{\displaystyle \frac{i}{2}}\overline{\mathrm{\Psi }}\epsilon D^\mu \xi +{\displaystyle \frac{i}{2}}D_\nu \overline{\mathrm{\Psi }}\mathrm{\Gamma }^{\mu \nu }\epsilon \xi .`$
We will consider the reduction of this system to two dimensions, which means that the field configurations are assumed to be independent of the space–like dimension $`x^2`$. In the resulting two dimensional system we will implement light–cone quantization, which means that initial conditions as well as canonical commutation relations will be imposed on a light–like surface $`x^+=const`$. In particular we construct the supercharge by integrating the current (4) over the light–like surface:
$`\overline{\epsilon }Q`$ $`=`$ $`{\displaystyle }dx^{}dx^2({\displaystyle \frac{i}{4}}\overline{\epsilon }\mathrm{\Gamma }^{\alpha \beta }\mathrm{\Gamma }^+\text{tr}\left(\mathrm{\Lambda }F_{\alpha \beta }\right)+{\displaystyle \frac{i}{2}}D_{}\xi ^{}\overline{\epsilon }\mathrm{\Psi }+{\displaystyle \frac{i}{2}}\xi ^{}\overline{\epsilon }\mathrm{\Gamma }^{+\nu }D_\nu \mathrm{\Psi }`$ (5)
$``$ $`{\displaystyle \frac{i}{2}}\overline{\mathrm{\Psi }}\epsilon D^+\xi +{\displaystyle \frac{i}{2}}D_\nu \overline{\mathrm{\Psi }}\mathrm{\Gamma }^{+\nu }\epsilon \xi ).`$
Since all fields are assumed to be independent of $`x^2`$, the integration over this coordinate gives just a constant factor, which we absorb by a field redefinition.
If we consider a specific representation for the Dirac matrices in three dimensions:
$$\mathrm{\Gamma }^0=\sigma _2,\mathrm{\Gamma }^1=i\sigma _1,\mathrm{\Gamma }^2=i\sigma _3,$$
(6)
then the Majorana fermion $`\mathrm{\Lambda }`$ can be chosen to be real, it is also convenient to write the fermions in the component form:
$$\mathrm{\Lambda }=(\lambda ,\stackrel{~}{\lambda })^T,\mathrm{\Psi }=(\psi ,\stackrel{~}{\psi })^T,Q=(Q^+,Q^{})^T$$
(7)
In terms of this decomposition the superalgebra has explicit $`(1,1)`$ form:
$$\{Q^+,Q^+\}=2\sqrt{2}P^+,\{Q^{},Q^{}\}=2\sqrt{2}P^{},\{Q^+,Q^{}\}=0.$$
(8)
The traditional way of solving the bound state problem is based on a simultaneous diagonalization of the momentum $`P^+`$ and the Hamiltonian $`P^{}`$, but as one can see from the structure of (8), the same problem can be solved by diagonalizing $`P^+`$ and $`Q^{}`$ instead .
In order to solve the bound state problem we impose the light cone gauge ($`A^+=0`$), then the supercharges are given by:
$`Q^+`$ $`=`$ $`2{\displaystyle 𝑑x^{}\left(\lambda _{}A^2+\frac{i}{2}_{}\xi ^{}\psi \frac{i}{2}\psi ^{}_{}\xi \frac{i}{2}\xi ^{}_{}\psi +\frac{i}{2}_{}\psi ^{}\xi \right)}`$ (9)
$`Q^{}`$ $`=`$ $`2{\displaystyle 𝑑x^{}\left(\lambda _{}A^{}+i\xi ^{}D_2\psi iD_2\psi ^{}\xi +\frac{i}{\sqrt{2}}_{}(\stackrel{~}{\psi }^{}\xi \xi ^{}\stackrel{~}{\psi })\right)}.`$
Note that apart from a total derivative these expressions involve only left–moving components of fermions ($`\lambda `$ and $`\psi `$). In fact in the light–cone formulation only these components are dynamical. To see this we consider the equations of motion that follow from the action (1), in the light cone gauge. Three of them serve as constraints rather than as dynamical statements:
$$_{}^2A^{}=gJ,J=i[A^2,_{}A^2]+\frac{1}{\sqrt{2}}\{\lambda ,\lambda \}ih_{}\xi \xi ^{}+ih\xi _{}\xi ^{}+\sqrt{2}h\psi \psi ^{},$$
(11)
$`_{}\stackrel{~}{\lambda }={\displaystyle \frac{ig}{\sqrt{2}}}\left([A^2,\lambda ]+ih\xi \psi ^{}ih\psi \xi ^{}\right),`$ (12)
$`_{}\stackrel{~}{\psi }={\displaystyle \frac{ig^{}}{\sqrt{2}}}A^2\psi +{\displaystyle \frac{g}{\sqrt{2}}}\lambda \xi .`$ (13)
We introduced the relative coupling for the fundamental matter: $`h=g^{}/g`$. Apart from the zero mode problem , one can invert the first constraint to write the auxiliary field $`A^{}`$ in terms of physical degrees of freedom. Substituting the result into the expression for the supercharge and omitting the boundary term, we get:
$$Q^{}=2𝑑x^{}\left(gJ\frac{1}{_{}}\lambda +g^{}\xi ^{}A^2\psi +g^{}\psi ^{}A^2\xi \right).$$
(14)
This supercharge gives rise to a whole family of supersymmetric theories. Namely one can see that the expression (14) is meaningful even if we exclude some of the fields from the theory. As soon as we have at least one fermion and at least one field in the adjoint representation, (14) defines an interacting theory with supersymmetry. Some of these theories were studied before (namely the pure adjoint systems with or without bosons), but many of them are new. In this paper we will study the mesonic spectrum of all these models, their field content is summarized in the table 1.
In order to solve the bound state problem we apply the methods of Supersymmetric DLCQ. Namely we compactify the two dimensional theory on a light–like circle ($`L<x^{}<L`$), and impose periodic boundary conditions on all physical fields. This leads to the following mode expansions:
$`A_{ij}^2(0,x^{})`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{4\pi }}}{\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{\sqrt{k}}}\left(a_{ij}(k)e^{ik\pi x^{}/L}+a_{ji}^{}(k)e^{ik\pi x^{}/L}\right),`$ (15)
$`\lambda _{ij}(0,x^{})`$ $`=`$ $`{\displaystyle \frac{1}{2^{\frac{1}{4}}\sqrt{2L}}}{\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}\left(b_{ij}(k)e^{ik\pi x^{}/L}+b_{ji}^{}(k)e^{ik\pi x^{}/L}\right),`$ (16)
$`\xi _i(0,x^{})`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{4\pi }}}{\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{\sqrt{k}}}\left(c_i(k)e^{ik\pi x^{}/L}+\stackrel{~}{c}_i^{}(k)e^{ik\pi x^{}/L}\right),`$ (17)
$`\psi _i(0,x^{})`$ $`=`$ $`{\displaystyle \frac{1}{2^{\frac{1}{4}}\sqrt{2L}}}{\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}\left(d_i(k)e^{ik\pi x^{}/L}+\stackrel{~}{d}_i^{}(k)e^{ik\pi x^{}/L}\right).`$ (18)
We drop the zero modes of the fields. Including them could lead to new and interesting effects (see , for example), but this is beyond the scope of this work. In the light–cone formalism one treats $`x^+`$ as a time direction, thus the commutation relations between fields and their momenta are imposed on the surface $`x^+=0`$. For the system under consideration this means:
$`[A_{ij}^2(0,x^{}),_{}A_{kl}^2(0,y^{})]`$ $`=`$ $`i\left(\delta _{il}\delta _{kj}{\displaystyle \frac{1}{N}}\delta _{ij}\delta _{kl}\right)\delta (x^{}y^{}),`$ (19)
$`\{\lambda _{ij}(0,x^{}),\lambda _{kl}(0,y^{})\}`$ $`=`$ $`\sqrt{2}\left(\delta _{il}\delta _{kj}{\displaystyle \frac{1}{N}}\delta _{ij}\delta _{kl}\right)\delta (x^{}y^{}),`$ (20)
$`[\xi _i(0,x^{}),_{}\xi _j(0,y^{})]`$ $`=`$ $`i\delta _{ij}\delta (x^{}y^{}),`$ (21)
$`\{\psi _i(0,x^{}),\psi _j(0,y^{})\}`$ $`=`$ $`\sqrt{2}\delta _{ij}\delta (x^{}y^{}).`$ (22)
These relations can be rewritten in terms of creation–annihilation operators:
$$[a_{ij},a_{kl}^{}]=\left(\delta _{il}\delta _{kj}\frac{1}{N}\delta _{ij}\delta _{kl}\right),\{b_{ij},b_{kl}^{}\}=\left(\delta _{il}\delta _{kj}\frac{1}{N}\delta _{ij}\delta _{kl}\right),$$
(23)
$$[c_i,c_j^{}]=\delta _{ij},[\stackrel{~}{c}_i,\stackrel{~}{c}_j^{}]=\delta _{ij},\{d_i,d_j^{}\}=\delta _{ij}\{\stackrel{~}{d}_i,\stackrel{~}{d}_j^{}\}=\delta _{ij}.$$
(24)
In this paper we will discuss numerical results obtained in the large $`N`$ limit, i.e. we neglect $`1/N`$ terms in the above expressions. Although $`1/N`$ corrections may lead to interesting effects , they are beyond the scope of this work. In the presence of fundamental matter, however, one can also define different limits by making the number of flavors comparable with number of colors . While this is an interesting direction for a future exploration, in the current paper we concentrate on models with one flavor and an infinite number of colors.
## 3 Mesons involving static (s)quarks.
We begin our consideration with the simplest models involving only adjoint matter. While the glueball spectrum of such models was studied before , it might also be interesting to look at the meson–like states. In the string interpretation of these theories such states would correspond to open strings with freely moving endpoints. In the QCD language the model corresponds to a system of interacting gluons and gluinos which is bounded by nondynamical (s)quark and anti-(s)quark. In the large $`N`$ limit we will have to consider only a single (s)quark—anti-(s)quark pair, thus the Fock space is constructed from the states of the following type:
$$\overline{f}_{i_1}a_{i_1i_2}^{}(k_1)\mathrm{}b_{i_ni_{n+1}}^{}(k_n)\mathrm{}f_{i_p}|0.$$
(25)
Here $`|0`$ is a vacuum defined by annihilation operators $`a_{ij}`$ and $`b_{ij}`$ and $`\overline{f}_i`$ and $`f_{i_p}`$ are sets of c–numbers.
The supercharges for such a system can be constructed by eliminating fundamental matter from the expressions (9) and (14):
$`Q^+`$ $`=`$ $`2{\displaystyle 𝑑x^{}\lambda _{}A^2},`$
$`Q^{}`$ $`=`$ $`2g{\displaystyle 𝑑x^{}\left(i[A^2,_{}A^2]+\frac{1}{\sqrt{2}}\{\lambda ,\lambda \}\right)\frac{1}{_{}}\lambda }.`$ (26)
The supercharge $`Q^{}`$ has an interesting property which can be seen from its mode expansion . Acting on states (25), this operator changes the numbers of bosons $`a^{}`$ by an even number ($`0`$ or $`\pm 2`$), thus one can perform the diagonalization of $`P^{}`$ on two separate spaces: those containing either even or odd number of bosons. The second supercharge $`Q^+`$ makes the situation even more interesting: it maps one of these spaces into another and thus leads to the same massive spectra in both sectors (see ). We have studied the mesonic mass spectrum of $`\left(Q^{}\right)^2`$ at low values of the harmonic resolution $`k`$ and found no massless states, thus the spectra in two different sectors are completely identical. Note that this property is not satisfied for the “glueball spectrum”: there we found a lot of exact massless states .
In figure 1a we present the results of the numerical diagonalization of the Hamiltonian. Note that we considered only one of two equivalent sectors. Although we have not seen any massless states at any finite value of resolution (we considered $`K=4\mathrm{}7`$), it appears that at least the two lowest states converge to $`M=0`$ in the continuum limit. This conclusion is supported by the quadratic fit to the data. The third lowest state is also well–defined and an extrapolation of its mass gives a value of $`1.83`$ for the continuum limit. To reveal the structure of higher states we need some additional information about their wavefunctions, this will lead to a clear distinction between nearly degenerate states.
While this direction is definitely worth pursuing, we already can formulate some interesting properties of the mesonic mass spectrum. Unlike the glueball case, there are no massless mesons at any finite value of resolution. In the continuum limit, however, the massless mesons appear, but the number of such massless states is still an open question. We also see an example of a meson state converging to a finite mass, and the data presented in Figure 1a suggests that there are many such states in the continuum limit. We should mention that the convergence properties of this model are very good, which seems to be a general properties of supersymmetric DLCQ as opposed to the traditional DLCQ approach.
As we mentioned in the previous section, the supercharge $`Q^{}`$ is well–defined as soon as we have at least one fermion. This suggests an interesting truncation of the supersymmetric adjoint model : one can eliminate the bosonic field $`A^2`$ from the theory. This reduces the number of supersymmetries to $`(1,0)`$ and the remaining supercharge is given by:
$$Q^{}=\frac{g}{\sqrt{2}}𝑑x^{}\{\lambda ,\lambda \}\frac{1}{_{}}\lambda .$$
(27)
While the “closed string” sector of this theory has well–defined continuum bound states , all mesonic masses appear to be pushed to infinity in the continuum limit. In fact if one looks at the spectrum of DLCQ Hamiltonian $`P^{}`$, rather than the mass operator $`M^2=2KP^{}`$, a finite limit can be found. The result, presented in figure 1b, has a peculiar property: an extrapolation of the four lowest eigenvalues gives $`3.08`$, $`3.01`$, $`3.09`$ and $`3.06`$ accordingly. Thus it seems that the spectrum of the Hamiltonian has some kind of a threshold. This interesting observation, however, does not undermine the fact that this model does not have any sensible continuum limit.
## 4 Models without adjoint scalar.
As we already saw in the previous section, the simplest supersymmetric model can be constructed by truncating all fields in the supercharge (14), except for gluino $`\lambda `$. In this section we add the dynamics of fundamental matter to the model. There are three different ways of doing this, and they give rise to the systems which we call $`\lambda \mathrm{\Psi }`$, $`\lambda \xi `$ and $`\lambda \mathrm{\Psi }\xi `$. We will show that all three systems have a well–defined mass spectrum and their bound states exhibit similar behavior under a variation of the coupling constant $`h`$.
We begin with the pure fermionic system $`\lambda \mathrm{\Psi }`$. This system has $`(1,0)`$ supersymmetry and the supercharge has the form:
$$Q^{}=\frac{g}{\sqrt{2}}𝑑x^{}\left(\{\lambda ,\lambda \}+2h\psi \psi ^{}\right)\frac{1}{_{}}\lambda .$$
(28)
After substituting the expansions (16), (18) one gets the mode decomposition of the supercharge:
$`Q^{}={\displaystyle \frac{i2^{1/4}g\sqrt{L}}{\pi }}{\displaystyle \underset{k_1=1}{\overset{\mathrm{}}{}}}{\displaystyle \underset{k_2=1}{\overset{\mathrm{}}{}}}\{{\displaystyle \frac{h}{k_1}}[\stackrel{~}{d}_i^{}(k_2)b_{ij}^{}(k_1)\stackrel{~}{d}_j(k_1+k_2)+`$ (29)
$`\stackrel{~}{d}_j^{}(k_1+k_2)\stackrel{~}{d}_i(k_2)b_{ij}(k_1)+b_{ij}^{}(k_1)d_j^{}(k_2)d_i(k_1+k_2)+d_i^{}(k_1+k_2)b_{ij}(k_1)d_j(k_2)]+`$
$`({\displaystyle \frac{1}{k_1}}+{\displaystyle \frac{1}{k_2}}{\displaystyle \frac{1}{k_1+k_2}})[b_{ik}^{}(k_1)b_{kj}^{}(k_2)b_{ij}(k_1+k_2)+b_{ij}^{}(k_1+k_2)b_{ik}(k_1)b_{kj}(k_2)]\}`$
The mass spectrum of this theory is presented in figure 2a (we truncated it at $`M^2=25`$ and chose $`h=1`$). There is one massless state and all other masses converge to finite continuum limits. We illustrate this convergence for the five lowest states in the Table 2. Using the structure of the supercharge (29), one can easily calculate the wavefunction of the massless state for any finite resolution as well as its continuum limit. This state appears to have only two partons in it and the wavefunction is equal to a constant:
$$|M=0,K=C\underset{n=1}{\overset{K1}{}}\stackrel{~}{d}_i^{}(n)d_i^{}(Kn)|0_0^{P^+}𝑑p\stackrel{~}{D}_i^{}(p)D_i^{}(P^+p)|0.$$
(30)
Here we have introduced the continuum modes $`D^{}`$ and $`\stackrel{~}{D}^{}`$. One can see that the state with the wavefunction (30) stays massless for an arbitrary values of coupling $`h`$. It interesting to note that a massless state with constant wavefunction was also observed in a non–supersymmetric adjoint QCD in two dimensions .
It is also interesting to look at other masses as functions of the coupling constant. Figure 2b shows this dependence for the resolution $`K=4`$. We should note that this behavior is typical for all values of resolution and for all of the systems we are considering in this section. In particular one can see that the lowest states stays near $`M^2=0`$ for a wide range of negative couplings. A closer look at this state at resolutions $`4`$ and $`5`$ is presented in figure 3. Looking at higher resolutions, we observe, that this state become massless for some value of coupling at all resolutions except $`K=5`$, thus the graph 3a is typical, while 3b is an artifact of the DLCQ. Interestingly, such odd behavior at $`k=5`$ is observed for all three systems we are studying here. The values of the critical coupling and the extrapolation to the continuum limit is presented in Table 3. The wavefunction of this massless state is concentrated in the two–parton sector: at resolutions $`4`$, $`6`$, $`7`$ and $`8`$ the two–particle sector contains $`82\%`$, $`92\%`$, $`93\%`$ and $`95\%`$ of wavefunction, however we believe that in continuum limit the wavefunction has small, but nonzero contributions from sectors with an arbitrarily large number of partons. This property is common for massless states in other supersymmetric theories . While we were not able to solve for the continuum wavefunction in two–parton sector, the SDLCQ results presented in figure 4 points to a linear behavior in the continuum limit. Note that for the state we are looking at:
$$|M=0,K,h\underset{n=1}{\overset{K1}{}}f(n,Kn)\stackrel{~}{d}_i^{}(n)d_i^{}(Kn)|0,$$
(31)
the wavefunction is antisymmetric: $`f(p,q)=f(q,p)`$, so in figure 4 we present only the region $`pq`$.
Let us now discuss two other theories without adjoint scalars. As we already mentioned, their properties are similar to the ones of the $`\lambda \psi `$ system, so we will consider both $`\lambda \xi `$ and $`\lambda \psi \xi `$ models only briefly. The $`\lambda \xi `$ system has $`(1,0)`$ supersymmetry and its supercharge reads:
$$Q^{}=\frac{g}{\sqrt{2}}𝑑x^{}\left(\{\lambda ,\lambda \}+i\sqrt{2}h\xi _{}\xi ^{}i\sqrt{2}h_{}\xi \xi ^{}\right)\frac{1}{_{}}\lambda .$$
(32)
The mass spectrum of this model is presented in figure 5a (again, we put $`h=1`$) and the extrapolation of the masses to the continuum limit is given in Table 4. The masses exhibit the same coupling dependence as in the pure fermionic model (the only exception is the absence of the coupling–independent massless state), in particular there exists a critical coupling at which one of the states becomes massless. The values of this critical coupling are plotted in figure 5b and for the highest resolution we have $`h_{cr}=1.2`$.
Finally we analyze the system which includes everything except for the adjoint scalar ($`\lambda \psi \xi `$ in our notation). This model has two supercharges, but we look only at one of them:
$$Q^{}=\frac{g}{\sqrt{2}}𝑑x^{}\left(\{\lambda ,\lambda \}+i\sqrt{2}h\xi _{}\xi ^{}i\sqrt{2}h_{}\xi \xi ^{}+2h\psi \psi ^{}\right)\frac{1}{_{}}\lambda .$$
(33)
The second supercharge can be formally constructed:
$$Q^+=2𝑑x^{}\left(\frac{i}{2}_{}\xi ^{}\psi \frac{i}{2}\psi ^{}_{}\xi \frac{i}{2}\xi ^{}_{}\psi +\frac{i}{2}_{}\psi ^{}\xi \right),$$
(34)
but its square does not give the canonical $`P^+`$ but rather:
$$\{Q^+,Q^+\}=2\sqrt{2}P^+2\sqrt{2}𝑑x^{}\lambda _{}\lambda .$$
(35)
Moreover, from this expression one concludes that
$$[\{Q^+,Q^+\},Q^{}]0,$$
(36)
thus the two supercharges do not anticommute and they cannot be diagonalized simultaneously. Our formulation of the bound state problem is based on diagonalization of $`P^+`$ and $`\left(Q^{}\right)^2`$, which still commute, and the $`Q^+`$ operator must be abandoned.
In the large $`N`$ calculations there are four different sectors to consider. One can start from either one of the four types of states:
$$\stackrel{~}{d}^{}b^{}\mathrm{}b^{}d^{}|0,\stackrel{~}{c}^{}b^{}\mathrm{}b^{}c^{}|0,\stackrel{~}{c}^{}b^{}\mathrm{}b^{}d^{}|0,\stackrel{~}{d}^{}b^{}\mathrm{}b^{}c^{}|0,$$
(37)
then $`Q^{}`$ acts only inside the corresponding subspace. The first and second sectors reproduce the results we just obtained for $`\lambda \psi `$ and $`\lambda \xi `$ models. Two remaining models are mapped into each other under the Z2 transformation:
$$b_{ij}b_{ji},\stackrel{~}{c}_ic_i,\stackrel{~}{d}_id_i,$$
(38)
which is a symmetry of $`Q^{}`$. The spectrum for one of these sectors is presented in figure 6a and in the table 5 and the critical coupling is plotted in figure 6b. Note that we don’t have a coupling–independent massless state in this sector. The critical coupling for resolutions $`K>6`$ is relatively constant at $`h_{cr}=2.0`$.
## 5 Models containing adjoint bosons.
Let us now discuss the models with an adjoint scalar $`A^2`$. We will see that the general property of such theories is the presence of a large number of low energy bound states which complicates the extrapolation to the continuum limit. We begin with the system which does not have a gluino as a dynamical field. Although this model has two supercharges:
$`Q^{}=2g^{}{\displaystyle 𝑑x^{}\left(\xi ^{}A^2\psi +\psi ^{}A^2\psi \right)},`$ (39)
$`Q^+=2{\displaystyle 𝑑x^{}\left(\frac{i}{2}_{}\xi ^{}\psi \frac{i}{2}\psi ^{}_{}\xi \frac{i}{2}\xi ^{}_{}\psi +\frac{i}{2}_{}\psi ^{}\xi \right)},`$ (40)
they do not commute, thus we will ignore the supercharge $`Q^+`$ in our consideration (this situation is analogous to the case of $`\lambda \xi \psi `$ model). We constructed mesonic spectrum of this model and the lowest masses are presented in figure 7a. One can see that the number of light states grows with resolution, in fact we will argue that in the continuum limit this model has a continuous mass spectrum. But first let us look at the heaviest bound state one can construct at a given value of resolution. The masses of such states are plotted in figure 7b, and one can see that this graph has a good linear approximation:
$$M_{max}^2(K)=\frac{(g^{})^2N}{\pi }\left(1.08K1.86\right),$$
(41)
where $`K`$ is a value of resolution. The negative constant in the above expression is not important, since one cannot consider $`K<4`$ and the interesting limit is $`K\mathrm{}`$. For the total number of bosonic states at a particular resolution, simple combinatorics gives:
$$N_{total}(K)=2^K.$$
(42)
In order to analyze the continuum limit of the spectrum we will study the density of states $`dN/N_{total}(K)`$ as a function of the reduced mass: $`dM^2/M_{max}^2(K)`$. In particular we will plot the distribution function defined as
$$F(M^2/M_{max}^2(K))=\frac{N(\text{Mass}<M)}{N_{total}}.$$
(43)
Such plots for different resolutions are presented in color in figure 8a, which gives a convincing argument for the convergence of the function $`F(x)`$ in the continuum limit. In figure 8b we present this function for resolution $`10`$, which is a good approximation to the continuum limit.
Let us now discuss the models that include both fields in the adjoint representation. First we consider the $`A\lambda \psi `$ system. It has two supercharges which don’t commute, so we look only at one of them:
$$Q^{}=2g𝑑x^{}\left(i[A^2,_{}A^2]+\frac{1}{\sqrt{2}}\{\lambda ,\lambda \}+\sqrt{2}h\psi \psi ^{}\right)\frac{1}{_{}}\lambda .$$
(44)
The mass spectrum, obtained as the result of the diagonalization of Hamiltonian $`P^{}=\left(Q^{}\right)^2`$, is presented in the figure 9a. As usual we put $`h=1`$ and truncated the spectrum at some value of mass (in this case $`M^2=10`$). The states can be easily traced from one resolution to another, and a new low state appears at every even resolution and a new massless state appears at every odd resolution (thus there is one massless state at $`K=4`$, two massless states at $`K=5`$ and $`K=6`$, $`3`$ massless states at $`K=7,8`$ and so on). Such behavior was observed in other supersymmetric systems with adjoint scalars and it points to a continuum spectrum in the limit $`K=\mathrm{}`$. One can also look at the coupling dependence of the states, in particular the lowest state with nonzero mass at $`K=4`$ becomes massless at the coupling $`h=1.25`$ (see figure 9b). But since this state ultimately becomes a part of a continuous spectrum, its properties are not as interesting as its counterpart in the $`\lambda \psi `$ model.
As we saw in our study of the models without gauge fields, the theories which differ only by replacing $`\psi `$ by $`\xi `$, behave in the same fashion. The same is true for the models with adjoint scalar. The spectrum of the $`A\lambda \xi `$ model is presented in figure 10, it also converges to a continuous spectrum and has a critical coupling $`h`$ at any $`K5`$. The figure 10b illustrates that it is not only low mass states that appear at high resolution, but all values of mass seem to be filled in the continuum limit.
Finally we consider the system without truncation, i.e. we study the $`A\lambda \psi \xi `$ model. Unlike all other systems we studied in this section, it has a complete $`(1,1)`$ supersymmetry. Thus the two supercharges given by (9) and (14) anticommute and can be diagonalized simultaneously. Thus, apart from the massless sector, the spectrum is four–fold degenerate and we can look only at a quarter of the theory, while diagonalizing the mass operator. In particular, we consider only bosonic states with an even number of creation operators $`a^{}`$. The combined action of $`Q^+`$ and $`Q^{}`$ gives a boson with an odd number of creation operators $`a^{}`$, while the action of either one of the supercharges leads to the fermionic sectors. The low energy spectrum in a single sector at $`h=1`$ is presented in figure 11, there are two massless states for every even value of resolution. One can see that the presence of these states is not the only difference between the odd and even values of $`K`$, it seems that we are dealing with the SDLCQ of two different theories. Of course, in the large $`K`$ limit they should converge to the same result, and, as figure 11 demonstrates, the resulting theory has many light states and the possibility of a continuum spectrum is not ruled out. Note that the lowest states appear in almost degenerate pairs. The explanation of this doubling poses an interesting question, which might be answered by performing a careful study of the wavefunctions. We also looked at the lowest masses as functions of a coupling $`h`$, and the result for resolution $`4`$ is presented in figure 12. This shows two peculiar properties of this model. First, there is a smooth interchange between almost degenerate states, as opposed to the level crossing we observed in the other theories. In addition to this, the model exhibits two critical couplings at resolution $`4`$ and it is interesting to see whether this property persists at higher resolutions. Unfortunately this model is the hardest one to study since it has a maximal number of fields. Also additional calculational effort is require to handle the fact that the $`A\lambda \psi \xi `$ model (and only this model) has a complex supercharge. We hope to overcome both difficulties in the future work. In spite of these difficulties we can already conclude that the $`A\lambda \psi \xi `$ has many light states in the continuum limit and possibly it converges to a continuous spectrum as the other models with adjoint scalars.
## 6 Discussion.
In this work we studied the mesonic mass spectrum of various supersymmetric models. The calculations were performed in the framework of supersymmetric DLCQ, namely we compactified the light–like coordinate $`x^{}`$ on a finite circle and performed a numerical diagonalization of supercharge $`Q^{}`$. We found that the systems with adjoint scalars tend to give a continuous spectrum in the decompactification limit. For one of these systems ($`A\psi \xi `$ model) we found a limiting form of the distribution function for the mass. By contrast, the models without adjoint bosons have well defined spectra of bound states in the continuum limit. The only exception from this rule is the system containing only the gaugino field, which seems not to have any finite mass mesons in the decompactification limit. For the well defined systems we found the masses of the lightest mesons and demonstrated the fast convergence of SDLCQ approximation.
We also looked at the mass spectrum at different values of coupling constant. The nontrivial phase diagram is an essential property which distinguish the models we considered here from the two dimensional systems studied previously . The coupling constant of a pure gauge theory in two dimensions has a dimension of mass, thus all the bound state masses scale like $`g`$, leaving no space for nontrivial coupling constant dependence. One can avoid this problem by introducing the masses for a gauge field or its superpartner, but such terms usually lead to breaking of either a gauge invariance or supersymmetry<sup>1</sup><sup>1</sup>1Actually the pure fermionic system $`\lambda `$ is supersymmetric only at nonzero value of a particle mass $`m=g^2N/\pi `$, but this still does not leave a space for an adjustable mass parameter.. Another way of introducing a free parameter in the theory is to add a new supermultiplet with a different charge $`g^{}`$. Of course, this parameter is not completely free: the quantization of charge requires the value $`g^{}/g`$ to be a rational number, but formally we can study the bound states as functions of $`h=g^{}/g`$. We found an interesting property of the lowest nonzero mass: it vanishes at a particular negative value of $`h`$. The nature of this property is still unclear. One can also introduce the true free parameter by considering the massive matter supermultiplet (now there is no obstacle coming from gauge invariance) and we leave this possibility for a future investigation.
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# Fast and Accurate Computation Tools for Gravitational Waveforms from Binary Stars with any Orbital Eccentricity
## 1 INTRODUCTION
Gravitational wave detection experiments in space, including satellite Doppler-Tracking (Bertotti and Iess, 1999) and LISA (http://lisa.jpl.nasa.gov), will hopefully open a window on the low-frequency part of the gravitational wave (henceforth GW) spectrum of cosmic origin. In these frequency bands, binary stars are among the most promising continuous detectable source.
A substantial fraction of binaries are expected to have orbits with non negligible eccentricity (Barone et al., 1988; Hils et al., 1992; Pierro and Pinto, 1996c) resulting into the emission of several harmonics of the fundamental orbital frequency. The importance of this fact from the standpoint of signal detection and estimation has been already noted.
For coalescing binaries, Pierro and Pinto (1996b) and Martel and Poisson (1999) pointed out that neglecting residual (albeit very small) orbital eccentricities may seriously deteriorate matched-filter detection performance. Their results, obtained in the frame of the simplest (newtonian) Peters Mathews (henceforth PM) model (Peters and Mathews, 1963; Peters 1964; Pierro and Pinto 1996c), support the qualitative conclusion that residual orbital eccentricities cannot be bona fide disregarded in building templates for matched-filter detection of gravitational wave chirps from inspiraling binaries<sup>1</sup><sup>1</sup>1 The effect of a residual (tiny) orbital eccentricity on the radiation emitted from an inspiraling binary system was also considered in (Moreno et al., 1994), with special emphasis on the possible relevance of periastron advance. The results in (Moreno et al., 1994) are unfortunately affected by several errors and misprints..
For steady-state binaries with non-zero orbital eccentricity, on the other hand, using circular-orbit waveform templates, i.e. neglecting higher order harmonics, implies a potentially large loss of signal-to-noise ratio (henceforth SNR), leading to significantly worse detector’s performance, as will be shown in the sequel.
The main goals of the present paper are:
* to provide some quantitative hint for validating the applicability of the simple PM model to steady-state binaries;
* to gauge the loss in SNR due to the simple circular-orbit assumption and, more generally, to set some criteria for spectral waveform truncation;
* to introduce efficient (accurate and fast) computational tools for constructing gravitational waveform templates for (steady-state) binary sources with any orbital eccentricity.
The paper is organized as follows. In Sect. 2 we introduce some (dimensionless) parameters whereby the applicability of the PM model to specific sources can be assessed. In Sect.s 3a and 3b we review the GW spectra and waveforms in the frame of the PM model. In Sect.s 4a and 4b we show how to evaluate the total harmonic distortion due to spectral waveform truncation, and introduce a modified Carlini Meissel expansion tool for fast and accurate GW harmonics computation. The results in this section can be readily extended, in principle, to higher-order post-newtonian (henceforth PN) models. As an application, in Sect. 5 we apply our formalism to some paradigm eccentrical binary sources. Conclusions follow under Sect. 6. Technical developments are collected in Appendix A to C.
## 2 STEADY-STATE BINARIES: THE PETERS-MATHEWS MODEL
The PM model for gravitational wave emission from binary systems in a Keplerian orbit was introduced in the sixties (Peters and Mathews, 1963, Peters 1964), and recently re-examined (Pierro and Pinto, 1996a) . It relies on the following main assumptions: i) point mass, ii) weak field, iii) slow motion, and iv) adiabatic evolution (negligible change of the orbital parameters over each orbit). These conditions can be checked in terms of the following inequalities (Pierro and Pinto, 1996a):
$$\xi _1:=\frac{sourcegravitationalradius}{aphastralseparation}=2\chi ^{2/3}(1e)^11,$$
(1)
$$\xi _2:=\frac{aphastralvelocity}{velocityoflight}=\chi ^{1/3}\left(\frac{1+e}{1e}\right)^{1/2}1,$$
(2)
$$\xi _3:=sup\{\left|\frac{dT}{dt}\right|,\left|e^1\frac{de}{dt}T\right|\}=\frac{152\pi }{15}(1\mathrm{\Delta }^2)\chi ^{5/3}\left(1+\frac{121}{304}e^2\right)(1e^2)^{5/2}1,$$
(3)
where as already stated $`\chi =cT/\pi r_g`$, $`T`$ being the orbital period, $`r_g=2G(M_1+M_2)/c^2`$ is the source gravitational radius, $`M_{1,2}`$ are the companion masses, $`\mathrm{\Delta }=|M_1M_2|/(M_1+M_2)`$ and $`e`$ is the eccentricity.
Tidal effects could be neglected provided neither companion star fills its Roche lobe. Following Eggleton (1983), this translates into:
$$\xi _1\mathrm{\Lambda }^1\frac{2}{1+\mathrm{\Delta }}\left\{0.6+\left(\frac{1\mathrm{\Delta }}{1+\mathrm{\Delta }}\right)^{2/3}\mathrm{ln}\left[1+\left(\frac{1\mathrm{\Delta }}{1+\mathrm{\Delta }}\right)^{1/3}\right]\right\}^1,$$
(4)
where $`\mathrm{\Lambda }`$ is the ratio between the physical and gravitational companion radius<sup>2</sup><sup>2</sup>2Typical values of $`\mathrm{\Lambda }`$ range from $`10^4`$ for white dwarfs down to $`3`$ for hadronic stars. Further departures from the standard model are expected due to the possible occurrence of mass-transfer phenomena, which would be present in closely-orbiting classical stars, as well as in binaries where one companion is an accreting collapsed object..
For most steady-state binary systems, i.e. long before coalescence, $`\xi _1`$ to $`\xi _3`$ above are fairly small (see e.g. Sect. 5), and the PM model turns out to be perfectly adequate.
## 3a STEADY-STATE BINARIES: SPECTRA
According to the PM model, the GW power $`\overline{}_n^{+,\times }`$ radiated at the $`n^{th}`$ harmonic of the orbital frequency by a steady binary source can be conveniently cast into the following universal form (Barone et al., 1988):
$$\overline{}_n^{+,\times }=\frac{2G}{5c^5}\chi ^{10/3}(1\mathrm{\Delta }^2)^2G_{max}(e)\overline{g}^{+,\times }(n,e)$$
(5)
where the superscripts $`+,\times `$ refer to the fundamental GW polarization states. The spectral power distribution is embodied in the universal dimensionless functions $`\overline{g}^{+,\times }(n,e)`$ shown in fig.s 1.1-1.10 for $`e=0(0.1)0.9`$. For circular orbits ($`e=0`$) only the second harmonic is emitted. The function $`G_{max}(e)`$ plotted in fig. 2 is the ratio between the total luminosity (sum over both polarizations) of the brightest GW spectral line, and the total luminosity of a circular-orbit binary having the same $`\chi `$ and $`\mathrm{\Delta }`$. The brightest spectral line is the $`N_{max}`$-th harmonic of the orbital frequency, where $`N_{max}`$ is a function of $`e`$ only, displayed in fig. 3.
It is seen that for non circular orbits, several spectral lines with comparable intensities are emitted. Thus, use of the circular orbit waveform templates implies a potentially sizeable loss in the available signal power and hence in the SNR, which can spoil the detector’s performance.
### 3b WAVEFORMS
The far-field metric deviation (TT gauge) in the PM model is<sup>3</sup><sup>3</sup>3 The GW field can also be obtained by inverting the Keplerian integral of motion relating time to the true anomaly, and exploiting the simple dependance of the radiated waveforms on this latter (Wahlquist, 1987). The referred procedure is purely numerical and, to the best of our knowledge, its generalization to higher PN order models is not immediate.:
$$h_\times =\frac{\mathrm{cos}\vartheta }{\sqrt{2}}\left[2h_{xy}\mathrm{cos}2\phi (h_{xx}h_{yy})\mathrm{sin}2\phi \right],$$
(6)
$$h_+=\frac{1}{\sqrt{2}}\left\{\frac{3+\mathrm{cos}2\vartheta }{4}\left[2h_{xy}\mathrm{sin}2\phi +(h_{xx}h_{yy})\mathrm{cos}2\phi \right]\frac{1\mathrm{cos}2\vartheta }{4}(h_{xx}+h_{yy})\right\},$$
(7)
where the coordinates $`\vartheta `$ and $`\phi `$ specify the direction of the observer in a spherical polar system where the orbit lies in the equatorial plane and the binary center of mass is at the origin.
The metric components in (6), (7) can be expanded into Fourier series under the adiabatic assumption that the orbital parameters do not change appreciably over each orbit. Hence<sup>4</sup><sup>4</sup>4 The unknown irrelevant phase at $`t=0`$ has been set to zero.:
$$h_{xy}=\underset{n=1}{\overset{\mathrm{}}{}}h_{xy}^{(n)}\mathrm{sin}\left(n\frac{2\pi }{T}t\right),$$
(8)
$$h_{x\pm y}=\underset{n=1}{\overset{\mathrm{}}{}}h_{x\pm y}^{(n)}\mathrm{cos}\left(n\frac{2\pi }{T}t\right),$$
(9)
where $`h_{x\pm y}`$ is a shorthand for $`h_{xx}\pm h_{yy}`$ (see Appendix A),
$$h_{xy}^{(n)}=h_0n(1e^2)^{1/2}\left[J_{n2}(ne)+J_{n+2}(ne)2J_n(ne)\right],$$
(10)
$$h_{xy}^{(n)}=2h_0n\left\{J_{n2}(ne)J_{n+2}(ne)2e\left[J_{n1}(ne)J_{n+1}(ne)\right]+(2/n)J_n(ne)\right\},$$
(11)
$$h_{x+y}^{(n)}=4h_0J_n(ne),$$
(12)
and<sup>5</sup><sup>5</sup>5Note that for $`n=1`$ eq.s (10) and (11) contain Bessel functions of order $`1`$, for which $`J_1(x)=J_1(x)`$.
$$h_0=\frac{cT}{4\pi r}\frac{1\mathrm{\Delta }^2}{\chi ^{5/3}}.$$
(13)
For circular orbits one has simply:
$$h_{xy}^{(n)}=2h_{xy}^{(n)}=4h_0\delta _{n2},h_{x+y}^{(n)}=0,$$
(14)
where $`\delta _{pq}`$ is the Kronecker symbol.
For steady state binaries the (Robertson) periastron advance<sup>6</sup><sup>6</sup>6The relativistic periastron advance was heuristically (i.e., inconsistenlty, from the post-newtonian expansion view point) included in (Moreno et al., 1995). does not produce sensible effects on the waveforms, and is thus deliberately ignored. Inclusion of the periastron advance amounts to splitting each GW spectral line into a doublet at $`(2\pi /T)(1\pm 6\chi ^{2/3})`$, which cannot be resolved unless the signal is Fourier-transformed over a timespan $`\chi ^{2/3}T\text{sec}`$. This time is, e.g., $`510^5`$ years and $`2.810^5`$ years for $`PSR1534+12`$ and $`PSR1913+16`$, respectively.
## 4a SPECTRAL TRUNCATION AND APPROXIMATION ERROR
In order to discuss the effect of spectral truncation of (8) and (9) on the available SNR it is convenient to introduce the total harmonic distortion (henceforth THD):
$$THD=\frac{h\stackrel{}{h}}{h}=\left(\frac{{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}(\stackrel{(n)}{\stackrel{}{h}}\stackrel{(n)}{\stackrel{}{h}})^2}{{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}|h^{(n)}|^2}\right)^{1/2},$$
(15)
where $`h`$, $`\stackrel{}{h}`$ represent the exact and approximate values of the metric tensor, $`h^{(n)}`$, $`\stackrel{(n)}{\stackrel{}{h}}`$ are the Fourier coefficients of $`h`$, $`\stackrel{}{h}`$, respectively, and the $`L^2`$-norms are computed by taking the time average over one orbital period of the square of the argument, within the spirit of the adiabatic approximation. If only $`N_T`$ harmonics are included, then
$$\stackrel{(n)}{\stackrel{}{h}}=\{\begin{array}{c}h^{(n)},nN_T,\hfill \\ 0,n>N_T\hfill \end{array},THD=(1\frac{{\displaystyle \underset{n=1}{\overset{N_T}{}}}\left|h^{(n)}\right|^2}{{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}\left|h^{(n)}\right|^2})^{1/2}.$$
(16)
It is readily recognized that $`THD^2`$ represents the fraction of signal power which is lost as an effect of truncation<sup>7</sup><sup>7</sup>7 The THD is closely related to the fitting factor FF (Apostolatos, 1996) between the exact and spectral-truncated (template) waveform. From the very definitions one gets:
$$FF1\frac{THD^2}{2}+𝒪(THD{}_{}{}^{3}).$$
..
In the most general case, where besides spectral truncation, the Fourier coefficients are computed in approximate form (as e.g. in the next subsection), one has:
$$THD=\left(1\frac{2{\displaystyle \underset{n=1}{\overset{NT}{}}}h^{(n)}\stackrel{(n)}{\stackrel{}{h}}{\displaystyle \underset{n=1}{\overset{NT}{}}}|\stackrel{(n)}{\stackrel{}{h}}|^2}{{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}|h^{(n)}|^2}\right)^{1/2},$$
(17)
The harmonic distortions $`\text{THD}_{x\pm y}`$, $`\text{THD}_{xy}`$ due to the spectral truncation of (8), (9) can be computed for any given $`N_T`$ using Kapteyn’s theory (Watson, 1966, ch. 17) to evaluate in closed form the infinite sums in (17). After some lengthy but simple algebra, one obtains (see Appendix B):
$$h_{x+y}^2=\underset{n=1}{\overset{\mathrm{}}{}}\left|h_{x+y}^{(n)}\right|^2=8\left[(1e^2)^{1/2}1\right],$$
(18)
$$h_{xy}^2=\underset{n=1}{\overset{\mathrm{}}{}}\left|h_{xy}^{(n)}\right|^2=e^4\left\{4(1e^2)^{1/2}(812e^2+9e^4)8(e^22)^2\right\},$$
(19)
$$h_{xy}^2=\underset{n=1}{\overset{\mathrm{}}{}}\left|h_{xy}^{(n)}\right|^2=e^2(1e^2)^{1/2}\left\{12+e^2+8e^2\left[(1e^2)^{3/2}1\right]\right\},$$
(20)
The corresponding harmonic distortions for the $`TT`$ metric components $`h_+`$, $`h_\times `$ can be conveniently written as follows:
$$THD_\times =\left\{\left[4THD_{xy}^2h_{xy}^2\mathrm{cos}^22\phi +THD_{xy}^2h_{xy}^2\mathrm{sin}^22\phi \right]\left[4h_{xy}^2\mathrm{cos}^22\phi +h_{xy}^2\mathrm{sin}^22\phi \right]^1\right\}^{1/2},$$
(21)
and:
$$THD_+=\{[(\stackrel{}{3}+\mathrm{cos}2\vartheta )^2(4THD_{xy}^2h_{xy}^2\mathrm{sin}^22\phi +THD_{xy}^2h_{xy}^2\mathrm{cos}^22\phi )+(1\mathrm{cos}2\vartheta )^2THD_{x+y}^2h_{x+y}^2+$$
$$+2(1\mathrm{cos}2\vartheta )(3+\mathrm{cos}2\vartheta )\mathrm{cos}2\phi (h_{xy}\underset{xy}{\overset{}{h}}),(h_{x+y}\underset{x+y}{\overset{}{h}})][(\stackrel{}{3}+\mathrm{cos}2\vartheta )^2(4h_{xy}^2\mathrm{sin}^22\phi +h_{xy}^2\mathrm{cos}^22\phi )+$$
$$+(1\mathrm{cos}2\vartheta )^2h_{x+y}^2+2(1\mathrm{cos}2\vartheta )(3+\mathrm{cos}2\vartheta )\mathrm{cos}2\phi \stackrel{}{h_{xy}},\stackrel{}{h_{x+y}}]^1\}^{1/2},$$
(22)
where $`,`$ is the scalar product in $`L_{[0,T]}^2`$. In order to evaluate (22) the further infinite sum:
$$h_{xy},h_{x+y}=\underset{n=1}{\overset{\mathrm{}}{}}h_{xy}^{(n)}h_{x+y}^{(n)}=8(1e^2)^{1/2}\left\{1+(12e^2)\left[1(1e^2)^{1/2}\right]\right\}.$$
(23)
is needed, which is also readily obtained as explained in Appendix B.
The harmonic distortions (21) and (22) can be sensible even at very low eccentricities ($`e.1)`$. Expanding (21) and (22) to lowest order in $`e`$ yields:
$$THD_\times =\frac{3\sqrt{10}}{4}e+𝒪(e^3),$$
(24)
$$THD_+=\frac{\sqrt{4(1\mathrm{cos}2\phi )^2+12(1\mathrm{cos}2\phi )(3+\mathrm{cos}2\phi )\mathrm{cos}2\vartheta +90(3+\mathrm{cos}2\phi )^2}}{4(3+\mathrm{cos}2\phi )}e+𝒪(e^3).$$
(25)
The above simple expressions are fairly accurate for $`e.1`$, as seen, e.g., from fig. 4, where the angular averages of the approximate and exact harmonic distortion are drawn, and seen to be almost indistinguishable and non-negligible. The $`(\vartheta ,\phi )`$dependent factor in (25) is plotted in fig. 5. Its average value over the sphere is exactly equal to the $`(\vartheta ,\phi )`$independent factor in (24).
The obvious question is how many terms should be included in (8) and (9) so as to keep both $`THD_+`$ and $`THD_\times `$ below some specified level, for any $`(\vartheta ,\phi )`$.
To answer this question one may resort to the following inequalities:
$$\underset{(\vartheta ,\phi )}{\mathrm{max}}THD_\times max(THD_{xy},THD_{xy}),\underset{(\vartheta ,\phi )}{\mathrm{max}}THD_+max(THD_{xy},THD_{xy},THD_{x+y})\underset{(\vartheta ,\phi )}{\mathrm{max}}[Q(\vartheta ,\phi ,e)].$$
(26)
where:
$$Q(\vartheta ,\phi ,e)=\{[(\stackrel{}{3}+\mathrm{cos}2\vartheta )^2(4h_{xy}^2\mathrm{sin}^22\phi +h_{xy}^2\mathrm{cos}^22\phi )+(1\mathrm{cos}2\vartheta )^2h_{x+y}^2+$$
$$+2|(1\mathrm{cos}2\vartheta )(3+\mathrm{cos}2\vartheta )\mathrm{cos}2\phi |h_{xy}h_{x+y}\frac{}{}][(\stackrel{}{3}+\mathrm{cos}2\vartheta )^2(4h_{xy}^2\mathrm{sin}^22\phi +h_{xy}^2\mathrm{cos}^22\phi )+$$
$$+(1\mathrm{cos}2\vartheta )^2THD_{x+y}^2h_{x+y}^2+2(1\mathrm{cos}2\vartheta )(3+\mathrm{cos}2\vartheta )\mathrm{cos}2\phi \stackrel{}{h_{xy}},\stackrel{}{h_{x+y}}]^1\}^{1/2},$$
(27)
The first of (26) follows immediately from (21); the second one is obtained from (22) using Schwartz inequality. The supremum of the function $`Q(\vartheta ,\phi ,e)`$ occurs at $`\vartheta =\pi /2`$, $`\phi =m\pi `$, $`e`$, where $`Q(\pi /2,m\pi ,e)\stackrel{<}{}1.5`$ (see fig. 6).
The truncation orders required to keep $`THD_{\times ,+}0.01`$, deduced from (26) are collected in Table-I.
| $`e_0`$ | $`N_T`$ |
| --- | --- |
| .1 | 4 |
| .2 | 6 |
| .3 | 8 |
| .4 | 11 |
| .5 | 15 |
| .6 | 22 |
| .7 | 36 |
| .8 | 68 |
| .9 | 206 |
Table I - Truncation orders needed to keep $`THD_{+,\times }.1`$.
## 4a A GENERALIZED CARLINI-MEISSEL FORMULA
A key issue for an efficient computation of waveform-templates based on (10), (11) and (12) involves clever evaluation of terms like:
$$J_n(ne),J_{n\pm 1}(ne),J_{n\pm 2}(ne).$$
(28)
It is well known that, in general, whenever the argument and the order are close (here, in fact they are proportional through the orbital eccentricity $`e`$), numerical computation of Bessel functions either by series summation (Abramowitz and Stegun, 1968, ch. X), or by (re-normalized, downward) recurrence (Press et al, 1992, Sect. 6.5) is inefficient. As a convenient alternative, we suggest the following generalization of the well-known (see Watson, 1976, ch. XVII) Carlini-Meissel (henceforth CM) expansion:
$$J_{n\pm k}(ne)J_n^{(CM)}(ne)\mathrm{\Psi }_{\pm k}(n,e),$$
(29)
where (see Appendix C for the detailed deduction):
$$J_n^{CM}(ne)=\frac{\left({\displaystyle \frac{ne}{2}}\right)^n}{n!}\left(\frac{1+\sqrt{1e^2}}{2}\right)^n(1e^2)^{1/4}\mathrm{exp}\left\{n\left[\sqrt{1e^2}1\right]+n^1\left[\frac{3e^22}{24(1e^2)^{3/2}}+\frac{1}{12}\right]\right\},$$
(30)
$$\mathrm{\Psi }_{\pm k}(n,e)=\frac{n!}{(n\pm k)!}\left(\frac{ne}{1+\sqrt{1e^2}}\right)^{\pm k}exp\left\{\frac{1}{n}\left[\frac{k}{2}\frac{e^2}{1e^2}+\frac{k^2}{2}\left(1\frac{1}{\sqrt{1e^2}}\right)\right]\right\}.$$
(31)
Using (31) to evaluate the Fourier coefficients $`\stackrel{(n)}{\stackrel{}{h}}`$ does not significantly spoil the accuracy of the waveforms. Indeed, spectral truncation according to Table-I still yields THD values below $`0.1`$.
## 5 PROTOTYPE SOURCES
As an application of the above, we wrote a code for waveform template construction, and used it to compute the waveforms for several prototype sources. Taylor et al. (1993) provide data for 24 binary pulsars. In Table II we quote PSR 1913+16 and PSR 1534+12, as possible paradigm sources for space detectors, being respectively the most popular and closest known binary pulsars.
| Binary | 1534+12 | 1913+16 |
| --- | --- | --- |
| Right ascension $`B1950`$ | 15:34:47.686 | 19:13:12.46769 |
| Declination $`B1950`$ | +12:05:45.23 | +16:01:08.0323 |
| Orbital inclination $`i`$ \[degrees\] | 74 | 45 |
| Distance \[kpc\] | 0.68 | 7.13 |
| Projected semimajor axis $`a_i\mathrm{sin}i[lights]`$ | 3.729468 | 2.3417592 |
| Eccentricity $`e`$ | 0.2736779 | 0.6171308 |
| Orbital period $`P_b`$ \[d\] | 0.4207372998 | 0.322997462736 |
| Companion masses $`[M_{}]`$ | 1.34, 1.34 | 1.42, 1.41 |
| $`\xi _1`$ ($`10^6`$) | 3.4849 | 4.3102 |
| $`\xi _2`$ ($`10^3`$) | 1.3200 | 1.4680 |
| $`\xi _3`$ ($`10^{14}`$) | 7.6549 | 13.023 |
| $`\mathrm{\Delta }`$ | 0 | $`3.5336\times 10^3`$ |
| $`\chi `$ | 434777882.4767 | 316085232.7313 |
| $`h_0`$ ($`10^{23}`$) | 16.518 | 2.0575 |
Table II - Paradigm compact binary sources
The gravitational waveforms at $`\vartheta =\phi =0,45,90deg`$, computed using $`8`$ harmonics for PSR1534+12 and $`22`$ harmonics for PSR1913+16 (consistent with Table-I) are displayed in fig.s 7.1-7.15, and fig.s 8.1-8.15, respectively. By comparison, the waveforms corresponding to $`e=0`$ are also drawn.
## 6 CONCLUSIONS
The main results in this paper can be summarized as follows. Orbital eccentricity should not be neglected in detecting gravitational waves from steady-state binaries, for which the simple Peters Mathews model has been shown to be accurate enough. GW spectral truncation criteria have been discussed, and computationally efficient tools/techniques have been introduced for constructing reliable templates. We stress that the above tools/techniques could be readily extended, to higher order PN models with relative ease.
## ACKNOWLEDGEMENTS
V. Pierro has been a Visiting Scientist at the European Space Research & Technology Centre ESTEC-ESA, under a grant from the University of Salerno; A.D.A.M. Spallicci, formerly staff at ESTEC-ESA, has been a Visiting Professor at the University of Salerno in 1996. Both wish to express their appreciation to the hosting Institutions.
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## APPENDIX A: RELEVANT TO EQ.S (10) TO (16).
In the weak-field slow-motion approximation, the cartesian far-field harmonic-gauge metric tensor deviation components in $`(8)`$, $`(9)`$ are simply related to the source quadrupole tensor $`I_{ij}`$ through:
$$h_{xy}=\frac{2G}{c^4r}\frac{d^2I_{xy}}{dt^2},h_{xx}=\frac{2G}{c^4r}\frac{d^2I_{xx}}{dt^2},h_{yy}=\frac{2G}{c^4r}\frac{d^2I_{yy}}{dt^2},$$
$`(A1)`$
where:
$$I_{xx}=\mu \rho ^2\mathrm{cos}^2(\varphi ),I_{yy}=\mu \rho ^2\mathrm{sin}^2(\varphi ),I_{xy}=\mu \rho ^2\mathrm{cos}(\varphi )\mathrm{sin}(\varphi ),$$
$`(A2)`$
$`\rho `$ being the companion star separation, $`e`$ the eccentricity, $`\varphi `$ the true anomaly and $`\mu `$ the reduced mass.
The relevant terms of the (reduced) quadrupole moment can be conveneintly rewritten:
$$I_{xx}=\mu a^2\xi ^2,I_{yy}=\mu a^2\eta ^2,I_{xy}=\mu a^2\xi \eta ,$$
$`(A3)`$
where $`a`$ is the orbit semimajor axis,
$$\xi =\left(\frac{\rho \mathrm{cos}\varphi }{a}\right),\eta =\left(\frac{\rho \mathrm{sin}\varphi }{a}\right).$$
$`(A4)`$
Then, using the well known Keplerian equations (see, e.g., Watson, 1976, ch. XVII)
$$\frac{\rho \mathrm{cos}\varphi }{a}=\mathrm{cos}Ee,\frac{\rho \mathrm{sin}\varphi }{a}=(1e^2)^{1/2}\mathrm{sin}E,$$
$`(A5)`$
where $`E`$ is the eccentric anomaly, and the relation between the latter and the mean anomaly $`M`$,
$$M=\frac{2\pi t}{T}=Ee\mathrm{sin}E,$$
$`(A6)`$
one can expand $`\xi ^2`$ ,$`\eta ^2`$ and $`\xi \eta `$ into Fourier series of argument $`M`$, taking properly into account their parities, viz.:
$$\xi ^2=\frac{\gamma _0}{2}+\underset{n=1}{\overset{\mathrm{}}{}}\gamma _n\mathrm{cos}(nM),$$
$`(A7)`$
$$\eta ^2=\frac{\delta _0}{2}+\underset{n=1}{\overset{\mathrm{}}{}}\delta _n\mathrm{cos}(nM),$$
$`(A8)`$
$$\xi \eta =\underset{n=1}{\overset{\mathrm{}}{}}\eta _n\mathrm{sin}(nM).$$
$`(A9)`$
The relevant Fourier coefficients are readily found. Hence, using $`(A5)`$ and $`(A6)`$:
$$\gamma _n=\frac{2}{n\pi }_0^\pi \mathrm{sin}[n(Ee\mathrm{sin}E)](\mathrm{sin}2E2e\mathrm{sin}E)𝑑E,$$
$`(A10)`$
$$\delta _n=\frac{2}{n\pi }(1e^2)_0^\pi \mathrm{sin}[n(Ee\mathrm{sin}E)]\mathrm{sin}2EdE,$$
$`(A11)`$
$$\eta _n=\frac{2}{n\pi }(1e^2)^{1/2}_0^\pi \mathrm{cos}[n(Ee\mathrm{sin}E)](cos2Ee\mathrm{cos}E)𝑑E.$$
$`(A12)`$
Upon repeated use of trivial trigonometric identities, and in view of the integral definition of the Bessel function of the $`1`$st kind,
$$J_\nu (\alpha )=\frac{1}{\pi }_0^\pi cos[\nu x\alpha \mathrm{sin}x]𝑑x,$$
$`(A13)`$
the Fourier coefficients $`(A10)`$ to $`(A12)`$ can be written:
$$\gamma _n=\frac{1}{n}\left[J_{n2}(ne)J_{n+2}(ne)\right]\frac{2e}{n}\left[J_{n1}(ne)J_{n+1}(ne)\right],$$
$`(A14)`$
$$\delta _n=\frac{1}{n}(1e^2)\left[J_{n2}(ne)J_{n+2}(ne)\right],$$
$`(A15)`$
$$\eta _n=\frac{1}{n}(1e^2)^{1/2}\left[J_{n2}(ne)+J_{n+2}(ne)2J_n(ne)\right].$$
$`(A16)`$
Using $`(A14)`$ to $`(A16)`$ and $`(A7)`$ to $`(A9)`$ in $`(A1)`$ to $`(A3)`$ gives equations $`(10)`$ to $`(16)`$.
## APPENDIX B: RELEVANT TO EQUATIONS $`(22)(26)`$
In order to establish eq.s $`(22)`$ to $`(26)`$ one may repeatedly use the recurrency formula:
$$J_{n\pm 1}(z)=\frac{n}{z}J_n(z)\pm J_n^{}(z),$$
$`(B1)`$
so as to reduce the sought series to combinations of the following (generalized) Kapteyn’s expansions of the second kind:
$$\underset{n=1}{\overset{\mathrm{}}{}}n^2[J_n^{^{}}(ne)]^2,$$
$`(B2)`$
$$\underset{n=1}{\overset{\mathrm{}}{}}[J_n^{^{}}(ne)]^2,$$
$`(B3)`$
$$\underset{n=1}{\overset{\mathrm{}}{}}\frac{[J_n^{^{}}(ne)]^2}{n^2},$$
$`(B4)`$
$$\underset{n=1}{\overset{\mathrm{}}{}}n^2J_n^2(ne),$$
$`(B5)`$
$$\underset{n=1}{\overset{\mathrm{}}{}}J_n^2(ne),$$
$`(B6)`$
$$\underset{n=1}{\overset{\mathrm{}}{}}nJ_n(ne)J_n^{^{}}(ne).$$
$`(B7)`$
These latter can be summed as follows. From the Fourier analysis of Kepler motion, the following equations are readily established (see, e.g., Watson, 1966), ch. 17.2 :
$$\mathrm{cos}E=\frac{e}{2}+2\underset{n=1}{\overset{\mathrm{}}{}}\frac{J_n^{^{}}(ne)}{n}\mathrm{cos}nM,$$
$`(B8)`$
$$\mathrm{sin}E=\frac{2}{e}\underset{n=1}{\overset{\mathrm{}}{}}\frac{J_n(ne)}{n}\mathrm{sin}nM,$$
$`(B9)`$
$$\frac{dE}{dM}=(1e\mathrm{cos}E)^1.$$
$`(B10)`$
Differentiating eq. $`(B8)`$ w.r.t. $`M`$, and using $`(B10)`$, one gets:
$$\frac{\mathrm{sin}E}{1e\mathrm{cos}E}=2\underset{n=1}{\overset{\mathrm{}}{}}J_n^{^{}}(ne)\mathrm{sin}nM,$$
$`(B11)`$
$$\frac{\mathrm{cos}Ee}{(1e\mathrm{cos}E)^3}=2\underset{n=1}{\overset{\mathrm{}}{}}nJ_n^{^{}}(ne)\mathrm{cos}nM.$$
$`(B12)`$
Similarly, from $`(B9)`$:
$$\frac{\mathrm{cos}E}{1e\mathrm{cos}E}=\frac{2}{e}\underset{n=1}{\overset{\mathrm{}}{}}J_n(ne)\mathrm{cos}(nM),$$
$`(B13)`$
$$\frac{\mathrm{sin}E}{(1e\mathrm{cos}E)^3}=\frac{2}{e}\underset{n=1}{\overset{\mathrm{}}{}}nJ_n(ne)\mathrm{sin}(nM),$$
$`(B14)`$
where $`E`$ is the eccentric anomaly, $`M`$ the mean anomaly, and $`e`$ the eccentricity. The following procedure can be then applied to eq.s $`(B8)`$ and $`(B11)`$-$`(B14)`$: i) squaring; ii) taking the average in $`M`$ over $`(0,2\pi )`$, using again eq. $`(B10)`$; iii) using the well known (Euler) transformations:
$$\mathrm{cos}E=(z+z^1)/2,\mathrm{sin}E=i(zz^1)/2,dE=iz^1dz,$$
so as to express the sought series as contour integrals on $`|z|=1`$ of rational functions of $`z`$, which are trivially computed in terms of residues.
As an example, applying the above procedure to eq. $`(B13)`$, one gets:
$$\frac{4}{e^2}\underset{n=1}{\overset{\mathrm{}}{}}J_n^2(ne)=\frac{1}{2\pi }_0^{2\pi }\frac{\mathrm{cos}^2E}{1e\mathrm{cos}E}𝑑E=\frac{1}{2\pi i}_{|z|=1}\frac{(1+z^2)^2}{z^2[4z2e(1+z^2)]}𝑑z=$$
$$=\underset{|z_i|<1}{}Res\left[\frac{(1+z^2)^2}{z^2[4z2e(1+z^2)]}\right]_{z=z_i}.$$
$`(B15)`$
The integrand function on the r.h.s. of $`(B15)`$ has a double pole at $`z=0`$ and two simple ones at $`z=(2e)^1[1(1e^2)^{1/2}]`$. Only two poles above fall within $`|z|<1`$, and $`(B15)`$ gives:
$$\underset{n=1}{\overset{\mathrm{}}{}}J_n^2(ne)=\frac{1}{2}\left[(1e^2)^{1/2}1\right]$$
$`(B16)`$
in agreement with Watson, ch. 17.6, eq. (2). Similarly, starting from $`(B8)`$, $`(B11)`$, $`(B13)`$ and $`(B14)`$ one gets, respectively <sup>8</sup><sup>8</sup>8Note that equation (3) in Watson ch. 17.6, is in error, as seen by comparison with $`(B20)`$, and by direct numerical check. For this (erroneous) result Watson quotes (Schott, 1912). The same error appears in (Prudnikov et al., 1986, sect. 5.7.31).:
$$\underset{n=1}{\overset{\mathrm{}}{}}\frac{[J_n^{^{}}(ne)]^2}{n^2}=\frac{1}{2}\left(1\frac{e^2}{4}\right),$$
$`(B17)`$
$$\underset{n=1}{\overset{\mathrm{}}{}}[J_n^{^{}}(ne)]^2=\frac{1}{2e^2}\left[1(1e^2)^{1/2}\right],$$
$`(B18)`$
$$\underset{n=1}{\overset{\mathrm{}}{}}n^2[J_n^{^{}}(ne)]^2=\frac{4+3e^2}{8\left(1e^2\right)^{5/2}},$$
$`(B19)`$
$$\underset{n=1}{\overset{\mathrm{}}{}}n^2J_n^2(ne)=\frac{e^2(4+e^2)}{16\left(1e^2\right)^{7/2}}.$$
$`(B20)`$
The series $`(B7)`$ can be summed by differentiating both sides of eq. $`(B16)`$ w.r.t. $`e`$. Hence:
$$\underset{n=1}{\overset{\mathrm{}}{}}nJ_n(ne)J_n^{^{}}(ne)=\frac{e}{4\left(1e^2\right)^{3/2}}.$$
$`(B21)`$
## APPENDIX C: GENERALIZED CARLINI-MEISSEL EXPANSIONS
To obtain the generalized Carlini Meissel expansion for $`J_{n\pm k}(ne)`$ we start from Bessel equation for $`J_{n\pm k}(ne)`$:
$$\frac{d^2J_{n\pm k}(ne)}{de^2}+\frac{1}{e}\frac{dJ_{n\pm k}(ne)}{de}+\left[n^2\frac{(n\pm k)^2}{e^2}\right]J_{n\pm k}(ne)=0,$$
$`(C1)`$
and let<sup>9</sup><sup>9</sup>9This formula is suggested by the well-known McLaurin expansions of Bessel functions.:
$$J_n(ne)=J_{n\pm k}(ne)=\frac{n^{(n\pm k)}}{(n\pm k)!}\mathrm{exp}\left[_0^eu_{n\pm k}(z)\right].$$
$`(C2)`$
On letting eq.s $`(C2)`$ into $`(C1)`$, we get:
$$\dot{u}_{n\pm k}+u_{n\pm k}^2+e^1u_{n\pm k}+n^2\frac{(n\pm k)^2}{e^2}=0,$$
$`(C3)`$
then, following Carlini and Meissel, we assume that the following asymptotic representation for $`u_{n\pm k},(k=0,1,2)`$ holds:
$$u_{n\pm k}(z)\frac{u_{n\pm k}^{}(z)}{n}+u_{n\pm k}^0(z)+nu_{n\pm k}^+(z).$$
$`(C4)`$
Substituting $`(C4)`$ into $`(C3)`$, and equating like powers of $`n`$ (as required by consistency), we get:
$$u_{n\pm k}^+=\frac{\sqrt{1z^2}}{z},$$
$`(C5)`$
$$u_{n\pm k}^0=\frac{\pm 2keu_{n\pm k}^+e^2\dot{u}_{n\pm k}^+}{2e^2u_{n\pm k}^+},$$
$`(C6)`$
$$u_{n\pm k}^{}=\frac{k^2eu_{n\pm k}^0e^2\left[\dot{u}_{n\pm k}^0+(u_{n\pm k}^0)^2\right]}{2e^2u_{n\pm k}^+}.$$
$`(C7)`$
Hence:
$$u_{n\pm k}^0=\frac{z}{2(1z^2)}\pm k\frac{1}{z\sqrt{1z^2}},$$
$`(C8)`$
$$u_{n\pm k}^{}=\frac{z^34z}{8(1z^2)^{5/2}}k\frac{z}{(1z^2)^2}+k^2\frac{z^3z}{2(1z^2)^{5/2}}.$$
$`(C9)`$
Carrying out the integrations in $`(C2)`$, and taking into account that $`J_{n\pm k}(0)=\delta _{n\pm k,0}`$ we get:
$$^eu_{n\pm k}^+(z)𝑑z=\sqrt{1e^2}+log\left[\frac{e}{1+\sqrt{1e^2}}\right]+C_{n\pm k}^+,$$
$`(C10)`$
$$^eu_{n\pm k}^0(z)𝑑z=\frac{1}{4}log(1e^2)\pm klog\left[\frac{e}{1+\sqrt{1e^2}}\right]+C_{n\pm k}^0,$$
$`(C11)`$
$$^eu_{n\pm k}^{}(z)𝑑z=\frac{3e^22}{24(1e^2)^{3/2}}k\frac{1}{2(1e^2)}+k^2\frac{e^21}{2(1e^2)^{3/2}}+C_{n\pm k}^{}.$$
$`(C12)`$
Plugging the last three eq.s into eq. $`(C2)`$ we obtain:
$$J_{n\pm k}(ne)=\frac{(ne)^{n\pm k}}{(n\pm k)!}(1+\sqrt{1e^2})^{(n\pm k)}(1e^2)^{1/4}exp\{n\sqrt{1e^2}+n^1[\frac{3e^22}{24(1e^2)^{3/2}}+$$
$$k\frac{1}{2(1e^2)}k^2\frac{1e^2}{2(1e^2)^{3/2}}]+nC_{n\pm k}^++C_{n\pm k}^0+n^1C_{n\pm k}^{}\}.$$
$`(C13)`$
The unknown integration constants can be found by enforcing the following obvious asymptotic equality, valid for all $`n`$:
$$J_{n\pm k}(e0)\frac{(ne/2)^{(n\pm k)}}{(n\pm k)!}.$$
$`(C14)`$
Hence:
$$\{\begin{array}{c}1+C_{n\pm k}^+=0,\hfill \\ \\ C_{n\pm k}^0=0,\hfill \\ \\ 1/12k/2k^2/2+C_{n\pm k}^{}=0.\hfill \end{array}$$
$`(C15)`$
Hence, from $`(C13)`$:
$$J_{n\pm k}(ne)\frac{\left({\displaystyle \frac{ne}{2}}\right)^{n\pm k}}{(n\pm k)!}\left(\frac{1+\sqrt{1e^2}}{2}\right)^{(n\pm k)}(1e^2)^{1/4}exp\{n[\sqrt{1e^2}1]+n^1[\frac{3e^22}{24(1e^2)^{3/2}}+$$
$$k\frac{1}{2(1e^2)}k^2\frac{1e^2}{2(1e^2)^{3/2}}+\frac{1}{12}\pm \frac{k}{2}+\frac{k^2}{2}]\}.$$
$`(C16)`$
The r.h.s. of eq. $`(C16)`$ above will be henceforth denoted as $`J_{n\pm k}^{CM}(ne)`$, and can be more conveniently written as in $`(29)`$ to $`(31)`$.
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# “Embedded solitons”: solitary waves in resonance with the linear spectrum
## 1 Introduction
Recent studies have revealed a novel type of solitary waves (which we will loosely called “solitons”, without assuming integrability of the underlying models) that are embedded into the continuous spectrum, i.e., the soliton’s internal frequency is in resonance with linear (radiation) waves. Generally, such a soliton should not exist, one finding instead a quasi- (delocalized) soliton with nonvanishing oscillatory tails (radiation component) . Nevertheless, bona fide (exponentially decaying) solitons can exist as codimension-one solutions if, at discrete values of the (quasi-)soliton’s internal frequency, the amplitude of the tail exactly vanishes, while the soliton remains embedded into the continuous spectrum of the radiation modes. This requires the spectrum of the corresponding linearized system to consist of (at least) two branches, one corresponding to exponentially localized solutions, and the other to oscillatory radiation modes. In terms of the corresponding ordinary differential equations (ODEs) for the traveling-wave solutions, the origin must be a saddle-centre, that is its linearisation gives rise to both real and pure imaginary eigenvalues.
Examples of such embedded solitons (ESs) were found in water-wave models and in several nonlinear-optical ones, including a Bragg-grating model with the wave-propagation (second-order-derivative) linear terms taken into regard , and a model of the second-harmonic generation (SHG) in the presence of a Kerr nonlinearity . The term “ES” was proposed in the latter work.
ESs are interesting for several reasons, firstly because they frequently appear when higher-order (singular) perturbations are added to the system, which may completely change its soliton spectrum (see, e.g., ). Secondly, optical ESs have considerable potential for applications, just because they are isolated solitons, rather than members of continuous families. Finally, and most crucial for physical applications, ESs are semi-stable objects. That is, as argued in Ref. analytically, in a general form applicable to ESs in any system, and checked numerically for the SHG model with the additional defocusing Kerr nonlinearity, ESs are stable in the linear approximation, but do have a slowly growing (sub-exponential) one-sided nonlinear instability (see Section 4 below).
In the next section we discuss the existence and stability of ESs in four different nonlinear partial-differential-equation (PDE) models. Mathematically speaking, in each case the reduced traveling-wave or steady-state ODEs have the structure of a fourth-order, reversible Hamiltonian system. In accord with what was said above, a parameter region we focus on is where the origin (trivial fixed point) in these ODE systems is a saddle-centre. That is, after diagonalizing the system, one two-dimensional (2D) component of the dynamical system gives rise to imaginary eigenvalues $`\pm i\omega `$ (corresponding to a continuous radiation branch in the linear spectrum of the PDE system), and the other to real eigenvalues $`\pm \lambda `$ (corresponding to a gap in the radiation spectrum).
All four models considered in this paper share the feature that ESs exist as non-generic, codimension-one solutions. Section 3 then presents three different views as to why this should be, which together provide for general insight into the existence and multiplicity of ESs. Section 4 goes on to discuss their stability, arguing that ESs, in general, may be neutrally stable linearly, but suffer from a one-sided sub-exponential instability. This situation has been termed semi-stability in . Section 5 treats “moving” embedded solitons and, finally, Section (6) draws conclusions and briefly discusses potential physical applications of ESs in optical memory devices.
## 2 Physical Examples
### 2.1 An extended 5th-order KdV equation
One of the first systems in which ESs were found (although without being given that name) is an extended 5th-order Korteweg - de Vries (KdV) equation ,
$$u_t=\left[(2/15)u_{xxxx}bu_{xx}+au+(3/2)u^2+\mu \left((1/2)(u_x)^2+(uu_x)_x\right)\right]_x,$$
(1)
which with $`\mu =0`$ reduces to the usual 5th-order KdV equation studied by a number of authors, see, e.g., Refs. . The extended form (1) may be derived via a regular Hamiltonian perturbation theory from an exact Euler-equation formulation for water-waves with surface tension . Looking for traveling-wave solutions $`u(xct)`$, integrating once, setting the constant of integration to be zero, and absorbing the linear term $`u_x`$ by redefining $`a`$, one arrives at the following ODE (the prime stands for $`d/d(xct)`$),
$$\frac{2}{15}u^{\prime \prime \prime \prime }bu^{\prime \prime }+au+\frac{3}{2}u^2+\mu \left[\frac{1}{2}(u^{})^2+(uu^{})^{}\right]=0.$$
(2)
When $`a<0`$, Eq. (2) is in the ES regime, since the linearization around the origin, $`u=0`$, yields both real and imaginary eigenvalues. Note that, for the particular case $`\mu =0`$, it has been proved that there are no dynamical orbits homoclinic to the origin . Nevertheless, at least in the limit $`a0`$, the system does possess a large family of homoclinic connections to periodic orbits, rather than to fixed points . The latter generic solutions correspond to the above-mentioned delocalized solitary waves .
However, ESs do exist in the case $`\mu =1`$ , where one can find the explicit family of homoclinic-to-zero solutions,
$$u(t)=3\left(b+\frac{1}{2}\right)\text{sech}^2\left(\sqrt{\frac{3(2b+1)}{4}}t\right),a=\frac{3}{5}(2b+1)(b2),b1/2.$$
(3)
Moreover, this curve (family) of ESs was found in to be only the first in a countable set of curves that appear to bifurcate from $`a=0`$ at a discrete set of negative values of $`b`$, see Fig. 1. Note that there are numerical difficulties in computing up to the limit point $`a=0`$ for $`b<0`$; this is because, as will be motivated in Section 3.3 below, the bifurcation of these solutions from $`a=0`$ is a ‘beyond-all-orders’ effect
We remark that Grimshaw and Cook found a similar bifurcation-type phenomenon in a system of two coupled KdV equations, although they did not explicitly identify the vanishing of the tail amplitude of the delocalized solitary waves as what we now call ESs. Also, Fujioka and Espinosa found a single explicit ES in a higher-order NLS equation with a quintic nonlinearity.
Note also the following feature of this family of ES states. The first member of the family, viz., the explicit solution (3)) is a fundamental “ground state”, while all other solution branches represent ‘ “excited states”, having the form of the ground-state soliton with small-amplitude “ripples” superimposed on it. As yet unpublished numerical results suggest that the ground-state soliton appears to be dynamically stable (in the sense described in Section 4 below).
### 2.2 A generalized Massive Thirring model
The first example of ESs in nonlinear optics was found in a generalized Thirring model (GTM), introduced long ago in . ESs appear in this model when additional wave-propagation (second-derivative) terms are included ($`k`$ being the carrier wavenumber), so that the model takes the form
$$\begin{array}{ccc}\hfill iu_t+iu_x+(2k)^1\left(u_{xx}u_{tt}\right)+\left[(1/2)|u|^2+|v|^2\right]u+v& =& 0,\hfill \\ \hfill iv_tiv_x+(2k)^1\left(v_{xx}v_{tt}\right)+\left[(1/2)|v|^2+|u|^2\right]v+u& =& 0.\hfill \end{array}$$
(4)
Here $`u(x,t)`$ and $`v(x,t)`$ are right- and left- traveling waves coupled by resonant reflections on the grating. Soliton solutions are sought for as $`u(x,t)=\mathrm{exp}(i\mathrm{\Delta }\omega t)U(\xi )`$, $`v(x,t)=\mathrm{exp}(i\mathrm{\Delta }\omega t)V(\xi )`$, where $`\xi xct`$, $`c`$ and $`\mathrm{\Delta }\omega `$ being velocity and frequency shifts. The functions $`U(\xi )`$ and $`V(\xi )`$ satisfy ODEs
$$\begin{array}{ccc}\hfill \chi U+i(1C)U^{}+DU^{\prime \prime }+\left[(1/2)|U|^2+|V|^2\right]U+V& =& 0,\hfill \\ \hfill \chi Vi(1+C)V^{}+DV^{\prime \prime }+\left[(1/2)|V|^2+|U|^2\right]V+U& =& 0,\hfill \end{array}$$
(5)
where $`\chi \mathrm{\Delta }\omega +\left(\mathrm{\Delta }\omega \right)^2/2k`$, the effective velocity is $`C(1+\mathrm{\Delta }\omega /k)c`$, and an effective dispersion coefficient is $`D\left(1c^2\right)/2k`$. The same equations were derived in for: (i) temporal-soliton propagation in nonlinear fiber gratings, including spatial-dispersion effects; and (ii) spatial solitons in a planar waveguide with a Bragg grating in the form of parallel scores, taking diffraction into regard, with $`t`$ replaced by the propagation coordinate $`z`$, while $`x`$ is the transverse coordinate. If we set $`C=0`$ then Eqs. (5) admit the further invariant reduction $`U=V^{}`$, leading to a single ODE ,
$$DU^{\prime \prime }+iU^{}+\chi U+(3/2)|U|^2U+U^{}=0,$$
(6)
which is equivalent to a (Hamiltonian and reversible) system of four first-order equations for the real and imaginary parts of $`U`$ and $`U^{}`$. In this system, the origin is a saddle-centre, provided that $`D>0`$ and $`|\chi |<1`$. In it was found numerically that Eq. (6) admits exactly three branches of the fundamental ESs (in contrast to Eq. (1) where countably many ES states have been found). Also moving ($`c0`$) ESs, satisfying the 8th-order system $`(\text{5})`$, have also recently been found in Ref. ; see Section 5 below for more details.
### 2.3 A three-wave interaction model
ESs can be found in far greater abundance in a model for spatial solitons, assuming a planar waveguide with a quadratic ($`\chi ^{(2)}`$) nonlinearity , where two fundamental-harmonic (FH) waves $`v_{1,2}`$ are coupled by the Bragg reflections from a set of parallel scores. These two waves then interact nonlinearly, and generate a third wave, the second-harmonic (SH), with its wave-vector equal to the sum of those of the two FH components. The set of equations are:
$$\begin{array}{ccc}\hfill i(v_{1,2})_z\pm i(v_{1,2})_x+v_{2,1}+v_3v_{2,1}^{}& =& 0,\hfill \\ \hfill 2i(v_3)_zqv_3+D(v_3)_{xx}+v_1v_2& =& 0.\hfill \end{array}$$
(7)
Here $`v_3`$ is the SH field, $`x`$ is the normalized transverse coordinate, $`q`$ is a mismatch parameter, and $`D`$ is an effective diffraction coefficient. Solutions to Eq. (7) are sought in the form $`v_{1,2}(x,z)=\mathrm{exp}(ikz)u_{1,2}(\xi )`$, $`v_3(x,z)=\mathrm{exp}(2ikz)u_3`$, with $`\xi xcz`$, $`c`$ being the slope of the soliton’s axis relative to the propagation direction $`z`$. In Ref. , many ESs of the zero-walkoff type ($`c=0`$) were found (summarized in Fig. 2), in which case the ODEs reduce, as in Eqs. (5), to a fourth-order real system. When $`c0`$, one cannot assume all the amplitudes to be real. In this case, one finds “moving” ESs as solutions to an eighth-order real ODE system (see Section 5 below).
### 2.4 A second-harmonic-generation system
The model we shall study in most detail is that for which the term ‘ES’ was first proposed in , viz., a nonlinear optical medium with competing quadratic and cubic nonlinearities
$$\begin{array}{ccc}\hfill iu_z+(1/2)u_{tt}+u^{}v+\gamma _1(|u|^2+2|v|^2)u& =& 0,\hfill \\ \hfill iv_z(1/2)\delta v_{tt}+qv+(1/2)u^2+2\gamma _2(|v|^2+2|u|^2)v& =& 0.\hfill \end{array}$$
(8)
Here, $`u`$ and $`v`$ are FH and SH amplitudes, $`\delta `$ is a relative SH/FH dispersion coefficient, $`q`$ is a phase-velocity mismatch, and $`\gamma _{1,2}`$ are cubic (Kerr) nonlinear coefficients. In the absence of the Kerr nonlinearities, these equations are the same as those use by Karamzin and Sukhorukov in 1974 to obtain their famous $`\chi ^{(2)}`$ soliton solution, in which both the FH and SH fields are proportional to sech<sup>2</sup>. A detailed analysis of the higher-order soliton solutions in that model with purely quadratic nonlinearity can be found in Ref. . However the solutions that we will consider here are in a different class, in that the FH field will be more like a sech than a sech<sup>2</sup>.
The particular case of Eqs. (8) with $`\delta =1/2`$ is specially important, as it corresponds, with $`t`$ replaced by the transverse coordinate $`x`$, to a second-harmonic-generation model in the spatial domain (in fact, in a nonlinear planar optical waveguide). In this special case, the model (8) is Galilean invariant, which allows one to generate a whole family of “moving” solitons from the single zero-walkoff one . At all other values of $`\delta `$, construction of a “moving” (nonzero-walkoff) soliton is a nontrivial problem.
Stationary solutions to Eq. (8) are sought for in the form $`u=U(t)\mathrm{exp}(ikz)`$, $`v=V(t)\mathrm{exp}(2ikz)`$, where $`k`$ is real, and $`U,V`$ satisfy ODEs
$$\begin{array}{ccc}\hfill (1/2)U^{\prime \prime }kU+U^{}V+\gamma _1(|U|^2+2|V|^2)U& =& 0,\hfill \\ \hfill (1/2)\delta V^{\prime \prime }+(q2k)V+(1/2)U^2+2\gamma _2(|V|^2+2|U|^2)V& =& 0.\hfill \end{array}$$
(9)
In Ref. , ES solutions to these equations were found for $`\delta >0`$ and $`\gamma _{1,2}<0`$ (which implies anomalous and normal dispersions, respectively, at FH and SH, and self-defocusing Kerr nonlinearity. Alternatively, the same case may be physically realized as the normal and anomalous dispersions at FH and SH and self-focusing Kerr nonlinearity). In the same work, stability of these ESs was studied in detail.
Here we shall present new results for $`\delta <0`$, which corresponds to a more common case where the dispersion has the same sign at both harmonics. The results can be naturally displayed in the form of curves of ESs in the $`(\delta ,q)`$ or $`(q,k)`$ parameter planes, see Figs. 3 and 4 below.
## 3 Existence
We now give three distinct explanations of why and how an ES may exist.
### 3.1 Nonlinearizablity
Consider first the ODE system (9). It is important to note that the system gets fully decoupled in the linear approximation, the linearization of its second equation immediately telling one that no true soliton (with exponentially decaying tails) can exist inside the continuous (radiation-mode) SH spectrum. However, it may happen that the tail of the soliton’s SH component decays at the same rate as the square of the tail of the FH component. In that case, the second equation of the system (9) is nonlinearizable, which opens the way for the existence of truly localized solitons inside the continuous spectrum. Note from the profiles of the solutions in Figs. 3 and 4, that the $`V`$-component appears to decay to zero much faster than $`U`$, in accordance with this nonlinearizability property.
Let us remark on some qualitative features of the ES branches presented in Fig. 3. First, the computed branches appear to end in “mid air”. At the low $`q`$ end, we have a boundary. The condition for the origin of the ODE system (9) to be a saddle-centre is $`k<q/2`$. For $`q`$ below this limit, one may linearize (9), and find regular (non-embedded) solitons, for any $`q`$, provided $`k>0`$. These latter solitons continuously match into the embedded solitons at $`q=k/2`$. Thus, since we took $`k=0.3`$, we have that both branches of ESs will end at $`q=0.6`$. At the high $`q`$ end, there were numerical difficulties in continuing the branches, using the software AUTO , for precisely the same reasons why the computations presented in Fig. 1 were difficult at $`a=0^{}`$ for $`b<0`$. As one can see in the figure, one does have the amplitudes of both $`U`$ and $`V`$ tending to zero (note also that the $`V`$-component is much smaller, in accordance with the nonlinearizability principle).
Second, note that for high $`q`$, both ES solutions appear to be single-humped and fundamental. But upon going to smaller $`q`$, where the amplitudes of both $`U`$ and $`V`$ become larger, it is apparent that the right-hand branch is a higher-order state, with a superimposed ripple in the $`V`$ component (akin to the second branch in Figs. 1 and 2).
Finally, we remark that only this second branch passes through the physically significant value of $`\delta =1/2`$. For the spatial ESs, corresponding to $`\delta =1/2`$, one is interested in how the solutions will vary as a function of $`k`$ as well. That is shown in Fig. 4. Again, the low-$`q`$ limit corresponds to the boundary of the saddle-centre region is $`k=q/2`$.
### 3.2 Homoclinic orbits to saddle-centres
Next, let us try to understand why ESs should be of codimension-one. We start from a general fourth-order Hamiltonian and reversible (invariant with respect to reflections of time and ‘velocity’ variables) ODE system. Linearizing about the zero solution (the origin), we assume the system to have a saddle-centre equilibrium. Thus, we have
$`\dot{x}=f(x),xIR^4,R,`$ $`R^2=\text{id},Rf(Rx)=f(x),`$ (10)
$`H(x)=\text{const. along solutions},`$ $`\text{eigs}(Df(0))=\{\lambda ,\lambda ,i\omega ,i\omega \}.`$
For example, for the ODE (2) the reversibility operator $`R`$ is given by
$$(u,u^{},u^{\prime \prime },u^{\prime \prime \prime })(u,u^{},u^{\prime \prime },u^{\prime \prime \prime }),$$
and for the system (9) by
$$R:(U,U^{},V,V^{})(U,U^{},V,V^{}).$$
A homoclinic orbit connecting such a saddle-centre equilibrium to itself is formed by a trajectory simultaneously belonging to the one-dimensional unstable and stable manifolds of the origin. Both of these manifolds lie in a 3D phase space $`H(0)`$. Therefore, were it not for the reversibility, such homoclinic orbits would be of codimension-two in general , since we require the coincidence of two lines in the three-dimensional space. But, reversible homoclinic solutions (i.e., solutions that somewhere intersect the fixed-point set of the reversibility) are of codimension-one. This is because the unstable manifold and $`\text{fix}(R)H(0)`$ are both one-dimensional, and we only require a point intersection between them. Hence, varying two parameters, we should expect to see ESs occurring along lines in the corresponding two-dimensional parameter plane. Moreover, the solutions themselves must be reversible; asymmetric ESs would be of a higher codimension still.
Mielke, Holmes and O’Reilly proved a general theorem valid in the neighborhood of such a curve in the parameter plane of reversible saddle-centre homoclinic orbits: under a sign condition, essentially governing how reversibility and the Hamiltonian interact, they showed that there will be an accumulation of infinitely many curves of $`N`$-pulse “bound states” of the primary homoclinic orbit, for each $`N>1`$. Note that, for systems that are reversible and also have odd symmetry (such as (6)), the sign condition is always satisfied by virtue of the system’s admitting both reversibilities $`R`$ and $`R`$ (i.e., the model (6) is symmetric too under $`UU`$, and hence is also invariant under the reversibilities $`(U,U^{})(U^{},U_{}^{}{}_{}{}^{})`$ and $`(U,U^{})(U^{},U^{})`$). Here, the sign condition determines whether there are “up-up” or “up-down” bound states. In Ref. , a large number of the bound states of the “up-down” type were found for the generalized massive Thirring model (6) in agreement with this theory. We also remark that Buryak and co-workers (see ) found a similar discrete sequence of bound states of ‘nonexistent’ dark solitons in a SHG model and a higher-order NLS equation, however none of these solutions were linearly stable.
So far, in every case that we have investigated, we have found the higher-order ESs to be unstable against linear perturbations. A distinction should be stressed between what one would call “bound-states” — which are like several copies of a fundamental soliton placed end to end — and the higher-order solitons, such as those displayed in Figs. 1 and 2, which are like a fundamental with internal ripples. At the moment, there seems to be no connection between these states, although it is still may happen that, as some parameter is varied, a continuous branch may connect solutions of the two different types.
### 3.3 A singular limit
We now look at a mechanism which explains how fundamental ESs (possibly with ripples) may appear from the singular limit $`\lambda 0^+`$. Such a limit for the general class of systems (10) has been studied using the normal-form theory by Lombardi , incorporating careful estimation of various exponentially small terms (cf. related results obtained using exponential asymptotics, e.g. ). A crucial additional ingredient we shall add to the Lombardi’s work is that $`\lambda `$ and $`\omega `$ are assumed to play the role of two independent parameters. We provide here only an oversimplified sketch, more details will appear elsewhere .
The appropriate normal form is
$$\begin{array}{cccc}\hfill \dot{x}_1& =& x_2,\hfill & \\ \hfill \dot{x}_2& =& \lambda ^2x_1(3/2)x_1^2b_1(x_3^2+x_4^2)\hfill & +\rho N_2(x;\lambda ),\hfill \\ \hfill \dot{x}_3& =& x_4(\omega +b_2x_1)\hfill & +\rho N_3(x;\lambda ),\hfill \\ \hfill \dot{x}_4& =& x_3(\omega +b_2x_1)\hfill & +\rho N_4(x;\lambda )\hfill \end{array}$$
(11)
where $`b_{1,2}`$ are $`\omega `$-dependent constants to be determined for a particular system, and $`N_i`$ are higher-order (remainder) terms that break (for nonzero $`\rho `$) the completely-integrable structure of the truncated normal form. It is not difficult to see that the truncated system possess a $`\text{sech}^2`$-like homoclinic connection to the origin, whose amplitude is $`O(\lambda )`$. The key question is whether this homoclinic orbit persists under the inclusion of the remainder terms. This question can be posed in terms of the vanishing of a certain Melnikov integral (whose vanishing measures the splitting distance between the stable and unstable manifolds), which, after a lengthy calculation , can be written as
$$I=\frac{\rho }{\lambda ^2}\mathrm{exp}(\omega \pi /\lambda )(\mathrm{\Lambda }(N_3,N_4,\omega )(1+O(\lambda )+O(\rho )),$$
where $`\mathrm{\Lambda }(N_3,N_4,\omega )`$ can be computed explicitly for each monomial which is pure in $`x_1`$ and $`x_2`$ in the Taylor-series expansion of either $`N_3`$ or $`N_4`$.
Now, something beautiful happens because, for each such monomial, $`\mathrm{\Lambda }`$ turns out to be a (single-signed $`\omega `$-dependent) constant multiple of either a Bessel function, $`\mathrm{\Lambda }J_n(4\sqrt{\omega b_2})`$ for $`b_2>0`$, or modified Bessel function $`\mathrm{\Lambda }I_n(4\sqrt{\omega |b_2|})`$ for $`b_2<0`$, for some integer $`n`$. Recall the basic properties of the Bessel functions, according to which $`J_n(x)`$ has infinitely many zeros for $`x>0`$, whereas the modified Bessel function $`I_n(x)=i^nJ_n(ix)`$ has no zeros. Hence, notice the crucial role played by the coefficient $`b_2`$ (in the truncated, scaled normal form (11)) in the case when the remainder $`N`$ consists of a single monomial in $`(x_1,x_2)`$. If $`b_2>0,`$ there will be an infinite number of $`\omega `$ values corresponding to zeros of $`\mathrm{\Lambda }`$ and hence homoclinic solutions will exist for small $`\rho `$ and $`\lambda `$. However, if $`b_2<0,`$ then $`\mathrm{\Lambda }`$ is strictly non-zero and hence there are no homoclinic solutions to the origin.
The coefficient $`b_2`$ is easy to calculate in the examples with the quadratic nonlinearities, such as the extended 5th-order KdV model (2) with $`\mu =1`$ (see for details). There it is found $`b_2>0`$ and this entirely explains the approximately periodic sequence of points on the negative $`b`$-axis at which the homoclinic solutions bifurcate from $`a=0`$ with the zero amplitude, as shown in Fig. 1. In contrast, for $`\mu =1`$, $`b_2`$ is negative and no ES bifurcates from $`a=0`$. For models with purely cubic nonlinearity, such as the generalized Thirring model (6), it can be shown that all quadratic coefficients in the normal form (11) vanish, hence a new, odd-symmetric normal form would need to be studied. This is left for future work. We mention also that for the SHG model (9), although there are quadratic terms in it, the coefficient $`b_2`$ is identically zero, so this analytic technique gives no information.
## 4 Stability
The stability of the embedded solitons in the SHG model has been studied both numerically and analytically in . It was shown that the fundamental (single-hump) ES is linearly stable, but nonlinearly semi-unstable, while all the multi-humped ES are linearly unstable. In the semi-stability analysis of the fundamental ES, a crucial role is played by the energy, as ESs are isolated solutions with uniquely determined values of the energy, with the adjacent delocalized soliton states, on either side, having an infinite energy. It was argued in (see also a similar argument given by Buryak ) that perturbations which slightly increases the fundamental ES’s energy can be safe, while perturbation which decreases the energy inevitably triggers a slow (sub-exponential) decay of the soliton into radiation. Thus, the weak instability of an ES is one-sided. This fact also follows from a simple argument that the usual exponential instability is always dual-sided. Equivalently, the usual instability is linear, while the weak one-sided instability of ESs must be nonlinear.
The situation is the same in the present case where $`\delta <0`$. A study of the linearized equation around the ES shows that the fundamental ES branch (the left one in Fig. 3) is linearly stable, while the branch of multi-humped ES (the right one in Fig. 3) is linearly unstable. The semi-stability argument for the fundamental ES branch also applies here. Positive energy perturbations can be safe, while negative energy perturbations trigger decay of ES. However, as we shall see below, for $`\delta <0`$, this decay seems to be significantly slower than that found in Ref. for $`\delta >0`$. As was done there, we numerically simulated the system $`(\text{8})`$, with the initial data
$$u(0,t)=U(t)+\alpha _1\text{sech}2t,v(0,t)=V(t)+\alpha _2\text{sech}2t,$$
(12)
where $`U(t)`$ and $`V(t)`$ is an ES solution on the left-hand branch of Fig. 3 (the values are $`\delta =2.9292`$, $`q=6.0556`$, $`k=0.2936251`$ and $`\alpha _1=\alpha _2=0.05`$). Figs. 5 and 6 depict, respectively, the effects of the positive-energy and negative-energy perturbations. In both cases, we observe fast oscillations on top of slow ones in the evolution of $`|u|`$ and $`|v|`$ at $`t=0`$. These oscillations are very similar to those reported in for perturbed solitons in the standard SHG system with no cubic nonlinearity (note that those solitons are ordinary ones, rather than ES). In that case, the fast oscillations were attributed to an intrinsic mode, while the slow oscillations were attributed to beatings between the intrinsic mode and a quasi-mode. The latter one is localized in the FH component, but resonates with the continuous spectrum in the SH component. Both oscillations could last for a very long time, even though they were expected to eventually decay due to a very weak radiation damping.
We believe that similar mechanisms are also at work in our model. However, there are important differences because of the fact that the solitons in our model are embedded, and those in the model considered in Ref. were not. On the other hand, there are important similarities because perturbations of an ES naturally resonate with the continuous spectrum of the SH component, and we do see such slow oscillations as well. Although according to the semi-stability argument, a negative perturbation of an ES, as shown in Fig. 6 would eventually decay, it is a very slow decay. Detailed examination of the numerical solutions shows that the central pulse (the $`v`$ component) in Fig. 6 keeps shedding oscillating tails into the far field. However, the tail amplitudes are extremely small (about 0.001 or smaller). This is why the expected decay of the perturbed ES is not obvious in that figure. Because of this, the actually observed evolution is dominated by the beating and internal oscillations, just as in Fig. 5 with a positive perturbation, and as in Ref. . It will take an extremely long time for the pulse in Fig. 6 to show considerable decay. In fact, it is clear that the ESs in the present model with $`\delta <0`$ are virtually stable, when compared to the previously considered case , $`\delta >0`$, where the semi-instability was a really observed feature. The relative stability of an ES for $`\delta <0`$ is a new result reported in the present paper, and its importance is quite obvious.
If the ES were linearly unstable, the semi-stability argument would not apply. This is the case for the solutions investigated in , and also for the right-hand branch depicted in Fig. 3, which corresponds to the case of a nonlinear planar optical waveguide in the spatial domain. To verify this, we have performed a time integration of the PDE, the results of which are presented in Fig. 7. One can see the onset of a violent exponential instability. We have verified that this solution does have exponentially unstable eigenmodes, from a numerical study of the linearization of Eq. (8) about this solution.
In view of these results we conjecture the following. Fundamental ESs are, in general, linearly neutrally stable but semi-stable nonlinearly. The higher-order ESs (which usually have internal ripples in their profiles) are, generally, linearly unstable. Preliminary results for the extended 5th-order KdV (1) indicate qualitatively the same properties. This would also accord with previous numerical results for higher-order NLS equations that an isolated fundamental ES is semi-stable whereas multihumped bound-states are exponentially unstable.
To conclude this section, we notice that multihumped solitons (which are not a subject of the present paper) of the ordinary (nonembedded) type were found in many models, see, e.g., early works dealing with the Langmuir waves in plasmas, and a recent work on the resonant three-wave interactions
## 5 Moving embedded solitons
We will now discuss a possibility that ESs in the generalized Massive Thirring model (5) may be moving at a non-zero velocity $`c`$ (this discussion applies also to ESs in the three-wave system (7) with a non-zero walkoff). In this case, the reduction of the 8th-order ODE to a 4th-order one is no longer possible. Moreover, since the 8th-order system is obtained by separating the real and imaginary parts of a 4th-order complex system, the spectrum will have a double degeneracy when $`c=0`$. This means that, in the parameter region of interest in the $`(\chi ,c,D)`$-space, the linearization yields four pure imaginary eigenvalues, plus two with positive real parts and two with negative ones. A similar counting argument, as in Section (3.2), shows that reversible homoclinic orbits to such equilibria are of codimension two. Hence ESs lie on one-dimensional curves in the three-parameter space. Alternatively, this property can be explained as follows: in addition to the energy $`E`$, the full system also preserves the momentum $`P`$, so we can view a moving ES as being isolated in both invariants, i.e., the ES solution family is described by curves $`E(D)`$ and $`P(D)`$. Finally, we find that such curves can be found naturally as bifurcation points at $`c=0`$ from curves of the zero-velocity ESs.
Fig. 8 shows one such bifurcation occurring from one of the three branches of fundamental ESs for the Thirring model. The ES branches terminate (as at point “4”) and beyond that, ES become ordinary (non-embedded) solitons (i.e., where the saddle-centre equilibrium becomes a pure saddle with only real eigenvalues).
Similarly, Fig. 9 shows two branches of moving ESs bifurcating from the ground-state ES (the nearly vertical line) in the three-wave model. Note that the two bifurcating curves become globally connected on the left, through a regular tangent (fold) bifurcation.
Finally, let us turn to the second-harmonic model (8). As it was mentioned above, precisely at the value $`\delta =1/2`$, this model gives rise to moving solitons in a trivial way, via the Galilean transformation. Search for moving ES in this model with $`\delta 1/2`$ is the subject of ongoing work.
## 6 Conclusion
In this paper, we have given a brief overview of recent results that establish the existence of isolated (codimension-one) solitons embedded into the continuous spectrum of radiation modes. A necessary condition for the existence of the embedded solitons is the presence of (at least) two different branches in the spectrum of the corresponding linearized system, so that one branch can correspond to purely imaginary eigenvalues, and another to purely real ones. The fundamental (single-humped) embedded solitons are always stable in the linear approximation, being subject to a weak sub-exponential one-sided instability. Moving embedded solitons may also exist as codimension-two solutions. Moreover, bound states in the form of multi-humped embedded solitons exist too, but they are linearly unstable.
We should remark that the work presented and reviewed here does not necessarily represent the very first time that the existence of ESs have been established in any physical model. For example, as stated earlier, multi-humped dark ESs were already observed in a SHG model and for higher-order NLS equations where they were found to be unstable. Fundamental ESs were also found for the latter with competing cubic and quintic nonlinearities and their existence was also suggested for a coupled KdV system , without evidence to suggest their stability. Also recent work has shown the existence of front solutions to higher-order Frekel-Kontorova models which are in resonance with the linear spectrum; so called ‘embedded kinks’. Nevertheless, we believe the results presented in this paper — the new stability results for fundamental ESs in Section 4, together with the three distinct views given in Section 3 as to why ESs should exist, and the similarities found in Section 2 between their existence properties for four distinct models — provide a new theory for ESs as a phenomenon in their own right.
Moreover, the very fact that ESs are isolated states suggest that they may find potential application in photonics, such as in all-optical switching. For example, taking the examples in Sections 2.2 and 2.3, switching from one ES state to a neighboring one with a smaller energy might be easily initiated by a small perturbation, in view of the semi-stability inherent to ESs. Switching between two branches of moving ESs with $`c0`$ might be quite easy to realize too, due to the small energy-flux and walkoff differences between them. There remains much work to be done in investigating these potential aplications further.
Many issues concerning the embedded solitons remain open and are a subject of ongoing investigations. In particular, immediate questions arise concerning moving ESs, and interactions caused by collisions between them.
We appreciate a useful discussion with A. Buryak. The research of JY was supported in part by the NSF and the AFOSR. The research of DJK was supported in part by the AFOSR. A collaboration between ARC and BAM was supported by a fellowship granted by the Benjamin Meaker Foundation through the University of Bristol. ARC is suppoted by an Advanced Fellowship from the EPSRC.
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# A Phenomenological Expression for Deuteron Electromagnetic Form Factors Based on Perturbative QCD PredictionsThe Project Supported by the National Science Foundation of China (NSFC)and Grant of Academia Sinica.
## 1 Introduction
It was found a long time ago that the traditional meson-nucleon picture can not explain the form factors of the deuteron as the momentum transfer $`Q^2>1`$ GeV$`^{2[1]}`$. It means that the fundamental degrees of freedom of QCD, the quark and gluon degrees of freedom, must be taken into account. However, A pure perturbative QCD (PQCD) calculation<sup></sup> shows that the theoretical prediction is much smaller than the data at currently accessible energies, although it may be correct in very large $`Q^2`$. To explain the deuteron form factors in the intermediate energy region, we have suggested a QCD-inspired model for the helicity-zero to zero matrix elements $`G_{00}^+`$ in the light-cone frame<sup></sup>, which should be the dominant amplitude from PQCD predictions<sup></sup>. This model can explain the data of the deuteron electromagnetic structure function $`A(Q^2)`$ and shows that $`G_{00}^+`$ is already dominant at $`Q^2`$ of 1 GeV<sup>2</sup>. Furthermore, it was found that $`G_{+0}^+`$ can not be neglected in the form factor $`G_M(Q^2)^{[5,6]}`$ and $`G_+^+`$ plays an important role in $`G_Q(Q^2)^{[6]}`$. Neglecting them will result in contradictions with both the data and the conventional meson-nucleon picture in the low energy region.
Perturbative QCD (PQCD) predicts that $`G_{00}^+`$ becomes the dominant helicity amplitude at large $`Q^2`$ and that $`G_{+0}^+`$ and $`G_+^+`$ are suppressed by factors $`\mathrm{\Lambda }_{\mathrm{QCD}}/Q`$ and $`\mathrm{\Lambda }_{\mathrm{QCD}}^2/Q^2`$, respectively. neglecting $`G_{+0}^+`$ and $`G_+^+`$ contributions at large $`Q^2`$, we have the relation approximately,
$$G_C:G_M:G_Q=(1\frac{2}{3}\eta ):2:1,$$
$`(1)`$
where $`\eta =Q^2/4M^2`$ and $`M`$ is the mass of the deuteron. However, the helicity-flip amplitude $`G_{+0}^+`$ and $`G_+^+`$ contribute to $`G_M`$ and $`G_Q`$ because of the kinematic enhancement in the intermediate energy region. In order to explore the role of $`G_{+0}^+`$ and $`G_{+0}^+`$, we have tried to expand them to the second order in $`\mathrm{\Lambda }_{\mathrm{QCD}}/M^{[6]}`$, according to the QCD predictions at large $`Q^2`$. This expansion can connects smoothly PQCD predictions in the high energy region with traditional nuclear physics predictions in the low energy region. It was shown that the second order contribution strongly affects the behavior of $`G_Q`$ in the intermediate energy region. At large $`Q^2`$, the ratio of form factors (1) is slightly modified.
Following this approach, $`G_{+0}^+`$ and $`G_+^+`$ will be expanded to higher orders (beyond the second order) in $`\mathrm{\Lambda }_{\mathrm{QCD}}/M`$ according to the PQCD prediction at large $`Q^2`$. In order to explore the role of higher order contributions we discuss the possibility to interpolate an expression for $`G_{+0}^+`$ and $`G_+^+`$ to the intermediate energy region in this paper. It is worthwhile to unify the predictions for the deuteron form factors from the low energy to large energy region.
A general consideration based on the PQCD predictions for the deuteron from factors is given in Sec. 2. As an example, a phenomenological analytic form including the higher order corrections in $`\mathrm{\Lambda }_{\mathrm{QCD}}/M`$ is suggested in Sec. 3. The numerical results and summary are presented in Sec. 4 and Sec. 5, respectively.
## 2 A General Consideration Based on the PQCD Predictions
For the deuteron case, the matrix element of the electromagnetic current $`J^\mu `$ is defined as
$$G_{\lambda ^{}\lambda }^\mu =P^{}\lambda ^{}J^\mu P\lambda ,$$
$`(2)`$
where $`Q^2=(P^{}P)^2`$ and $`|P\lambda `$ is an eigenstate of the deuteron with momentum $`P`$ and helicity $`\lambda `$. In the standard light-cone frame (LCF), defined by Ref. $`q^+=0,q_y=0`$, and $`q_x=Q`$, the charge, magnetic, and quadrapole form factors can be obtained from the plus component of three helicity matrix elements:
$$G_C=\frac{1}{2p^+(2\eta +1)}\left[(1\frac{2}{3}\eta )G_{00}^++\frac{8}{3}\sqrt{2\eta }G_{+0}^++\frac{2}{3}(2\eta 1)G_+^+\right],$$
$`(3a)`$
$$G_M=\frac{1}{2p^+(2\eta +1)}\left[2G_{00}^++\frac{2(2\eta 1)}{\sqrt{2\eta }}G_{+0}^+2G_+^+\right]$$
$`(3b)`$
and
$$G_Q=\frac{1}{2p^+(2\eta +1)}\left[G_{00}^++\sqrt{\frac{2}{\eta }}G_{+0}^+\frac{\eta +1}{\eta }G_+^+\right].$$
$`(3c)`$
In terms of $`G_C`$, $`G_M`$, and $`G_Q`$, the Rosenbluth cross section and the tensor polarization $`T_{20}`$ for elastic $`ed`$ scattering can be expressed as
$$\frac{d\sigma }{d\mathrm{\Omega }}=\left(\frac{d\sigma }{d\mathrm{\Omega }}\right)_{\mathrm{Mott}}\left[A(Q^2)+B(Q^2)\mathrm{tan}^2(\frac{\theta }{2})\right],$$
$`(4)`$
and
$$T_{20}=\frac{\frac{8}{9}\eta ^2G_Q^2+\frac{8}{3}\eta G_CG_Q+\frac{2}{3}\eta G_M^2\left[\frac{1}{2}+(1+\eta )\mathrm{tan}^2(\frac{\theta }{2})\right]}{\sqrt{2}\left[A+B\mathrm{tan}^2(\frac{\theta }{2})\right]},$$
$`(5)`$
where $`A(Q^2)`$ and $`B(Q^2)`$ are given by
$$A(Q^2)=G_C^2+\frac{2}{3}\eta G_M^2+\frac{8}{9}\eta ^2G_Q^2$$
$`(6)`$
and
$$B(Q^2)=\frac{4}{3}\eta (1+\eta )G_M^2.$$
$`(7)`$
It was shown <sup></sup> that, in LCF, the helicity-zero to zero matrix element $`G_{00}^+`$ would be the dominant helicity amplitude at large $`Q^2`$ for elastic $`ed`$ scattering from the PQCD predictions. It means the $`G_{00}^+`$ dominance in the structure function $`A(Q^2)`$. It was also argued<sup></sup> that the dominance of $`G_{00}^+`$ begins at $`Q^22M\mathrm{\Lambda }_{\mathrm{QCD}}0.8`$ GeV<sup>2</sup> but not $`\eta >>1`$. Thus $`2M\mathrm{\Lambda }_{\mathrm{QCD}}`$ is a scale of validity of PQCD predictions and the quark and gluon degrees of freedom in the deuteron should be taken into account to solve the problem that the experimental results of $`A(Q^2)`$ are in sharp disagreement with the meson exchange calculations for $`Q^2>0.8`$ GeV<sup>2</sup>. To make detailed prediction for $`G_{00}^+`$, we have suggested a QCD-inspired model<sup></sup> in the region of $`Q^2>1`$ GeV<sup>2</sup> based on the reduced form factor method<sup></sup>, which fit the data well.
PQCD predicts that $`G_{+0}^+`$ and $`G_+^+`$ are suppressed by factors $`\mathrm{\Lambda }_{\mathrm{QCD}}/Q`$ and $`\mathrm{\Lambda }_{\mathrm{QCD}}^2/Q^2`$, respectively. However, in the intermediate energy region, $`G_{00}^+`$ dominates the charge form factor $`G_C`$, but not $`G_M`$ and $`G_Q`$. As shown in Eq.(3), while $`\eta <\frac{1}{2}`$, the $`G_{+0}^+`$ contribution to $`G_M`$ and $`G_Q`$ are enhanced by a factor $`\frac{1}{\sqrt{2\eta }}`$ and $`G_+^+`$ contribution to $`G_Q`$ is enhanced by a factor $`\frac{1}{2\eta }`$. Although $`G_{+0}^+`$ and $`G_+^+`$ are suppressed for dynamic reason, they contribute significantly to $`G_M`$ and $`G_Q`$ because of the kinematic enhancement. Without these contributions, the predicted form factors, except for $`G_C`$, are in sharp disagreement with the data. In order to explore the role of $`G_{+0}^+`$ and $`G_+^+`$, we interpolate a general expression based on perturbative QCD predictions,
$$G_{+0}^+=\frac{1}{\sqrt{2\eta }}g_{+0}(\eta )G_{00}^+$$
$`(8a)`$
$$G_+^+=\frac{1}{2\eta }g_+(\eta )G_{00}^+,$$
$`(8b)`$
where $`g_{+0}(\eta )`$ and $`g_+(\eta )`$ are any functions of $`\eta `$ with $`\eta \frac{Q^2}{4M^2}`$. Obviously $`G_{+0}^+`$ and $`G_+^+`$ are suppressed by factors $`\mathrm{\Lambda }_{QCD}/Q`$ and $`\mathrm{\Lambda }_{QCD}^2/Q^2`$ as long as $`g_{+0}(\eta )`$ and $`g_+(\eta )`$ satisfy the following condition,
$$g_{+0}(\eta ),g_+(\eta )O(1)as\eta \mathrm{}.$$
$`(9)`$
Substituting Eqs.(8) into Eqs.(3), one can get
$$G_c=\frac{1}{2p^+(2\eta +1)}[(1\frac{2}{3}\eta )+\frac{8}{3}g_{+0}(\eta )+\frac{2}{3}(2\eta 1)\frac{1}{2\eta }g_+(\eta )]G_{00}^+,$$
$`(10a)`$
$$G_M=\frac{1}{2p^+(2\eta +1)}[2+\frac{2\eta 1}{\eta }g_{+0}(\eta )\frac{1}{\eta }g_+(\eta )]G_{00}^+$$
$`(10b)`$
and
$$G_Q=\frac{1}{2p^+(2\eta +1)}[1+\frac{1}{\eta }g_{+0}(\eta )\frac{\eta +1}{2\eta ^2}g_+(\eta )]G_{00}^+.$$
$`(10c)`$
Now we discuss the constriants on $`g_{+0}(\eta )`$ and $`g_+(\eta )`$ from the experimental data. As we know, $`G_M(Q^2)`$ changes sign at $`Q^2=Q_0^22GeV^{2[10]}`$ (or $`\eta =\eta _0=\frac{Q_0^2}{4M^2}0.13)`$. The dominance of $`G_{00}^+`$ can not explain this point (see Eq.(3b)) and at least the second term in Eq.(3b) should be in the same order of the first term to cancel it in order to fit data of $`G_M`$. That means $`g_{+0}(\eta )`$ is nonzero. If we first keep $`g_+(\eta )=0`$ in Eqs.(10), then $`g_{+0}(\eta _0)`$ can be determined by the zero at $`G_M(\eta _0)`$, which turns out to be $`g_{+0}(\eta _0)=\frac{2\eta _0}{12\eta _0}`$. In this case, $`G_Q`$ is negative at $`\eta =\eta _0`$ where PQCD begins to be valid. Thus there must be a node in the region $`Q^2<1GeV^2`$ since $`G_Q`$ is positive at the origin experimentally. The theoretical prediction is contrary to the experimental data without $`G_+^+`$ contribution. Therefore the predicted form factors are in sharp disagreement with the data without $`G_{+0}^+`$ and $`G_+^+`$ contributions and the existence of non-zero $`g_{+0}(\eta )`$ and $`g_+(\eta )`$ is necessary. In addition to the constrant (9), $`g_{+0}(\eta )`$ and $`g_+(\eta )`$ should astisfy
$$g_{+0}(\eta _0)=\frac{2\eta _0g_+(\eta _0)}{12\eta _0}$$
$`(11)`$
to ensure $`G_M(\eta _0)=0`$ and
$$2\eta g_{+0}(\eta )2\eta ^2<(\eta +1)g_+(\eta )$$
$`(12)`$
to keep $`G_Q`$ positive at any momentum transfer. In particular, as $`\eta =\eta _0`$ Eqs.(11) and (12) give the constraints on $`g_{+0}(\eta _0)`$ and $`g_+(\eta _0)`$,
$$g_+(\eta _0)>\frac{2\eta _0^2}{1\eta _0}$$
$`(13)`$
and
$$g_{+0}(\eta _0)<\frac{2\eta _0}{1\eta _0}$$
$`(14)`$
Eqs.(9), (11) and (12) are three constraints on the functions $`g_{+0}(\eta )`$ and $`g_+(\eta )`$.
## 3 A Phenomenological Example
As mentioned in the section 2, $`g_+(\eta )0`$ is important to keep $`G_Q>0`$ although $`G_+^+=\frac{1}{2\eta }g_+(\eta )G_{00}^+`$ is suppressed by the higher order factor $`\mathrm{\Lambda }_{QCD}^2/Q^2(=\frac{\mathrm{\Lambda }_{QCD}^2}{2M^2}\frac{1}{2\eta })`$. However the $`G_{+0}^+`$ is the first order correction which makes the zero in $`G_M`$ at $`Q_0^21.85GeV^2`$. We have expanded $`G_{+0}^+`$ and $`G_+^+`$ to the second order in $`\mathrm{\Lambda }_{QCD}/M`$ in Ref. and numerical results show that the second order contribution to $`G_{+0}^+`$ plays an important role in the intermediate energy region. In this paper we introduce an exponential form phenomenologically as an example,
$$g_{+0}(\eta )=fexp(\frac{bf}{\sqrt{2\eta }})$$
$`(15a)`$
and
$$g_+(\eta )=f^2exp(\frac{cf}{\sqrt{2\eta }})$$
$`(15b)`$
to interpolate the higher order corrections. Obviously, the exponential form (15) satisfies Eq.(9) and is consistent with the perturbative QCD prediction at the large transfer momentum region. Thus Eq.(15) is enable us to make an analytical evaluation to the deuteron form factors.
## 4 Numerical Results
We input the parameters $`b`$ and $`c`$, and determine $`f`$ by the zero in $`G_M`$. For a certain $`c`$, the obtained $`G_Q`$ increases proportionally with $`b`$. We can fix $`b`$ to connect our predictions with the data smoothly. To retain good convergence, we constrain $`b\frac{f}{\sqrt{2\eta }}`$ and $`c\frac{f}{\sqrt{2\eta }}`$ being smaller than unity. On the other hand, we restrict $`b`$ and $`c`$ to be positive to keep the exponential damping as $`Q^2`$ goes to infinity and it is a resonable assumption, after taking into account the higher order corrections. For $`c=0.0,0.5`$ and -0.5, the predicted $`B(Q^2)`$, $`G_Q(Q^2)`$ and $`T_{20}(Q^2)`$ are shown in figs. (1-3). Experimental data are taken from Refs.\[1,10-12\]. To smoothly connect with the data of $`G_Q`$, the parameter $`b`$ should be 1.1, 1.3, and 0.8, $`f=0.37,0.51`$, and 0.30, respectively. As argued above, $`c=0.5`$ should be abandaned. While $`c=0.5`$, the corresponding $`f`$ is 0.51, which is too large to keep $`b\frac{f}{\sqrt{2\eta }}`$ smaller than unity as $`Q^21`$. The parameter $`c=0.0,b=1.1,f=0.37`$ is an appropriate choice.
## 5 Summary
Based on perturbative QCD predictions at large momentum transfers we have tried to discuss the corrections to deuteron form factors beyond the second order in $`\mathrm{\Lambda }_{QCD}/M`$.A general consideration is given by introducing functions $`g_{+0}(\eta )`$ and $`g_+(\eta )`$ and the data at the present energy region put three constrants on the functions $`g_{+0}(\eta )`$ and $`g_+(\eta )`$. In order to explore the role of higher order contributions we suggest an exponential form for $`g_{+0}(\eta )`$ and $`g_+(\eta )`$ as an example. We conclude that (1) the helicity-zero to zero matrix element $`G_{00}^+`$ dominates the gross structure function $`A(Q^2)`$ in both of the large and intermediate energy regions. A QCD-inspired model can describe this matrix element well<sup></sup>. (2) $`G_{+0}^+`$ and $`G_+^+`$ contributions are important in determining form factors $`G_M`$ and $`G_Q`$; (3) By fitting the data, we get a set of parameters in $`g_{+0}(\eta )`$ and $`g_+(\eta )`$, $`c=0.0,b=1.1,f=0.37`$, which can describe $`G_M,G_Q`$ and $`T_{20}`$ appropriately. the parameters $`c=0.0`$ indicates that it is good for $`G_+^+`$ in the intermediate energy region to take the asymptotic behavior which were predicted by pQCD, and the higher order (beyond the second order) contributions are negligible; (4) The higher order corrections to $`G_{+0}^+`$ should be taken into account and they make sizeable contributions in the intermediate energy region.
## Figure Captions
Structure function $`B(Q^2)`$. The dashed dotted line corresponds to the Paris potential calculation. Experimental data are taken from Ref..
The form factor $`G_Q`$. Experimental data are taken from Ref..
The tensor polarization $`T_{20}`$ with scattering angle $`\theta =70^{}`$. The dashed dotted line corresponds to the calculation with Paris potential. Experimental data are taken from Ref.
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# On the determination of the last stable orbit for circular general relativistic binaries at the third post-Newtonian approximation
## I Introduction
The study of the late dynamical evolution of binaries made of compact objects (neutron stars or black holes) is important because such systems are the most promising candidate sources for interferometric gravitational-wave detectors such as LIGO and VIRGO. In particular, the global structure of the gravitational waveform emitted by such a binary sensitively depends on the frequency at which the system’s orbital evolution changes from a gravitational-radiation-driven inspiral phase to a plunge phase followed by coalescence .
In the test-mass limit $`(\mu M)`$ the orbital dynamics is that of a test particle (of mass $`\mu `$) moving in a Schwarzschild background (of mass $`M`$). A very important qualitative feature of circular orbits in such a background is the existence of a Last Stable Orbit (LSO) located at the (area) radius $`R_{\text{LSO}}=6GM/c^2`$. When considering the effect of gravitational-radiation reaction, one expects the test-particle motion to change abruptly near $`R_{\mathrm{LSO}}`$ from a slow inspiral to a fast plunge. By analogy, one believes that the motion of a compact binary made of comparable masses, $`m_1`$ and $`m_2`$, will exhibit a similar transition from inspiral to plunge, with the location of the transition being mainly determined by the existence of a Last Stable Orbit in the conservative part of the two-body Hamiltonian.
Several authors have tried to estimate the location of the LSO in comparable-mass (compact) binaries. An early analytical estimate was made by Clark and Eardley using the first post-Newtonian (1PN) Hamiltonian. Some authors tried to use initial value formalisms to locate the LSO. However, the initial value approaches used in these works assume a conformally flat metric, and therefore do not correctly incorporate the well known second post-Newtonian (2PN) dynamics . In this work we shall take the view that a correct incorporation of all 2PN effects is a necessary (if maybe not sufficient) prerequisite for an accurate determination of the location of the LSO. Previous treatments that used the full 2PN dynamics to try to analytically determine the location of the LSO include Refs.. The aim of the present work is to extend these PN-based analytic determinations of the LSO to the third post-Newtonian (3PN) level. The conservative part of the 3PN Hamiltonian for two point masses has been obtained in 1998 by Jaranowski and Schäfer , though with some remaining ambiguity due to the need to regularize the divergent integrals entailed by the use of point-like sources. The determination of the 3PN dynamics has been recently completed by deriving the Hamiltonian in a non-mass-centered frame, and by fixing a certain momentum-dependent regularization ambiguity . Recently, we have extracted from the 3PN Hamiltonian of Ref. all its dynamical invariants, i.e. all the functions linking dynamical quantities which do not depend on the choice of coordinates in phase-space . In this paper, we shall use several of these 3PN invariants to determine the location of the LSO. Some of the methods of LSO determination that we shall use below generalize previous works , but others are new.
## II Methods directly based on dynamical invariants
Before embarking on the discussion of the methods we shall use here to extract the LSO from PN expansions, let us stress that the basic theme underlying our endeavours is the following: Our problem is to extract some semi-non-perturbative information from (badly convergent) perturbation expansions. We shall do that by using several types of “resummation methods”. The basic idea of all resummation methods is simply the following: to complete the information contained in the first few terms of a perturbative expansion $`f(z)=c_0+c_1z+\mathrm{}+c_nz^n+𝒪(z^{n+1})`$ by injecting some non-perturbative information about the global behaviour (if possible in the complex plane) of the exact function $`f(z)`$. The amount of global information one has about the function $`f(z)`$ determines the best type of resummation method to use. For instance, if the only information at our disposal is that the function $`f(z)`$ is (probably) meromorphic in the complex $`z`$-plane, then the best, all-purpose method is to use Padé approximants. If we knew more about the location of the singularities of $`f(z)`$ in the complex plane one might contemplate to use other methods (e.g. change of independent variable, Borel transform, $`\mathrm{}`$).
In this paper, we shall use two distinct classes of methods for extracting the (invariant) location of the LSO from the knowledge of the PN-expanded dynamics. The first class of methods (which was introduced in Ref. ) is discussed in the present Section. The second class will be discussed in the next Section. The first class of methods is based on the combination of three ideas:
(i) to work only with invariant functions;
(ii) to make a maximal use of the known, exact functional form of invariant functions in the test-mass limit $`(\nu m_1m_2/(m_1+m_2)^20`$) and to assume some structural stability when the parameter $`\nu `$ is turned on;
(iii) to use Padé approximants to represent the invariant functions which are (because of (ii)) expected to be meromorphic functions of their argument.
This method was applied in Ref. to the two invariant functions which play an essential role in the gravitational-damping-driven inspiral of a binary system: the binding energy $`E(x)`$ of a circular orbit, and the gravitational-wave flux $`F(x)`$ emitted by a circular orbit, both being considered as functions of the (dimensionless) invariant parameter
$$x\left(\frac{GM\omega }{c^3}\right)^{2/3},$$
(1)
where $`Mm_1+m_2`$ denotes the total mass of the binary, and $`\omega `$ the orbital angular frequency along a circular orbit.
Before introducing some variations on this method, let us motivate the interest of using the three ideas (i)–(iii) above. First, we recall that the PN expansions of non-invariant, i.e. gauge-dependent, functions can have (and do have, in some gauges) worse convergence properties than the PN expansions of invariant functions. For instance, the PN expansion of the gravitational-wave flux, from a test mass in circular orbit, say $`F_{\mathrm{TM}}`$, considered as a function of the harmonic-coordinate parameter $`\gamma GM/(c^2r_{\mathrm{harmonic}})`$, is such that the 1PN “correction” to the leading “quadrupole” result becomes fractionally larger than 100% (while being negative!) for a radius $`r_{\mathrm{harmonic}}`$ larger than the LSO (which is at $`r_{\mathrm{LSO}}=5GM/c^2`$ in harmonic coordinates). More precisely, if one formally writes down the expansion of $`F_{\mathrm{TM}}(\gamma )`$, in powers of $`\gamma `$, near the LSO (i.e. for $`\gamma `$ near $`1/5`$) one gets a series numerically proportional to: $`11.74(5\gamma )+1.12(5\gamma )^{3/2}+1.29(5\gamma )^2+\mathrm{}`$, i.e. a series whose first terms do not exhibit any convergence near the LSO. By contrast, the flux $`F_{\mathrm{TM}}`$ expanded in terms of the invariant parameter $`x`$, Eq. (1), ($`x_{\mathrm{LSO}}=1/6`$ in the test-mass limit) has a more reasonable expansion proportional to $`10.619(6x)+0.855(6x)^{3/2}0.137(6x)^2+\mathrm{}`$. Still, it is clear that one needs some resummation technique for summing a series as slowly converging as $`F_{\mathrm{TM}}(x)`$. It was shown in great detail in Ref. , by making use both of the known analytical results on high-order (5.5PN) terms in the post-Newtonian expansion of the test-mass flux function $`F_{\mathrm{TM}}(x)`$ , and of the existence of a pole-like blow up of $`F_{\mathrm{TM}}(x)`$ at $`x=x_{\text{light ring}}=1/3`$, that one could considerably speed up the convergence of the straightforward (Taylor-like) PN series, by replacing it by a suitably defined sequence of Padé approximants (see, notably, Fig. 3 in ).
To further illustrate the idea (ii), and the need for an acceleration of convergence, let us recall the treatment of the energy function $`E(x)`$ introduced in . In this paper we follow in defining the dimensionless energy function
$$E\frac{^RMc^2}{\mu c^2},$$
(2)
where $`^R`$ is the total (“relativistic”) energy of the binary system (including the rest-mass contribution), and where
$$Mm_1+m_2,\mu \frac{m_1m_2}{m_1+m_2},\nu \frac{\mu }{M}=\frac{m_1m_2}{(m_1+m_2)^2}.$$
(3)
For notational simplicity we shall henceforth use units such that $`c=1`$. Note that the energy function used in was $`E^{\mathrm{DIS}}(^RM)/M=\nu E^{\mathrm{here}}`$.
Because the argument $`x`$ is, from its definition (1), of formal order $`𝒪(c^2)`$, the knowledge of the dynamics at, say, the $`n`$th post-Newtonian ($`n`$PN) order entails the knowledge of the expansion of the ratio $`E(x)/x`$ up to $`x^n`$:
$$E(x;\nu )=\frac{1}{2}x\left[1+E_1(\nu )x+E_2(\nu )x^2+\mathrm{}+E_n(\nu )x^n+𝒪(x^{n+1})\right].$$
(4)
\[The term $`\frac{1}{2}x`$ corresponds to the Newtonian binding energy $`^{\mathrm{NR}}=\frac{1}{2}\mu v^2`$. Then the term $`E_1(\nu )x`$, for instance, represents the fractional 1PN effects, etc.\] The symmetric mass ratio $`\nu `$, Eq. (3), enters the expansion coefficients $`E_n(\nu )`$ as a parameter. At present, only $`E_1(\nu )`$, $`E_2(\nu )`$ and $`E_3(\nu )`$ are known (with some ambiguity for the 3PN coefficient $`E_3(\nu )`$). They were written down in and will be repeated below. On the other hand, in the test-mass limit $`\nu 0`$, the coefficients $`E_n(0)`$ are known (in principle) for any $`n`$. What is more, the exact expression of $`E(x;\nu =0)`$ is known:
$$E(x;\nu =0)=\frac{12x}{\sqrt{13x}}1.$$
(5)
It is known (see, e.g., ) that the location of the LSO corresponds just to the minimum of the function $`E(x)`$ ($`(dE(x)/dx)_{x_{\text{LSO}}}=0`$). Therefore, it would seem that the most straightforward way of locating the LSO is to consider the successive “Taylor approximants” of $`E(x)`$, say $`E_{T_n}(x)`$ (defined for each integer $`n`$ as the R.H.S. of Eq. (4), without the $`𝒪(x^{n+1})`$ error term), and to solve the equations $`dE_{T_n}(x)/dx=0`$. Let us see what this gives in the test-mass limit where the Taylor expansion of Eq. (5) is known:
$$E(x;\nu =0)=\frac{1}{2}x\left(1\frac{3}{4}x\frac{27}{8}x^2\frac{675}{64}x^3\frac{3969}{128}x^4\mathrm{}\right).$$
(6)
The successive “Taylor” estimates of $`x_{\mathrm{LSO}}`$, or better $`x_{\mathrm{LSO}}^{T_n}/x_{\mathrm{LSO}}^{\mathrm{exact}}6x_{\mathrm{LSO}}^{T_n}`$, are found to be (in the test-mass case):
$$6x^{T_1}=4;6x^{T_2}=1.49284;6x^{T_3}=1.17565;6x^{T_4}=1.07680.$$
(7)
\[To avoid confusion note that we use here the convention that ‘$`T_n`$’ corresponds to the $`n`$PN approximation, i.e. a $`(v/c)^{2n}`$-accurate result, while in Ref. $`T_n`$’ referred to $`(v/c)^n`$-accuracy, i.e. to the $`\frac{n}{2}`$PN approximation.\] From Eq. (1) the corresponding values of the orbital frequency at the LSO, scaled to the exact value, $`\widehat{\omega }\omega /\omega ^{\mathrm{exact}}=(6x)^{3/2}`$, read
$$\widehat{\omega }^{T_1}=8;\widehat{\omega }^{T_2}=1.82398;\widehat{\omega }^{T_3}=1.27472;\widehat{\omega }^{T_4}=1.11739.$$
(8)
As the value of the frequency at the LSO is the most important observable one wishes to know (for data analysis purposes), one should reject any method which does not have the prospect of determining it to, say, better than about 10%. The test-mass results<sup>*</sup><sup>*</sup>*The logic here is to use the known convergence properties of the $`\nu =0`$ limit to estimate the convergence when $`\nu 0`$. It is, indeed, unlikely that turning on $`\nu `$ will drastically improve the convergence properties. (8) suggest that, if the dynamics is known only up to the 3PN level, straightforward Taylor approximants of the energy function $`E(x)`$ do not converge fast enough to determine satisfactorily the location of the LSO. \[The 4PN level might barely suffice, but it seems anyway excluded that one will be able to analytically derive the 4PN dynamics.\]
This preliminary discussion motivates the necessity of boosting up the convergence of the series (4) or (6). This is here that the ideas (ii) and (iii) enter. The exact test-mass result (5) suggests (under the assumption that the $`\nu 0`$ case represents a structurally stable deformation of the $`\nu =0`$ limit) that the function $`E(x;\nu )`$, for $`\nu 0`$, will have a branch-cut singularity (at some point $`x_0=\frac{1}{3}+𝒪(\nu )`$) in the complex $`x`$ plane. It is, a priori, much better to work with functions which are meromorphic in the complex plane, because we can then make use of Padé approximants, which are efficient tools for accurately representing meromorphic functions. This suggests to work with some (to be defined) new invariant energy function that is known to be meromorphic in the test-mass limit.
Looking at Eq. (5) one is tempted to consider the square of the function $`1+E`$. However, Ref. remarked that the usual test-mass limit definition of this function, namely ($`m_2m_1`$ being the test-mass orbiting $`m_1`$)
$$1+E\frac{_2^R}{m_2}$$
(9)
with
$$_{\mathrm{tot}}^R=m_1+_2^R+𝒪(m_2^2)=m_1\frac{p_1^\mu p_{2\mu }}{m_1}+𝒪(m_2^2),$$
(10)
looks unnaturally asymmetric with respect to the labels 1 and 2. \[Essentially, this asymmetry comes from the fact, hidden in the symmetric definition (2), that $`_2^R`$ represents the sum of the rest-mass energy of $`m_2`$ alone and of the binding energy of the binary system.\] Ref. therefore suggested to introduce the more symmetric function
$$\phi (s)\frac{sm_1^2m_2^2}{2m_1m_2}\frac{(^R)^2m_1^2m_2^2}{2m_1m_2}$$
(11)
of the Mandelstam invariant $`s(p_1^\mu +p_2^\mu )^2(^R)^2`$. Indeed, in the test-mass limit
$$\phi (s)\frac{p_1^\mu p_{2\mu }}{m_1m_2}\frac{_2^R}{m_2}=1+E=\frac{12x}{\sqrt{13x}}.$$
(12)
It is then natural to define the function $`e(x)`$ (or rather $`1+e(x)`$), by setting
$$1+e(x)\mathbf{\left(}\phi (s(x))\mathbf{\right)}^2=\left(\frac{(^R)^2m_1^2m_2^2}{2m_1m_2}\right)^2.$$
(13)
Then, $`1+e(x)`$, being equal to $`(12x)^2/(13x)`$ in the test-mass limit, i.e. (after subtraction of the trivial constant 1)
$$e(x;\nu =0)=x\frac{14x}{13x},$$
(14)
one expects the function $`e(x;\nu )`$ to be meromorphic in $`x`$ when $`\nu 0`$ (with a pole located at some $`x_{\mathrm{pole}}=\frac{1}{3}+𝒪(\nu )`$).
This finally leads to the following “$`P`$-approximant”-improved method for locating the LSO: starting from the Taylor approximants ($`T`$-approximants) of the original $`E(x)`$ function, Eq. (4), compute first the corresponding $`T`$-approximants of the new $`e(x)`$ function, say
$$e(x)=x\left[1+e_1(\nu )x+e_2(\nu )x^2+e_3(\nu )x^3+𝒪(x^4)\right].$$
(15)
Then construct a sequence of Padés of the Taylor-expanded $`e(x)`$:
$$e_{P_n}(x)P_{\mathrm{}}^k\left[T_n[e(x)]\right],$$
(16)
where $`k+\mathrm{}=n`$. Here $`k`$ and $`\mathrm{}`$ are the degrees of the polynomials $`N_k(x)`$ and $`D_{\mathrm{}}(x)`$ entering the Padé $`P_{\mathrm{}}^k(x)=N_k(x)/D_{\mathrm{}}(x)`$. It is known that, generically, the Padé improvements are best when one is near the “diagonal”, i.e. when $`|k\mathrm{}|`$ is “small” compared to $`k`$ and $`\mathrm{}`$. When dealing with a function $`f(x)`$ that is expected to have a pole at some $`x_00`$, one imposes the constraint $`\mathrm{}>0`$. At 1PN order this uniquely fixes the values of $`k`$ and $`\mathrm{}`$, namely $`k=0`$ and $`\mathrm{}=1`$. At 2PN order the Padé closest to the diagonal is that with $`k=1`$ and $`\mathrm{}=1`$. At 3PN order there are two possible Padés near the diagonal, namely $`k=1`$ and $`\mathrm{}=2`$, or $`k=2`$ and $`\mathrm{}=1`$. In this work we shall use the latter one ($`P_1^2`$) because we found it to be more robust under variations of the coefficients of the Taylor expansion of the Padéed function. (One aspect of this robustness is that the existence of a real pole in this Padé is always ensured, while this is not the case for the other 3PN possibility: $`P_2^1`$.)
Note also that Padés are originally defined only for series of the regular type $`\sigma (x)=c_0+c_1x+\mathrm{}`$ with $`c_00`$. When dealing with a function of the type $`f_p(x)=x^p\sigma (x)`$ with some relative integer $`p`$, we shall, by convention, define any Padé of $`f_p(x)`$ as being $`(k+\mathrm{}=n)`$
$$P_{\mathrm{}}^k[T_n[f_p(x)]]x^pP_{\mathrm{}}^k[T_n[x^pf_p(x)]].$$
(17)
The “$`e_{P_n}`$-estimate” of the location of the LSO is then defined as the value $`x_{P_n}`$ at which $`e_{P_n}(x)`$ reaches a minimum. \[It is easily seen that $`e(x)`$ follows the variations of $`E(x)`$, and in particular that it reaches a minimum at the same place as $`E(x)`$.\]
We spent some time explaining in detail on one example our general Padé-improved–test-mass-limit-motivated approach because we are going to extend it to several other invariant functions. In this (and the next) section, we present our various methodologies. The 3PN results obtained by them will be presented in a later section.
The introduction of the function $`e(x)`$ has two defects. First, it is not unique because we do not know for sure whether the function $`(\phi (s))^2`$, Eq. (13), is a “better” invariant than $`(1+E)^2=(1+(\sqrt{s}M)/\mu )^2`$. Second, the Padeing of $`e(x)`$ starts giving meaningful results only at the 2PN level. Indeed the 1PN expansion of $`e(x)`$, in the test-mass limit, $`e(x;\nu =0)=x(1x+𝒪(x^2))`$, yields a Padé $`e_{P_1}(x;\nu =0)=x/(1+x)`$ which contains no pole on the positive real axis, and which formally predicts an LSO (minimum of $`e_{P_1}(x)`$) located at $`x=+\mathrm{}`$.
In this paper, we propose to consider another invariant function, which is more uniquely defined, and which gives sensible results already at the 1PN level. Let us consider the reduced angular momentum
$$j\frac{𝒥}{\mu GM}=\frac{𝒥}{Gm_1m_2},$$
(18)
where $`𝒥`$ denotes the total angular momentum of the system. In the test-mass limit the invariant function giving the (dimensionless) quantity $`j`$ in terms of the quantity $`x`$, Eq. (1), reads
$$j(x;\nu =0)=\frac{1}{\sqrt{x(13x)}}.$$
(19)
This motivates the consideration of the squared (reduced) angular momentum $`j^2(x)`$ which is expected, when $`\nu 0`$, to be a meromorphic function of $`x`$, with a pole at the “light ring” $`x_{\mathrm{pole}}=\frac{1}{3}+𝒪(\nu )`$. Therefore we propose to work with the Padeed form of $`j^2`$:
$$j_{P_n}^2(x;\nu )P_{\mathrm{}}^k[T_n[j^2(x;\nu )]],$$
(20)
with $`k+\mathrm{}=n`$, and the choice of $`\mathrm{}>0`$ discussed above. (We use in Eq. (20) the convention (17), i.e. the factor $`x^1`$ in $`j^2(x)`$ is factored before taking a Padé.) Note that, in the test-mass limit, if we knew only the 1PN approximation to the function $`j^2(x)`$, i.e. $`j^2(x;\nu =0)=x^1\mathbf{\left(}1+3x+𝒪(x^2)\mathbf{\right)}`$, the procedure (20) would reconstruct the exact result: $`j_{P_1}^2(x;\nu =0)=P_1^0[T_1[j^2(x;\nu =0)]]=[x(13x)]^1`$.
It is important to note that the perturbative information contained in the PN expansion of the function $`j(x)`$ (or equivalently $`j^2(x)`$) is totally equivalent (at any PN accuracy) to the information contained in, either the original energy function $`E(x)`$, or the new one $`e(x)`$, Eq. (13). Indeed, the generic Hamiltonian equation $`\dot{\theta }_i=\omega _i=H/I_i`$ in action-angle variables $`(I_i,\theta _i)`$ yields
$$\omega _{\mathrm{circular}}=\frac{d^R}{d𝒥}=\frac{1}{GM}\frac{dE}{dj},$$
(21)
which implies the identity
$$\frac{dE(x)}{dx}=x^{3/2}\frac{dj(x)}{dx}.$$
(22)
The identity (22) proves the assertion just made about the identical information content in $`E(x)`$ and $`j(x)`$. It also proves several interesting facts. First, the location of the minimum of (the exact) $`j(x)`$ coincides with that of the minimum of (the exact) $`E(x)`$ (both of them equivalently defining the LSO). Second, the existence of a branch cut singularity $`(x_0x)^{1/2}`$ in either $`j(x)`$ or $`E(x)`$ necessarily implies the presence of a similar singularity $`(x_0x)^{1/2}`$ (at the same location $`x_0`$) in the other function ($`E(x)`$ or $`j(x)`$, respectively). This can be viewed as a confirmation of our generic assumption of structural stability. However, this argument also shows the ambiguities present when trying to work with the energy function. Indeed, if we assume that, near $`x_0`$, $`j(x)`$ can be expanded as $`\phi (x)(x_0x)^{1/2}`$, where $`\phi (x)`$ is a smooth function, one finds from Eq. (22) that $`E(x)=\psi (x)(x_0x)^{1/2}+c`$, with some unknown constant $`c`$. The lack of knowledge of the constant $`c`$ (which a priori depends on the parameter $`\nu `$) implies that we do not know which function $`(Ec(\nu ))`$ we would square to get a meromorphic function with a simple pole. The same reasoning shows another defect of the proposal to consider the (new) function $`e(x)`$. Indeed, when $`E(x)\mathrm{}`$ the leading term in $`e(x)`$, Eq. (13), will be $`e(x)\frac{1}{4}\nu ^2E^4(x)`$, which will have a double pole $`\nu ^2(x_0x)^2`$, if $`j^2(x)`$ has a simple pole $`(x_0x)^1`$. For all these reasons, we consider that the “$`j`$-method”, Eq. (20), appears as the best way of locating the LSO, within the class of methods dealt with in this section.
To conclude this section, let us, however, mention another invariant function one might wish to consider. This function is the fourth power of the dimensionless periastron parameter $`1+k=\mathrm{\Phi }/(2\pi )`$ considered as a function of the reduced angular momentum. Indeed, in the test-mass limit, and for circular orbits, one knows that
$$(1+k)^4=\left(1\frac{12}{j^2}\right)^1.$$
(23)
We recall that $`j^2=12`$ is the location of the LSO. Therefore, if we define $`K(1+k)^4=\mathbf{\left(}\mathrm{\Phi }/(2\pi )\mathbf{\right)}^4`$ and $`y1/j^2`$, we might consider
$$K_{P_n}(y;\nu )P_{\mathrm{}}^k[T_n[K(y;\nu )]],$$
(24)
with $`k+\mathrm{}=n`$, and our canonical choice of $`\mathrm{}>0`$. Then, we can take the pole of $`K_{P_n}(y;\nu )`$ as estimate of the value of $`1/j^2`$ at the LSO when $`\nu 0`$. We consider, however, that the $`j`$-method (or any energy method for that matter) has a better chance of accurately locating the LSO because it incorporates not only the information that something special (a minimum) takes place at the LSO, but also the information that the location of this minimum is a corollary of the presence of a blow up ($`j^2\mathrm{}`$, $`E\mathrm{}`$) at a point, $`x_{\mathrm{pole}}`$, further away on the real axis. Therefore, even if the location of $`x_{\mathrm{pole}}`$ is not known very accurately, one can hope that $`x_{\mathrm{LSO}}`$ will be more robustly determined.
## III Effective one-body method
In this section, we shall turn to a rather different method, though one which also makes use of the three ideas (i)–(iii) listed at the beginning of the previous section. This new method incorporates a fourth idea, which has been recently put forward by Buonanno and Damour . This fourth idea consists in mapping (through the use of invariant functions) the real two-body problem we are interested in (two masses $`m_1`$, $`m_2`$ orbiting around each other) onto an “effective one-body problem” (one mass $`m_0`$ moving in some background metric, $`g_{\alpha \beta }^{\mathrm{effective}}(x^\gamma )`$). At the 2PN level, Ref. has shown the possibility of mapping the real two-body problem onto geodesic motion in some spherically symmetric metric $`g_{\alpha \beta }^{\mathrm{effective}}(x^\gamma ;\nu )`$. It was found that, when $`\nu 0`$, $`g_{\alpha \beta }^{\mathrm{effective}}(x^\gamma ;\nu )`$ is a smooth deformation of the Schwarzschild metric. The “mapping rules” between the two problems were motivated by quantum considerations:
(i) the adiabatic invariants $`I_i=p_i𝑑q_i`$ (which are quantized in units of $`\mathrm{}`$) were identified in the two problems, and
(ii) the energies are mapped through a function $`_{\mathrm{effective}}=f[_{\mathrm{real}}]`$ which is, a priori, arbitrary, and which is determined in the process of matching the two problems.
In other words, the idea is to determine a metric $`g_{\alpha \beta }^{\text{effective}}`$ such that the “energy levels” $`_{\text{effective}}[I_i]`$ (i.e. the Hamiltonian expressed in action variables, or “Delaunay Hamiltonian”) of the bound states of a particle moving in $`g_{\alpha \beta }^{\text{effective}}`$, are in one-to-one correspondence with the two-body bound states:
$$_{\mathrm{effective}}[I_i]=f[_{\mathrm{real}}[I_i]].$$
(25)
The unknowns of the problem are the numerical coefficients entering a spherically symmetric metric
$$ds_{\mathrm{eff}}^2=A(R_{\mathrm{eff}})dt_{\mathrm{eff}}^2+\frac{D(R_{\mathrm{eff}})}{A(R_{\mathrm{eff}})}dR_{\mathrm{eff}}^2+R_{\mathrm{eff}}^2(d\theta _{\mathrm{eff}}^2+\mathrm{sin}^2\theta _{\mathrm{eff}}d\phi _{\mathrm{eff}}^2),$$
(26)
namely,
$`A(R)`$ $`=`$ $`1+a_1{\displaystyle \frac{GM_0}{R}}+a_2\left({\displaystyle \frac{GM_0}{R}}\right)^2+a_3\left({\displaystyle \frac{GM_0}{R}}\right)^3+a_4\left({\displaystyle \frac{GM_0}{R}}\right)^4+\mathrm{},`$ (28)
$`D(R)`$ $`=`$ $`1+d_1{\displaystyle \frac{GM_0}{R}}+d_2\left({\displaystyle \frac{GM_0}{R}}\right)^2+d_3\left({\displaystyle \frac{GM_0}{R}}\right)^3+\mathrm{},`$ (29)
and the coefficients entering the energy-map $`f`$ (here written for the “non-relativistic” energies $`_{\mathrm{eff}}^{\mathrm{NR}}_{\mathrm{eff}}m_0`$, $`_{\mathrm{real}}^{\mathrm{NR}}_{\mathrm{real}}M`$):
$$\frac{_{\mathrm{eff}}^{\mathrm{NR}}}{m_0}=\frac{_{\mathrm{real}}^{\mathrm{NR}}}{\mu }\left[1+\alpha _1\frac{_{\mathrm{real}}^{\mathrm{NR}}}{\mu }+\alpha _2\left(\frac{_{\mathrm{real}}^{\mathrm{NR}}}{\mu }\right)^2+\alpha _3\left(\frac{_{\mathrm{real}}^{\mathrm{NR}}}{\mu }\right)^3+\mathrm{}\right].$$
(30)
As discussed in Ref. it is natural to require that the effective mass $`m_0`$ be exactly equal to the usual non-relativistic effective mass $`\mu m_1m_2/(m_1+m_2)`$. One can also (by convention) choose the mass $`M_0`$ entering the effective metric to be $`M_0M=m_1+m_2`$. With these choices the Newtonian limit tells us that the first coefficient in $`A(R)=g_{00}`$ is simply $`a_1=2`$. The 1PN level then involves the coefficients $`a_2`$, $`d_1`$, and $`\alpha _1`$, while the 2PN and 3PN levels involve $`(a_3,d_2,\alpha _2)`$ and $`(a_4,d_3,\alpha _3)`$, respectively. In other words, at each PN level, we only have three arbitrary coefficients to play with. This seems to be quite a small number of degrees of freedom, compared to the many possible coefficients which can enter the PN-expansion of the Delaunay Hamiltonian. \[The Delaunay Hamiltonian was determined at the 2PN level in Ref. , and at the 3PN level in Ref. .\] In order to clarify the number of independent equations to be satisfied when mapping the real problem onto the effective one, let us consider a genericWe use the information that the leading kinetic terms in a PN Hamiltonian are given by the expansion of the free Hamiltonian $`\sqrt{𝐩_1^2+m_1^2}+\sqrt{𝐩_2^2+m_2^2}`$, so that, at the $`n`$PN level, they are $`𝐩^{2(1+n)}`$ without dependence on $`𝐧𝐪/q`$. PN-expanded Hamiltonian, with the symbolic structure
$$\widehat{H}_{n\mathrm{PN}}^{\mathrm{NR}}(𝐪,𝐩)=p^{2(n+1)}+\frac{1}{q}\left[p^{2n}+p^{2n2}(np)^2+\mathrm{}+(np)^{2n}\right]+\frac{1}{q^2}\left[p^{2(n1)}+\mathrm{}+(np)^{2(n1)}\right]+\mathrm{}+\frac{1}{q^{n+1}}.$$
(31)
Here, we consider the reduced Hamiltonian $`\widehat{H}^{\mathrm{NR}}H^{\mathrm{NR}}/\mu `$, in the center-of-mass frame, as a function of the reduced canonical variables $`𝐩𝐩_1/\mu =𝐩_2/\mu `$, $`𝐪(𝐱_1𝐱_2)/(GM)`$ ($`(np)`$ denotes $`𝐧𝐩`$ with $`𝐧𝐪/q`$). See Ref. for details. Note that we follow here in denoting by $`𝐪`$, $`𝐩`$ the original (ADM-like) coordinates. We have suppressed all coefficients in Eq. (31) to display the structure. What is important for our present purpose is the total number of coefficients in the $`n`$PN-level Hamiltonian (31). This is easily checked to be:
$$C_H(n)=\frac{(n+1)(n+2)}{2}+1.$$
(32)
As explained in Ref. , one way (and the only explicit one) to map the real Hamiltonian (31) onto an effective Hamiltonian, while keeping the action variables invariant, is to apply a canonical transformation, with generating function $`\stackrel{~}{G}(q,p^{})=q^ip_i^{}+G(q,p^{})`$. The most generic generating function that we need to consider has the symbolic structure (at the $`n`$PN level)
$$G_{n\mathrm{PN}}(𝐪,𝐩)=(𝐪𝐩)\left\{p^{2n}+\frac{1}{q}\left[p^{2(n1)}+\mathrm{}+(np)^{2(n1)}\right]+\mathrm{}+\frac{1}{q^n}\right\}.$$
(33)
Correspondingly to the pure $`𝐩`$-dependence of the leading kinetic term in (31) we have written here the leading term in $`G_{n\mathrm{PN}}`$ as $`(qp)p^{2n}`$. \[It is easily shown that any term of the form $`(qp)p^{2(nk)}(np)^{2k}`$ must have a vanishing coefficient.\] The number of arbitrary coefficients in the $`n`$PN-level generating function (33) is easily seen to be
$$C_G(n)=C_H(n1)=\frac{n(n+1)}{2}+1.$$
(34)
Finally, the difference $`\mathrm{\Delta }(n)`$ between the number of equations to satisfy, and the number of unknowns (including the 3 basic parameters ($`a_{n+1}`$, $`d_n`$, $`\alpha _n`$) appearing in the effective metric and the energy-map) reads, at the $`n`$PN level
$$\mathrm{\Delta }(n)=C_H(n)C_G(n)3=n2.$$
(35)
In particular: $`\mathrm{\Delta }(1)=1`$, which means that requiring a 1PN matching leaves one degree of freedom unrestricted. \[This freedom was used in to impose the natural condition $`d_1=0`$, i.e. that the linearized effective metric coincides with Schwarzschild.\] Then $`\mathrm{\Delta }(2)=0`$, which means that there will (barring any degeneracy) be a unique solution at the 2PN level. \[This was indeed the result of .\] But $`\mathrm{\Delta }(3)=+1`$, which means that, at the 3PN level, there is one more equation to satisfy than the number of free parameters. Then the situation would get worse and worse at higher PN levels.
By explicitly performing the matching between the canonically-transformed Hamiltonian and (modulo the energy map (30)) the effective Hamiltonian of a point particle moving in some $`g_{\mu \nu }^{\mathrm{eff}}`$, we have established (details will be given below) that, if we follow Ref. in imposing the natural condition $`d_1=0`$ at the 1PN level, there are, indeed, $`C_H(3)=11`$ linearly independent equations, for $`C_G(3)+3=10`$ unknowns, to be satisfied at the 3PN level. \[No miracle occurred!\]
At face value, this looks like bad news for the idea of the “effective one-body approach”. However, we think that there are acceptable ways to rescue this approach. A first cure would be to take advantage of the fact that the total number of equations at the three first PN levels (1, 2, and 3) is exactly equal to the number of unknowns (in other words $`\mathrm{\Delta }(1)+\mathrm{\Delta }(2)+\mathrm{\Delta }(3)=1+0+1=0`$). Therefore if we relax the (not really necessary) constraint $`d_1=0`$, there will be a unique 3PN-accurate effective metric $`g_{\mu \nu }^{\mathrm{eff}}`$ (and a unique energy mapping (30)) satisfying the necessary constraints. We have verified that this is true and, for completeness, we give this unique solution in Appendix A. But, we do not want to take this solution too seriously for the following reasons: (i) it does not look natural to have to wait to know the 3PN Hamiltonian to determine the 1PN and 2PN effective metrics; (ii) this solution looks more complicated than the 3PN Hamiltonian itself; and (iii) this trick is not expected to be sufficient to ensure the existence of solutions at higher PN levels. (Indeed, $`\mathrm{\Delta }(n)=n2`$ continues to increase.)
We propose therefore to consider a second (more radical, and simpler) way to cure the problem. Indeed, we have to face the fact that there is (probably) nothing deep in the effective-one-body approach. After all, it is just a somewhat artificial way of mapping the complicated two-body dynamics on a simpler one-body dynamics. There is no reason to assume that the one-body dynamics can, to all orders, be considered as equivalent to a simple geodesic motion. Let us recall that, in quantum mechanicsWe recall that the origin of the effective-one-body approach lies in the quantum electrodynamics of two-charge systems, see ., geodesic motion means a simple Klein-Gordon Lagrangian
$$_{\mathrm{eff}}=\sqrt{g_{\mathrm{eff}}}\left((\phi )^2+m_0^2\phi ^2\right)\text{with}(\phi )^2=g_{\mathrm{eff}}^{\alpha \beta }_\alpha \phi _\beta \phi .$$
(36)
It is well-known that effective actions generally include, at higher orders, some higher-derivative terms: for instance of the type $`(\mathrm{}_g\phi )^2`$, $`(\phi )^4`$, etc. Coming back to the classical limit, i.e. to the Hamilton-Jacobi equation (obtained by considering that $`\phi (x)=\mathrm{exp}(iS(x)/\mathrm{})`$ with $`\mathrm{}0`$), we should correspondingly expect that, at higher orders, the effective one-body “Hamilton-Jacobi” equation be of the generalized form
$$0=m_0^2+g_{\mathrm{eff}}^{\alpha \beta }(x)p_\alpha p_\beta +A^{\alpha \beta \gamma \delta }(x)p_\alpha p_\beta p_\gamma p_\delta +\mathrm{}$$
(37)
with $`p_\alpha =S(x)/x^\alpha `$. If we were to use a Lagrangian formulation, the general form (37) would correspond to an action $`S=m_0𝑑s_{\mathrm{eff}}\left[1+A_{\alpha \beta \gamma \delta }(x)u^\alpha u^\beta u^\gamma u^\delta +\mathrm{}\right]`$ with $`u^\alpha =dx^\alpha /ds_{\mathrm{eff}}`$, i.e. to a general (perturbative) Finsler structure.
If we use perturbatively the lowest-order “on shell” condition (i.e. $`m_0^2+g_{\mathrm{eff}}^{\alpha \beta }p_\alpha p_\beta 0`$) we can (for instance) eliminate the presence of the energy $`_{\mathrm{eff}}=p_0`$ in the quartic (and higher) terms in (37). In other words, we can restrict ourselves to considering purely spatial higher-order tensors $`A^{\alpha \beta \gamma \delta }(x)=A^{ijk\mathrm{}}(x)`$, etc. Dimensional analysis shows that the quartic terms $`(𝒪(𝐩^4))`$ in Eq. (37) which can enter at the 3PN level must have a $`q`$-dependence of the type $`A^{ijk\mathrm{}}(x)q^2`$. Finally, we are led to consider, at the 3PN accuracy, (after solving Eq. (37) with respect to the effective energy $`_{\mathrm{eff}}=p_0`$) a generalized effective Hamiltonian of the form
$$\widehat{H}_{\mathrm{eff}}^\mathrm{R}(𝐪^{},𝐩^{})=\sqrt{A(q^{})\left[1+𝐩^2+\left(\frac{A(q^{})}{D(q^{})}1\right)(𝐧^{}𝐩^{})^2+\frac{1}{q^2}\mathbf{\left(}z_1(𝐩^2)^2+z_2𝐩^2(𝐧^{}𝐩^{})^2+z_3(𝐧^{}𝐩^{})^4\mathbf{\right)}\right]},$$
(38)
where the quartic (in $`𝐩^{}`$) terms come from the $`A^{\alpha \beta \gamma \delta }`$ coupling and modify the normal “geodesic” Hamiltonian $`\sqrt{g_{00}^{\mathrm{eff}}(1+g_{\mathrm{eff}}^{ij}p_i^{}p_j^{})}`$. In anticipation of the need to transform (via a canonical transformation) the original coordinates $`(𝐪,𝐩)`$ of the real (reduced) Hamiltonian into the coordinates of the effective dynamics, we have denoted the latter by $`(𝐪^{},𝐩^{})`$.
This procedure introduces three new arbitrary degrees of freedom at the 3PN level: $`z_1`$, $`z_2`$, and $`z_3`$. It is clear that it now becomes possible to map the real dynamics on the generalized effective dynamics (38). This becomes, in fact, possible in many ways. As we are primarily interested in (quasi-)circular orbits we shall find convenient to consider only the simple case where $`z_1=z_2=0`$, i.e. to use only $`z_30`$ as new degree of freedom. \[This degree of freedom then disappears in the discussion of circular orbits, which can then be considered as following essentially from a “geodesic” dynamics.\]
An important feature of our proposal (37) is that it is clearly general enough to allow for the existence of solutions at arbitrary PN orders. For instance, at 4PN we would have the freedom to introduce either arbitrary sextic terms, $`q^2[p^6+p^4(n^{}p^{})^2+\mathrm{}]`$, or a modification of the quartic terms: $`q^3[p^4+\mathrm{}]`$. This is more freedom than is needed to compensate for $`\mathrm{\Delta }(4)=42=+2`$, Eq. (35). We can also clearly arrange things so that circular orbits always follow from a “geodesic” dynamics. The only somewhat unsatisfactory feature of the proposal (37) is that we cannot offer any principle for determining a priori the structure of the “post geodesic” terms. In particular, it would have been nice to say that (as is often the case in effective actions) the additional terms are somehow generated by the leading term, i.e. by the effective metric. We have in mind here relations of the type:
$$A_{\alpha \beta \gamma \delta }=\lambda _1R_{\alpha \beta }R_{\gamma \delta }+\lambda _2_{\alpha \beta }R_{\gamma \delta }+\lambda _3_{\alpha \beta \gamma \delta }R+\mathrm{},$$
where $`R_{\alpha \beta }`$ is the Ricci tensor of $`g_{\alpha \beta }^{\mathrm{eff}}`$. However, we have checked that such “geometrical” tensors do not yield possible 3PN corrections, but only much smaller contributions (starting at 5PN).
We shall give in the next section the details of the computation of the coefficients entering $`A(q^{})`$, $`D(q^{})`$, the generating function $`G_{3\mathrm{P}\mathrm{N}}`$, the energy-mapping function and $`z_3`$. For the time being, let us only stress some conceptual points. First, it is remarkable that the unique solution of the 3PN real $``$ effective matching problem leads to the simple value
$$\alpha _3=0,$$
(39)
for the 3PN parameter entering the energy-mapping (30). Ref. had already found that $`\alpha _2=0`$ at 2PN, and $`\alpha _1=\nu /2`$ at 1PN. These values correspond exactly to
$$\frac{_{\mathrm{eff}}^R}{m_0}=\phi (s_{\mathrm{real}})=\frac{(_{\mathrm{real}}^R)^2m_1^2m_2^2}{2m_1m_2}.$$
(40)
We find remarkable that the simplest, symmetric function of the Mandelstam invariant $`s_{\mathrm{real}}=(_{\mathrm{real}}^R)^2`$ which reduces to $`_0/m_0`$ in the test-mass limit turns out to define the unique energy-map needed to link the real 3PN dynamics to the effective dynamics. We interpret this as a good sign for our generalized dynamics (37). \[By contrast, the other proposal of relaxing the natural constraint $`d_1=0`$ leads to a very complicated energy-map with $`\alpha _1\nu /2`$, $`\alpha _20`$, and $`\alpha _30`$, see Appendix A.\]
Let us now consider the consequences of the effective-one-body approach for the determination of the LSO. For the case of circular orbits the effective-one-body approach boils down to saying that the real energy is the following function of the effective-one,
$$_{\mathrm{real}}^R=M\sqrt{1+2\nu \frac{_{\mathrm{eff}}^R\mu }{\mu }}$$
(41)
(as obtained by inverting Eq. (40)) and that the effective energy along circular orbits is the square-root of the minimum value of a certain “effective radial potential”:
$$\frac{_{\mathrm{eff}}^R}{\mu }=\sqrt{[W_j(q^{})]_{\mathrm{min}}},$$
(42)
where $`W_j(q^{})`$ is obtained from (the square of) Eq. (38) by setting $`𝐧^{}𝐩^{}=0`$ (and $`𝐩^2=(𝐧^{}\times 𝐩^{})^2+(𝐧^{}𝐩^{})^2=j^2/q^2`$)
$$W_j(q^{})=A(q^{})\left(1+\frac{j^2}{q^2}+z_1\frac{j^4}{q^6}\right).$$
(43)
As said above, we shall assume (as is always possible) that $`z_1=0`$, so that the effective potential has the usual “Schwarzschild-like” value $`g_{00}(q^{})(1+j^2/q^2)`$. The value of the metric coefficient $`A(q^{})=g_{00}(q^{})`$ is obtained, at the 3PN level, as a perturbative expansion in $`1/q^{}=GM/R_{\mathrm{eff}}`$:
$$A(q^{})=1\frac{2}{q^{}}+\frac{2\nu }{q^3}+\frac{a_4(\nu )}{q^4}+𝒪\left(\frac{1}{q^5}\right).$$
(44)
Note that, finally, in this approach the entire effect of the 3PN dynamics (for circular orbits) is contained in the sole coefficient $`a_4(\nu )`$ (whose value will be discussed in the next section).
In Ref. the metric coefficient $`A(q^{})`$ was used (at the 2PN level) as a simple truncated Taylor expansion: $`A_{2\mathrm{P}\mathrm{N}}(q^{})=12/q^{}+2\nu /q^3`$. This simple-minded approach is not adequate for dealing with the 3PN level. Indeed, we shall see in next section that the coefficient $`a_4(\nu )`$ is positive and can be relatively large. In keeping with the spirit of the present work where we systematically try to use adequate resummation methods to improve the convergence of PN expansions, we shall define our effective-one-body radial potential, at the $`n`$PN accuracy (expressed in terms of the convenient variable $`u1/q^{}`$)
$$W_j^{P_n}(u)=A_{P_n}(u)\left(1+j^2u^2\right)$$
(45)
by using a suitable Padé approximant of Eq. (44):
$$A_{P_n}(u)P_{\mathrm{}}^k\left[T_{n+1}[A(u)]\right],$$
(46)
with $`k+\mathrm{}=n+1`$ (because it is $`q^{n1}`$ which corresponds to the $`n`$PN level) and, now the constraint $`k>0`$ (instead of $`\mathrm{}>0`$ as above), because we want to factor a zero of $`A(u)`$ (and no longer a pole). The Padé improvement of $`A(u)`$ is really needed (and makes a difference) only at the 3PN level. We have found that the most robust Padés (under variation of the Taylor coefficients) are defined by taking $`k=1`$ and $`\mathrm{}=n`$.
Summarizing the present method: The effective radius $`q^{}`$ of circular orbits is obtained as a function of the reduced angular momentum $`j`$ by looking for the minimum of the radial potential (45), where $`u1/q^{}`$ and where the Padé-improved function $`A`$ is given by Eq. (46), with $`k=1`$ and $`\mathrm{}=n`$. For each value of $`j`$ above some threshold $`j_{\mathrm{min}}`$, $`W_j(u)`$ admits a (unique) minimum $`u_{}(j)`$. From this one then determines the effective-energy (42), and then the real one (41), namely
$$_{\mathrm{real}}^R(j)=M\sqrt{1+2\nu \left[\sqrt{W_j(u_{}(j))}1\right]}.$$
(47)
The real circular orbital frequency corresponding to $`j`$ is then given by using the identity (21). This yields
$$GM\omega _{\mathrm{real}}(j)=\frac{ju_{}^2(j)\sqrt{A_{P_n}(u_{}(j))}}{\sqrt{1+j^2u_{}^2(j)}\sqrt{1+2\nu \left[\sqrt{W_j(u_{}(j))}1\right]}}.$$
(48)
Note that in this approach $`jj_{\mathrm{real}}j_{\mathrm{effective}}`$. Finally, the LSO is invariantly defined as the minimum value of $`j`$, $`j_{\mathrm{LSO}}=j_{\mathrm{min}}`$, for which $`W_j(u)`$ admits a local minimum. When $`j<j_{\mathrm{LSO}}`$, $`W_j(u)`$ has no local minimum and there are no (stable or unstable) circular orbits (see, e.g., Fig. 1 of Ref. ). If one is only interested in locating the LSO (as a function of $`\nu `$) it suffices to look for the existence of an inflection point of $`W_j(u)`$, i.e. to solve the simultaneous equations $`W_j(u)/u=0`$ and $`^2W_j(u)/u^2=0`$.
## IV Results
Let us now give the details of the application, at the 3PN level, of the methods explained above. We follow the order of presentation given in the previous two sections. In what follows, we use the results of Ref. for the dynamical invariants of the 3PN two-body Hamiltonian. We recall that the 3PN Hamiltonian derived in Ref. contained two ambiguous parameters $`\omega _{\text{static}}`$ and $`\omega _{\text{kinetic}}`$ (these ambiguities arise because of the need to regularize badly divergent integrals ), and that all the dynamical invariants involve only the combination $`\sigma (\nu )\omega _{\text{static}}\nu +\omega _{\text{kinetic}}\nu ^2`$ .
Recently the ‘kinetic’ ambiguity parameter $`\omega _{\text{kinetic}}`$ was uniquely determined (see also ) to be $`\omega _{\text{kinetic}}=41/24`$, so that $`\sigma (\nu )=\omega _{\text{static}}\nu +\frac{41}{24}\nu ^2`$ and the remaining 3PN ambiguity is embodied in the product $`\omega _{\text{static}}\nu `$. As discussed in the Appendix A of one expects (both by judging from the other coefficients, and by looking at some of the sources of ambiguity) that $`\omega _{\text{static}}`$ varies in the range
$$10\omega _{\text{static}}+10.$$
(49)
### A $`e`$-method
For self-containedness let us quote the results obtained in our previous paper for the link between the energy and the $`x`$-variable, Eq. (1). The original energy function $`E`$, Eq. (2), admits the PN expansion (4) with coefficients
$`E_1(\nu )`$ $`=`$ $`{\displaystyle \frac{1}{12}}(9+\nu ),`$ (51)
$`E_2(\nu )`$ $`=`$ $`{\displaystyle \frac{1}{24}}(8157\nu +\nu ^2),`$ (52)
$`E_3(\nu )`$ $`=`$ $`{\displaystyle \frac{10}{3}}\mathbf{\left(}w_1(\nu )\omega _{\text{static}}\nu \mathbf{\right)},`$ (53)
where
$$w_1(\nu )\frac{405}{128}+\frac{1}{64}\left(41\pi ^2\frac{6889}{6}\right)\nu +\frac{31}{64}\nu ^2+\frac{7}{3456}\nu ^3.$$
(54)
Correspondingly to this expansion, the new energy function $`e`$, Eq. (13), admits the expansion (15) with coefficients
$`e_1(\nu )`$ $`=`$ $`\left(1+{\displaystyle \frac{1}{3}}\nu \right),`$ (56)
$`e_2(\nu )`$ $`=`$ $`\left(3{\displaystyle \frac{35}{12}}\nu \right),`$ (57)
$`e_3(\nu )`$ $`=`$ $`{\displaystyle \frac{10}{3}}\mathbf{\left(}w_2(\nu )\omega _{\text{static}}\nu \mathbf{\right)},`$ (58)
where
$$w_2(\nu )\frac{27}{10}+\frac{1}{16}\left(\frac{41}{4}\pi ^2\frac{4309}{15}\right)\nu +\frac{103}{120}\nu ^2\frac{1}{270}\nu ^3.$$
(59)
The 2PN and 3PN Padés of $`e(x)`$ are given by (see Ref. for the 2PN case)
$`e_{P_2}(x)`$ $``$ $`P_1^1\left[T_2[e(x)]\right]=x{\displaystyle \frac{1+\frac{1}{3}\nu \left(4\frac{9}{4}\nu +\frac{1}{9}\nu ^2\right)x}{1+\frac{1}{3}\nu \left(3\frac{35}{12}\nu \right)x}},`$ (61)
$`e_{P_3}(x)`$ $``$ $`P_1^2\left[T_3[e(x)]\right]=x{\displaystyle \frac{1\mathbf{\left(}1+\frac{1}{3}\nu +w_3(\nu )\mathbf{\right)}x\mathbf{\left(}3\frac{35}{12}\nu \left(1+\frac{1}{3}\nu \right)w_3(\nu )\mathbf{\right)}x^2}{1w_3(\nu )x}},`$ (62)
where
$$w_3(\nu )\frac{40}{3635\nu }\mathbf{\left(}w_2(\nu )\omega _{\text{static}}\nu \mathbf{\right)}.$$
(63)
The corresponding $`x`$-location of the $`e`$-LSO (minimum of $`e(x)`$) can be written down analytically only at the 2PN level :
$$6x_{\mathrm{LSO}}^{e_{P_2}}(\nu )=\frac{1+\frac{1}{3}\nu }{1\frac{35}{36}\nu }\left(2\frac{1+\frac{1}{3}\nu }{\sqrt{1\frac{9}{16}\nu +\frac{1}{36}\nu ^2}}\right).$$
(64)
Note that $`6x_{\mathrm{LSO}}^{e_{P_2}}\left(\frac{1}{4}\right)=1.1916`$, which means that the $`e_{P_2}`$-predicted value of the orbital frequency at the LSO differs from the “Schwarzschild value”, $`GM\omega _{\mathrm{LSO}}^{\mathrm{Schw}}=(x_{\mathrm{LSO}}^{\mathrm{Schw}})^{3/2}`$ with $`x_{\mathrm{LSO}}^{\mathrm{Schw}}=1/6`$, by the factor
$$\widehat{\omega }_{\mathrm{LSO}}\frac{\omega _{\mathrm{LSO}}}{\omega _{\mathrm{LSO}}^{\mathrm{Schw}}}=(6x_{\mathrm{LSO}})^{3/2},$$
(65)
which is about $`1.3007`$ in the present case. We will use $`\widehat{\omega }_{\mathrm{LSO}}`$ as our main tracer of the “observable” location of the LSO. It is important to note from the start that (as emphasized in ) the $`e`$-method predicts (at 2PN) that the orbital frequency at the LSO is larger than the “Schwarzschild value”. The corresponding results, at 3PN, are exhibited in Table I. Let us only note here that the tendency to get $`\widehat{\omega }_{\mathrm{LSO}}>1`$ seems confirmed, at the 3PN level, rather independently of the value of the ambiguity parameter $`\omega _{\text{static}}`$.
Once the value of $`x_{\mathrm{LSO}}(\nu )`$ is determined (analytically or numerically) one can compute the corresponding value of the (real) reduced binding energy $`E`$, Eq. (2). It is obtained by solving Eq. (13) in terms of $`^RM+\mu E`$. The solution is explicitly given by
$$E(x)=\frac{1}{\nu }\left[\sqrt{1+2\nu \left(\sqrt{1+e(x)}1\right)}1\right].$$
(66)
Then, knowing $`E(x)`$ we can also compute the value of the reduced angular momentum $`j`$ by integrating the identity (22). Integrating Eq. (22) by parts yields
$$j(x)=2x^{1/2}\frac{dE(x)}{dx}+2_0^x𝑑\overline{x}\overline{x}^{1/2}\frac{d^2E(\overline{x})}{d\overline{x}^2},$$
(67)
where we have incorporated the information that $`j(x)x^{1/2}`$ when $`x0`$ (i.e. in the limit of very wide circular orbits, described by Newtonian dynamics). By applying this result to $`x=x_{\mathrm{LSO}}(\nu )`$, one finally gets the value of $`j_{\mathrm{LSO}}(\nu )`$. The results so obtained by the $`e`$-method are shown in Table I.
### B $`j`$-method
In this approach the basic PN expansion we consider is that of $`1/j^2(x)`$. It reads (cf. Eq. (5.11) in Ref. )
$$\frac{1}{j^2(x)}=x\left[1\frac{1}{3}(9+\nu )x+\frac{25}{4}\nu x^2\frac{16}{3}\mathbf{\left(}w_4(\nu )\omega _{\text{static}}\nu \mathbf{\right)}x^3\right],$$
(68)
where
$$w_4(\nu )\frac{1}{64}\left(41\pi ^2\frac{5269}{6}\right)\nu +\frac{61}{64}\nu ^2\frac{1}{432}\nu ^3.$$
(69)
From Eq. (68) one gets
$$j^2(x)=\frac{1}{x}\left[1+\frac{1}{3}(9+\nu )x+\frac{1}{36}(36\nu )(94\nu )x^2+\frac{16}{3}\mathbf{\left(}w_5(\nu )\omega _{\text{static}}\nu \mathbf{\right)}x^3\right],$$
(70)
where
$$w_5(\nu )\frac{81}{16}+\frac{1}{64}\left(41\pi ^2\frac{7321}{6}\right)\nu +\frac{23}{64}\nu ^2+\frac{1}{216}\nu ^3.$$
(71)
As explained above we construct the following sequence of near-diagonal Padés of $`j^2(x)`$:
$`j_{P_1}^2(x)`$ $``$ $`P_1^0\left[T_1[j^2(x)]\right]={\displaystyle \frac{1}{x\mathbf{\left(}1\left(3+\frac{1}{3}\nu \right)x\mathbf{\right)}}},`$ (73)
$`j_{P_2}^2(x)`$ $``$ $`P_1^1\left[T_2[j^2(x)]\right]={\displaystyle \frac{1+\frac{1}{9}\nu +\frac{25}{12}\nu x}{x\mathbf{\left(}1+\frac{1}{9}\nu \left(3\frac{17}{12}\nu +\frac{1}{27}\nu ^2\right)x\mathbf{\right)}}},`$ (74)
$`j_{P_3}^2(x)`$ $``$ $`P_1^2\left[T_3[j^2(x)]\right]={\displaystyle \frac{1+\mathbf{\left(}3+\frac{1}{3}\nu w_6(\nu )\mathbf{\right)}x+\mathbf{\left(}9\frac{17}{4}\nu +\frac{1}{9}\nu ^2\left(3+\frac{1}{3}\nu \right)w_6(\nu )\mathbf{\right)}x^2}{x\mathbf{\left(}1w_6(\nu )x\mathbf{\right)}}},`$ (75)
where
$$w_6(\nu )\frac{192}{(36\nu )(94\nu )}\mathbf{\left(}w_5(\nu )\omega _{\text{static}}\nu \mathbf{\right)}.$$
(76)
The corresponding $`x_{\text{LSO}}`$ (now defined as the location of the minimum of $`j(x)`$, or, equivalently, $`j^2(x)`$) can be written down analytically at 1PN and 2PN
$`6x_{\mathrm{LSO}}^{j_{P_1}}(\nu )`$ $`=`$ $`{\displaystyle \frac{1}{1+\frac{1}{9}\nu }},`$ (78)
$`6x_{\mathrm{LSO}}^{j_{P_2}}(\nu )`$ $`=`$ $`{\displaystyle \frac{8(9+\nu )}{25\nu }}\left[{\displaystyle \frac{2(9+\nu )}{\sqrt{(36\nu )(94\nu )}}}1\right].`$ (79)
Note that while $`6x_{\mathrm{LSO}}^{j_{P_1}}(\nu )`$ is very slightly smaller than 1, $`6x_{\mathrm{LSO}}^{j_{P_2}}(\nu )`$ is (like for the $`e_{P_2}`$ estimate) larger than 1. In particular, $`6x_{\mathrm{LSO}}^{j_{P_2}}\left(\frac{1}{4}\right)=1.1121`$, and the corresponding dimensionless frequency is $`\widehat{\omega }_{\mathrm{LSO}}^{j_{P_2}}\left(\frac{1}{4}\right)=1.1728`$. This tendency to get “larger than Schwarzschild” frequency at the LSO is confirmed by the (numerical) 3PN results which are exhibited in Table I and Fig. 1.
Within the present $`j`$-method, once the value of $`x_{\mathrm{LSO}}(\nu )`$ is determined one can compute the corresponding value of the (real) reduced binding energy $`E`$ by integrating the identity (22). Indeed, one can write
$$E(x)=_0^x𝑑\overline{x}\overline{x}^{3/2}\frac{dj(\overline{x})}{d\overline{x}},$$
(80)
where one has incorporated the boundary condition that $`E(x)0`$ when $`x0`$. The results so obtained are shown in Table I and Fig. 1.
### C $`k`$-method
For completeness, let us (though it is not among our preferred methods) mention some results obtained by using the $`k`$-method, Eq. (24). The PN expansion of the function $`K(y)`$, where $`K(1+k)^4`$ and $`y1/j^2`$, reads (cf. Eq. (5.27) in Ref. )
$$K(y)=1+12y+24(6\nu )y^2+24\mathbf{\left(}w_7(\nu )\omega _{\text{static}}\nu \mathbf{\right)}y^3,$$
(81)
where
$$w_7(\nu )72+\left(\frac{41}{64}\pi ^2\frac{128}{3}\right)\nu +\frac{1}{2}\nu ^2.$$
(82)
As explained above, we construct the following sequence of near-diagonal Padés of $`K(y)`$:
$`K_{P_1}(y)`$ $``$ $`P_1^0\left[T_1[K(y)]\right]={\displaystyle \frac{1}{112y}},`$ (84)
$`K_{P_2}(y)`$ $``$ $`P_1^1\left[T_2[K(y)]\right]={\displaystyle \frac{1+2\nu y}{12(6\nu )y}},`$ (85)
$`K_{P_3}(y)`$ $``$ $`P_1^2\left[T_3[K(y)]\right]={\displaystyle \frac{1+\mathbf{\left(}12w_8(\nu )\mathbf{\right)}y+12\mathbf{\left(}2(6\nu )w_8(\nu )\mathbf{\right)}y^2}{1w_8(\nu )y}},`$ (86)
where
$$w_8(\nu )\frac{1}{6\nu }\mathbf{\left(}w_7(\nu )\omega _{\text{static}}\nu \mathbf{\right)}.$$
(87)
Then we take the poles of $`K_{P_n}`$ as estimates of the value of $`1/j^2`$ at the LSO. The results for equal-mass binaries ($`\nu =1/4`$) are given in Table I.
### D Effective one-body method
Let us explain in detail how we implemented the effective one-body method. Two methods of implementation were presented in Ref. . We could have used the first one by starting from the 3PN Delaunay Hamiltonian given in Ref. . However, we found finally as convenient (given the existence of good algebraic manipulators) to use the second method, which has the advantage of being more informative. This second method consists of writing explicitly the equations that have to be satisfied by the looked for generating function $`G(q,p^{})`$ and solving them. Indeed, we look for a canonical transformation between the original (quasi-ADM) coordinates $`(q,p)`$ of the real problem (i.e. the phase-space coordinates in which was obtained the order-reduced Hamiltonian $`H(q,p)`$ in ), and the “effective” coordinates $`(q^{},p^{})`$ (i.e. the coordinates used in Eq. (38)). The effect of the generating function $`G(q,p^{})`$ reads
$$q^i=q^i+\frac{G(q,p^{})}{p_i^{}},p_i=p_i^{}+\frac{G(q,p^{})}{q^i}.$$
(88)
Note that, as is well known, the canonical transformation is defined only in implicit form: $`q^{}`$ and $`p`$ being given as functions of $`q`$ and $`p^{}`$. But there is, in fact, no need to solve for, e.g., $`(q^{},p^{})`$ as functions of $`(q,p)`$. As the basic idea is anyway to identify the numerical value of, say, $`H_{\mathrm{eff}}(q^{},p^{})`$ with the numerical value of some (to be determined) function $`f(H_{\mathrm{real}}(q,p))`$, we can do this identification by expressing both sides in terms of any set of common variables, say $`q`$ and $`p^{}`$. Finally we write (using Eq. (30))
$$\left[\widehat{H}_{\mathrm{eff}}^\mathrm{R}\mathbf{(}q^{}(q,p^{}),p^{}\mathbf{)}\right]^2=\left\{1+\widehat{H}_{\mathrm{real}}^{\mathrm{NR}}\mathbf{(}q,p(q,p^{})\mathbf{)}\left[1+\alpha _1\widehat{H}_{\mathrm{real}}^{\mathrm{NR}}+\alpha _2\left(\widehat{H}_{\mathrm{real}}^{\mathrm{NR}}\right)^2+\alpha _3\left(\widehat{H}_{\mathrm{real}}^{\mathrm{NR}}\right)^3\right]\right\}^2,$$
(89)
in which the L.H.S. is given by the R.H.S. of Eq. (38) (without the square root, because we work with the squared equation), while $`\widehat{H}_{\mathrm{real}}^{\mathrm{NR}}(q,p)`$ on the R.H.S. is the order-reduced Hamiltonian of , obtained from the higher-order 3PN Hamiltonian derived in . Both sides of Eq. (89) are written in terms of $`q`$ and $`p^{}`$ by using Eq. (88). This procedure has been used at the 2PN level in , so that we know the values of $`\alpha _1=\frac{1}{2}\nu `$, $`\alpha _2=0`$, the 2PN values of the metric functions $`A(q^{})`$ and $`D(q^{})`$, as well as the 2PN-accurate generating function $`G(q,p^{})`$ (see e.g. Eqs. (6.24)–(6.27) in ). The new unknowns entering at the 3PN level are: $`\alpha _3`$, $`a_4`$, Eq. (28), $`d_3`$, Eq. (29), $`z_1`$, $`z_2`$, $`z_3`$, Eq. (38), and the 7 arbitrary coefficients $`c_1`$, …, $`c_7`$ entering the generic form of $`G_{3\mathrm{P}\mathrm{N}}`$:
$$G_{3\mathrm{P}\mathrm{N}}(q,p^{})=(𝐪𝐩^{})\left[c_1𝐩^6+\frac{1}{q}\mathbf{\left(}c_2𝐩^4+c_3𝐩^2(𝐧𝐩^{})^2+c_4(𝐧𝐩^{})^4\mathbf{\right)}+\frac{1}{q^2}\mathbf{\left(}c_5𝐩^2+c_6(𝐧𝐩^{})^2\mathbf{\right)}+\frac{c_7}{q^3}\right].$$
(90)
The basic input for writing these equations is the explicit value of the 3PN-accurate Hamiltonian
$$\widehat{H}_{\text{real}}^{\text{NR}}(𝐪,𝐩)=\widehat{H}_\mathrm{N}(𝐪,𝐩)+\widehat{H}_{1\mathrm{P}\mathrm{N}}(𝐪,𝐩)+\widehat{H}_{2\mathrm{P}\mathrm{N}}(𝐪,𝐩)+\widehat{H}_{3\mathrm{P}\mathrm{N}}(𝐪,𝐩),$$
(91)
where
$`\widehat{H}_\mathrm{N}(𝐪,𝐩)={\displaystyle \frac{𝐩^2}{2}}{\displaystyle \frac{1}{q}},`$ (93)
$`\widehat{H}_{1\mathrm{P}\mathrm{N}}(𝐪,𝐩)={\displaystyle \frac{1}{8}}(3\nu 1)(𝐩^2)^2{\displaystyle \frac{1}{2}}\left[(3+\nu )𝐩^2+\nu (𝐧𝐩)^2\right]{\displaystyle \frac{1}{q}}+{\displaystyle \frac{1}{2q^2}},`$ (94)
$`\widehat{H}_{2\mathrm{P}\mathrm{N}}(𝐪,𝐩)={\displaystyle \frac{1}{16}}\left(15\nu +5\nu ^2\right)(𝐩^2)^3+{\displaystyle \frac{1}{8}}\left[\left(520\nu 3\nu ^2\right)(𝐩^2)^22\nu ^2(𝐧𝐩)^2𝐩^23\nu ^2(𝐧𝐩)^4\right]{\displaystyle \frac{1}{q}}`$ (95)
$`+{\displaystyle \frac{1}{2}}\left[(5+8\nu )𝐩^2+3\nu (𝐧𝐩)^2\right]{\displaystyle \frac{1}{q^2}}{\displaystyle \frac{1}{4}}(1+3\nu ){\displaystyle \frac{1}{q^3}},`$ (96)
$`\widehat{H}_{3\mathrm{P}\mathrm{N}}(𝐪,𝐩)={\displaystyle \frac{1}{128}}\left(5+35\nu 70\nu ^2+35\nu ^3\right)(𝐩^2)^4`$ (97)
$`+{\displaystyle \frac{1}{16}}\left[\left(7+42\nu 53\nu ^25\nu ^3\right)(𝐩^2)^3+(23\nu )\nu ^2(𝐧𝐩)^2(𝐩^2)^2+3(1\nu )\nu ^2(𝐧𝐩)^4𝐩^25\nu ^3(𝐧𝐩)^6\right]{\displaystyle \frac{1}{q}}`$ (98)
$`+\left[{\displaystyle \frac{1}{16}}\left(27+136\nu +109\nu ^2\right)(𝐩^2)^2+{\displaystyle \frac{1}{16}}(17+30\nu )\nu (𝐧𝐩)^2𝐩^2+{\displaystyle \frac{1}{12}}(5+43\nu )\nu (𝐧𝐩)^4\right]{\displaystyle \frac{1}{q^2}}`$ (99)
$`+\left\{\left[{\displaystyle \frac{25}{8}}+\left({\displaystyle \frac{1}{64}}\pi ^2{\displaystyle \frac{335}{48}}\right)\nu {\displaystyle \frac{23}{8}}\nu ^2\right]𝐩^2+\left({\displaystyle \frac{85}{16}}{\displaystyle \frac{3}{64}}\pi ^2{\displaystyle \frac{7}{4}}\nu \right)\nu (𝐧𝐩)^2\right\}{\displaystyle \frac{1}{q^3}}`$ (100)
$`+\left[{\displaystyle \frac{1}{8}}+\left({\displaystyle \frac{109}{12}}{\displaystyle \frac{21}{32}}\pi ^2+\omega _{\text{static}}\right)\nu \right]{\displaystyle \frac{1}{q^4}}.`$ (101)
As explained in Refs. and at the beginning of the present section the 3PN Hamiltonian contains one dimensionless ambiguity parameter $`\omega _{\text{static}}`$.
When written explicitly, the constraint equation (89) (truncated at 3PN accuracy) yields a system of 11 equations for the $`10+3`$ unknowns $`(\alpha _3,a_4,d_3,c_1,\mathrm{},c_7;z_1,z_2,z_3)`$ ($`\omega _{\text{static}}`$ being assumed to be known). This system can be decomposed into three subsystems. The first subsystem consists of 5 equations:
$`\alpha _3+16c_1`$ $`=`$ $`\nu 3\nu ^25\nu ^3,`$ (103)
$`\alpha _3+2c_12c_2`$ $`=`$ $`{\displaystyle \frac{1}{8}}\nu ^22\nu ^3,`$ (104)
$`6c_1+c_23c_3`$ $`=`$ $`{\displaystyle \frac{17}{16}}\nu +{\displaystyle \frac{19}{4}}\nu ^2+{\displaystyle \frac{27}{16}}\nu ^3,`$ (105)
$`3c_35c_4`$ $`=`$ $`{\displaystyle \frac{3}{2}}\nu {\displaystyle \frac{27}{2}}\nu ^2{\displaystyle \frac{81}{16}}\nu ^3,`$ (106)
$`c_4`$ $`=`$ $`{\displaystyle \frac{3}{2}}\nu ^2+{\displaystyle \frac{7}{16}}\nu ^3.`$ (107)
The second subsystem contains 3 equations:
$`3\alpha _3+2c_22c_5+z_1`$ $`=`$ $`{\displaystyle \frac{3}{2}}\nu {\displaystyle \frac{9}{4}}\nu ^2+{\displaystyle \frac{19}{8}}\nu ^3,`$ (109)
$`8c_2+6c_3+4c_56c_6+z_2`$ $`=`$ $`{\displaystyle \frac{61}{8}}\nu +{\displaystyle \frac{11}{2}}\nu ^2{\displaystyle \frac{11}{2}}\nu ^3,`$ (110)
$`4c_3+10c_4+8c_6+z_3`$ $`=`$ $`{\displaystyle \frac{79}{6}}\nu {\displaystyle \frac{55}{3}}\nu ^2+{\displaystyle \frac{39}{8}}\nu ^3,`$ (111)
and the third one consists also of 3 equations:
$`2\alpha _3+c_5c_7`$ $`=`$ $`\left({\displaystyle \frac{271}{48}}+{\displaystyle \frac{1}{64}}\pi ^2\right)\nu +{\displaystyle \frac{5}{8}}\nu ^2{\displaystyle \frac{5}{8}}\nu ^3,`$ (113)
$`d_3+4c_5+6c_6+6c_7`$ $`=`$ $`\left({\displaystyle \frac{35}{8}}{\displaystyle \frac{3}{32}}\pi ^2\right)\nu {\displaystyle \frac{57}{4}}\nu ^2+2\nu ^3,`$ (114)
$`2\alpha _3+a_4+2c_7`$ $`=`$ $`\left({\displaystyle \frac{221}{12}}{\displaystyle \frac{21}{16}}\pi ^2+2\omega _{\text{static}}\right)\nu +{\displaystyle \frac{3}{4}}\nu ^2+{\displaystyle \frac{1}{4}}\nu ^3.`$ (115)
The first subsystem, Eqs. (IV D), yields 5 linear equations for the 5 unknowns $`c_1`$, $`c_2`$, $`c_3`$, $`c_4`$, and $`\alpha _3`$. It is easily found to have a unique solution, namely
$`\alpha _3`$ $`=`$ $`0,`$ (117)
$`c_1`$ $`=`$ $`{\displaystyle \frac{1}{16}}(1+3\nu +5\nu ^2)\nu ,`$ (118)
$`c_2`$ $`=`$ $`{\displaystyle \frac{1}{16}}(1+2\nu 11\nu ^2)\nu ,`$ (119)
$`c_3`$ $`=`$ $`{\displaystyle \frac{1}{24}}(12+48\nu +23\nu ^2)\nu ,`$ (120)
$`c_4`$ $`=`$ $`{\displaystyle \frac{1}{16}}(24+7\nu )\nu ^2.`$ (121)
As already mentioned above, note the remarkably simple result $`\alpha _3=0`$ (which confirms that the energy map takes the nice form (40)). It is also remarkable that the result $`\alpha _3=0`$ holds independently of any assumption about the “quartic” parameters $`z_1`$, $`z_2`$, and $`z_3`$.
The second subsystem (IV D) can be viewed (after inserting the unique solution of the first subsystem) as an overdetermined system for the two unknowns $`c_5`$, $`c_6`$. It is then easily seen that it will admit a solution if and only if the parameters $`z_1`$, $`z_2`$, and $`z_3`$ satisfy the following linear constraint:
$$8z_1+4z_2+3z_3=6(43\nu )\nu .$$
(122)
This linear constraint forbids us to consider the simplest “geodesic” case where $`z_1=z_2=z_3=0`$. We can, however, continue to impose the natural conditions $`z_1=0=z_2`$ which simplify very much the 3PN effective dynamics of circular orbits. With this choice, the general constraint (122) yields
$$z_3=2(43\nu )\nu .$$
(123)
Having fixed the values of $`z_1`$, $`z_2`$, and $`z_3`$, the system (IV D) uniquely determines $`c_5`$ and $`c_6`$ to be
$`c_5`$ $`=`$ $`{\displaystyle \frac{1}{16}}(1316\nu +6\nu ^2)\nu ,`$ (125)
$`c_6`$ $`=`$ $`{\displaystyle \frac{1}{48}}(115+116\nu 26\nu ^2)\nu .`$ (126)
Finally, the subsystem (IV D) gives 3 equations for the remaining 3 unknowns $`c_7`$, $`d_3`$, and $`a_4`$. The unique solution of this system reads:
$`c_7`$ $`=`$ $`\left({\displaystyle \frac{1}{64}}\pi ^2+{\displaystyle \frac{155}{24}}\right)\nu +{\displaystyle \frac{3}{8}}\nu ^2+{\displaystyle \frac{1}{8}}\nu ^3,`$ (128)
$`d_3`$ $`=`$ $`2(3\nu 26)\nu ,`$ (129)
$`a_4`$ $`=`$ $`\left({\displaystyle \frac{94}{3}}{\displaystyle \frac{41}{32}}\pi ^2+2\omega _{\text{static}}\right)\nu .`$ (130)
Note how simple the structure of the coefficient $`a_4`$, Eq. (130), is. Indeed, the right-hand-sides of all the equations (IV D), (IV D), (IV D), were polynomials of the third degree in $`\nu `$. Therefore, one would have a priori expected the 3PN coefficient $`a_4`$ to have the same structure: $`a_4=a_{41}\nu +a_{42}\nu ^2+a_{43}\nu ^3`$. It is remarkable that the coefficients of $`\nu ^2`$ and $`\nu ^3`$ happen to vanish in $`a_4`$ (such a simplification does not occur in the 3PN-level coefficients appearing in the functions $`E(x)`$, $`e(x)`$ and $`j^2(x)`$ considered above). This simple structure of $`a_4`$ can be brought out by defining the following quantity,
$$\omega _{\text{static}}^{}\frac{47}{3}+\frac{41}{64}\pi ^2=9.3439\mathrm{},$$
(131)
in terms of which the value of $`a_4`$ can be written as
$$a_4=2(\omega _{\text{static}}\omega _{\text{static}}^{})\nu .$$
(132)
The presence of already two cancellations in $`a_4`$ ($`a_{42}=0=a_{43}`$) suggests that the yet undetermined value of $`\omega _{\text{static}}`$ might be precisely $`\omega _{\text{static}}^{}`$, so that $`a_4`$, Eq. (130), simply vanishes. We shall see below that this conjecture is indirectly supported by the fact that a numerical value $`\omega _{\text{static}}9`$ is selected by the requirement that the various methods discussed in this paper agree in their predictions for LSO quantities. Note also the remarkable fact that if $`\omega _{\text{static}}=\omega _{\text{static}}^{}`$ all the $`\pi ^2`$ terms cancell in all the 3PN-level coefficients, so that they all become rational (as were the 2PN ones). \[Stated in reverse, if one could a priori prove that all the 3PN coefficients are rational, this would support the conjecture that $`\omega _{\text{static}}=\omega _{\text{static}}^{}`$, though it would be also compatible with having $`\omega _{\text{static}}=\omega _{\text{static}}^{}+`$ a rational number.\]
The coefficient $`a_4`$ enters the PN expansion of $`A(u)g_{00}(q^{})`$ (with $`u1/q^{}`$):
$$A(u)=12u+2\nu u^3+a_4(\nu )u^4+𝒪(u^5).$$
(133)
As mentioned above, we improve the behaviour of the PN expansion of $`A(u)`$ by Padeeing it: $`A_{P_n}(u)=P_{\mathrm{}}^k[T_{n+1}[A(u)]]`$ with $`k+\mathrm{}=n+1`$. We impose the constraint $`k>0`$ to inject the information that $`A(u)`$ should qualitatively look like $`A_{\mathrm{Schw}}(u)=12u`$, i.e. have a zero at some $`u=\frac{1}{2}+𝒪(\nu )`$. As said above, the most robust (for our purpose) Padés of $`A(u)`$ are the ones with $`k=1`$ and $`\mathrm{}=n`$. Finally, we get the following sequence of Padé-improved $`A`$:
$`A_{P_1}(u)`$ $``$ $`P_1^1\left[T_2[A(u)]\right]=12u,`$ (135)
$`A_{P_2}(u)`$ $``$ $`P_2^1\left[T_3[A(u)]\right]={\displaystyle \frac{1\left(2\frac{1}{2}\nu \right)u}{1+\frac{1}{2}\nu u+\nu u^2}},`$ (136)
$`A_{P_3}(u)`$ $``$ $`P_3^1\left[T_4[A(u)]\right]={\displaystyle \frac{2(4\nu )+\mathbf{\left(}a_4(\nu )16+8\nu \mathbf{\right)}u}{2(4\nu )+\mathbf{\left(}a_4(\nu )+4\nu \mathbf{\right)}u+2\mathbf{\left(}a_4(\nu )+4\nu \mathbf{\right)}u^2+4\mathbf{\left(}a_4(\nu )+\nu ^2\mathbf{\right)}u^3}}.`$ (137)
To extract the LSO quantities from these Padéed $`A`$’s, we must consider the effective radial potential
$$W_j^{P_n}(u)=A_{P_n}(u)(1+j^2u^2).$$
(138)
The value of $`j`$ for which this radial potential has an inflection point defines $`j_{\mathrm{LSO}}(\nu )`$; the corresponding value of $`u`$ being $`u_{\mathrm{LSO}}=u_{}(j_{\mathrm{LSO}})`$. As explained in Eqs. (47) and (48) above one then deduces the energy and the orbital frequency of the LSO. Note that the 2PN Padéed $`A_{P_2}`$ that we use here differs from the straightforward Taylor approximant $`A_{T_2}`$ used in Ref. . However, this difference is essentially negligible (as shown by comparing the lines “eff. method” and “BD” in Table I). The Padé improvement is, however, rather important at 3PN in the case where $`\omega _{\text{static}}`$ is significantly larger than $`\omega _{\text{static}}^{}`$. Indeed, in this case, Padéeing allows one to tame the effect of a largish (positive) $`a_4(\nu )`$: $`a_4\left(\frac{1}{4}\right)4.67+\frac{1}{2}\omega _{\text{static}}`$. For instance, when $`\nu =\frac{1}{4}`$ and $`\omega _{\text{static}}1.2`$ the radial potential built from a straight Taylor-approximated $`A_{T_4}(u)`$ would not give rise to an inflection point continuously connected to the test-mass limit. We attribute this lack of structural stability to the known bad properties of high-order PN expansions, and not to the effective-one-body approach. The (numerical) results obtained by the effective one-body approach are exhibited in Table I.
## V Discussion
Before discussing the meaning of the results obtained above, let us state what we would a priori expect. First, we recall that the study in Ref. has shown that the sequence of Padé approximants of the invariant function $`F(v)`$, giving the gravitational wave flux in terms of $`v(GM\omega /c^3)^{1/3}=x^{1/2}`$, had very good (and very monotonic) convergence properties toward the exact result. \[By contrast, the sequence of Taylor approximants was badly convergent, and unstable when $`vv_{\mathrm{LSO}}^{\mathrm{Schw}}=0.40825`$; see Figs. 3a and 3b of .\] In our case, as one can meaningfully (at least for the $`j`$\- and effective-one-body methods) consider the 1PN, 2PN, and 3PN approximations, we would expect that a good resummation technique would ensure that any LSO quantity, say $`Q_{\mathrm{LSO}}`$, be determined with increasing accuracy, when using higher PN information, and, more precisely, that
$$Q_{\mathrm{LSO}}^{P_n}Q_{\mathrm{LSO}}^X+a(bx_{\mathrm{LSO}})^{n+1},$$
(139)
with (hopefully) coefficients $`a`$ and $`b`$ small enough to ensure a visible convergence (when $`x_{\mathrm{LSO}}x_{\mathrm{LSO}}^{\mathrm{Schw}}=1/6`$). As a minimum test of improved convergence we hope that $`|Q_{3\mathrm{P}\mathrm{N}}Q_{2\mathrm{P}\mathrm{N}}|`$ would be significantly smaller than $`|Q_{2\mathrm{P}\mathrm{N}}Q_{1\mathrm{P}\mathrm{N}}|`$, i.e. that the addition of the 3PN information would have only slightly refined the previous 2PN estimates of LSO quantities .
Independently of this expectation, we had also hoped, when starting this investigation, that the LSO quantities might be “robust” under the lack of precise knowledge of a sole ambiguous coefficient ($`\omega _{\text{static}}`$) among many others in $`H_{3\mathrm{P}\mathrm{N}}`$. \[Given that the amplitude of this coefficient would have some plausible upper bound; as discussed in the Appendix A of .\] However, the results exhibited in Table I and Fig. 1 show that, in spite of our use of resummation techniques, the LSO quantities appear to be quite sensitive to the exact value of $`\omega _{\text{static}}`$. A first conclusion of our work is therefore that it is quite important to resolve the problem of static ambiguity, arising at 3PN when using delta functions to represent compact (but extended) objects. Until this problem is unambiguously solved, it will not be clear whether (as proposed in ) it is possible to trust suitably resummed versions of PN-expanded results.
In the meantime, however, we wish to point out several remarkable features of the dependence of our various results on $`\omega _{\text{static}}`$. In Fig. 2 we plot (for the equal-mass case, $`\nu =1/4`$) our various predictions, at the 3PN level, and using various methods, as a function of the 3PN ambiguity parameter $`\omega _{\text{static}}`$. It is quite interesting to note that two a priori independent things happen:
(i) there is a value of $`\omega _{\text{static}}`$, namely
$$\omega _{\text{static}}^{\text{best}}9,$$
(140)
for which the three different methods give, at 3PN, nearly coincident LSO predictions.
(ii) For this “best fit” value $`\omega _{\text{static}}^{\text{best}}`$ the 3PN LSO predictions exhibit the expected convergence property that $`|Q_{3\mathrm{P}\mathrm{N}}Q_{2\mathrm{P}\mathrm{N}}|`$ is significantly smaller than $`|Q_{2\mathrm{P}\mathrm{N}}Q_{1\mathrm{P}\mathrm{N}}|`$ (see Table II below).
We have checked that these two remarkable properties hold for all values of the parameter $`\nu \frac{1}{4}`$. Actually there are several other ways of selecting the approximate value (140), i.e. of understanding why it plays a special role. First, we have seen above that the precise value (131) (which is near the “best fit” value (140)), played a special role in simplifying not only $`a_4`$ but also all the other 3PN coefficients. Second, it seems natural to expect that the true value of $`\omega _{\text{static}}`$ will be such that the full Taylor expansions of most of the invariant functions will be smooth deformations of their test-mass limits. A minimum requirement for this property of “structural stability” under the turning-on of the parameter $`\nu `$ seems to be that the Taylor coefficients of the functions $`e(x)`$, $`j^2(x)`$, $`K(y)`$, and $`1/A(u)`$ do not change sign as $`\nu `$ varies from 0 to 1/4 (we only consider functions with infinitely many non-zero Taylor coefficients in the test-mass limit). \[One could actually impose a more restricted bound on the $`\nu `$variation of the 3PN coefficients, especially given the information that the 2PN coefficients are found to vary by a smallish fractional amount.\] We find that this minimum requirement is satisfied only if $`\omega _{\text{static}}<0.62`$ (the consideration of the expansion of $`e(x)`$ gives the most stringent bound). Another natural requirement for “structural stability” under $`\nu 0`$ would be to impose that all the near-diagonal Padés (without restricting oneself, as above, to the most “robust” ones) of the functions above exhibit real poles that are smoothly connected to their test-mass counterparts. This requirement is most stringent when considering $`P_2^2(1/A(u))`$, and yields the limit $`\omega _{\text{static}}<8.35`$. Combining this with the general limits Eq. (49) suggests that $`\omega _{\text{static}}`$ lies within the small range $`10<\omega _{\text{static}}<8.35`$.
Let us quote a last way of selecting the value (140). It consists in comparing the 3PN Taylor coefficients of the invariant functions that contain only a finite number of terms in the test-mass limit. For instance, consider $`A(u)`$ and $`1/j^2(x)`$. In the test-mass limit $`A(u)=12u`$ and $`1/j^2(x)=x3x^2`$. When $`\nu 0`$ there will appear further powers of $`u`$ or $`x`$ with coefficients vanishing with $`\nu `$. At the 3PN level there is a term $`a_4u^4`$ in $`A(u)`$, and a term $`\frac{8}{3}b_4x^4`$ in $`1/j^2(x)`$. \[The factor 8/3 is introduced to have the same (linear) dependence on $`\omega _{\text{static}}`$ in $`a_4`$ and $`b_4`$.\] These two coefficients read
$`a_4(\nu )`$ $`=`$ $`\left({\displaystyle \frac{94}{3}}{\displaystyle \frac{41}{32}}\pi ^2\right)\nu +2\omega _{\text{static}}\nu `$ (141)
$``$ $`18.6879\nu +2\omega _{\text{static}}\nu `$ (142)
and
$`b_4(\nu )`$ $`=`$ $`\left({\displaystyle \frac{5269}{192}}{\displaystyle \frac{41}{32}}\pi ^2\right)\nu {\displaystyle \frac{61}{32}}\nu ^2+{\displaystyle \frac{1}{216}}\nu ^3+2\omega _{\text{static}}\nu `$ (143)
$``$ $`14.7973\nu 1.9063\nu ^2+0.00463\nu ^3+2\omega _{\text{static}}\nu .`$ (144)
Let us first note that the terms $`\nu ^2`$ and $`\nu ^3`$ in Eq. (143) are numerically nearly negligible. Forgetting about them (i.e. working with $`b_4^{}(\frac{5269}{192}\frac{41}{32}\pi ^2)\nu `$), we then see, by comparing Eqs. (141) and (143), that the coefficients $`a_4(\nu )`$ and $`b_4^{}(\nu )`$ are approximately identical. In particular, this means that there will be a small range of values of $`\omega _{\text{static}}`$ for which $`a_4`$ and $`b_4^{}`$ will be simultaneously small. The existence of this range explains why the $`j`$\- and effective one-body methods can give numerically similar results. We can then make an analytical estimate of the ‘best’ value of the ambiguity parameter $`\omega _{\text{static}}`$ by looking for the value of $`\omega _{\text{static}}`$ which simultaneously minimizes (in a least square sense) $`a_4`$ and $`b_4^{}`$. It is easily seen that the expression $`[a_4(\nu )]^2+[b_4^{}(\nu )]^2`$ attains its minimal value (as function of $`\omega _{\text{static}}`$) for
$$\omega _{\text{static}}^{\text{min}}=\frac{41}{64}\pi ^2\frac{11285}{768}8.37.$$
(145)
This numerical value is not very far the special values selected by the other arguments discussed above.
Summarizing: several (partially) independent arguments suggest that the true value of $`\omega _{\text{static}}`$ lies in the range $`10\omega _{\text{static}}8`$. For definiteness, and for the purpose of the following discussion, we shall henceforth assume that the “correct” value of $`\omega _{\text{static}}`$ is
$$\omega _{\text{static}}=\omega _{\text{static}}^{}.$$
(146)
In Fig. 2 we have included vertical lines corresponding to $`\omega _{\text{static}}=\omega _{\text{static}}^{}`$, to show visually that Eq. (146) is well compatible with our argument based on the convergence of the various methods.
One can also see in Fig. 2 that the curves related to the $`e`$\- and $`j`$-methods have a second intersection point, besides the one around $`\omega _{\text{static}}^{\text{best}}9`$. However, we have checked that the value of $`\omega _{\text{static}}`$ at this point strongly depends on the value of the parameter $`\nu `$. For this reason, and also for the fact that this point does not give an agreement with the “effective” method, we do not take this second intersection point as evidence for a different value of $`\omega _{\text{static}}`$.
Admitting (for the sake of the following argument) Eq. (146) we wish to propose a further way of improving the accuracy of the predictions of LSO observables. Indeed, if one has at one’s disposal three successive approximations, namely the 1PN, 2PN, and 3PN estimates of some quantity $`Q_{\mathrm{LSO}}`$, one can combine this information to refine the estimate of the (unknown) exact value $`Q_{\mathrm{LSO}}^X`$. The rationale for this is to assume that the approach to the limit, when the order $`n`$ of the approximant increases, is approximately described by Eq. (139) (i.e. essentially that the inaccuracy of the $`n`$th estimate decreases proportionally to the $`(n+1)`$th power of some constant $`cbx_{\mathrm{LSO}}<1`$). Then, under this assumption the knowledge of three (successive) approximants, say $`Q_{n1}`$, $`Q_n`$, and $`Q_{n+1}`$, gives three equations ($`Q_m=Q_X+ac^{m+1}`$) for the three unknowns $`(Q_X,a,c)`$. One can solve this system of equations and deduce, in particular, the value of the looked for $`n\mathrm{}`$ limit $`Q_X`$ in terms of $`Q_{n1}`$, $`Q_n`$, and $`Q_{n+1}`$. The result defines the so-called “Shanks transformation” , namely
$$Q_XS_n[Q]\frac{Q_{n+1}Q_{n1}Q_n^2}{Q_{n+1}+Q_{n1}2Q_n}.$$
(147)
When one disposes of more than three $`Q_m`$’s, the Shanks transformation associates to the original (truncated) sequence $`(Q_1,Q_2,\mathrm{},Q_N)`$ a shorter, but hopefully faster converging sequence $`(S_2[Q],\mathrm{},S_{N1}[Q])`$. In our case, the Shanks procedure associates to any triplet of LSO quantities $`(Q_1,Q_2,Q_3)(Q_{1\mathrm{P}\mathrm{N}}^{\mathrm{LSO}},Q_{2\mathrm{P}\mathrm{N}}^{\mathrm{LSO}},Q_{3\mathrm{P}\mathrm{N}}^{\mathrm{LSO}})`$ a single number,
$$Q_S^{\mathrm{LSO}}\frac{Q_{3\mathrm{P}\mathrm{N}}^{\mathrm{LSO}}Q_{1\mathrm{P}\mathrm{N}}^{\mathrm{LSO}}(Q_{2\mathrm{P}\mathrm{N}}^{\mathrm{LSO}})^2}{Q_{3\mathrm{P}\mathrm{N}}^{\mathrm{LSO}}+Q_{1\mathrm{P}\mathrm{N}}^{\mathrm{LSO}}2Q_{2\mathrm{P}\mathrm{N}}^{\mathrm{LSO}}},$$
(148)
which is a (hopefully better) estimate of the (unknown) exact value $`Q_{\mathrm{LSO}}^X`$. We shall refer to (148) as the $`S`$-estimate of $`Q^{\mathrm{LSO}}`$.
In Table II (see also Fig. 3) we apply this procedure to our two best methods: the $`j`$-method and the effective-one-body one, under the assumption (146) (which is needed to exhibit a visible convergence among the first three PN approximations).
Given our present (incomplete) knowledge we consider that the $`S`$-estimates exhibited in Table II represent our best estimates of LSO observables. To verify the plausiblity of these estimates one should resolve the issue of the ambiguous coefficient $`\omega _{\text{static}}`$ in the 3PN dynamics. \[In principle, this can be done by implementing the matching method described in and used there at the 2PN level.\] If this resolution approximately confirms the estimate (146) the $`S`$-estimates will be confirmed. If a very different value of $`\omega _{\text{static}}`$ is found, it might still be compatible with a less evidently convergent PN sequence. And hopefully, the $`S`$-estimate (148) of this new sequence will give an improved 3PN-accurate estimate of LSO observables.
Under the assumption that the $`S`$-estimates are accurate, there are several interesting conclusions that we can draw. First, we remark that the final estimates are quite near the 2PN-level predictions of the effective one-body approach, see Fig. 3. Although this may seem disappointing (an enormous, not yet completed, 3PN work leading to a confirmation of 2PN estimates), this would be a scientifically very useful conclusion. Indeed, this would (in our minds at least) establish the soundness of the philosophy advocated in Refs. and here, namely that resummation methods can be meaningfully employed to make analytical predictions concerning physics near the Last Stable Orbit. This would then also give support to the recent work of Buonanno and Damour in which the 2PN effective one-body Hamiltonian has been used, together with Padé-resummed estimates of gravitational-radiation damping, to study the transition between the inspiral motion and the final plunge of a binary system. \[Let us note, in passing, that this work shows that, though it is crucial to have good initial estimates of the LSO quantities defined by the Hamiltonian, the final observable effects linked to the presence of an LSO are blurred by radiation-reaction effects.\]
## Acknowledgments
This work was supported in part by the KBN Grant No. 2 P03B 094 17 (to P.J.) and the Max-Planck-Gesellschaft Grant No. 02160-361-TG74 (to G.S.). P.J. and G.S. thank the Institut des Hautes Études Scientifiques for hospitality during crucial stages of the collaboration.
## A 3PN effective “geodesic” one-body dynamics
We consider here an effective ‘relativistic’ one-body Hamiltonian $`\widehat{H}_{\text{eff}}^\mathrm{R}`$ of the simple “geodesic” form
$$\widehat{H}_{\mathrm{eff}}^\mathrm{R}(𝐪^{},𝐩^{})=\sqrt{A(q^{})\left[1+𝐩^2+\left(\frac{A(q^{})}{D(q^{})}1\right)(𝐧^{}𝐩^{})^2\right]}.$$
(A1)
The Hamiltonian $`\widehat{H}_{\text{eff}}^\mathrm{R}`$ is related to the real ‘non-relativistic’ Hamiltonian $`\widehat{H}_{\mathrm{real}}^{\mathrm{NR}}`$, Eq. (91), through the constraint equation
$$\left[\widehat{H}_{\mathrm{eff}}^\mathrm{R}\mathbf{(}q^{}(q,p^{}),p^{}\mathbf{)}\right]^2=\left\{1+\widehat{H}_{\mathrm{real}}^{\mathrm{NR}}\mathbf{(}q,p(q,p^{})\mathbf{)}\left[1+\alpha _1\widehat{H}_{\mathrm{real}}^{\mathrm{NR}}+\alpha _2\left(\widehat{H}_{\mathrm{real}}^{\mathrm{NR}}\right)^2+\alpha _3\left(\widehat{H}_{\mathrm{real}}^{\mathrm{NR}}\right)^3\right]\right\}^2.$$
(A2)
As in the text both sides of Eq. (A2) are written in terms of the variables $`q`$ and $`p^{}`$ by means of Eqs. (88) with a generating function $`G`$ of the form
$$G(q,p^{})=G_{1\mathrm{P}\mathrm{N}}(q,p^{})+G_{2\mathrm{P}\mathrm{N}}(q,p^{})+G_{3\mathrm{P}\mathrm{N}}(q,p^{}),$$
(A3)
where
$`G_{1\mathrm{P}\mathrm{N}}(q,p^{})`$ $`=`$ $`(𝐪𝐩^{})\left(g_1𝐩^2+{\displaystyle \frac{g_2}{q}}\right),`$ (A5)
$`G_{2\mathrm{P}\mathrm{N}}(q,p^{})`$ $`=`$ $`(𝐪𝐩^{})\left[b_1𝐩^4+{\displaystyle \frac{1}{q}}\mathbf{\left(}b_2𝐩^2+b_3(𝐧𝐩^{})^2\mathbf{\right)}+{\displaystyle \frac{b_4}{q^4}}\right],`$ (A6)
$`G_{3\mathrm{P}\mathrm{N}}(q,p^{})`$ $`=`$ $`(𝐪𝐩^{})\left[c_1𝐩^6+{\displaystyle \frac{1}{q}}\mathbf{\left(}c_2𝐩^4+c_3𝐩^2(𝐧𝐩^{})^2+c_4(𝐧𝐩^{})^4\mathbf{\right)}+{\displaystyle \frac{1}{q^2}}\mathbf{\left(}c_5𝐩^2+c_6(𝐧𝐩^{})^2\mathbf{\right)}+{\displaystyle \frac{c_7}{q^3}}\right].`$ (A7)
Written explicitly, the constraint equation (A2) is equivalent to a sytem of 23 equations for the 23 unknowns: $`a_1`$, …, $`a_4`$; $`d_1`$, $`d_2`$, $`d_3`$; $`\alpha _1`$, $`\alpha _2`$, $`\alpha _3`$; $`g_1`$, $`g_2`$; $`b_1`$, …, $`b_4`$; $`c_1`$, …, $`c_7`$. The unique solution to these equations reads
$`a_1`$ $`=`$ $`2,`$ (A9)
$`a_2`$ $`=`$ $`{\displaystyle \frac{(3\nu 4)\nu }{2(52\nu )}},`$ (A10)
$`a_3`$ $`=`$ $`{\displaystyle \frac{(16001576\nu +392\nu ^29\nu ^3)\nu }{16(52\nu )^2}},`$ (A11)
$`a_4`$ $`=`$ $`{\displaystyle \frac{(42803349\nu +692\nu ^29\nu ^3)\nu }{6(52\nu )^2}}{\displaystyle \frac{41\pi ^2}{32}}\nu +2\omega _{\text{static}}\nu ,`$ (A12)
$`d_1`$ $`=`$ $`{\displaystyle \frac{(3\nu 4)\nu }{2(52\nu )}},`$ (A13)
$`d_2`$ $`=`$ $`{\displaystyle \frac{(2400+1936\nu 408\nu ^2+9\nu ^3)\nu }{16(52\nu )^2}},`$ (A14)
$`d_3`$ $`=`$ $`{\displaystyle \frac{(486400+703680\nu 383904\nu ^2+93704\nu ^38580\nu ^427\nu ^5)\nu }{64(52\nu )^3}},`$ (A15)
$`\alpha _1`$ $`=`$ $`{\displaystyle \frac{(4\nu 3)\nu }{2(52\nu )}},`$ (A16)
$`\alpha _2`$ $`=`$ $`{\displaystyle \frac{(8032\nu +7\nu ^2)(3\nu 4)\nu }{4(52\nu )^2}},`$ (A17)
$`\alpha _3`$ $`=`$ $`{\displaystyle \frac{(36503660\nu +1829\nu ^2421\nu ^3+44\nu ^4)(3\nu 4)\nu }{8(52\nu )^3}},`$ (A18)
$`g_1`$ $`=`$ $`{\displaystyle \frac{(\nu 6)\nu }{4(52\nu )}},`$ (A19)
$`g_2`$ $`=`$ $`{\displaystyle \frac{5(42\nu +\nu ^2)}{4(52\nu )}},`$ (A20)
$`b_1`$ $`=`$ $`{\displaystyle \frac{(210178\nu +50\nu ^23\nu ^3)\nu }{16(52\nu )^2}},`$ (A21)
$`b_2`$ $`=`$ $`{\displaystyle \frac{(350+419\nu 149\nu ^2+13\nu ^3)\nu }{8(52\nu )^2}},`$ (A22)
$`b_3`$ $`=`$ $`{\displaystyle \frac{(12081\nu +38\nu ^26\nu ^3)\nu }{8(52\nu )^2}},`$ (A23)
$`b_4`$ $`=`$ $`{\displaystyle \frac{200+200\nu 816\nu ^2+360\nu ^349\nu ^4}{32(52\nu )^2}},`$ (A24)
$`c_1`$ $`=`$ $`{\displaystyle \frac{(590012700\nu +8270\nu ^22082\nu ^3+192\nu ^411\nu ^5)\nu }{64(52\nu )^3}},`$ (A25)
$`c_2`$ $`=`$ $`{\displaystyle \frac{(25850+44320\nu 24876\nu ^2+5157\nu ^3283\nu ^4+13\nu ^5)\nu }{32(52\nu )^3}},`$ (A26)
$`c_3`$ $`=`$ $`{\displaystyle \frac{(25200+21360\nu 4364\nu ^21077\nu ^3+408\nu ^447\nu ^5)\nu }{96(52\nu )^3}},`$ (A27)
$`c_4`$ $`=`$ $`{\displaystyle \frac{(21601834\nu +900\nu ^2228\nu ^3+23\nu ^4)\nu ^2}{32(52\nu )^3}},`$ (A28)
$`c_5`$ $`=`$ $`{\displaystyle \frac{(247000407080\nu +225416\nu ^249240\nu ^3+3730\nu ^4151\nu ^5)\nu }{128(52\nu )^3}},`$ (A29)
$`c_6`$ $`=`$ $`{\displaystyle \frac{7(1030016360\nu +7136\nu ^21372\nu ^3+154\nu ^47\nu ^5)\nu }{192(52\nu )^3}},`$ (A30)
$`c_7`$ $`=`$ $`{\displaystyle \frac{(962800+1472880\nu 813024\nu ^2+206500\nu ^325110\nu ^4+1479\nu ^5)\nu }{384(52\nu )^3}}{\displaystyle \frac{\pi ^2}{64}}\nu .`$ (A31)
As said in the text, in view of the complexity of these results, we do not take this possibility seriously. We prefer to it the non-minimal (“non-geodesic”) Hamiltonian given in Sec. IV.
## B 2PN results for the conformally-flat truncation of general relativity
By contrast, let us note that other approximation philosophies are, in our opinion, less reliable to make predictions concerning the LSO. We have in mind here: (i) the use of non-resummed (or only partially resummed) PN expansions, and (ii) the “Wilson-Mathews”-type truncation of Einstein’s theory, in which the spatial metric is taken to be conformally flat. As an example of the first philosophy, let us consider the proposal of Kidder, Will, and Wiseman to partially resum the Damour-Deruelle equations of motion by separating out (and resumming) the “Schwarzschild” $`(\nu =0)`$ terms. This approach led to the prediction (at 2PN) that the LSO is significantly less bound (when $`\nu =1/4`$) than the “Schwarzschild” limit. In terms of orbital frequency at the LSO, Ref. predicts $`\widehat{\omega }_{\mathrm{LSO}}(1/4)0.891<1`$. This contrasts very much with our 2PN and 3PN estimates above which consistently indicate that $`\widehat{\omega }_{\mathrm{LSO}}(\nu )`$ is larger than one (and that the LSO is more bound than its Schwarzschild limit: $`E_{\mathrm{LSO}}/|E_{\mathrm{LSO}}^{\mathrm{Schw}}|<1`$). Independently of this (biassed) argument, we think that both the robustness and the consistency of the “hybrid” approach of are seriously in doubt. Indeed, Refs. and have shown that the hybrid approach was robust neither under a change of formulation (Hamiltonian versus equations-of-motion), nor under a change of coordinate system. Moreover, Ref. has questioned the consistency of this approach by pointing out that the non-resummed “$`\nu `$-corrections” represent, in several cases, a very large (larger than 100%) modification of the corresponding $`\nu `$-independent terms.
Regarding the conformally-flat truncation it was noted by Rieth that this implies significant deviations from the Einstein dynamics already at the 2PN level. We have investigated this question further. In Ref. we gave the invariant functions defined (at 2PN accuracy) by the Wilson-Mathews-type truncation. Applying now our $`j`$-method, we find (at 2PN)
$$j_{\mathrm{WM}}^2(x)=\frac{1}{x}\left[1+\frac{1}{3}\left(9+\nu \right)x+\frac{1}{36}\left(324333\nu 50\nu ^2\right)x^2\right],$$
(B1)
whose Padeed form is
$$j_{\mathrm{WM}P_2}^2(x)P_1^1\left[T_2[j_{\mathrm{WM}}^2(x)]\right]=\frac{1+\frac{1}{9}\nu +\left(\frac{15}{4}\nu +\frac{1}{2}\nu ^2\right)x}{x\mathbf{\left(}1+\frac{1}{9}\nu \left(3\frac{37}{12}\nu \frac{25}{54}\nu ^2\right)x\mathbf{\right)}}.$$
(B2)
This leads to a prediction for the 2PN $`x_{\mathrm{LSO}}`$ which can be written down analytically
$$6x_{\mathrm{LSO}}^{j_{\mathrm{WM}P_2}}(\nu )=\frac{8(9+\nu )}{3(15+2\nu )\nu }\left[\frac{2(9+\nu )}{\sqrt{324333\nu 50\nu ^2}}1\right].$$
(B3)
The corresponding dimensionless orbital frequency, for $`\nu =1/4`$, equals $`\widehat{\omega }_{\mathrm{LSO}}=1.4378`$, the reduced binding energy $`E_{\mathrm{LSO}}/|E_{\mathrm{LSO}}^{\mathrm{Schw}}|=1.2253`$, and the reduced angular momentum $`j_{\mathrm{LSO}}/j_{\mathrm{LSO}}^{\mathrm{Schw}}=0.9293`$.
We have also studied, at the 2PN level, the effective one-body method for the Wilson-Mathews dynamics. Using the procedure decribed in Sec. IV D above, imposing at the 1PN level the condition $`d_1=0`$, we have found that the effective-metric function $`A_{\mathrm{WM}}(u)`$ at the 2PN accuracy reads
$$A_{\mathrm{WM}}(u)=12u+a_3(\nu )u^3+𝒪(u^4),$$
(B4)
where
$$a_3(\nu )=\frac{1}{4}(185\nu )\nu .$$
(B5)
We have improved the behaviour of the 2PN expansion of $`A_{\mathrm{WM}}(u)`$ by Padeeing it:
$$A_{\mathrm{WM}P_2}(u)P_2^1\left[T_3[A_{\mathrm{WM}}(u)]\right]=\frac{1\left(2\frac{9}{8}\nu +\frac{5}{16}\nu ^2\right)u}{1+\left(\frac{9}{8}\nu \frac{5}{16}\nu ^2\right)u+\left(\frac{9}{4}\nu \frac{5}{8}\nu ^2\right)u^2}.$$
(B6)
To extract the LSO quantities from this Padéed $`A`$, we have considered the inflection point of the effective radial potential $`A_{\mathrm{WM}P_2}(u)(1+j^2u^2)`$, which defines the angular momentum $`j_{\mathrm{LSO}}(\nu )`$ and the location $`u_{\mathrm{LSO}}(\nu )`$ of the LSO; then using Eqs. (47) and (48) one calculates the energy and the orbital frequency of the LSO. The results, for $`\nu =1/4`$, are: dimensionless orbital frequency $`\widehat{\omega }_{\mathrm{LSO}}=1.1482`$, reduced binding energy $`E_{\mathrm{LSO}}/|E_{\mathrm{LSO}}^{\mathrm{Schw}}|=1.0972`$, and reduced angular momentum $`j_{\mathrm{LSO}}/j_{\mathrm{LSO}}^{\mathrm{Schw}}=0.9647`$. Let us also mention that Cook, using another conformally flat approximation , obtained, for $`\nu =1/4`$, the following LSO parameters: the dimensionless orbital frequency $`\widehat{\omega }_{\mathrm{LSO}}=2.528`$, the reduced binding energy $`E_{\mathrm{LSO}}/|E_{\mathrm{LSO}}^{\mathrm{Schw}}|=1.579`$, and the reduced angular momentum $`j_{\mathrm{LSO}}/j_{\mathrm{LSO}}^{\mathrm{Schw}}=0.8591`$. Similar results have been obtained in Ref. .
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# Singular Projective Varieties and Quantization
## Introduction
Compact Kähler manifolds which are quantizable, i.e. which admit a holomorphic line bundle with curvature form equal to the Kähler form (a so called quantum line bundle) are projective algebraic manifolds. This means that with the help of the global holomorphic sections of a suitable tensor power of the quantum line bundle they can be embedded into a projective space of certain dimension. Submanifolds of the projective space are always projective varieties, i.e. can be given as zero sets of finitely many homogeneous polynomials. As will be explained in this contribution the basic objects in the set-up of Berezin-Toeplitz (or equivalently geometric) quantization of quantizable compact Kähler manifolds can be completely described inside this algebraic-geometric context. For example the quantum Hilbert space will be essentially the projective coordinate ring of the variety.
By definition Kähler manifolds are nonsingular hence the varieties obtained are nonsingular. But from the point of view of varieties the singular ones are on equal footing. Hence one might expect that it is possible to find a direct way towards quantization of singular spaces by exploiting the theory of varieties.
In this contribution I do not present a solution for the quantization of singular spaces. I will only explain the above mentioned path from compact quantizable Kähler manifolds to projective varieties. The quantization procedure I am considering is the Berezin-Toeplitz quantization, resp. the Berezin-Toeplitz deformation quantization. This quantization procedure is adapted to the complex structure which is a requirement for the fact that it can be formulated in terms of complex algebraic geometry. I recall the results on this quantization scheme in Section 1.
The rest of the contribution is considered to be tutorial. There is nothing new there, and everything is well-known for researchers working in algebraic geometry. But I hope that the collection of concepts and results will be useful for researcher in quantization. Some concepts used elsewhere in this volume are explained. In Section 2 and in Section 3 basic concepts of algebraic geometry are introduced. First projective varieties are defined. Compactified moduli spaces are candidates for projective varieties. The projective (homogeneous) coordinate ring is discussed. It will turn out to be the quantum Hilbert space of the theory. It incorporates the vector space of global holomorphic sections of all tensor power of the quantum line bundle at once. On this quantum Hilbert space the total Berezin-Toeplitz quantization operator operates. It is used to show that the quantization scheme has the correct semi-classical limit and to prove the existence of an associated deformation quantization. As already pointed out, quantizable Kähler manifolds can be embedded with the help of the very ample quantum line bundle into projective space (as complex manifolds, not necessarily as Kähler manifolds). Such embeddings are discussed in detail in Section 2.
Projective varieties are not necessarily smooth, they can have singularities. After giving some examples of singular varieties in Section 2 (e.g. the singular cubic curves) singularities are treated in more detail in Section 3. Beside the definition of a singular point using the rank of the Jacobi matrix of the defining equations for the variety, a more intrinsic definition in terms of the local ring $`𝒪_{V,\alpha }`$ of a point $`\alpha `$ on the variety $`V`$ and the Zariski tangent space at the point $`\alpha `$ is given. In terms of algebraic properties of the local ring a hierarchy of types of singularities can be introduced. As special examples normal singularities are discussed. Whereas on an arbitrary singular variety the set of singular points can have codimension one, on a normal variety (i.e. on a variety where all local rings are normal rings, see the definition below) this subset has codimension at least two. The singularities of moduli spaces are very often normal singularities.
Typically, moduli spaces are obtained by dividing out a group action on a nonsingular variety. The main question is whether it is possible to define a geometric structure on the orbit space, i.e. whether there exists some algebraic geometric quotients. The “Geometric Invariant Theory (GIT)” as developed by Mumford gives a powerful tool how to deal with such quotients. If one considers only certain suitable subsets of points of the variety the group is acting on (i.e. the subset of semi-stable, or stable points) one obtains a good quotient (which is also a categorical quotient), resp. a geometric quotient. They will carry a compatible structure of a projective variety, resp. of an open subset of a projective variety. This will be explained in Section 4. Also there the results on the relation with the symplectic quotients obtained via moment maps and symplectic reduction due to Kirwan, Kempf and Ness will be explained. These results are taken from the appendix to , written by Kirwan. Roughly speaking, the geometric quotient and the symplectic quotient coincides on the regular points of the symplectic reduction (see Theorem 4.5 for a precise statement). But in general the singularity structure will differ.
## 1. From quantizable compact Kähler manifolds to projective varieties
Let $`(M,\omega )`$ be a Kähler manifold, i.e. $`M`$ a complex manifold and $`\omega `$ a Kähler form on $`M`$. In this contribution I will only consider compact Kähler manifolds. If nothing else is said compactness is assumed. A further data we need is the triple $`(L,h,)`$, with a holomorphic line bundle $`L`$ on $`M`$, a hermitian metric $`h`$ on $`L`$ (with the convention that it is conjugate linear in the first argument) and a connection $``$ compatible with the metric on $`L`$ and the complex structure. With respect to local holomorphic coordinates of the manifold and with respect to a local holomorphic frame for the bundle the metric $`h`$ can be given as
$$h(s_1,s_2)(x)=\widehat{h}(x)\overline{\widehat{s}_1}(x)\widehat{s}_2(x),$$
(1)
where $`\widehat{s}_i`$ is a local representing function for the section $`s_i`$ ($`i=1,2`$) and $`\widehat{h}`$ is a locally defined real-valued function on $`M`$. The compatible connection is uniquely defined and is given in the local coordinates as $`=+(\mathrm{log}\widehat{h})+\overline{}`$. The curvature of $`L`$ is defined as the two-form
$$curv_{L,}(X,Y):=_X_Y_Y_X_{[X,Y]},$$
(2)
where $`X`$ and $`Y`$ are vector fields on $`M`$. In the local coordinates the curvature can be expressed as $`curv_{L,}=\overline{}\mathrm{log}\widehat{h}=\overline{}\mathrm{log}\widehat{h}`$.
A Kähler manifold $`(M,\omega )`$ is called *quantizable* if there exists such a triple $`(L,h,)`$ which obeys
$$curv_{L,}(X,Y)=\mathrm{i}\omega (X,Y).$$
(3)
The condition (3) is called the (pre)quantum condition. The bundle $`(L,h,)`$ is called a (pre)quantum line bundle. Usually we will drop $``$ and sometimes also $`h`$ in the notation.
For the following we assume $`(M,\omega )`$ to be a quantizable Kähler manifold with quantum line bundle $`(L,h,)`$.
There is an important observation. If $`M`$ is a compact Kähler manifold which is quantizable then from the prequantum condition (3) we obtain for the Chern form of the line bundle the relation
$$c(L):=\frac{\mathrm{i}}{2\pi }curv_{L,}=\frac{\omega }{2\pi }.$$
(4)
This implies that $`L`$ is a positive line bundle. In the terminology of algebraic geometry it is an ample line bundle, see Definition 2.5 for the definition of ampleness. By the Kodaira embedding theorem $`M`$ can be embedded (as algebraic submanifold) into projective space $`^N()`$ using a basis of the global holomorphic sections $`s_i`$ of a suitable tensor power $`L^{m_0}`$ of the bundle $`L`$
$$z(s_0(z):s_1(z):\mathrm{}:s_N(z))^N().$$
(5)
These algebraic submanifolds can be described as zero sets of homogeneous polynomials, i.e. they are projective varieties. Note that the dimension of the space $`\mathrm{\Gamma }_{hol}(M,L^{m_0})`$ consisting of the global holomorphic sections of $`L^{m_0}`$, can be determined by the Theorem of Grothendieck-Hirzebruch-Riemann-Roch, see , . By passing to the Kähler form $`m_0\omega `$ and to the associated quantum line bundle $`L^{m_0}`$ we might assume that the sections of our quantum line bundle do already the embedding (i.e that it is already very ample).
So even if we start with an arbitrary Kähler manifold the quantization condition will force the manifold to be an algebraic manifold and we are in the realm of algebraic geometry. This should be compared with the fact that there are “considerable more” Kähler manifolds than algebraic manifolds.
In Section 2 I will explain what projective varieties are. But first I like to introduce the quantum operator we are dealing with. We take $`\mathrm{\Omega }=\frac{1}{n!}\omega ^n`$ as volume form on $`M`$. On the space of $`C^{\mathrm{}}`$ sections of the bundle $`L`$ we have the scalar product
$$\phi ,\psi :=_Mh(\phi ,\psi )\mathrm{\Omega },\phi :=\sqrt{\phi ,\phi }.$$
(6)
Let $`\mathrm{L}^2(M,L)`$ be the L<sup>2</sup>-completion of the space of $`C^{\mathrm{}}`$-sections of the bundle $`L`$ and $`\mathrm{\Gamma }_{hol}(M,L)`$ be its (due to compactness of $`M`$) finite-dimensional closed subspace of holomorphic sections. Let $`\mathrm{\Pi }:\mathrm{L}^2(M,L)\mathrm{\Gamma }_{hol}(M,L)`$ be the projection.
###### Definition 1.1.
For $`fC^{\mathrm{}}(M)`$ the Toeplitz operator $`T_f`$ is defined to be
$$T_f:=\mathrm{\Pi }(f):\mathrm{\Gamma }_{hol}(M,L)\mathrm{\Gamma }_{hol}(M,L).$$
In words: We multiply the holomorphic section with the differentiable function $`f`$. This yields only a differentiable section. To obtain a holomorphic section again, we project it back to the subspace of global holomorphic sections. From the point of view of Berezin’s approach , $`T_f`$ is the operator with contravariant symbol $`f`$.
The linear map
$$T:C^{\mathrm{}}(M)\mathrm{End}\left(\mathrm{\Gamma }_{hol}(M,L)\right),fT_f,$$
is called the Berezin-Toeplitz quantization. Recall that $`(C^{\mathrm{}}(M),,\{.,.\})`$ is a Poisson algebra. To define the Poisson bracket (i.e. a compatible Lie algebra structure) on $`C^{\mathrm{}}(M)`$ we use the Kähler form $`\omega `$ as symplectic form and define $`\{f,g\}:=\omega (X_f,X_g)`$ where $`X_f`$ is the Hamiltonian vector field assigned to $`fC^{\mathrm{}}(M)`$ given by $`\omega (X_f,.)=df(.)`$. The Berezin-Toeplitz quantization map is neither a Lie algebra homomorphism nor an associative algebra homomorphism, because in general
$$T_fT_g=\mathrm{\Pi }(f)\mathrm{\Pi }(g)\mathrm{\Pi }\mathrm{\Pi }(fg)\mathrm{\Pi }.$$
Due to the compactness of $`M`$ this defines a map from the commutative algebra of functions to a noncommutative finite-dimensional (matrix) algebra. A lot of information will get lost. To recover this information one should consider not just the bundle $`L`$ alone but all its tensor powers $`L^m`$ for $`m_0`$ and apply all the above constructions for every $`m`$. In this way one obtains a family of matrix algebras and maps
$$T_f^{(m)}:C^{\mathrm{}}(M)\mathrm{End}\left(\mathrm{\Gamma }_{hol}(M,L^m)\right),fT_f^{(m)}.$$
This infinite family should in some sense “approximate” the algebra $`C^{\mathrm{}}(M)`$.(See for a definition of such an approximation.)
If we group all $`T_f^{(m)}`$ together we obtain a map
$$C^{\mathrm{}}(M)\underset{m_0}{}\mathrm{End}(\mathrm{\Gamma }_{hol}(M,L^m))\mathrm{End}(\underset{m_0}{}\mathrm{\Gamma }_{hol}(M,L^m)),$$
(7)
$$fT_f^{()}:=(T_f^{(m)})_{m_0}.$$
(8)
We will see later on that $`_m\mathrm{\Gamma }_{hol}(M,L^m)`$ with a slight modification (i.e. $``$ is replaced by $``$) is the projective coordinate ring of the embedded $`M`$. The operator $`T_f^{()}`$ is called the total Berezin-Toeplitz operator. It operates on the projective coordinate ring.
It was shown by Bordemann, Meinrenken and Schlichenmaier that this quantization scheme has the correct semi-classical behavior and yields an associated star product (a deformation quantization). Denote by $`f_{\mathrm{}}`$ the sup-norm of $`f`$ on $`M`$ and by $`T_f^{(m)}=sup_{s\mathrm{\Gamma }_{hol}(M,L^m),s0}\frac{T_f^{(m)}s}{s}`$ the operator norm on $`\mathrm{\Gamma }_{hol}(M,L^m)`$.
###### Theorem 1.2.
\[Bordemann, Meinrenken, Schlichenmaier\]
(a) For every $`fC^{\mathrm{}}(M)`$ there exists $`C>0`$ such that
$$f_{\mathrm{}}\frac{C}{m}T_f^{(m)}f_{\mathrm{}}.$$
(9)
In particular, $`lim_m\mathrm{}T_f^{(m)}=f_{\mathrm{}}`$.
(b) For every $`f,gC^{\mathrm{}}(M)`$
$$m\mathrm{i}[T_f^{(m)},T_g^{(m)}]T_{\{f,g\}}^{(m)}=O(\frac{1}{m}).$$
(10)
(c) For every $`f,gC^{\mathrm{}}(M)`$
$$T_f^{(m)}T_g^{(m)}T_{fg}^{(m)}=O(\frac{1}{m}).$$
(11)
Let me recall the definition of a star product. Let $`𝒜=C^{\mathrm{}}(M)[[\nu ]]`$ be the algebra of formal power series in the variable $`\nu `$ over the algebra $`C^{\mathrm{}}(M)`$. A product $``$ on $`𝒜`$ is called a (formal) star product if it is an associative $`[[\nu ]]`$-linear product such that
1. $`𝒜/\nu 𝒜C^{\mathrm{}}(M)`$, i.e. $`fgmod\nu =fg`$,
2. $`{\displaystyle \frac{1}{\nu }}(fggf)mod\nu =\mathrm{i}\{f,g\}`$,
where $`f,gC^{\mathrm{}}(M)`$. We can also write
$$fg=\underset{j=0}{\overset{\mathrm{}}{}}\nu ^jC_j(f,g),$$
(12)
with $`C_j(f,g)C^{\mathrm{}}(M)`$. The $`C_j`$ should be $``$-bilinear in $`f`$ and $`g`$. The conditions 1. and 2. can be reformulated as
$$C_0(f,g)=fg,\text{and}C_1(f,g)C_1(g,f)=\mathrm{i}\{f,g\}.$$
(13)
###### Theorem 1.3.
There exists a unique (formal) star product on $`C^{\mathrm{}}(M)`$
$$fg:=\underset{j=0}{\overset{\mathrm{}}{}}\nu ^jC_j(f,g),C_j(f,g)C^{\mathrm{}}(M),$$
(14)
in such a way that for $`f,gC^{\mathrm{}}(M)`$ and for every $`N`$ we have with suitable constants $`K_N(f,g)`$ for all $`m`$
$$T_f^{(m)}T_g^{(m)}\underset{0j<N}{}\left(\frac{1}{m}\right)^jT_{C_j(f,g)}^{(m)}K_N(f,g)\left(\frac{1}{m}\right)^N.$$
(15)
See , and .
It has a couple of nice properties, i.e. (i) $`1f=f1=f`$, (ii) the selfadjointness $`\overline{fg}=\overline{g}\overline{f}`$, and (iii) it admits a naturally defined trace (see ).
As is shown in the star product is a differential star product, i.e. the $`C_j`$ are bidifferential operators and it has the property of “separation of variables” (resp. it is of Wick type ). This says that it respects the holomorphic structure. In more precise terms: if the star product is restricted to open subsets the star multiplication from the right with local holomorphic functions is pointwise multiplication, and the star multiplication from the left with local anti-holomorphic functions is pointwise multiplication.
Let me close this section with two remarks.
###### Remark.
More traditionally one considers the operator $`Q`$ of geometric quantization (with Kähler polarization) defined as $`Q=\mathrm{\Pi }P`$ with
$$P:C^{\mathrm{}}(M)\mathrm{End}(\mathrm{\Gamma }_{\mathrm{}}(M,L)),fP_f:=_{X_f}+\mathrm{i}fid,$$
where $`\mathrm{\Gamma }_{\mathrm{}}(M,L)`$ is the space of $`C^{\mathrm{}}`$ sections of the bundle $`L`$ and $`\mathrm{\Pi }`$ is the projection onto the space of global holomorphic sections. Now $`Q_f\mathrm{End}(\mathrm{\Gamma }_{hol}(M,L))`$. Again one should consider $`Q_f^{(m)}`$ for all $`m_0`$. For compact Kähler manifolds both quantization procedures are related via the Tuynman relation. It reads as
$$Q_f^{(m)}=\mathrm{i}T_{f\frac{1}{2m}\mathrm{\Delta }f}^{(m)}=\mathrm{i}\left(T_f^{(m)}\frac{1}{2m}T_{\mathrm{\Delta }f}^{(m)}\right).$$
(16)
Hence, the $`T_f^{(m)}`$ and the $`Q_f^{(m)}`$ have the same asymptotic behavior.
###### Remark.
There is another kind of embedding of the manifold $`M`$ into projective space. It is the embedding using the coherent states of Berezin-Rawnsley. This embedding turns out to be a special case of the embedding considered at the beginning of this section where one uses a orthogonal basis of the sections, resp. (depending on the conventions) the conjugate of it, see for details.
## 2. Projective varieties
### 2.1. The definition of a projective variety
Let $`𝕂`$ be an algebraically closed field and let us assume for simplicity that its characteristic is zero. Without any harm the reader might even assume $`𝕂=`$. The projective space $`^n=^n(𝕂)`$ is given as the space of lines through the origin in $`𝕂^{n+1}`$, i.e. as the equivalence classes of points in $`𝕂^{n+1}\{0\}`$ where two points $`\alpha `$ and $`\beta `$ are equivalent if $`\lambda 𝕂\{0\}`$ with $`\beta =\lambda \alpha `$. The point $`[\alpha ]`$ in projective space defined by the point $`\alpha =(\alpha _0,\alpha _1,\mathrm{},\alpha _n)`$, $`\alpha 0`$ can be given by its (non-unique) homogeneous coordinates $`[\alpha ]=(\alpha _0:\alpha _1:\mathrm{}:\alpha _n)`$.
Let $`f𝕂[X_0,X_1,\mathrm{},X_n]`$ be a homogeneous polynomial of degree $`k`$. As usual we obtain an associated $`𝕂`$-valued function on $`𝕂^{n+1}`$ by assigning to the point $`\alpha 𝕂^{n+1}`$ the value $`f(\alpha )`$ obtained by “setting” $`X_i`$ to be $`\alpha _i`$. If $`\beta =\lambda \alpha `$ with $`\lambda 𝕂`$, $`\lambda 0`$, is another point in the same equivalence class as $`\alpha `$, then we obtain $`f(\lambda \alpha )=\lambda ^kf(\alpha )`$. In particular, the induced function is only well-defined on the whole projective space if $`k=0`$, i.e. $`f`$ is a constant. But we also see that if $`\alpha `$ is a zero of $`f`$ then any other element $`\beta =\lambda \alpha `$ will be a zero too. Hence the zero-set
$$𝒵(f):=\{[\alpha ]^nf(\alpha )=0\}$$
(17)
is a well-defined subset of $`^n`$. The Zariski topology is the coarsest topology in which the sets $`𝒵(f)`$ are closed subsets for all polynomials $`f`$, or equivalently for which the complements $`D_f=^n𝒵(f)`$ are open sets. Because the zero-sets of polynomials are also closed in the “usual” topology if the base field is $``$ the sets which are closed (open) in the Zariski topology are closed(open) in the “usual” topology. The Zariski topology has a number of quite unusual properties. For example, it is not separated, i.e. two distinct points do not necessarily have disjoint open neighborhoods. Even more is true: every non-empty Zariski open set $`U`$ is automatically dense in $`^n`$.
###### Definition 2.1.
(a) A subset $`W`$ of $`^n`$ is called a (projective) variety if it is the set of common zeros of finitely many homogeneous polynomials $`f_1,f_2,\mathrm{},f_m`$ (which are not necessarily of the same degree)
$$W=𝒵(f_1,f_2,\mathrm{},f_m):=\{[\alpha ]^nf_i(\alpha )=0,i=1,\mathrm{},m\}.$$
(18)
(b) A variety is called a linear variety if it can be given as the zero-set of linear polynomials.
(c) A variety is called irreducible if every decomposition
$$W=V_1V_2$$
(19)
with varieties $`V_1`$ and $`V_2`$ implies that
$$V_1V_2\text{ or }V_2V_1.$$
(20)
A variety which is not irreducible is called reducible.
(d) A Zariski open set of a projective variety is called a quasiprojective (or sometimes just algebraic) variety.
Note that some authors reserve the term variety for irreducible ones.
###### Definition 2.2.
Let $`V`$ be an irreducible variety, then its dimension $`dimV`$ is defined as the maximal length $`n`$ of chains of strict subvarieties which are irreducible
$$\mathrm{}V_0V_1\mathrm{}V_{n1}V_n=V.$$
(21)
For arbitrary varieties the dimension is defined to be the maximum of the dimensions of its irreducible subvarieties.
Subvarieties of dimension 0 are called points, subvarieties of dimension 1 curves, etc.
Let $`V`$ be a projective variety i.e. $`V=𝒵(f_1,f_2,\mathrm{},f_m)`$. Take $`I=(f_1,f_2,\mathrm{}f_m)`$ to be the ideal generated by the polynomials $`f_1,f_2,\mathrm{}f_m`$, i.e.
$$I=\{\underset{i=1}{\overset{m}{}}g_if_ig_i𝕂[X_0,X_1,\mathrm{}X_n],i=1,\mathrm{},k\}.$$
(22)
Obviously $`V=𝒵(I)`$. Ideals which can be generated by homogeneous elements are called homogeneous ideals. Hence, projective varieties can always be given as zero-sets of homogeneous ideals. The converse is also true. Clearly
$$𝒵(I):=\{x^nf(x)=0,fI\}.$$
(23)
is by the homogeneity of the generators a well-defined subset of $`^n`$. Because the polynomial ring is a Noetherian ring, i.e. every ideal $`I`$ can be generated (as ideal in the sense of (22)) by finitely many elements, e.g. $`I=(g_1,g_2,\mathrm{},g_s)`$, we get $`𝒵(I)=𝒵(g_1,g_2,\mathrm{},g_s)`$ and hence $`𝒵(I)`$ is a projective variety in the sense of Definition 2.1. Any other set of generators of the ideal will define the same zero-set. Even the ideal is not fixed uniquely by $`V`$. As a simple example one might consider the hyperplane $`H=𝒵(I)`$ with $`I=(X_0)`$. The same variety might be defined as $`H=𝒵(I^{}))`$ with $`I^{}=(X_0^2)`$. But note that $`I^{}I`$. One might expect that for a given variety $`V`$ there is a largest ideal which still defines $`V`$. This is indeed true.
###### Definition 2.3.
Let $`V`$ be a projective variety, i.e. $`V=𝒵(I)`$ for some ideal $`I`$. The vanishing ideal $`(V)`$ is defined to be
$$(V):=\{f𝕂[X_0,\mathrm{},X_n]f(x)=0,xV\}.$$
(24)
The subset $`(V)`$ is a homogeneous ideal and contains any other defining ideal $`I`$ for $`V`$. It can completely be described in algebraic terms. For this we define for any ideal $`I`$ its radical ideal
$$Rad(I):=\{f𝕂[X_0,\mathrm{},X_n]n:f^nI\}.$$
(25)
If $`I`$ is homogeneous $`Rad(I)`$ will again be homogeneous. We obtain
$$(𝒵(I))=Rad(I),$$
(26)
with only one exception in the case when $`𝒵(I)=\mathrm{}`$. Note that $`\mathrm{}`$ corresponds to two homogeneous radical ideal, the full ring $`𝕂[X_0,\mathrm{},X_n]`$ and the ideal $`I_0:=(X_0,X_1,\mathrm{},X_n)`$. Note that the only possible zero of $`I_0`$ is the point $`0`$ which is not an element of projective space.
There is another warning necessary. One might think that the dimension $`r`$ of a variety is exactly $`nk`$ if $`k`$ is the minimal number of necessary polynomials to generate its vanishing ideal $`I`$. Unfortunately this is not true. The only information one has is that $`rnk`$, with equality if $`k=1`$. A variety is called a complete intersection if indeed $`r=nk`$.
Projective varieties are not always manifolds (of course not all manifolds are projective varieties either). Varieties have not necessarily to be smooth. They might have singularities. In Section 3 I will deal with singularities in more detail. Here I would like to show some non-trivial examples of singular varieties. For this I give a first definition of a singular point. Further definitions will follow in Section 3.
###### Definition 2.4.
Let $`V=𝒵(f_1,f_2,\mathrm{},f_m)`$ be a projective variety of dimension $`r`$ in $`^n`$ with vanishing ideal $`(V)`$ generated by the polynomials $`f_1,f_2,\mathrm{},f_m`$. Consider the $`m\times (n+1)`$-matrix (the Jacobi matrix)
$$J(X)=\left(\begin{array}{cccc}\frac{f_1}{X_0}& \frac{f_1}{X_1}& \mathrm{}& \frac{f_1}{X_n}\\ \frac{f_2}{X_0}& \frac{f_2}{X_1}& \mathrm{}& \frac{f_2}{X_n}\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ \frac{f_m}{X_0}& \frac{f_m}{X_1}& \mathrm{}& \frac{f_m}{X_n}\end{array}\right).$$
(27)
A point $`x`$ on $`V`$ with
$$\mathrm{rank}J(x)<nr$$
(28)
is called a singular point of the variety. A point on a variety which is not a singular point is called a regular point. If $`V`$ has no singular points it is called a non-singular (or smooth, or regular) variety.
For a variety which is the union of two different subvarieties the points where the subvarieties meet are always singular points. As a typical example one might take the variety $`V=𝒵(X_0X_1)`$ in $`^2`$. Then $`V=𝒵(X_0)𝒵(X_1)`$ and the singular point is the point $`(0:0:1)=𝒵(X_0)𝒵(X_1)`$. But even irreducible varieties can have singularities. As an example let me consider the varieties $`Y`$ in $`^2`$ defined by irreducible cubic polynomials. These polynomials can be written (after a suitable change of coordinates) as
$$f(X,Y,Z)=Y^2Z4X^3+g_2XZ^2+g_3Z^3,$$
(29)
with certain elements $`g_2,g_3𝕂`$. The variety $`𝒵(f)`$ is non-singular if and only if the coefficients $`g_2`$ and $`g_3`$ are such that $`g_2^327g_3^20`$. One obtains in this way the elliptic curves. See the Figure 1. Here the curve is defined over $``$ and the real-valued points are plotted. The elliptic curves correspond to 1-dimensional complex tori, see Section 2.2 below.
In the singular cases we obtain two different types of curves. The first one is the nodal cubic $`𝒵(Y^2Z4X^2(X+Z))`$, see Figure 2. The only singular point is the point $`(0:0:1)`$. Moving along the curve we pass through the singular point twice, each time with a different tangent direction. The second one is the cuspidal cubic $`𝒵(Y^2Z4X^3)`$, see Figure 3. Again the point $`(0:0:1)`$ is a singular point. But now there is only one tangent direction at this point.
### 2.2. Embeddings into Projective Space
Let us now take the complex numbers $``$ as base field $`𝕂`$. With the complex topology $`^n`$ is a compact, $`n`$-dimensional complex manifold. A coordinate covering is given by the “affine” sets
$$D_{X_i}:=𝒵(X_i)=\{(\alpha _0:\alpha _1:\mathrm{}:\alpha _n)\alpha _i0\}^n$$
(30)
for $`i=0,1,\mathrm{},n`$. Every nonsingular projective variety is closed in the Zariski topology and hence also in the complex topology and hence a compact submanifold of $`^n`$. An abstract compact complex manifold $`M`$ is called a projective algebraic manifold if there exists an injective holomorphic embedding
$$\mathrm{\Phi }:M^n,$$
(31)
such that $`\mathrm{\Phi }(M)M`$ as complex manifolds. The Theorem of Chow \[7, p.166\] says that in this case $`\mathrm{\Phi }(M)`$ is a nonsingular projective variety, i.e. it can be given as the zero-set of finitely many homogeneous polynomials. This is even true in the strong sense that every meromorphic function on $`\mathrm{\Phi }(M)`$ is a rational function (i.e. it can be expressed as quotient of homogeneous polynomials of the same degree in $`(n+1)`$ variables), every meromorphic differential is a rational differential, and every holomorphic map between two embedded complex manifolds is an algebraic map, i.e. can be given locally as a set of rational functions without poles.
Let me illustrate this in the case of the above introduced elliptic curves. Let $`T=/\mathrm{\Gamma }`$ be the one-dimensional complex torus defined as the quotient of $``$ by the lattice $`\mathrm{\Gamma }:=\{m+n\tau m,n\}`$ for fixed $`\tau `$ with $`\mathrm{Im}\tau >0`$. The associated Weierstraß $`\mathrm{}`$-function and its derivative $`\mathrm{}^{}`$ are doubly-periodic meromorphic functions with respect to the lattice $`\mathrm{\Gamma }`$, i.e.
$$\mathrm{}(z+\omega )=\mathrm{}(z),\mathrm{}^{}(z+\omega )=\mathrm{}^{}(z),\text{for all}\omega \mathrm{\Gamma },z.$$
Hence they are meromorphic functions on $`T`$. The function $`\mathrm{}`$ fulfills the famous differential equation
$$(\mathrm{}^{})^2=4\mathrm{}^3g_2\mathrm{}g_3$$
(32)
with the Eisenstein series
$$g_2:=60\underset{\omega \mathrm{\Gamma },\omega 0}{}\frac{1}{\omega ^4},g_3:=140\underset{\omega \mathrm{\Gamma },\omega 0}{}\frac{1}{\omega ^6}.$$
An embedding of the torus into the projective plane is given by
$$\mathrm{\Psi }:T^2,[z]\{\begin{array}{cc}(\mathrm{}(z):\mathrm{}^{}(z):1),\hfill & [z]0\hfill \\ (0:1:0),\hfill & [z]=0.\hfill \end{array}$$
(33)
Here $`[z]=zmod\mathrm{\Gamma }`$ denotes the point on the torus represented by $`z`$. If one compares the differential equation (33) with the polynomial (29) one sees that $`\mathrm{\Psi }(T)`$ is a cubic curve hence a projective variety (indeed it is nonsingular). Via $`\mathrm{\Psi }`$ the meromorphic function $`\mathrm{}`$ corresponds to the rational function $`X/Z`$ and $`\mathrm{}^{}`$ corresponds to $`Y/Z`$. Note that the field of meromorphic functions on the torus consists of rational expressions in $`\mathrm{}`$ and $`\mathrm{}^{}`$. For more details see, , p.34 and p.62 ff.
After this excursion let me return to the situation discussed in the Section 1. Let $`M`$ be a compact complex manifold and $`\pi :LM`$ a holomorphic line bundle (not necessarily a quantum line bundle). Choose a basis of the global holomorphic sections $`s_0,s_1,\mathrm{},s_n\mathrm{\Gamma }_{hol}(M,L)`$. For every point $`xM`$ there exists an open neighborhood $`U`$ of $`x`$ such that $`L`$ can be locally trivialized over $`U`$, i.e. that there is an (holomorphic) bundle map $`\rho :L_U:=\pi ^1(U)U\times `$. With respect to this trivialization the section $`s_i`$ can be given by a local holomorphic function $`\widehat{s}_i:U`$ defined by $`\rho (s_i(x))=(x,\widehat{s}_i(x))`$. The map
$$U^{n+1},y\stackrel{~}{\mathrm{\Phi }}(y):=(\widehat{s}_0(y),\widehat{s}_1(y),\mathrm{},\widehat{s}_n(y))$$
(34)
is a holomorphic map. It depends not only on the basis chosen, but also on the trivialization. If $`\rho ^{}`$ is a different trivialization defined over the open set $`U`$ where $`\rho `$ is defined (or a subset of it) then $`\rho ^{}\rho ^1(x,\lambda )=(x,g(x)\lambda )`$ with a holomorphic function $`g:U`$ nowhere vanishing on $`U`$. The map $`\stackrel{~}{\mathrm{\Phi }}^{}:U^{n+1}`$ corresponding to $`\rho ^{}`$ fulfills $`\stackrel{~}{\mathrm{\Phi }}^{}(y)=g(y)\stackrel{~}{\mathrm{\Phi }}(y)`$. Hence, $`[\mathrm{\Phi }(y)]:=\stackrel{~}{\mathrm{\Phi }}(y)`$ will be well-defined, i.e. not depend on the trivialization chosen if we assure that $`\stackrel{~}{\mathrm{\Phi }}(y)0`$. But $`\stackrel{~}{\mathrm{\Phi }}(y)=0`$ if and only if $`s(y)=0`$ for all sections $`s\mathrm{\Gamma }_{hol}(M,L)`$. Hence we obtain a well-defined holomorphic map
$$\mathrm{\Phi }:M\{yMs(y)=0,s\mathrm{\Gamma }_{hol}(M,L)\}\stackrel{}{}^n,$$
(35)
obtained by glueing together the local maps $`\stackrel{~}{\mathrm{\Phi }}`$. A change of basis of $`\mathrm{\Gamma }_{hol}(M,L)`$ is given by an element of $`\mathrm{GL}(n+1,)`$. The images of the two mappings obtained by the two set of basis elements are related by the corresponding $`\mathrm{PGL}(n+1,)`$ action. Note that if there exists a nontrivial section $`s`$ (i.e. $`s0`$) then the map $`(\text{35})`$ is defined on a dense open subset of $`M`$.
###### Definition 2.5.
(a) A line bundle $`L`$ is called very ample if the map $`\mathrm{\Phi }`$ (with respect to one and hence to all set of basis elements) is an embedding.
(b) A line bundle $`L`$ is called ample if there exists $`m`$ such that $`L^m`$ is very ample.
It follows that a compact complex manifold is projective algebraic if it admits an ample line bundle. The converse is also true. To see this we first study $`^n`$. Here we have the tautological line bundle whose fiber over the point $`[z]`$ is the complex line trough $`0`$ and the point $`z^{n+1}`$. The hyperplane section bundle $`H`$ is the dual of the tautological line bundle. Its space of global sections is generated by the coordinate functions $`X_0,X_1,\mathrm{},X_n`$, i.e. it can be identified with the space of linear polynomials in $`(n+1)`$ variables. All line bundles over $`^n`$ are given as $`H^m`$ where this denotes for $`m>0`$ the $`m`$-th tensor power of $`H`$, for $`m<0`$ the $`|m|`$-th tensor power of the dual bundle of $`H`$, and for $`m=0`$ the trivial bundle $`𝒪`$. The space of global holomorphic sections of $`H^m`$ can canonically be identified with the space of homogeneous polynomials of degree $`m`$ in $`(n+1)`$ variables. In particular, there exists no nontrivial sections for $`m<0`$. If $`\mathrm{\Phi }:M^n`$ is a holomorphic map then the pullback $`\mathrm{\Phi }^{}H`$ is a holomorphic line bundle on $`M`$. The space of global sections of $`\mathrm{\Phi }^{}H`$ is generated by the pullback $`\mathrm{\Phi }^{}(X_i)=X_i\mathrm{\Phi }`$ of the global sections $`X_i`$, $`i=0,1,\mathrm{},n+1`$. If $`\mathrm{\Phi }`$ is a holomorphic embedding than $`\mathrm{\Phi }^{}H`$ is a very ample line bundle. If the pull-backs of the $`(n+1)`$ sections $`X_i`$ stay linearly independent then $`\mathrm{\Phi }`$ is exactly given by the embedding defined via the bundle $`\mathrm{\Phi }^{}H`$. If not, then the embedding defined via $`\mathrm{\Phi }^{}H`$ goes into a linear subvariety of $`^n`$ of lower dimension, hence in a $`^k`$ for $`k<n`$.
Altogether we see that the embeddings of $`M`$ into projective space correspond to very ample line bundles over $`M`$. The pair $`(M,L)`$ where $`M`$ is a compact complex manifold and $`L`$ is a very ample line bundle is called a polarized projective algebraic manifold.
Note that the same manifold considered with different $`L`$ may “look” quite differently. As a simple example take $`M=^1`$ and $`L=H`$ then $`\mathrm{\Phi }:^1^1`$ is the identity. Now consider $`L=H^2`$, which gives an embedding into $`^2`$. Let $`X_0,X_1`$ be the basis of the sections of $`H`$ then $`X_0^2,X_0X_1,X_1^2`$ is a basis of $`H^2`$. If $`(\alpha _0:\alpha _1)`$ are homogeneous coordinates on $`^1`$ the image of $`^1`$ in $`^2`$ is given as
$$\mathrm{\Phi }(^1)=\{(\alpha _0^2:\alpha _0\alpha _1:\alpha _1^2)^2\alpha _0,\alpha _1\}=𝒵(X_1^2X_0X_3).$$
The obtained subvariety is not linear anymore. Nevertheless it is algebraically isomorph to the linear variety $`^1`$.
Let us come back to the quantization condition. Recall that the quantization condition says that the Chern form of the quantum line bundle $`L`$ is essentially the Kähler form. But the Kähler form is a positive form, hence $`L`$ is a positive line bundle. Kodaira s embedding theorem says that a certain positive tensor power of $`L`$ will give an embedding into projective space. Hence $`L`$ is an ample line bundle. This implies that quantizable compact Kähler manifolds are always projective algebraic. In Section 2.3 we will see that the converse is also true.
### 2.3. The projective coordinate ring
Let $`V`$ be a projective variety in $`^n`$ and $`I=(V)`$ its vanishing ideal (24). Recall that it is a homogeneous ideal.
###### Definition 2.6.
The projective (or homogeneous) coordinate ring is the graded ring
$$𝕂[V]:=𝕂[X_0,X_1,\mathrm{}X_n]/(V).$$
(36)
For $`V=^n`$, we have $`(^n)=(0)`$, hence
$$𝕂[^n]:=𝕂[X_0,X_1,\mathrm{}X_n]=\underset{m0}{}H^0(^n,H^m)$$
is the full polynomial ring.
Inside $`𝕂[V]`$ the whole geometry of the variety $`V`$ is encoded. For example the points correspond to maximal homogeneous ideals $`MI`$ which are not identical to the ideal $`(X_0,X_1,\mathrm{}X_n)`$. Note that the only element of $`𝕂^{n+1}`$ which is a zero of all $`X_i`$ is $`0`$, which is not an element of projective space.
###### Definition 2.7.
The Krull dimension $`dimR`$ of a ring $`R`$ is defined to be the maximal length $`k`$ of strict chains of prime ideals $`P_i`$
$$P_0P_1\mathrm{}P_kR.$$
(37)
Recall that an ideal $`P`$ is called a prime ideal if from $`fgP`$ it follows that $`fP`$ or $`gP`$. Clearly, for prime ideals $`P`$ we have $`Rad(P)=P`$, hence $`(𝒵(P))=P`$. Moreover, the variety $`𝒵(P)`$ is always irreducible. Any chain (21) of irreducible subvarieties of an irreducible variety gives a chain of homogeneous prime ideals
$$𝕂[X_0,\mathrm{},X_n](V_0)(V_1)\mathrm{}(V_n)=(V).$$
(38)
lying between the vanishing ideal of $`V`$ and the whole ring. Passing to the quotient, i.e. to the coordinate ring one obtains a chain of prime ideals of the coordinate ring $`𝕂[V]`$. This works also in the opposite direction with the one exception that to both the whole ring $`𝕂[V]`$ and to the ideal $`(X_0,X_1,\mathrm{},X_n)modI(V)`$ corresponds the empty set. This implies <sup>1</sup><sup>1</sup>1 Note that for homogeneous coordinate rings to determine the Krull dimension it is enough to consider chains of homogeneous prime ideals.
$$dimV=dim𝕂[V]1.$$
Now let $`\mathrm{\Phi }:M^n`$ be the embedding obtained via the quantum line bundle $`L`$, which we assume already to be very ample. Let $`I:=(𝒵(\mathrm{\Phi }(M))`$ be the vanishing ideal of $`\mathrm{\Phi }(M)`$. We obtain $`\mathrm{\Phi }^{}H=L`$, $`i^{}X_i=s_i`$ for $`i=0,1,\mathrm{},n`$ for the sections $`s_i`$ used for the embedding, and $`\mathrm{\Phi }^{}(H^m)=(\mathrm{\Phi }^{}H)^m=L^m`$. In particular, the pull-backs of the global sections of $`H^m`$ generate the space of global sections of $`L^m`$. But in general they will not be a basis. The algebraic relations between them are exactly given by the elements of the ideal I. The projective coordinate ring $`[V]`$ can be identified with $`_{m0}H^0(M,L^m)`$.
In Section 1 we have defined the Berezin-Toeplitz quantization map
$$C^{\mathrm{}}(M)\mathrm{End}\left(\underset{m_0}{}H^0(M,L^m)\right),fT_f^{()}=(T_f^{(m)})_{m_0}.$$
(39)
Due to the fact that $`T_f^{()}`$ respects the grading given by $`m`$, it can also be considered as an element of
$$\mathrm{End}\left(\underset{m_0}{}H^0(M,L^m)\right)$$
and is fixed by this restriction. Hence, $`T_f^{()}`$ is an algebraic object operating on an algebraic vector space which coincides with the coordinate ring. The coordinate ring should be considered as the quantum Hilbert space. Note that this set-up makes perfect sense also for singular projective varieties.
Clearly, there is also a metric aspect in the theory. Our line bundle comes with a hermitian metric. On $`^n`$ we have the Fubini-Study Kähler form $`\omega _{FS}`$ induced by the standard metric in $`^{n+1}`$. This defines a metric on the tautological bundle and by taking the inverse metric a hermitian metric $`h_{FS}`$ on the hyperplane section bundle $`H`$. Suitable normalized it turns out that $`H`$ with $`h_{FS}`$ is the quantum line bundle of the Kähler manifold $`(^n,\omega _{FS})`$. If $`N`$ is a closed submanifold of $`^n`$, i.e. a nonsingular projective variety and $`i:N^n`$ is the embedding then the pair $`(N,i^{}\omega _{FS})`$ is a Kähler manifold with associated quantum line bundle $`(i^{}H,i^{}h_{FS})`$. In particular, nonsingular projective varieties are always quantizable. But note that if we start with a fixed Kähler manifold $`(M,\omega _M)`$ with very ample quantum line bundle $`(L,h)`$ and induced embedding $`\mathrm{\Phi }:M^n`$ then $`(M,\mathrm{\Phi }^{}\omega _{FS})`$ is again a quantizable Kähler manifold with quantum line bundle $`(L\mathrm{\Phi }^{}H,\mathrm{\Phi }^{}h_{FS})`$. But in general we have for the two Kähler forms defined on the same complex manifold $`\mathrm{\Phi }^{}\omega _{FS}\omega _M`$. We only know that they are cohomologous because they are representatives of the Chern class of the same bundle $`L`$. The question whether they coincide as forms has to do with the question whether the embedding is a Kähler embedding. This is related to Calabi’s diastatic function, respectively to Rawnsley’s epsilon function. I will not discuss this matter here, but see for a discussion and references to further results.
Via the metric the projective coordinate ring $`_{m0}H^0(M,L^m)`$ carries also a metric structure. To have a full description of the quantization also in the singular case the metric should be studied in more detail.
## 3. Singularities
In the last section a point on a projective variety was called a singular point if the rank of the matrix (27) is less than expected (see Definition 2.4). In this section I will give a different characterization of singular points. In particular, it will turn out, that there exist singularities which are better than others.
Clearly, the definition of a singular point as given in Definition 2.4 is a local one. Hence it is enough to study the local situation. For the local situation it is more convenient to consider affine varieties instead of projective varieties. If the projective space is replaced by an affine space the definitions work accordingly. After choosing coordinates the $`n`$-dimensional affine space is given as $`𝕂^n`$. A subset $`V`$ of $`𝕂^n`$ is called an affine variety if there exists finitely many polynomials $`f_1,f_2,\mathrm{},f_m𝕂[X_1,X_2,\mathrm{},X_n]`$ such that
$$V=𝒵(f_1,f_2,\mathrm{},f_m)=\{\alpha 𝕂^nf_i(\alpha )=0,i=1,\mathrm{},m\}.$$
If one replaces homogeneous polynomials, homogeneous ideals, etc. by arbitrary polynomials, arbitrary ideals, etc, the whole theory develops like in the projective case. Again, let $`I=(f_1,\mathrm{},f_m)`$ be the ideal generated by the above polynomials then $`V=𝒵(I)`$. Vice versa, given a variety $`V`$ in $`𝕂^n`$ we can define its vanishing ideal
$$(V):=\{f𝕂[X_1,X_2,\mathrm{},X_n]f(\alpha )=0,\alpha V\}.$$
(40)
With the same definition (25) of the radical ideal we obtain $`(𝒵(I))=Rad(I)`$. The affine coordinate ring of the variety $`V`$ is defined to be
$$𝕂[V]:=𝕂[X_1,X_2,\mathrm{},X_n]/(V).$$
(41)
The subset
$$U^{(i)}:=^n𝒵(X_i)=\{(\alpha _0:\alpha _1:\mathrm{}:\alpha _n)\alpha _i0\}$$
(42)
of $`^n`$ is a Zariski open (and hence dense) subset of $`^n`$. It can be identified with the affine space $`𝕂^n`$ via the map
$$\mathrm{\Phi }_i((\alpha _0:\alpha _1::\alpha _n))(\frac{\alpha _0}{\alpha _i},\mathrm{},\frac{\alpha _{i1}}{\alpha _i},\frac{\alpha _{i+1}}{\alpha _i},\mathrm{},\frac{\alpha _n}{\alpha _i}).$$
(43)
In this way $`^n`$ is covered by $`(n+1)`$ copies of affine $`n`$-space, i.e. $`^n=_{i=0}^nU^{(i)}`$. Every projective variety can be covered by affine varieties. Let $`f_l(X_0,X_1,\mathrm{},X_n)`$ for $`l=1,\mathrm{},m`$ be defining homogeneous polynomials for the projective variety $`V`$. Fix $`i`$ with $`0in`$ and let $`f_l^{(i)}`$ be the polynomials in $`n`$ variables obtained from the $`f_l`$ by setting the variable $`X_i`$ to $`1`$. Then $`V^{(i)}=𝒵(f_1^{(i)},\mathrm{},f_m^{(i)})`$ defines an affine variety. Via the map (43) we can identify $`V^{(i)}=VU^{(i)}`$. Again $`V=_{i=0}^nV^{(i)}`$. In particular every point of the projective variety lies at least in one of these affine varieties $`V^{(i)}`$.
In the following let $`V`$ be an affine variety. Again the dimension $`dimV`$ can be defined by Definition 2.2. This coincides with the Krull dimension of the coordinate ring $`𝕂[V]`$, i.e. $`dimV=dim𝕂[V]`$. In the affine case there is no subtraction of 1 necessary, because in this case there is a complete 1:1 correspondence between prime ideals and irreducible subvarieties. Note that if $`Y`$ is a irreducible projective variety all covering affine varieties $`Y^{(i)}`$ will be irreducible affine varieties and vice versa. Additionally we have $`dimY=dimY^{(i)}`$ for non-empty $`Y^{(i)}`$.
Singular points of affine varieties can be defined according to Definition 2.4 (of course now only $`n`$ variables will appear, hence we get an $`m\times n`$ matrix) using generators $`f_1,f_2,\mathrm{},f_m`$ of the vanishing ideal of the variety. If the affine variety $`V^{(i)}`$ comes from a projective variety $`V`$ as described above then $`xV`$ corresponding to $`\mathrm{\Phi }(x)`$ will be a singular point of $`V`$ if and only if $`\mathrm{\Phi }(x)`$ is a singular point of $`V^{(i)}`$.
There are some problems with this definition of a singular point. First, it is not a priori clear that it does not depend on the chosen generators of the ideal $`(V)`$. Second, the starting point of the definition is a variety lying in some affine space. But the singularity should be something intrinsic to the variety and not depend on the affine space the variety is lying in. It can be shown that indeed the notion does not depend on these choices. Nevertheless, a more intrinsic definition of a singularity would be desirable.
There is such a definition which deals with the local ring $`𝒪_{V,\alpha }`$ of the point $`\alpha `$ on $`V`$. This local ring is defined as follows. Let $`\alpha =(\alpha _1,\mathrm{},\alpha _n)𝕂^n`$ be a point on $`V`$. The vanishing ideal of $`\alpha `$ in the polynomial ring is the ideal $`M_\alpha =(X_1\alpha _1,X_2\alpha _2,\mathrm{},X_n\alpha _n)`$. It is a maximal ideal. This says that every ideal which is strictly bigger than $`M_\alpha `$ is the whole polynomial ring. The condition $`\alpha V`$ is equivalent to $`M_\alpha I=(V)`$. From this it follows that $`M_\alpha modI`$ is a maximal ideal of $`𝕂[V]`$. The local ring $`𝒪_{V,\alpha }`$ of the variety $`V`$ at the point $`\alpha `$ is defined as the localization of the ring $`𝕂[V]`$ with respect to the maximal ideal $`M_\alpha modI`$ (for simplicity we will denote it also by $`M_\alpha `$)
$$𝒪_{V,\alpha }=𝕂[V]_{M_\alpha }.$$
(44)
The localization is the ring of fractions where the denominators are elements from the multiplicative set $`𝕂[V]M_\alpha `$. It is a Noetherian local ring. Noetherian means that every ascending chain of ideals becomes stationary (or terminates). Local means that the ring has only one maximal ideal. Here the unique maximal ideal is $`M_\alpha modI/1`$ (again simply denoted by $`M_\alpha `$).
###### Definition 3.1.
A local ring $`R`$ with maximal ideal $`M`$ is called a regular local ring if
$$dim_{R/M}M/M^2=dimR.$$
(45)
Due to the fact that $`M`$ is a maximal ideal of $`R`$ the quotient $`R/M`$ is a field, and $`M/M^2`$ is a vector space over $`R/M`$. On the left hand side of (45) the vector space dimension is meant, on the right hand side the Krull dimension is meant.
In our case where $`R`$ is the local ring coming from the coordinate ring of a variety over an algebraically closed field $`K`$ we have $`R/M𝕂`$.
###### Definition 3.2.
A point $`\alpha V`$ is called a non-singular (or regular) point of $`V`$ if its local ring $`𝒪_{V,\alpha }`$ is regular. If not it is called a singular point. A variety is called non-singular (or smooth, or regular) if all points are regular. The subset of singular points of $`V`$ is denoted by $`Sing(V)`$.
$`Sing(V)`$ is always an algebraic subvariety of codimension $`\mathrm{codim}_VSing(V)1`$. In particular the non-singular locus is a non-empty Zariski open subset of $`V`$. If $`V`$ is irreducible then $`Sing(V)`$ is dense. For irreducible $`V`$ the dimensions of all local rings $`𝒪_{V,\alpha }`$ are constant and equal to the dimension of $`V`$.
The $`𝕂`$-vector space $`M_\alpha /M_\alpha ^2`$ is also called the Zariski cotangent space, resp. its dual $`(M_\alpha /M_\alpha ^2)^{}`$ the Zariski tangent space. In general
$$dim_𝕂M_\alpha /M_\alpha ^2dim𝒪_{V,\alpha }=dimV,$$
(46)
where we assume for the last equality $`V`$ to be irreducible. Hence we can also define the singular points to be the points where the dimension of the Zariski tangent space is bigger than the dimension of the (irreducible) variety.
Let me illustrate this in the case of cubic curves. In an affine chart using the ideal
$$I=\left(Y^24X(Xa)(Xb)\right)$$
(47)
with $`a,b𝕂`$ we obtain the cubic curves as $`𝒵(I)`$. In this normalisation $`\alpha =(0,0)`$ lies on the cubic. The cotangent space at $`\alpha `$ is given as
$$M_\alpha /M_\alpha ^2=(X,Y)modI/(X^2,Y^2,XY)modI.$$
(48)
From the relations given by $`I`$ we calculate
$$Y^2=4abX8(a+b)X^2+4X^3modI.$$
(49)
Hence
$$4abX(X^2,Y^2,XY)modI.$$
(50)
If $`ab0`$ the element $`Y`$ will be enough to generate the quotient (48). The tangent space will be one-dimensional and $`(0,0)`$ will be a nonsingular point. If either $`a`$ or $`b`$ equals $`0`$ the element $`X`$ will also be necessary to generate the tangent space. Hence $`(0,0)`$ will be a singular point.
Given an arbitrary irreducible (projective or affine) variety $`V`$ then there exists a stratification of the singularity set $`Sing(V)`$ obtained in the following manner. Let $`U=VSing(V)`$ be the Zariski open set of regular points then $`VU`$ is a closed subvariety. It can be decomposed into finitely many irreducible varieties of dimension less than $`dimV`$
$$VU=V_1^{(1)}V_2^{(1)}\mathrm{}V_l^{(1)}.$$
Again from this complement $`Sing(VU)`$ can be determined. It is a subvariety of the variety $`V`$ from higher codimension. This process can be repeated as long as there are singularities. Because the codimension strictly increases it has to stop after finitely many steps.
Guided by the algebraic properties of the local rings we have an hierarchy for the types of singularities. If $`R`$ is a ring without zero divisors then $`Quot(R)`$ is the ring whose elements are the fractions of elements in $`R`$ with denominator $`0`$.
###### Definition 3.3.
A ring $`R`$ (without zero divisors) is called a normal ring if the elements of $`Quot(R)`$ which are solutions of algebraic equations with coefficients from $`R`$ and highest coefficient 1 lie already in $`R`$.
It is a classical result (Gauß Lemma) that $``$ is normal and also that polynomial rings over fields are normal.
###### Definition 3.4.
Let $`V`$ be an irreducible variety.
(a) A variety is normal at a point $`\alpha V`$ if the local ring $`𝒪_{V,\alpha }`$ is normal.
(b) A variety is called normal if it is normal at every point $`\alpha V`$.
(c) A singular point is called a normal singular point, and the singularity is called a normal singularity, if the variety is normal at this point.
A regular local ring is always normal. Hence, regular points are always normal. If $`V`$ is a normal variety it follows
$$\mathrm{codim}_VSing(V)2.$$
(51)
This says that normal singular varieties are “less singular” than generic singular varieties. A lot of the singular varieties which appear as moduli spaces are normal. Normal varieties behave from the point of functions defined on them similar to nonsingular varieties. For example, if $`V`$ is a variety of dimension $`2`$ and $`xV`$ a normal point then every regular function in $`V\{x\}`$ can be extended to a regular function on $`V`$. Additionally normality is necessary to have a well-behaved theory of (Weil-)divisors based on codimension 1 irreducible subvarieties.
For every irreducible affine variety $`V`$ with singularity set $`Sing(V)`$ there exists a normal affine variety $`\stackrel{~}{V}`$ and an algebraic morphism $`\pi :\stackrel{~}{V}V`$ such that
$$\pi ^1(VSing(V))VSing(V).$$
(52)
The variety $`\stackrel{~}{V}`$ is called the normalization of $`V`$. It is obtained by a purely algebraic process, i.e. by taking the normal closure of the coordinate ring in its quotient field. This can be extended to the projective case too.
For algebraic curves normal points are always regular (there is no space for codimension two subvarieties). Hence the normalization gives already a desingularization. In the case of the above discussed singular cubic curves the normalization is given by the line (affine, resp. projective).
The question arises whether it might be even possible to find for every projective variety $`V`$ with singularity set $`Sing(V)`$ a nonsingular projective variety $`Y`$ which coincides with $`V`$ outside $`Sing(V)`$, and is minimal in a certain sense. Such a $`Y`$ is called a desingularization and the whole process is called a resolution of singularities. It was shown by Hironaka (see also ) that there exists for projective varieties over fields of characteristic zero (and this is the case we are dealing with) a resolution of singularities. More precisely, for every projective variety $`V`$ there exists a nonsingular projective variety $`Y`$ and a proper <sup>2</sup><sup>2</sup>2 the algebraic equivalent of a compact map algebraic map $`f:YV`$ such that $`f`$ is an isomorphism over an open non-empty subset $`UV`$, i.e. $`f^1(U)U`$.
## 4. Quotients
In this section let us assume $`𝕂=`$ for the base field.
### 4.1. Quotients in algebraic geometry
Moduli spaces of geometric objects are very often varieties with singularities. Typically, they are obtained starting from a smooth variety classifying the objects with respect to a certain “presentation”. To obtain the moduli space one has to “divide out” the different presentations. Usually, one has a group operating on the presentations and a candidate of the moduli space is given by the quotient set under the group action, the orbit space. Unfortunately, it is not always possible to endow the quotient set with a compatible structure of a variety again. Even if we allow the quotient to be an algebraic scheme it will not be possible. In our context schemes will appear as “varieties with multiplicities”. It is quite reasonable that one should at least incorporate such objects in the theory. E.g. if we have two lines in the plane meeting at a point and we move one line with the intersection point fixed, we will nearly always have two lines. There is only one exception, if the moving line coincides with the fixed one. In this case the configuration consists of one line. But from the deformation point of view we should better count this special line twice, i.e. we should consider it as double line. The language of schemes deals with such objects and even with much more general ones. Nevertheless to avoid giving additional definitions I will still work on algebraic varieties (affine, projective, quasiprojective) in the following. But the reader should keep in mind that the language of algebraic schemes would be more appropriate for moduli problems. See for an introduction to this field.
Let $`X`$ be an algebraic variety and $`G`$ a reductive algebraic group acting algebraically on $`X`$. This means that $`G`$ is the complexification of a maximal compact subgroup $`K`$ of $`G`$. Of special importance (and this are the examples that the reader should keep in mind) are the groups $`\mathrm{GL}(n)`$, $`\mathrm{SL}(n)`$, and $`\mathrm{PGL}(n)`$. As indicated above it is important to study “quotients” of $`X`$ under actions of the group $`G`$. Mumford has given with his geometric invariant theory (GIT) the principal tool to deal with such quotients.
###### Definition 4.1.
A morphism of algebraic varieties $`f:XY`$ is called a good quotient if
* $`f`$ is surjective and $`G`$-invariant, i.e. $`f(gx)=f(x)`$, for all $`gG`$ and $`xX`$,
* $`\left(f_{}(𝒪_X)\right)^G=𝒪_Y`$,
* if $`V`$ is a $`G`$-invariant closed subset of $`X`$ then $`f(V)`$ is closed in $`Y`$, and if $`V_1`$ and $`V_2`$ are $`G`$-invariant closed subsets of $`X`$ then
$$V_1V_2=\mathrm{}f(V_1)f(V_2)=\mathrm{}.$$
In Condition (2) $`𝒪_X`$ and $`𝒪_Y`$ are the structure sheaves of the varieties $`X`$ and $`Y`$. They are essentially nothing else as the sheaves of local regular functions on $`X`$ and $`Y`$ respectively. Condition (2) states that the local regular functions on $`Y`$ can be given as those local regular functions on $`X`$ which are constant along the fiber and invariant under $`G`$.
A good quotient is a categorical quotient in the sense that
* $`f`$ is constant on the orbits of the action,
* for every algebraic variety $`Z`$ with a morphism $`g:XZ`$ which is constant on the orbits of the $`G`$-action on $`X`$ there exist a unique morphism
$`\overline{g}:YZ`$ with $`g=\overline{g}f`$.
###### Definition 4.2.
A morphism of algebraic varieties $`f:XY`$ is called a geometric quotient if
* $`f:XY`$ is a good quotient,
* for every $`yY`$ the fiber $`f^1(y)`$ consists exactly of one orbit under the group action.
If any of these quotients exists then they are unique.
For a good quotient there might exists fibers consisting of several orbits under the group action and these orbits are not necessarily closed (for a geometric quotient the orbits are always closed).
If we have a geometric quotient than the orbit space carries a structure of an algebraic variety. But this condition is very often to strong to be fulfilled. We have sometimes to assign several orbits to one geometric point to obtain a geometric structure and to end up (hopefully) with a good quotient. Mumford’s concept of stability will help to decide what to do. Let $`X^n`$ be a projective algebraic variety and $`G`$ a reductive algebraic group embedded into $`\mathrm{GL}(n+1)`$ with an action of $`G`$ on $`X`$ given by the standard linear action of $`\mathrm{GL}(n+1)`$ on the points in $`^n`$.
###### Definition 4.3.
(1) A point $`xX`$ is called semi-stable if and only if there exists a non-constant $`G`$-invariant homogeneous polynomial $`F[X_0,\mathrm{},X_n]`$ with $`F(x)0`$.
(2) A point $`xX`$ is called stable if and only if
* the dimension of the Orbit $`O(x)`$ under the $`G`$-action equals the dimension of the group and
* there exists a non-constant $`G`$-invariant homogeneous polynomial $`F[X_0,\mathrm{},X_n]`$ with $`F(x)0`$, and the action of $`G`$ on the zero set $`X_F:=\{yXF(y)=0\}`$ is closed, i.e. if for every $`y_0X_F`$ the orbit $`O(y_0)`$ is closed.
The set of stable <sup>3</sup><sup>3</sup>3“Stable” in the above introduced sense corresponds to “properly stable” in the definition of Mumford. Stability in his sense does not require the condition on the dimension of the orbit. points of $`X`$ under the above group action $`G`$ is denoted by $`X^s`$, the set of semi-stable points is denoted by $`X^{ss}`$. Both are open subset. Clearly, $`X^sX^{ss}X`$.
Let me point out that the notion of stability might depend on the embedding of the projective variety $`X`$ into projective space and a corresponding linearization of the action of $`G`$. Recall from Section 2.2 that for an abstract projective variety $`X`$ such an embedding is defined by the choice of a very ample line bundle $`L`$ on $`X`$ and a choice of basis of its global sections.
###### Theorem 4.4.
Assume that $`X^s`$ is non-empty then there exists a projective algebraic variety $`Y`$ and a morphism $`f_{ss}:X^{ss}Y`$ such that
* $`f_{ss}`$ is a good quotient of $`X^{ss}`$ by $`G`$,
* there exists an open subset $`UY`$ such that $`f_{ss}^1(U)=X^s`$
and $`f_s:=f_{ss}^{}{}_{|X^s}{}^{}:X^sU`$ is a geometric quotient of $`X^s`$ by $`G`$.
The good quotient is projective, but the geometric quotient is as an open subset of something projective in general only quasi-projective. If we interpret this in the opposite way, we see that we will need in general also non-stable (but still semi-stable) points to obtain projective (this means compact in the complex topology) moduli spaces. Clearly, even if the projective variety we started with was smooth there is no reason to expect that the quotient will be smooth.
For more details one might consult , or for a more leisurely reading .
### 4.2. The relation with the symplectic quotient
In this subsection I want to quote results on the relation between the quotients in algebraic geometry and the symplectic quotients. The results are taken from Francis Kirwan’s appendix to the third edition of Mumford’s book on GIT and are due to Kirwan, Kempf and Ness. More details and references can be found there.
Let $`X`$ be a nonsingular projective complex variety in $`^n`$, and $`G`$ a reductive group acting linearly on $`X`$ via $`\rho :G\mathrm{GL}(n+1)`$. If $`K`$ is any fixed maximal compact subgroup of $`G`$ then after a suitable choice of coordinates the subgroup $`K`$ acts unitarily on $`X`$, i.e. $`\rho _{|K}:K\mathrm{U}(n+1)`$. Let $`𝔨`$ be the Lie algebra of $`K`$, $`𝔨^{}`$ its dual, and $`\mu :X𝔨^{}`$ the standard moment map defined for all $`a𝔨`$ by
$$\mu (x)(a):=\frac{{}_{}{}^{t}\overline{\widehat{x}}\rho _{}(a)\widehat{x}}{2\pi \mathrm{i}\widehat{x}}.$$
(53)
Here $`\widehat{x}^{n+1}\{0\}`$ is any vector of homogeneous coordinates representing the element $`xX^n`$ and $`\rho _{}:𝔨𝔲(n+1)`$ is the tangent map of $`\rho _{|K}`$.
In this situation the symplectic quotient (or Marsden-Weinstein reduction) is defined as $`\mu ^1(0)/K`$. On the other hand we can define the good quotient (which is also a categorical quotient) of the semi-stable points $`X^{ss}`$ by $`G`$. By Theorem 4.4 it is a projective variety which is commonly also denoted by $`X//G`$. It contains as open subset the geometric quotient $`X^s/G`$. Immediately the following question arrises: How are these quotients related?
###### Theorem 4.5.
(a) The point $`xX`$ is semi-stable if and only if $`\overline{O_G(x)}\mu ^1(0)\mathrm{}`$, i.e. the closure of the orbit of $`x`$ under $`G`$ meets $`\mu ^1(0)`$.
(b) $`\mu ^1(0)X^{ss}`$.
(c) The inclusion under (b) induces a homeomorphism
$$\mu ^1(0)/KX//G.$$
(54)
(d) If we denote by $`\mu ^1(0)_{reg}`$ the set of the $`x\mu ^1(0)`$ for which the tangent map $`d\mu _x`$ of the moment map is surjective, then the homeomorphism under (c) restricts to a homeomorphism
$$\mu ^1(0)_{reg}/KX^s/G.$$
(55)
Hence we see that as topological spaces the symplectic quotient is isomorphic to the good quotient and the subspace of the regular points $`\mu ^1(0)_{reg}/K`$ is isomorphic to the geometric quotient.
Things get slightly more complicated if we consider also the complex structure. If $`X`$ is a compact Kähler manifold then the symplectic quotient $`\mu ^1(0)_{reg}/K`$ carries a structure of a complex Kähler manifold away from the singularities. But $`X^s/G`$ carries also a complex structure from the geometric quotient construction. These two structures coincide on the subset where the symplectic quotient has no singularities. Hence if $`\mu ^1(0)_{reg}/K`$ is a Kähler manifold the complex structure coincides with the complex structure of the geometric quotient. But there exists examples (e.g. given by Kirwan \[13, p.159\]) where the geometric quotient has no singularities but the symplectic quotient has singularities in the sense that the Kähler structure coming from the reduction process is singular at certain points. Further examples of the relation between the structure of the singularities of the two type of quotients are given in the contribution of J. Huebschmann to this volume.
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# Small cones of 𝑚-hemimetrics
## 1 Introduction
The notions of $`m`$-hemimetrics and $`m`$-partition hemimetrics are generalizations of the notions of metrics and cuts, which are well-known and central objects in Graph Theory, Combinatorial Optimization and, more generally, Discrete Mathematics.
Obviously the $`m`$-hemimetrics on an $`n`$-element set $`E`$ form a cone. What are its facets and extreme rays? The calculations are rather complex due to the high dimension, $`\left(\genfrac{}{}{0pt}{}{n}{m+1}\right)`$, of the cone even for small $`n`$; moreover, matrices can no longer be used.
Recall that a metric or a metric space is a pair $`(E,d)`$ where $`E`$ is a nonvoid set and $`d:E^2_+`$(the set of nonnegative reals) satisfies for all $`x,y,zE`$:
(d1) $`d(x,y)=0x=y`$,
(d2) $`d(x,y)=d(y,x)`$ (symmetry),
(d3) $`d(x,y)d(x,z)+d(z,y)`$ (the triangle inequality).
A basic example is $`(^2,d)`$, where $`d`$ is the Euclidean distance of $`x`$ and $`y`$; i.e., the length of the segment joining $`x`$ and $`y`$. An immediate extension is $`(^3,d)`$, where $`d(x,y,z)`$ is the area of the triangle with vertices $`x,y`$ and $`z`$. This leads to the following definition (see \[Me28, Bl53, Fr58, Gä63\]). A 2-metric is a pair $`(E,d)`$, where $`E`$ is a nonempty set and $`d:E^3_+`$ satisfies for all $`x,y,z,tE`$
(d1) $`d(x,x,y)=0`$,
(d1<sup>′′</sup>) $`xyd(x,y,u)>0`$ for some $`uE`$,
(d2) $`d(x,y,z)`$ is totally symmetric,
(d3) $`d(x,y,z)d(t,y,z)+d(x,t,z)+d(x,y,t)`$ (the tetrahedron inequality).
The axiom (d2) means that the value of $`d(x,y,z)`$ is independent of the order of $`x,y`$ and $`z`$. The axiom (d3) captures that fact that in $`^3`$ the area of a triangle face of a tetrahedron does not exceed the sum of the areas of the remaining three faces. A 2-metric allows the introduction of several geometrical and topological concepts – e.g. the betweenness, convexity, line and neighborhood – which lead to interesting results.
For finite 2-metrics (the object of this study), for their polyhedral aspects and for applications, the axiom (d1) and (d1<sup>′′</sup>) seem to be too restrictive and so we drop them. The definition given below is formulated for an arbitrary positive integer $`m`$. A map $`d:E^{m+1}`$ is totally symmetric if for all $`x_1,\mathrm{},x_{m+1}E`$ and every permutation $`\pi `$ of $`\{1,\mathrm{},m+1\}`$
$$d(x_{\pi (1)},\mathrm{},x_{\pi (m+1)})=d(x_1,\mathrm{},x_{m+1}).$$
Definition. Let $`m>0`$. An $`m`$-hemimetric is a pair $`(E,d)`$, where $`d:E^{m+1}`$ is totally symmetric and satisfies the simplex inequality: for all $`x_1,\mathrm{},x_{m+2}`$$`E`$
$$d(x_1,\mathrm{},x_{m+1})\underset{i=1}{\overset{m+1}{}}d(x_1,\mathrm{},x_{i1},x_{i+1},\mathrm{},x_{m+2}).$$
$`(1)`$
Call a $`m`$-hemimetric $`d`$ nonnegative if $`d`$ takes only nonnegative values. The notion of a $`m`$-hemimetric is new, but it is closely related to the notion of a $`m`$-metric, considered in about 200 references collected in \[Gä90\]. The study of $`m`$-hemimetrics is motivated (apart, of course, that it represents an extension of metrics) by applications in Statistics and Data Analysis (see \[DeRo99\] for some relevant references). Notice the following immediate:
Fact 1. If $`(E,d)`$ and $`(E,d^{})`$ are $`m`$-hemimetrics and $`a,b_+`$ then $`(E,ad+bd^{})`$ is an $`m`$-hemimetric.
Here, as usual, for all $`x_1,\mathrm{},x_{m+1}E`$
$$(ad+bd^{})(x_1,\mathrm{},x_{m+1}):=ad(x_1,\mathrm{},x_{m+1})+bd^{}(x_1,\mathrm{},x_{m+1}).$$
Given an $`(m+1)`$-partition $`S_1,\mathrm{},S_{m+1}`$ of $`V_n:=\{1,2,\mathrm{},n\}`$, a partition $`m`$-hemimetric $`\alpha (S_1,\mathrm{},S_{m+1})`$ is defined by setting $`\alpha (S_1,\mathrm{},S_{m+1})(i_1,\mathrm{},i_{m+1})`$ is equal to 1 if for no $`1j<lm+1`$ both $`i_j`$ and $`i_l`$ belong to the same $`S_k`$, and 0 otherwise. It is easy to see that $`\alpha (S_1,\mathrm{},S_{m+1})`$ is a nonnegative $`m`$-hemimetric and for $`m=1`$ it is the usual cut semimetric (see, for example, \[DeLa97\]).
For small values of $`n`$ and $`m`$ we consider the cone of all $`m`$-hemimetrics, the cone of all nonnegative $`m`$-hemimetrics and the cone, generated by all partition $`m`$-hemimetrics on $`V_n`$. Using computer search we list facets and generators for these cones and tables of their adjacencies and incidences. The different orbits were determined manually, using symmetries. We study two graphs, the $`1`$-skeleton and the ridge graph, of these polyhedra: the number of their nodes and edges, their diameters, conditions of adjacency, inclusions among the graphs and their restrictions on some orbits of nodes. In fact, we would like to describe two graphs $`G(C),G(C^{})`$ for our three cones as fully as possible, but in the cases, when it is too difficult, we will give some partial information on adjacencies in those graphs. Especially we are interested in the diameters of the graphs, in a good criterion of adjacency, in their local graphs (i.e. in the subgraphs induced by all neighbors of a given vertex) and in their restrictions on some orbits. Finally, we compare obtained results with similar results for metric case (see \[DeDe95, DDFu96, DeLa97\]) and quasi-metric case (see \[DePa99\]). All computation was done using the programs cdd of \[Fu95\].
The following notation will be used below:
* the (m+1)-simplex inequality (1) and, in particular, for $`m=2`$, the tetrahedron inequality
$`T_{ijk,l}:x_{ijl}+x_{ikl}+x_{jkl}x_{ijk}0`$;
* the nonnegativity inequality $`N_{i_1,\mathrm{},i_{m+1}}:x_{i_1,\mathrm{},i_{m+1}}0`$;
* the cone $`P_n^m`$ of partition m-hemimetrics, generated by all $`(m+1)`$-partitions of $`V_n`$;
* the cone $`NHM_n^m`$ of nonnegative m-hemimetrics, defined by all
$`(nm1)`$$`\left(\genfrac{}{}{0pt}{}{n}{m+1}\right)`$ $`(m+1)`$-simplex inequalities and all $`\left(\genfrac{}{}{0pt}{}{n}{m+1}\right)`$ nonnegativity inequalities on $`V_n`$;
* the cone of all m-hemimetrics, $`HM_n^m`$, defined by all $`(m+1)`$-simplex inequalities on $`V_n`$.
Clearly, $`P_n^mNHM_n^mHM_n^m`$ and these 3 cones are of full dimension $`\left(\genfrac{}{}{0pt}{}{n}{m+1}\right)`$ each. For $`m=1`$ the last two cones coincide and the first two cones are the cut cone $`CUT_n`$ and the semimetric cone $`MET_n`$, considered in detail in \[DeLa97\] and in the references listed there. The cone $`P_n^1=CUT_n`$ has $`\frac{n}{2}`$ orbits of extreme rays (represented by the cuts $`\alpha (1\mathrm{}i,(i+1)\mathrm{}n)`$).
To simplify the notation we keep $`n`$ fixed and denote by $`E_k`$ the family of all $`k`$-element subsets of $`V_n`$ ($`k=1,\mathrm{},n`$). Let $`d`$ be a semimetric on the set $`V_n`$. Because of symmetry (d2) and since $`d(i,i)=0`$ for all $`iV_n`$, we can view the semimetric $`d`$ as a vector $`(d_{i_1,i_2})_{}^{E_2}`$. In the same way, we can view the $`m`$-hemimetric $`d`$ on the set $`V_n`$ as a vector $`(d_{i_1,\mathrm{},i_{m+1}})_{}^{E_{m+1}}`$. In particular, each extreme ray of the cones, considered below, will be represented by an integer vector on the ray with relatively prime coordinates. We also will represent the facets of cones by such vectors.
Any such vector $`v=(v_{i_1,\mathrm{},i_{m+1}})^{E_{m+1}}`$ can be represented by the following vertex-labeled induced subgraph $`R(v)`$ of the Johnson graph $`J(n,m+1)`$. The vertices of $`R(v)`$ are all unordered $`(m+1)`$-tuples $`(i_1,\mathrm{},i_{m+1})`$ such that $`(v_{i_1,\mathrm{},i_{m+1}})`$ is not zero. This value will be the label of the vertex; we will omit the label when it is 1. Two vertices of the Johnson graph (and also of its induced subgraph $`R(v)`$) are adjacent if the corresponding $`(m+1)`$-tuples have $`m`$ common elements.
For example, $`R(\alpha (S_1,\mathrm{},S_{m+1}))`$ is the complement to the Hamming graph $`H(|S_1|,\mathrm{},|S_{m+1}|)`$, i.e. the direct (Cartesian) product of the cliques $`K_{|S_i|}`$, $`1im+1`$. Another example: the graph $`R`$ of the vector defining a nonnegativity facet is a vertex and, for a $`(m+1)`$-simplex facet, it is the complete graph $`K_{m+2}`$ with one vertex labeled $`1`$.
## 2 Partition $`m`$-hemimetrics and related polyhedra
Recall that for a partition $`S_1,S_2`$ of $`V_n`$ the cut semimetric $`\alpha (S_1,S_2)`$ satisfies $`\alpha (S_1,S_2)_{ij}=1`$ if $`\{i,j\}S_1`$ is a singleton and $`\alpha (S_1,S_2)_{ij}=0`$ otherwise. We extend it as follows. Let $`q2`$ be an integer and let $`S_1,\mathrm{},S_q`$ be pairwise disjoint nonvoid subsets of $`V_n`$, forming a partition of $`V_n`$. The multicut semimetric $`\delta (S_1,\mathrm{},S_q)`$ is the vector in $`^{E_2}`$, defined by $`\delta (S_1,\mathrm{},S_q)_{ij}=0`$, if $`i,jS_h`$ for some $`h`$, $`1hq`$, and $`\delta (S_1,\mathrm{},S_q)_{ij}=1`$, otherwise.
The connection between $`\delta (S_1,\mathrm{},S_q)`$ and $`\alpha (S_1,\mathrm{},S_q)`$ from Section 1 is given by
$`\alpha (S_1,\mathrm{},S_q)(i_1,\mathrm{},i_q)`$=$`_{1s<tq}\delta (S_1,\mathrm{},S_q)(i_s,i_t)`$=
$`\frac{_{1s<tq}\delta (S_1,\mathrm{},S_q)(i_s,i_t)}{\left(\genfrac{}{}{0pt}{}{q}{2}\right)}`$; compare it with the half-perimeter $`m`$-semimetric from \[DeRo99\].
The cone generated by all multicut semimetrics $`\delta (S_1,\mathrm{},S_q)`$ ($`q2`$) on $`V_n`$, is called the multicut cone and denoted by $`MCUT_n`$; it coincides with $`CUT_n`$ (see \[DeLa97\], Proposition 4.2.9). The convex hull of the cut semimetrics (multicut semimetrics) on $`V_n`$, is called the cut polytope ( multicut polytope) and is denoted by $`CUT_n^{\mathrm{}}`$ ($`MCUT_n^{\mathrm{}}`$); the two polytopes not coincide.
## 3 Facets, extreme rays and their orbits in polyhedra
We recall some terminology. Let $`C`$ be a polyhedral cone in $`^n`$. Given $`v^n`$, the inequality $`v^Tx0`$ is said to be valid for $`C`$, if it holds for all $`xC`$. Then the set $`\{xC|v^Tx=0\}`$ is called the face of $`C`$, induced by the valid inequality $`v^Tx0`$. A face of dimension $`dim(C)1`$ is called a facet of $`C`$; a face of dimension $`1`$ is called an extreme ray of $`C`$. A face of dimension $`dim(C)2`$ is called a ridge.
Two vertices $`x,y`$ of $`C`$ are said to be adjacent, if they generate a face of dimension 2 of $`C`$. Two facets of $`C`$ are said to be adjacent, if their intersection has dimension $`dim(C)2`$. The $`1`$-skeleton graph of $`C`$ is the graph $`G(C)`$ whose nodes are the extreme rays of $`C`$ and whose edges are the pairs of adjacent nodes. Denote by $`C^{}`$ the dual cone of $`C`$. The ridge graph of $`C`$ is the graph whose nodes are the facets of $`C`$ and with an edge between two facets if they are adjacent on $`C`$. So, the ridge graph of a cone $`C`$ is the $`1`$-skeleton $`G(C^{})`$ of its dual cone.
A mapping $`f:^n^n`$ is called a symmetry of a cone $`C`$ (or a polytope $`P`$), if it is an isometry, satisfying $`f(C)=C`$ (or $`f(P)=P`$). (An isometry of $`^n`$ is a linear mapping preserving the Euclidean distance.) Given a face $`F`$, the orbit $`\mathrm{\Omega }(F)`$ of $`F`$ consists of all faces, that can be obtained from $`F`$ by the group of all symmetries of $`C`$.
Clearly, all the faces of $`CUT_n`$ and $`CUT_n^{\mathrm{}}`$ are preserved by any permutation of $`V_n`$.
For $`m>1`$ all orbits of faces of $`m`$-hemimetric cones $`P_n^m,NHM_n^m,HM_n^m`$ on $`V_n`$ are also preserved under any permutation of the set $`V_n=\{1,\mathrm{},n\}`$. We conjecture that the symmetry group consists only of permutations of $`V_n`$, i.e. it is the group $`Sym(n)`$ of all permutations on $`V_n`$ ( see Theorem 3.3 in \[DGLu91\] stating that the symmetry group of a truncated multicut polytope is $`Sym(n)`$).
## 4 The case of $`n=m+2`$
The minimal $`n`$ for which the three cones are nontrivial is $`m+2`$; the dimension of the cones is also $`m+2`$ for $`n=m+2`$.
First, we present a complete linear description for case $`(m,n)=(2,4)`$.
It turns out that $`P_4^2=NHM_4^2`$. This cone has 6 extreme rays (all in the same orbit under $`Sym(4)`$): $`\alpha (S_1,S_2,S_3)`$ for the 3-partitions
$`(1,2,34),(1,3,24),(1,4,23),(2,3,14),(2,4,13),(3,4,12)`$.
There are $`8`$ facets, which form $`2`$ orbits: the orbit $`F_1`$ of all 4 tetrahedron facets and the orbit $`F_2`$ of all 4 nonnegativity facets.
The edge graph $`G(C)`$ is $`K_63K_2`$ (the octahedron); the 3 pairs of nonadjacent rays are of the form $`\alpha (a,b,cd),\alpha (c,d,ab)`$. Each extreme ray (say, $`\alpha (1,2,34)`$) is incident to 2 tetrahedron and to 2 nonnegativity facets (namely, to $`T_{123,4},T_{124,3}`$ and $`N_{134},N_{234}`$).
The ridge graph $`G(C^{})`$ is the cube. Adjacencies of facets of $`NHM_4^2`$ are shown in Table 1. For each orbit a representative and the number of adjacent facets from other orbits are given, as well as the total number of adjacent ones, the number of incident extreme rays and the cardinality of orbits.
More precisely, for the ridge graph of $`NHM_4^2`$ it holds:
(i) The tetrahedron facet $`T_{ijk,l}`$ is adjacent only to the facets $`N_{ijl},N_{ikl},N_{jkl}`$;
(ii) The nonnegativity facet $`N_{ijk}`$ is adjacent only to the facets $`T_{ijl,k},T_{ikl,j},T_{jkl,i}`$.
The cone $`HM_{m+2}^m`$ is a simplex $`(m+2)`$-dimensional cone; so $`G(C)=G(C^{})=K_{m+2}`$. Its facets are all $`(1,1)`$-valued $`(m+2)`$-vectors with only one $`1`$, its generators are all $`(1m,1)`$-valued $`(m+2)`$-vectors with only one $`1m`$. Notice that $`P_3^1=HM_3^1=CUT_3=MET_3`$.
In general, $`P_{m+2}^m=NHM_{m+2}^m`$ for any $`m2`$. This cone has $`\left(\genfrac{}{}{0pt}{}{m+2}{2}\right)`$ extreme rays, all in the same orbit, represented by $`\alpha (12,3,\mathrm{},m+2)`$, i.e. by any vector of length $`m+2`$, consisting of two ones and $`m`$ zeros. The skeleton of $`P_{m+2}^m`$ is the Johnson graph $`J(m+2,2)`$, called also the triangular graph $`T(m+2)`$, which is is the line graph $`L(K_{m+2})`$. It is also the skeleton of the $`(m+1)`$-polytope (called $`\mathrm{𝑎𝑚𝑏𝑜}\alpha _{m+1}`$), obtained from the $`(m+1)`$-simplex as the convex hull of the mid-points of all its edges; e.g. $`T(4)`$ is the skeleton of the octahedron, $`T(5)`$ is the complement of the Petersen graph. In general, $`T(m),m2,`$ has diameter 2; moreover, it is a strongly regular graph.
The cone $`P_{m+2}^m`$ has two orbits, $`F_1`$ and $`F_2`$, of facets, containing $`m+2`$ facets each and represented by the $`(m+1)`$-simplex facet $`T_{1\mathrm{}(m+1),(m+2)}`$ and by the nonnegativity facet $`N_{1\mathrm{}(m+1)}`$. The orbit $`F_1`$ consists of simplex cones, i.e. facets from this orbit are incident to $`m+1`$ linearly independent extreme rays. Any nonnegativity inequality $`N`$ defines the cone $`P_{m+1}^{m1}=NHM_{m+1}^{m1}`$, i.e. it becomes equality on this smaller cone. So, $`N`$ is non-facet only for $`m=1`$ and it is a simplex cone only for $`m=2`$; in general, $`N`$ is incident to $`\left(\genfrac{}{}{0pt}{}{m+1}{2}\right)`$ extreme rays. The ridge graph is $`\overline{K_{m+2}}`$ on $`F_1`$; on $`F_2`$ it is $`\overline{K_4}`$ for $`m=2`$ and $`K_{m+2}`$ for $`m3`$. Finally, the $`m+2`$ pairs $`(T_{\overline{i},i},N_{\overline{i}})`$ (of $`(m+1)`$-simplex and nonnegativity facets) are the only non-edges for pairs of facets from different orbits.
## 5 Small $`2`$-hemimetrics
### 5.1 The case of $`5`$ points
We present here the complete linear description of $`P_5^2`$, $`NHM_5^2`$ and $`HM_5^2`$. The cone $`P_5^2`$ has $`25`$ extreme rays, which form $`2`$ orbits with representatives $`\alpha (1,2,345)`$ (orbit $`O_1`$) and $`\alpha (1,23,45)`$ (orbit $`O_2`$). The skeleton and the ridge graph of $`P_5^2`$ has 270 and 1185 edges, respectively. The cone $`P_5^2`$ has $`120`$ facets divided into $`4`$ orbits, induced by the $`20`$ tetrahedron inequalities (orbit $`F_1`$), the $`10`$ nonnegativity inequalities (orbit $`F_2`$), the $`60`$ inequalities (orbit $`F_3`$), represented by
$$A:2x_{123}(x_{124}+x_{135})+(x_{134}+x_{125}+x_{245}+x_{345})0$$
and the $`30`$ inequalities (orbit $`F_4`$), represented by
$$B:2(x_{123}+x_{145}+x_{245}x_{345})+(x_{134}+x_{135}+x_{234}+x_{235}x_{124}x_{125})0$$
.
The above two inequalities are the 2-hemimetric analogs of the following 5-gonal inequality (the simplest inequality, different from the triangle inequality), appearing in the cone $`CUT_n`$ for $`n5`$:
$$(x_{13}+x_{14}+x_{15}+x_{23}+x_{24}+x_{25})(x_{12}+x_{34}+x_{35}+x_{45})0$$
.
This facet and $`B`$ have both the Petersen graph as their $`\overline{R}`$ (i.e. the complement of their graph $`R`$). Clearly, the graphs $`R`$ for partition 2-hemimetrics $`\alpha (1,2,34)`$, $`\alpha (1,2,345)`$, $`\alpha (1,23,45)`$ are the cycles $`C_2,C_3,C_4`$.
The graphs $`\overline{R(A)}`$, $`\overline{R(B)}`$ are given on the Figure 1.
Figure 1: $`\overline{R(A)}`$, $`\overline{R(B)}`$ in the cone $`P_5^2`$
The facets from the orbit $`F_4`$ are simplex-cones, i.e. the extreme rays on them are linearly independent. Among the 9 neighbors of $`B`$, 4 are from the orbit $`F_1`$, 4 are from the orbit $`F_3`$ and exactly one (actually, $`N_{123}`$) from the orbit $`F_2`$. The local graph of a facet from $`F_4`$ (i.e. the subgraph of the ridge graph of $`P_5^2`$, induced by all neighbors of $`B`$) is $`K_9C_4`$. In fact, all nonadjacencies in this local graph are the four edges of the 4-cycle of the 4 neighbors of $`B`$ from the orbit $`F_3`$.
Facets from orbits $`F_1,F_2,F_3,F_4`$ are incident, respectively, to 7,9; 7,9; 4,6; 3,6 extreme rays from orbits $`O_1,O_2`$ of $`P_5^2`$.
The skeleton and the ridge graph of $`NHM_5^2`$ have 420 and 355 edges, respectively. The adjacencies of 37 extreme rays and of 30 facets of this cone are given in Tables 4, 5. The extreme rays are divided into 3 orbits $`O_1,O_2,O_3`$, represented by (0,1)-valued vectors $`v_1,v_2,v_3`$ below; their $`R`$-graphs are $`C_3,C_4,C_5`$, respectively.
Each facet (from both orbits) of $`NHM_5^2`$ is incident to 7,9,6 extreme rays from orbits $`O_1,O_2,O_3`$, respectively. Each (tetrahedron) facet of $`HM_5^2`$ is incident to 7,9,6,6,9,15 extreme rays from orbits $`O_1,\mathrm{},O_6`$.
For any cone, let $`I_{O_i,F_j}`$ and $`I_{F_j,O_i}`$ denote the number of facets from the orbit $`F_j`$, incident to an extreme ray of the orbit $`O_i`$, and, respectively, the number of extreme rays from $`O_i`$, incident to a facet from $`F_j`$. Clearly, $`|O_i|I_{O_i,F_j}=|F_j|I_{F_j,O_i}`$.
The cone $`HM_5^2`$ has 92 extreme rays divided into 6 orbits. Below we give some representatives $`v_1,\mathrm{},v_6`$ of those orbits $`O_1,\mathrm{},O_6`$. The first two represent both orbits of $`P_5^2`$, the first 3 represent the 3 orbits of $`NHM_5^2`$.
$`x=(x_{123},x_{124},x_{125},x_{134},x_{135},x_{145},x_{234},x_{235},x_{245},x_{345})`$:
$`v_1=(1,1,1,0,0,0,0,0,0,0)`$;
$`v_2=(0,1,1,1,1,0,0,0,0,0)`$;
$`v_3=(1,0,1,0,0,1,1,0,0,1)`$;
$`v_4=(1,1,1,1,0,0,1,0,1,1)`$;
$`v_5=(1,1,1,1,1,1,1,1,1,1)`$;
$`v_6=(1,0,1,0,1,1,1,1,2,1)`$.
###### Proposition 1
The diameters of the skeleton graphs of $`P_5^2`$ and of $`NHM_5^2`$ are $`2`$.
In fact, each of the orbits $`O_1,O_2`$ of $`P_5^2`$ is a dominating clique. There is only one type of a non-edge, represented by $`\alpha (1,23,45),\alpha (2,3,145))`$, but $`\alpha (1,3,245)`$ is one of common neighbors. The complement of the skeleton of $`P_5^2`$ turns out to be the Petersen graph with a new vertex (corresponding to a member of the orbit $`O_2`$) on each of 15 edges. The result for $`NHM_5^2`$ comes also by finding out a common neighbor to each possible non-edge.
###### Proposition 2
For the ridge graphs of $`NHM_5^2`$ and $`HM_5^2`$ it holds:
(i) The diameter of the ridge graph of $`NHM_5^2`$ is $`2`$;
(ii) Its restriction on the orbits $`F_1`$ and $`F_2`$ is $`K_{4,4,4,4,4}`$ and the Petersen graph, respectively;
(iii) The ridge graph of $`HM_5^2`$ is $`K_{4,4,4,4,4}`$ (of diameter $`2`$).
### 5.2 The case of $`6`$ points
$`NHM_6^2`$ has exactly 12492 extreme rays, with $`(adjacency,incidence)`$ pairs being, respectively, (2278,64), (1321,56), (1030,40), (818,48), (731,48), (358,40), (270,36), (93,28), (66,28), (51,28), (47,28), (46,39), (37,31), (32,28), (30,27), (29,26), (27,23), (26,24), (26,23), (25,25), (23,22), (22,21), (21,21).
Three of the above pairs (1st, 2nd and 4th) are realized by orbits (say, $`0_1`$, $`O_2`$ and $`O_4`$), which are represented by 3-partition 2-hemimetrics $`\alpha (1,2,3456)`$, $`\alpha (1,23,456)`$, $`\alpha (12,34,56)`$ and have size 15, 60, 15, respectively. The graphs $`R`$ of members of orbits $`O_1,O_2`$ and $`O_4`$ are $`K_4,K_6C_6=K_3\times K_2`$ (the skeletons of the tetrahedron and 3-prism) and the skeleton of the cube. Three other orbits consist also of (0,1)-valued extreme rays: $`O_3`$ (with $`R`$ being the Petersen graph), $`O_5`$ (with $`R`$-graph being the skeleton of the simple polyhedron with $`p`$-vector $`p=(p_3=2,p_4=2,p_5=2)`$) and $`O_7`$ with graph $`R`$ (non-planar, non-regular), given on Figure 2, together with one for $`O_5`$. The extreme rays of the remaining orbits are (0,1,2)-valued and (0,1,2,3)-valued vectors.
Figure 2 : The graphs R of extreme rays from orbits $`O_7,O_5`$ of the cone $`NHM_6^2`$
The cone $`P_6^2`$ has more than 950.000 facets (computer stopped, by lack of memory, after 72, out of 90, iterations). Here are two examples of a $`(0,1,1)`$-valued facets of $`P_6^2`$; see also Figure 3 (for the facet W).
$`W:(x_{145}+x_{146}+x_{136}+x_{123}+x_{125})+(x_{245}+x_{234}+x_{346}+x_{356}+x_{256})(x_{235}+x_{236})0.`$
Figure 3 : $`R(W)`$ in the cone $`P_6^2`$
$`Z:x_{ijk}(x_{124}+x_{125}+x_{145})(x_{234}+x_{235}+x_{345})2(x_{146}+x_{156}+x_{456})2x_{236}0.`$
Remark that the triples with coefficients zero, in $`W`$ and $`Z`$, form the skeleton of 1- and 2-truncated tetrahedron, respectively; the triples with coefficient -1 form $`K_3+K_1`$ and $`K_2+K_1`$, respectively.
## 6 Small $`3`$-hemimetrics
The cone $`NHM_6^3`$ has 287 extreme rays divided into 5 orbits. Below we give representatives $`u_1,\mathrm{},u_5`$ of the orbits $`O_1,\mathrm{},O_5`$. These vectors are indexed by 4-subsets of the set $`\{1,\mathrm{},6\}`$; the 4-subsets are given as the complements of 2-subsets. The first four are $`(0,1)`$-valued; their $`R`$-graphs (in the Johnson graph $`J(6,4)`$ of all $`4`$-tuples) are the cycles $`C_3,C_4,C_5,C_6`$, respectively. The first two are partition $`3`$-hemimetrics; they represent both orbits of $`P_6^3`$. The graphs $`R(u_4)`$ and $`R(u_5)`$ are on Figure 4.
$`x=(x_{\overline{12}},x_{\overline{13}},x_{\overline{14}},x_{\overline{15}},x_{\overline{16}},x_{\overline{23}},x_{\overline{24}},x_{\overline{25}},x_{\overline{26}},x_{\overline{34}},x_{\overline{35}},x_{\overline{36}},x_{\overline{45}},x_{\overline{46}},x_{\overline{56}})`$:
$`u_1=(0,0,1,1,0,0,0,0,0,0,0,0,1,0,0)`$;
$`u_2=(0,0,1,1,0,0,0,0,0,1,1,0,0,0,0)`$;
$`u_3=(0,0,1,1,0,0,0,0,0,1,0,1,0,0,1)`$;
$`u_4=(0,0,1,1,0,0,0,1,1,1,0,1,0,0,0)`$;
$`u_5=(0,0,1,1,0,1,0,0,1,0,0,1,1,2,0)`$;
Figure 4 : $`R(u_4),R(u_5)`$ in the cone $`NHM_6^3`$
The cone $`P_6^3`$ has 4065 facets divided into at least 11 orbits. Below we give representatives $`f_1,\mathrm{},f_{11}`$ of the orbits $`F_1,\mathrm{},F_{11}`$. Their $`(adjacency,incidence)`$ pairs are, respectively, (1526,49), (703,41), (100,23), (37,19), (31,18), (30,18), (23,17), (23,15), (22,18), (18,16), (14,14). The facets $`f_1,f_2`$ are nonnegativity and 4-simplex facets; $`f_{11}`$ is a simplex cone. The $`R`$-graphs of the facets $`f_3`$ and $`f_4`$ are on Figure 5.
Figure 5 : $`\overline{R(f_3)},\overline{R(f_4)}`$ in the cone $`P_6^3`$
$`f_1=(0,0,0,0,0,1,0,0,0,0,0,0,0,0,0)`$;
$`f_2=(0,1,0,0,0,1,0,0,0,1,1,1,0,0,0)`$;
$`f_3=(1,1,0,1,1,2,1,0,0,1,0,0,1,1,0)`$;
$`f_4=(1,1,0,1,0,2,1,0,1,1,0,1,1,0,1)`$;
$`f_5=(1,1,1,2,1,2,2,1,0,2,1,0,1,2,1)`$;
$`f_6=(1,1,0,2,2,2,1,1,1,1,1,1,2,2,0)`$;
$`f_7=(1,1,1,3,2,2,2,2,1,2,2,1,2,3,1)`$;
$`f_8=(1,1,3,1,4,2,2,2,1,2,2,3,2,1,1)`$;
$`f_9=(1,1,1,2,2,2,2,1,1,2,1,1,1,1,2)`$;
$`f_{10}=(1,1,1,1,2,1,1,1,1,1,1,1,2,1,1)`$;
$`f_{11}=(1,1,1,2,0,2,2,1,1,2,1,1,1,1,2)`$.
###### Proposition 3
The diameter of the skeleton graph of $`P_6^3`$ is $`2`$. Moreover:
(i) $`G(O_1)=K_{20}`$, $`G(O_2)=K_{45}15K_3`$;
(ii) all non-edges are represented by $`\alpha (12,34,5,6)`$ that are nonadjacent to
$`\alpha (12,3,4,56),\alpha (1,2,34,56)`$ (from the same orbit $`O_2`$) and to
$`\alpha (125,3,4,6),\alpha (126,3,4,5),\alpha (345,1,2,6),\alpha (346,1,2,5)`$.
In fact, both non-neighbors of $`\alpha (12,34,5,6)`$ are in $`O_2`$. For both types of non-edges - $`\alpha (12,34,5,6)`$ with $`\alpha (12,3,4,56)`$ and $`\alpha (125,3,4,6)`$ \- the ray $`\alpha (13,24,5,6)`$ is a common neighbor. Also, all 9 non-neighbors of a ray from $`O_1`$, form $`K_9`$ in the skeleton graph.
Notice that the skeleton of $`P_6^3`$ is not an induced subgraph of the skeleton of $`NHM_6^3`$; the only difference is in their restriction $`G(O_2)`$ to the orbit of rays, represented by $`u_2`$. One can check that all neighbors of a partition hemimetric $`\alpha (a_1,a_2,b_1b_2,c_1c_2)`$ from the same orbit $`O_2`$ of $`NHM_6^3`$ are the 10 rays obtained by a transposition $`(xy)`$ and the 8 rays obtained by a product $`(a_1b_i)(a_2c_j)`$ or $`(a_1c_i)(a_2b_j)`$ of two transpositions. But in the skeleton of $`P_6^3`$, the ray $`\alpha (a_1,a_2,b_1b_2,c_1c_2)`$ is adjacent to all other members of $`O_2`$, except for the two rays, obtained from it by $`(a_1b_1)(a_2b_2)`$ or $`(a_1c_1)(a_2c_2)`$. The complement of the graph, induced by all 18 neighbors of the ray $`\alpha (a_1,a_2,b_1b_2,c_1c_2)`$ from the same orbit $`O_2`$ of $`NHM_6^3`$, is $`C_4+C_4`$ on 8 rays, obtained by a product $`(a_1b_i)(a_2c_j)`$ or a product $`(a_1c_i)(a_2b_j)`$, the skeleton of the cube on 8 rays, obtained by $`(a_ib_j)`$ or $`(a_ic_j)`$, and it is $`\overline{K_2}`$ on two rays obtained by $`(b_ic_j)`$.
The skeletons of $`P_6^3`$ and $`NHM_6^3`$ both contain a dominating clique $`O_1`$; so their diameters are 2 or 3. In order to see closer the skeleton of $`NHM_6^3`$, we now describe the local graph, denoted by $`H`$, of the ray $`u_5`$. All 26 neighbors are in orbits $`O_1,O_2,O_3,O_4`$ only. It will be easier to describe $`\overline{H}`$. The restrictions of $`\overline{H}`$ on them are $`\overline{K_6}`$, $`C_8`$, the skeleton of the cube and $`2K_2`$, respectively. Two vertices from $`O_1`$ (say $`15`$ and $`16`$) are isolated; so the diameter of $`H`$ is 2. Here we denote by $`ij`$ the $`j`$-th member of the orbit $`O_i`$ in $`H`$. All edges of $`\overline{H}`$ (without isolated vertices $`15`$ and $`16`$) are presented on Figure 6. On the right picture the members of $`O_1`$ are excluded while on the left one the members of $`O_2`$ are excluded. $`\overline{H}`$ does not contain cross-edges among orbits $`O_1`$ and $`O_2`$.
Figure 6 : A presentation of the local graph of a ray of the orbit $`O_5`$ of the cone $`NHM_6^3`$
###### Proposition 4
The ridge graph of $`NHM_6^3`$ has diameter $`2`$. Moreover:
(i) any 4-simplex facet $`T_{ijkl,m}`$ is adjacent to all but 5 facets: $`N_{ijkl}`$ and all 4 other 4-simplex facets with the same support;
(ii) the restrictions of the ridge graph to the orbits $`F_1`$ and $`F_2`$ are
$`K_{5,5,5,5,5,5}`$ and $`K_{15}`$, respectively.
## 7 Small 4-hemimetrics
The cone $`NHM_7^4`$ has 3692 extreme rays divided into 8 orbits. We give below representatives $`w_1,\mathrm{},w_8`$ of their orbits $`O_1,\mathrm{},O_8`$. These vectors are indexed by 5-subsets of the set $`\{1,\mathrm{},7\}`$; the 5-subsets are given as the complements of 2-subsets. The $`(adjacency,incidence)`$ pairs of those rays are, respectively, (985,48), (535,43), (315,38), (192,33), (126,28), (67,30), (43,25), (42,25). The first five vectors are $`(0,1)`$-valued; their graphs $`R`$ are $`C_3,C_4,C_5,C_6,C_7`$, respectively. The first two are partition $`4`$-hemimetrics; they represent both orbits of $`P_7^4`$. The vectors $`w_i,1i4,`$ and $`w_6`$ have same $`R`$-graphs as the members of orbits $`O_i,1i5,`$ of $`NHM_6^3`$, respectively; so, the graphs of Figure 5 represent also $`w_4`$ and $`w_6`$. The graphs $`R(w_7)`$ and $`R(w_8)`$ are on Figure 7.
$`(\overline{12},\overline{13},\overline{14},\overline{15},\overline{16},\overline{17},\overline{23},\overline{24},\overline{25},\overline{26},\overline{27},\overline{34},\overline{35},\overline{36},\overline{37},\overline{45},\overline{46},\overline{47},\overline{56},\overline{57},\overline{67})`$:
$`w_1=(0,0,0,1,1,0,0,0,0,0,0,0,0,0,0,0,0,0,1,0,0)`$;
$`w_2=(0,0,0,1,1,0,0,0,1,1,0,0,0,0,0,0,0,0,0,0,0)`$;
$`w_3=(0,0,0,1,1,0,1,0,0,1,0,0,1,0,0,0,0,0,0,0,0)`$;
$`w_4=(0,0,0,1,1,0,1,0,0,1,0,1,0,0,0,1,0,0,0,0,0)`$;
$`w_5=(0,0,0,1,1,0,1,0,0,1,0,1,0,0,0,0,0,1,0,1,0)`$;
$`w_6=(0,0,0,1,1,0,0,0,0,0,0,1,0,2,1,0,0,1,1,0,0)`$;
$`w_7=(0,0,0,2,2,0,0,0,0,1,1,1,1,0,0,1,0,0,0,0,1)`$;
$`w_8=(0,0,0,1,1,0,0,1,0,0,1,2,1,1,0,0,0,1,0,0,0)`$.
Figure 7 : $`R(W_6),R(W_7),R(W_8)`$ in the cone $`NHM_7^4`$
It is easy to check that the ridge graph of $`NHM_7^4`$ is $`K_{6,6,6,6,6,6,6}`$ on $`F_1`$ and $`K_{21}`$ on $`F_2`$. All non-edges among $`F_1,F_2`$ are of the form $`T_{i_1\mathrm{}i_5,i_6}`$ and $`N_{i_1\mathrm{}i_5}`$.
## 8 Comparison of the small cones
Now we compare some semimetric and m-hemimetric cones on $`n`$ points for small $`n`$. The triangle inequalities suffice to describe the cut cones for $`n4`$, but $`CUT_nMET_n`$ (strictly) for $`n5`$. The complete description of all the facets of the cut cone $`CUT_n`$ is known for $`n8`$, the complete description of the semimetric cone $`MET_n`$ is known for $`n7`$ (see, for example, the linear description of $`MET_7`$ in \[Gr92\]). Here the “combinatorial explosion” starts from $`n=8`$. The number of orbits of facets and of extreme rays of those and other cones, when it is known, is given in Table 12.
In fact, $`P_n^2=NHM_n^2`$ holds only for the smallest value $`n=4`$. For $`n=4,5`$ we computed all facets, extreme rays and their adjacencies and incidences for three cones $`P_n^2,NHM_n^2,HM_n^2`$. For 2-hemimetrics the “combinatorial explosion”(in terms of the amount of computation and memory) starts already for the cone $`P_6^2`$.
In the Table 12 we compare the small 2-hemimetric cones $`P_n^2,NHM_n^2`$ with the 1-hemimetric cones $`CUT_n,MET_n`$ and their generalization in another direction: the cones $`OMCUT_n,QMET_n`$. Last two cones consist of all quasi-semimetrics on $`V_n`$ and of those obtained from oriented multicuts; see \[DePa99\] for the notions and results for them given in the Table 12. The cones $`NHM_n^2`$ and $`QMET_n`$ have, besides of generalizations of the usual triangle inequality, only nonnegativity facets. In the Table 12, columns 3 and 4 give the number of extreme rays and facets, respectively; in parenthesis are given the numbers of their orbits. In column 5 are given the diameters of the skeleton and the ridge graphs of the cone specified in the row. In the Table 12, the number of orbits of extreme rays and the diameter for cones $`QMET_5`$, $`NHM_6^2`$, $`P_6^2`$, $`P_7^4`$ and dual $`P_6^3`$, $`OMCUT_5`$ are taken from recent work \[DuDe01\], as well as the exact value of the diameter for $`NHM_6^3`$, $`NHM_7^4`$ and for the duals of $`P_5^2`$, $`Cut_7`$.
Incidences (to the extreme rays) of facets $`T_{ijk,l}`$ and $`N_{ijk}`$ on the cones $`P_4^2=NHM_4^2`$, $`P_5^2,NHM_5^2`$ amount to 3, 14 and 22, respectively, but they are different (4001 and 3939) on $`HM_6^2`$. Incidences of similar facets $`T_{ij,k}`$ (oriented triangular inequality, i.e. $`d(x,y)d(x,z)+d(z,y)`$ for a quasimetric $`d`$ ), $`N_{ij}`$ (nonnegativity inequality) are equal (to 7, 43) on cones $`OMCUT_3=QMET_3,OMCUT_4`$, but they are different (78 and 80) on $`QMET_4`$.
For $`n=4,5`$ we observe that the ridge graphs of $`HM_n^2`$ and $`NHM_n^2`$ are induced subgraphs of the ridge graphs of $`NHM_n^2`$ and $`P_n^2`$, respectively. The similar property does not hold for the 1-skeletons of those cones. For example, any extreme ray of the orbit $`O_2`$ is adjacent to 14,6,2 members of the same orbit in the cones $`P_n^2,NHM_n^2,HM_n^2`$, respectively. Also, the ridge graph of $`QMET_4`$ is an induced subgraph of the ridge graph of $`OMCUT_4`$, but the skeleton of $`OMCUT_4`$ is not an induced subgraph of the skeleton of $`QMET_4`$ (see \[DePa99\]). On the other hand, the ridge graph of $`MET_n`$ and the skeleton of $`CUT_n`$ (for any $`n`$) have diameters 2 and 1, respectively, and those graphs are induced subgraphs of the ridge graph of $`CUT_n`$ and of the skeleton of $`MET_n`$, respectively (see Lemma 2.1 and Theorem 3.5 in \[DeDe94\])
## 9 Conjectures for general $`m,n`$
###### Conjecture 5
The two partition $`m`$-hemimetrics $`\alpha (S_1,\mathrm{},S_{m+1})`$ and $`\alpha (T_1,\mathrm{},T_{m+1})`$ on $`V_n`$ are nonadjacent in the skeleton of $`P_n^m`$ if and only if there exist six different subsets $`S_i,S_j,S_k`$ and $`T_i^{},T_j^{},T_k^{}`$, such that $`S_iS_j=T_k^{}`$ and $`S_k=T_i^{}T_j^{}`$.
The conjecture holds for $`m=1`$: all cut semimetrics are adjacent. It holds for $`nm=2`$: we have the graph $`J(m+2,2)`$. It also holds for $`(m,n)=(2,5)`$ and $`(3,6)`$.
###### Conjecture 6
The ridge graphs of $`HM_n^m`$ and of $`NHM_n^m`$ are induced subgraphs of the ridge graphs of $`NHM_n^m`$ and $`P_n^m`$, respectively.
Recall that the ridge graph of $`NHM_n^m`$ has two orbits of vertices: $`F_1,F_2`$, consisting of $`(nm1)`$$`\left(\genfrac{}{}{0pt}{}{n}{m+1}\right)`$ simplex and $`\left(\genfrac{}{}{0pt}{}{n}{m+1}\right)`$ nonnegativity inequalities.
###### Conjecture 7
The ridge graph $`NHM_n^m`$ satisfies:
(i) The $`(m+1)`$-simplex facet $`T_{i_1\mathrm{}i_{m+1},i_{m+2}}`$ is adjacent to all other facets, except the following $`m+2`$ facets:
all other $`(m+1)`$-simplex facets with the same support and $`N_{i_1\mathrm{}i_{m+1}}`$;
(ii) $`G(F_2)=\overline{J(n,3)}`$ for $`m=2`$ and $`G(F_2)=K_{\left(\genfrac{}{}{0pt}{}{n}{m+1}\right)}`$ for $`m3`$.
Clearly, (i) implies that the restriction of the ridge graph on $`F_1`$ is $`G(F_1)=K_{m+2,\mathrm{},m+2}`$. It is easy to see that Conjecture $`3`$ would imply that the diameter of the ridge graph of $`NHM_n^m`$ is 2 (it was proved in \[DeDe94\] that the diameter of the ridge graph of $`NHM_n^1=MET_n`$ is 2). In fact, to see it for $`m=2`$ consider all 3 types of pairs of nonadjacent vertices:
(i) let $`x,yF_1`$ have the same support, say, $`1234`$. Suppose that $`x_{124}=y_{124}=1`$. Then $`N_{123}`$ is a common neighbor for $`x`$ and $`y`$.
(ii) for $`N_{123}`$ and $`N_{124}`$, any tetrahedron facet $`T_{134,2}`$ is their common neighbor.
(iii) for $`N_{123}`$ and $`T_{123,4}`$, the facet $`N_{345}`$ is a common neighbor.
###### Conjecture 8
The extreme rays of $`NHM_n^m`$ include:
(i) any ray whose $`R`$-graph is an $`R`$-graph of an extreme ray of $`NHM_{n1}^{m1}`$;
(ii) every $`(0,1)`$-valued extreme ray of $`NHM_{m+3}^m`$ with $`R`$-graph $`C_i`$ ($`3im+3`$).
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# Half-quantum vortex and 𝑑̂-soliton in Sr2RuO4
## Abstract
Assuming that the superconductivity in Sr<sub>2</sub>RuO<sub>4</sub> is described by a planar p-wave order parameter, we consider possible topological defects in Sr<sub>2</sub>RuO<sub>4</sub>. In particular, it is shown that both of the $`\widehat{d}`$-soliton and half-quantum vortex can be created in the presence of the magnetic field parallel to the $`a`$-$`b`$ plane. We discuss how one can detect the $`\widehat{d}`$-soliton and half-quantum vortex experimentally.
PACS numbers: 74.20.-z, 74.80.-g, 74.25.-q
It has been suggested that the unconventional superconducting state of Sr<sub>2</sub>RuO<sub>4</sub> is described by the planar spin-triplet $`p`$-wave order parameter with broken time reversal symmetry in analogy to the <sup>3</sup>He A-phase. As known from the example of superfluid <sup>3</sup>He, one of the hallmarks of the triplet superconductivity is the presence of a manifold of topological defects. Thus, we expect that the creation and detection of topological defects in Sr<sub>2</sub>RuO<sub>4</sub> (or spin-triplet superconductor) will provide further insights about the nature of the unconventional superconducting state of Sr<sub>2</sub>RuO<sub>4</sub>. In this sense, the study of the topological defects in the planar $`p`$-wave superconducting state with broken time reversal symmetry is of great interest.
More specifically, it has been proposed that the superconducting order parameter in this system is described by the planar $`p`$-wave form written as,
$$\mathrm{\Delta }_{\alpha \beta }(\stackrel{}{k})=\stackrel{}{d}(\stackrel{}{k})(\stackrel{}{\sigma }i\sigma _2)_{\alpha \beta }$$
(1)
with
$$\stackrel{}{d}(\stackrel{}{k})=\mathrm{\Delta }\widehat{d}(\widehat{k}_1\pm i\widehat{k}_2),$$
(2)
where $`\sigma _\mu (\mu =1,2,3)`$ are Pauli matrices and $`\alpha ,\beta `$ represent spin $``$ or $``$. Here $`\widehat{k}_j(j=1,2)`$ represent the projection of the unit wave vector $`\widehat{k}`$ along two perpendicular directions $`\widehat{e}_1`$ and $`\widehat{e}_2`$ in two dimensional space. This order parameter describes the Cooper pair state with the zero spin projection on $`\widehat{d}`$ and with the unique projection of the pair orbital angular momentum given by $`\widehat{l}=\widehat{e}_1\times \widehat{e}_2`$. In Sr<sub>2</sub>RuO<sub>4</sub>, due to the spin-orbit coupling, $`\widehat{d}`$ is forced to be parallel to $`\pm \widehat{c}`$ and $`\stackrel{}{k}`$ is the quasi-particle momentum in the $`a`$-$`b`$ plane.
Indeed the spontaneous magnetization seen by muon spin relaxation experiment and flat Knight shift seen by NMR seem to be consistent with this picture. On the other hand, the origin of the spontaneous magnetization seen by muon spin relaxation experiment is somewhat mysterious since we do not expect such a magnetization in a homogeneous system.
It is important to notice that the superconducting ground state described by the order parameter of Eq.2 is doubly degenerate. We can designate these two ground states by the angular momentum $`l_z=\pm 1`$, where $`\widehat{l}`$ is parallel to the $`c`$-axis. Sigrist and Agterberg proposed recently that there will be in general a domain wall between $`l_z=1`$ and $`l_z=1`$ states, which we shall call $`\widehat{l}`$-soliton in analogy to the case of superfluid <sup>3</sup>He-A. It is likely that such a soliton is magnetically active, so it may be an origin of the spontaneous magnetization seen in muon spin relaxation experiment. In particular, in a magnetic field $`H||\widehat{c}`$, only one of these degenerate states is favored . Therefore, it is possible to control $`\widehat{l}`$-solitons by a magnetic field parallel to the $`c`$-axis. They also proposed that these $`\widehat{l}`$-solitons would provide very efficient barriers for the vortex motion and this effect is possibly related to the pinning of vortices observed in Sr<sub>2</sub>RuO<sub>4</sub> below $`T=30mK`$ . However, in this experiment the magnetic field is applied in a direction perpendicular to the $`c`$-axis. As we will see later, so-called $`\widehat{d}`$-solitons appear to be more appropriate than $`\widehat{l}`$-solitons in this configuration. Also, for the $`\widehat{l}`$-solitons, it is rather difficult to estimate the soliton energy and to make a further quantitative prediction.
The purpose of this paper is to propose an alternative model for the appearance of the spontaneous magnetization and the mechanism of the pinning of vortices; $`\widehat{d}`$-soliton and half-quantum vortex. The $`\widehat{d}`$-soliton is a domain wall between $`\widehat{d}\widehat{c}`$ and $`\widehat{d}\widehat{c}`$ as in superfluid <sup>3</sup>He-A. We believe that $`\widehat{d}`$ is parallel and antiparallel to $`\widehat{c}`$, because they are forced to be parallel to the angular momentum $`\widehat{l}`$ (or $`\widehat{l}`$) due to the spin-orbit coupling characterized by an energy scale $`\mathrm{\Omega }_d`$ . Therefore, if we use $`\mathrm{\Omega }_d`$ as a parameter, we can calculate the energy and shape of the $`\widehat{d}`$-soliton provided that $`\mathrm{\Omega }_d\mathrm{\Delta }(T)`$, where $`\mathrm{\Delta }(T)`$ is the superconducting gap. Unfortunately we do not know the precise value of $`\mathrm{\Omega }_d`$, but it may be about $`\frac{1}{10}\mathrm{\Delta }(T)`$. In this picture, moving $`\widehat{d}`$-soliton generates the local magnetization which can result in the spontaneous magnetization seen in muon spin relaxation experiment. One can also generate a large number of $`\widehat{d}`$-solitons by applying a burst of high frequency microwave with frequency $`\mathrm{\Omega }_d`$ sent parallel to the $`a`$-$`b`$ plane.
As in superfluid <sup>3</sup>He-A, each $`\widehat{d}`$-soliton is terminated by a pair of half-quantum vortices. We find that these pairs of half-quantum vortices are more stable than the usual single quantum vortex in the superconducting state in the presence of the magnetic field parallel to the $`a`$-$`b`$ plane. This means that the usual single quantum vortex would spilt into a pair of half-quantum vortices connected by the $`\widehat{d}`$-soliton. In this case, these objects would provide an extremely efficient pinning mechanism of vortices in Sr<sub>2</sub>RuO<sub>4</sub>. Also the half-quantum vortices should exhibit a clear electron spin resonance (ESR) signature. Further we believe that these objects are visible by the scanning tunneling microscopy (STM) imaging and by micromagnetometry developed by Kirtley et al used in high $`T_c`$ cuprate compounds.
Free energy of the conventional single vortex when the magnetic field is parallel to the $`a`$-$`b`$ plane
Let us assume that the magnetic field is parallel to the $`a`$-axis. Then the free energy of the conventional vortex with the flux quantum $`\varphi _0=hc/2e`$ is obtained within the London approximations as
$$f_v=(\frac{\varphi _0}{4\pi \lambda })^2\mathrm{ln}(\frac{\lambda }{\xi }),$$
(3)
where $`\xi `$ is the coherence length and $`\lambda `$ is the magnetic penetration depth. The magnetic penetration depth $`\lambda `$ is related to the superfluid density $`\rho _s(T)`$ by $`\lambda ^2=\frac{4\pi e^2}{mc^2}\rho _s(T)`$. When the magnetic field is parallel to the $`a`$-axis, in the anisotropic system like Sr<sub>2</sub>RuO<sub>4</sub>, $`\lambda `$ and $`\xi `$ should be reinterpreted as $`\lambda =\sqrt{\lambda _b\lambda _c}`$ and $`\xi =\sqrt{\xi _b\xi _c}`$. Here $`\lambda _{b,c}`$ and $`\xi _{b,c}`$ are the magnetic penetration depth and coherence length in the $`b`$ and $`c`$ directions respectively.
$`\widehat{d}`$-soliton and a pair of half-quantum vortices
There exists huge anisotropy in the in-plane and out-of-plane transport properties in Sr<sub>2</sub>RuO<sub>4</sub>. Thus Sr<sub>2</sub>RuO<sub>4</sub> may be regarded as an effectively two dimensional system. The large anisotropy or the effective two-dimensionality of the system forces the angular momentum of the Cooper pair to be parallel or antiparallel to the $`c`$-axis. In the $`p`$-wave superconducting state described by the order parameter described by Eq.1 and Eq.2, the $`\widehat{d}`$ vector is oriented along $`\pm \widehat{l}`$ in the presence of the spin-orbit coupling. Here we consider the case that the angular momentum $`\widehat{l}`$ is uniform in the entire system. We can assume, without loss of generality, that $`\widehat{l}\widehat{c}`$. We are interested in the deformation of the $`\widehat{d}`$ configuration from the uniform case; for example, $`\widehat{d}\widehat{c}`$. Any deviation from the uniform state would cost the energy associated with the spin-orbit coupling characterized by an energy scale $`\mathrm{\Omega }_d`$ . However, we will show that the so-called $`\widehat{d}`$-soliton (a particular form of the $`\widehat{d}`$ configuration) with a pair of half-quantum vortices can have lower energy than the conventional single vortex. Thus it is easier to excite a pair of half-quantum vortices with a $`\widehat{d}`$-soliton compared to single conventional vortex. In particular, a magnetic field parallel to the $`a`$-$`b`$ plane generates very likely pairs of half-quantum vortices rather than usual vortices when the formers are stable.
We consider the $`\widehat{d}`$-soliton that is a topological planar defect in the $`\widehat{d}`$ configuration. The orientation of $`\widehat{d}`$ changes by $`\pi `$ across the planar defect while $`\widehat{d}`$ vectors at far distances are still along the $`c`$-axis. Typical configurations of $`\widehat{d}`$-soliton in the $`y`$-$`z`$ plane can be found in Fig.1 and Fig.2 which we will explain later. We take $`y`$ and $`z`$ as the coordinates along $`b`$-axis and $`c`$-axis respectively. Now let us attach a pair of half-quantum vortices to the end points of the $`\widehat{d}`$-soliton of length $`R`$ in $`y`$-$`z`$ plane. In the case of an isolated half-quantum vortex, we have $`e^{i\pi }=1`$ factor in the order parameter due to phase winding around the half-quantum vortex. Therefore, an isolated half-quantum vortex cannot occur. On the other hand, if the half-quantum vortex is attached to the end points of the $`\widehat{d}`$-soliton, the disgyration in $`\widehat{d}`$ at the same point compensates the phase $`\pi `$ so that there is no net change in the overall phase of the order parameter.
In order to show that a pair of half-quantum vortices with the $`\widehat{d}`$-soliton is a lower energy excitation compared to single conventional vortex, we have to compare the free energies of two cases. The free energy required to create the $`\widehat{d}`$-soliton is obtained from
$$f_d=\frac{1}{2}\chi _NC^2d^3r[\underset{ij}{}|_i\widehat{d}_j|^2+\xi _d^2(1d_z^2)],$$
(4)
where $`\chi _N`$ is the spin susceptibility, $`\xi _d(T)=C(T)/\mathrm{\Omega }_d(T)`$ where $`C(T)`$ is the spin wave velocity, and $`\mathrm{\Omega }_d(T)`$ is the longitudinal spin resonance frequency.
On the other hand, $`\widehat{d}`$ vector of the $`\widehat{d}`$-soliton can be parametrized by the following expression.
$$\widehat{d}=\mathrm{cos}\psi \widehat{z}+\mathrm{sin}\psi \widehat{y},$$
(5)
where
$$\psi (y,z)=\frac{1}{2}\left(\mathrm{arctan}\frac{z+R/2}{y}\mathrm{arctan}\frac{zR/2}{y}\right),$$
(6)
where we put two half-quantum vortices at $`(y,z)=(0,R/2)`$ and $`(0,R/2)`$.
In the past, similar form of $`\psi `$ was also discussed in a different context in regard to <sup>3</sup>He. As one can see, there is a discontinuity in $`\psi `$ across the line defined by $`R/2<z<R/2`$ and $`y=0`$. The spatial configuration of the corresponding $`\widehat{d}`$ around a pair of half-quantum vortices is shown in Fig.1. As one can see from the figure, the planar defect is parallel to the $`\widehat{z}`$-direction or $`c`$-axis. One can also consider the planar defect lying along the $`\widehat{y}`$-axis given by
$$\psi =\frac{1}{2}\left(\mathrm{arctan}\frac{y+R/2}{z}\mathrm{arctan}\frac{yR/2}{z}\right),$$
(7)
where two half-quantum vortices are located at $`(y,z)=(R/2,0)`$ and $`(R/2,0)`$. The configuration of the $`\widehat{d}`$ vector using the above $`\psi `$ is shown in Fig.2. One can easily see that the free energies, $`f_d`$, associated with two possible $`\widehat{d}`$ configurations are the same.
The total free energy of the $`\widehat{d}`$-soliton and a pair of half-quantum vortices is given by
$`f_{pair}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\chi _NC^2{\displaystyle 𝑑y𝑑z[K(\mathrm{\Phi })^2i+\underset{ij}{}|_i\widehat{d}_j|^2+\xi _d^2\mathrm{sin}^2\psi ]}`$ (8)
$`=`$ $`{\displaystyle \frac{1}{2}}\chi _NC^2(\pi K\mathrm{ln}{\displaystyle \frac{\lambda }{R}}+I_1+I_2),`$ (9)
where $`\mathrm{\Phi }`$ represents the phase of the order parameter which couples to the external electromagnetic field. The parameter $`K`$ is defined by
$$K=\frac{\rho _s}{\rho _{sp}}=\frac{1+1/3F_1}{1+1/3F_1^a}\frac{1+1/3F_1^a(1\rho _s^0)}{1+1/3F_1(1\rho _s^0)},$$
(10)
where $`\rho _s`$ and $`\rho _{sp}`$ are the superfluid density and the spin superfluid density respectively. $`F_1`$ and $`F_1^a`$ are the Landau Parameters and $`\rho _s^0`$ ($`1Y(T)`$ and $`Y(T)`$ is the Yosida function) is the superfluid density without the Fermi liquid correction. Notice that $`K(T_c)=1`$ at $`T=T_c`$ and $`K(0)=\frac{1+1/3F_1}{1+1/3F_1^a}`$ at $`T=0`$. The temperature dependence of the parameter $`K`$ is shown in Fig.3 assuming that $`F_1=9`$ and $`F_1^a=0`$. This choice of the parameters will be explained later.
The first term in the first and second lines of Eq.9 is the contribution from two half-quantum vortices and represents the fact that these half-quantum vortices repel each other. $`I_{1,2}`$ are the contributions from the second and third terms in the first line of Eq.9. These contributions come from the disclination of the $`\widehat{d}`$-vector.
Using the form of $`\psi (y,z)`$ discussed above, Eq.6, $`I_1`$ and $`I_2`$ can be obtained as follows.
$`I_1`$ $`=`$ $`{\displaystyle \frac{1}{4}}{\displaystyle 𝑑y𝑑z\frac{R^2}{[y^2+(z+R/2)^2][y^2+(zR/2)^2]}}`$ (11)
$`=`$ $`\pi \mathrm{ln}{\displaystyle \frac{R}{\xi }},`$ (12)
$`I_2`$ $`=`$ $`{\displaystyle \frac{1}{2\xi _d^2}}{\displaystyle 𝑑y𝑑z\left(1\frac{y^2+z^2R^2/4}{\sqrt{(y^2+z^2R^2/4)^2+y^2R^2}}\right)}`$ (13)
$`=`$ $`\pi \left({\displaystyle \frac{R}{2\xi _d}}\right)^2\mathrm{ln}{\displaystyle \frac{4\xi _d}{R}},`$ (14)
where $`\xi _d`$ is the length scale associated with the spin-orbit coupling defined by $`\xi _d(T)=C(T)/\mathrm{\Omega }_d(T)`$.
By minimizing $`f_{pair}`$ with respect to $`R`$, we obtain the optimal $`R_0`$ for the lowest free energy configuration of a pair of half-quantum vortices and the $`\widehat{d}`$-soliton. The optimal $`R_0`$ is given by
$$R_0^2=\frac{(K1)2\xi _d^2}{\mathrm{ln}\frac{4\xi _d}{\sqrt{e}R_0}}>0.$$
(15)
Here we have assumed that the $`\xi _d>\sqrt{e}R_0/4`$. Notice that the half-quantum vortices with a $`\widehat{d}`$-soliton is possible only when $`K>1`$ in order to have $`R_0>0`$. Although we have no information about $`F_1^a`$, it is most likely that $`F_1^a0`$. The ratio between the effective mass and the bare mass, $`\frac{m^{}}{m}`$, is about $`4`$, which means that $`F_19`$. Therefore, $`K>1`$ in the superconducting state, as one can see from Eq.10. Thus this condition is always satisfied below $`T_c`$. However, the existence of the solution for $`R_0`$ depends on the value of $`K`$. We find that the solution exists only if $`1<K1.5`$. For example, for $`K=1.5`$, $`\xi _d/R=0.85`$. Since the parameter $`K`$ depends on temperature as shown in Fig.3, we find that a pair of half-quantum vortices with $`\widehat{d}`$-soliton exist only for $`0.78T/T_c<1`$.
Now the free energy of a pair of half-quantum vortices and the $`\widehat{d}`$-soltion at the optimal $`R_0`$ can be obtained as
$$f_{pair}=\frac{1}{2}\pi \chi _NC^2[K\mathrm{ln}\frac{\lambda }{\xi }+\frac{(K1)}{2}\mathrm{ln}\frac{\mathrm{\Lambda }\xi ^2}{2(K1)\xi _d^2}+\frac{K1}{2}]$$
(16)
where $`\mathrm{\Lambda }=\mathrm{ln}\frac{4\xi _d}{\sqrt{e}R_0}`$. In order to examine the stability of the half-quantum vortices, we have to compare $`f_{pair}`$ and the free energy of single vortex, $`f_v`$. The difference is given by
$$f_vf_{pair}=\frac{1}{2}\pi \chi _NC^2[\mathrm{ln}\frac{\lambda }{\xi }+\frac{(K1)}{2}\mathrm{ln}\frac{2(K1)\xi _d^2}{\mathrm{\Lambda }\lambda ^2}\frac{(K1)}{2}].$$
(17)
If $`f_vf_{pair}>0`$ for some values of $`K>1`$, a pair of half-quantum vortices are more stable than the conventional single vortex. This condition can be rewritten as
$$\frac{\lambda }{\xi }\left(\frac{\xi _d}{\lambda }\right)^{K1}>\frac{e^{(K1)/2}\mathrm{\Lambda }^{(K1)/2}}{2^{(K1)/2}(K1)^{(K1)/2}}.$$
(18)
Recalling that the solution of Eq. (15) exists if $`1<K1.5`$, one can investigate the stability condition given by Eq.18. One can see from Eq.15 and Eq.18 that, as long as $`K>1`$, a pair of half-quantum vortices can be stabilized over single vortex under certain conditions for the ratio between $`\xi _d`$ and $`\lambda `$. For example, for $`K=1.1`$ and $`K=1.5`$, the conditions for the stability of a pair of half-quantum vortices over single vortex are given by
$$\frac{\xi _d}{\lambda }>10^{11}\mathrm{and}\frac{\xi _d}{\lambda }>0.0094,$$
(19)
respectively. One can see that these conditions are easily satisfied. Here we use $`\lambda /\xi =\sqrt{(\lambda _b\lambda _c)/(\xi _b\xi _c)}=12.186`$ which is appropriate for Sr<sub>2</sub>RuO<sub>4</sub>. One can also see that the stability of a pair of half-quantum vortices with the $`\widehat{d}`$-soliton is determined by the value of $`K`$ which depends on temperatures as shown in Fig.3.
Now let us discuss the relation between $`\widehat{l}`$\- and $`\widehat{d}`$-solitons. It is difficult to estimate the energy of $`\widehat{l}`$-soliton in terms of the texture free energy given by Eq.4. However, it is likely that $`\widehat{l}`$-soliton costs much more energy because, if it exists, the order parameter given by Eq.2 should vanish inside the $`\widehat{l}`$-soliton. Therefore, if there is a natural passage for conversion of $`\widehat{l}`$-solitons to $`\widehat{d}`$-solitons, most of $`\widehat{l}`$-solitons will be converted into $`\widehat{d}`$-solitons.
In summary, assumming that the superconducting state of Sr<sub>2</sub>RuO<sub>4</sub> is characterized by the spin-triplet order parameter with broken time reversal symmetry, we investigated the existence of half-quantum vortices and associated topological defect; $`\widehat{d}`$-soliton. We showed that a pair of half-quantum vortices attached to a $`\widehat{d}`$-soliton can be created in the presence of the magnetic field parallel to the $`a`$-$`b`$ plane. It was found that a pair of half-quantum vortices with a $`\widehat{d}`$-soliton is more stable than the conventional single vortex for certain temperatures below $`T_c`$. As in superfluid <sup>3</sup>He-A, the presence of $`\widehat{d}`$-soliton may be detected as the deficit in the intensity of electron spin resonance signal at $`\omega =\mathrm{\Omega }_d`$. There should be a clear electron spin resonance signature due to the half-quamtum vortices. Detection of the half-quantum vortices by scanning tunneling microscopy (STM) would also provide a convincing evidence for the spin-triplet pairing state with time reversal symmetry breaking.
We thank Manfred Sigrist and Ying Liu for helpful discussions. The work of H.-Y. Kee was conducted under the auspices of the Department of Energy, supported (in part) by funds provided by the University of California for the conduct of discretionary research by Los Alamos National Laboratory. This work was also supported by NSF CAREER award grant No. DMR-9983783 (Y.B.K.) and Alfred P. Sloan Foundation (Y.B.K.).
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# X-rays and regions of star formation: a combined ROSAT-HRI/near-to-mid IR study of the 𝜌 Oph dark cloudTable A1, Fig. A1, Table B1 and Fig. B1 are only available in the on-line edition of the Journal (Table B1 is also available at the CDS).
## 1 Introduction
The $`\rho `$ Ophiuchi dark cloud complex is one of the nearest active site of low-mass star formation (see Wilking wilking92 (1992) for a review). It is composed of two main dark clouds, L1688 and L1689, from which filamentary dark clouds, called streamers, extend to the north-east over tens of parsecs (e.g., Loren loren89 (1989); de Geus et al. degeus90 (1990)). The main star formation activity is observed in the westernmost dark cloud, L1688, which shows a rich cluster of low mass young stellar objects (YSO) around two dense molecular cores, “core A” and “core F” in the terminology of Loren (loren89 (1989)) and Loren et al. (loren90 (1990)).
The distance to the molecular complex remains somewhat controversial (see Wilking wilking92 (1992)), with a usually adopted distance $`d160`$ pc from the Sun. From Hipparcos parallaxes and Tycho B–V colors of classes V and III stars, Knude & H$`ø`$g (knude98 (1998)) have detected at $`d=120`$ pc an abrupt rise of the reddening as expected from a molecular cloud. Based on the Hipparcos positions, proper motions, and parallaxes, de Zeeuw et al. (dezeeuw99 (1999)) gives $`d=145\pm 2`$ pc for the mean distance of the Upper Scorpius OB association. We adopt $`d=140`$ pc in this article, instead of 160 pc used in our previous work.
From infrared (IR) observations of star-forming regions, Lada and collaborators (e.g. Lada lada87 (1987); Wilking, Lada, & Wilking wilking89 (1989), hereafter WLY) introduced an IR classification and distinguished different stages of evolution of young stellar objects (YSO). This classification was subsequently revisited by André & Montmerle (AM (1994), hereafter AM) to incorporate results of millimeter continuum studies on circumstellar dust. The IR sources are classified in three classes, according to their spectral energy distributions (SEDs). This classification, initially defined empirically, is now well understood in terms of evolution of low-mass stars at their earliest stages. Submillimeter observations led to the discovery of cold objects, younger than the IR sources, and thus to the introduction of a fourth class named “Class 0” (André et al. andre93 (1993), andre00 (2000)). Class 0 sources are very young protostars, peaking in the submillimeter range, at the beginning of the main accretion phase. Class I sources are evolved IR protostars, optically invisible, in the late accretion phase. Class II sources are YSO surrounded by optically thick circumstellar disks. Class III sources are YSO with an optically thin circumstellar disk or no circumstellar disk. Studies of optically visible YSO, T Tauri stars, led to another classification based on the H<sub>α</sub> line, which separates “classical” T Tauri stars (CTTS) from “weak-line” T Tauri stars (WTTS) according to their equivalent width in emission, with a boundary at EW\[H<sub>α</sub>\]$`520`$ Å, depending on the spectral type (Martín 1997). CTTS and WTTS are usually taken to be identical to Class II and Class III sources respectively, on the basis of their IR SED (see AM for a discussion about these two classifications). We will associate in this article Class II (Class III) sources with CTTS (WTTS).
Several ground-based near-IR surveys (e.g.Wilking et al. wilking89 (1989); Greene et al. GWAYL (1994), hereafter GWAYL; Barsony et al. barsony97 (1997), hereafter BKLT; and references therein) discovered in a $``$1 square degree area around the densest regions (with survey completeness limit down to $`K14`$), $``$100 low-luminosity embedded sources. More recently, the ISOCAM camera on-board the Infrared Space Observatory satellite imaged a half square degree centered on L1688 in the mid-IR (LW2 and LW3 filters, respectively centered at 6.7 $`\mu `$m and 14.3 $`\mu `$m – ISOCAM central programme surveys by Nordh et al.; see Abergel et al. abergel96 (1996)), and recognized 68 new faint young stars with infrared excess (Bontemps et al. bontemps00 (2000)).
Near-IR spectroscopy has been used to determine spectral types of an increasingly large number of $`\rho `$ Oph YSO (see the pioneering works of Greene & Meyer greene95 (1995), and Greene & Lada greene96 (1996)). Recently, Luhman & Rieke (luhman99 (1999)) obtained $`K`$-band spectroscopy for $``$100 sources, combining a magnitude-limited sample in the cloud core ($`K12`$) with a representative population from the outer region of the cluster ($`K11`$).
The $`\rho `$ Oph dark cloud YSO have also been extensively studied in X-rays. Early observations with the Einstein Observatory satellite showed that at the T Tauri star stage YSO are bright and variable X-ray emitters in the 0.2–4 keV energy band (Montmerle et al. montmerle83 (1983)). When the S/N ratio is sufficient large, their X-ray spectra can be fitted by a thin thermal model, with temperatures $`1`$ keV and absorption column densities $`N_H10^{20}`$$`10^{22}`$ cm<sup>-2</sup>. Variability studies and modeling led to explain the X-ray emission in terms of bremsstrahlung from a hot ($`T_\mathrm{X}10^7`$ K) plasma trapped in very large magnetic loops, in other words in terms of an enhanced solar-like flare activity (see reviews by Montmerle et al. montmerle93 (1993); and Feigelson & Montmerle FM (1999), hereafter FM). Casanova et al. (CMFA (1995)) – hereafter CMFA – reported deep ROSAT Position Sensitive Proportional Counter (PSPC ) imaging of the $`\rho `$ Oph cloud dense cores A and F. They detected in the $`35\mathrm{}\times 35\mathrm{}`$ central portion of the field (the inner ring of the ROSAT detector entrance window support structure) 55 X-ray sources in the 1.0–2.4 keV energy band. For three X-ray sources, one or several Class I sources lie within the error boxes of X-ray peaks, but other counterparts are possible (unclassified IR sources, T Tauri stars). X-ray emission from one of these Class I sources, YLW15 (=IRS43 in WLY), was unambiguously confirmed with a follow-up ROSAT High Resolution Imager (HRI ) observation by Grosso et al. (grosso97 (1997)). The outer portion of the CMFA PSPC field, analyzed by Casanova (casanova94 (1994)), contains 36 X-ray sources. The optical spectroscopic classification of these X-ray sources and other X-ray selected stars in the $`\rho `$ Oph dark cloud vicinity, based on H<sub>α</sub> and Li I (670.8 nm) spectroscopy, was made by Martín et al. (martin98 (1998)), doubling the number of PMS stars spectroscopically classified in the $`\rho `$ Ophiuchi area.
Observations of harder X-ray ($`>`$4 keV) from the $`\rho `$ Oph dark cloud were initially only possible with non-imaging instruments. Tenma and Ginga revealed unresolved emission from the cloud core region, with a hard X-ray spectrum with $`kT_\mathrm{X}4`$ keV and $`N_H10^{22}`$ cm<sup>-2</sup> (Koyama koyama87 (1987); Koyama et al. koyama92 (1992)). Wide-energy band imaging observations became possible with ASCA in the range 0.5–10 keV. In the $`\rho `$ Oph dark cloud, Koyama et al. (koyama94 (1994)) detected hard X-rays from T Tauri stars, with $`kT_\mathrm{X}`$ up to $``$8 keV in the case of the WTTS DoAr21. There is also some evidence for unresolved hard X-ray emission from embedded young stars below the point source detection limit. From this ASCA observation, Kamata et al. (kamata97 (1997)), found additional T Tauri stars and detected three X-ray sources associated with Class I sources, but with large X-ray error boxes (15$`\mathrm{}`$–30$`\mathrm{}`$).
There is a deep connection between IR and X-ray observations of star-forming regions. Sensitive ground-based near-IR surveys penetrate dark clouds (except for dense cores) so that their source populations are frequently dominated by ordinary stars in the Galactic disk. Space-based mid-IR isolates YSO with significant circumstellar material and effectively eliminates the background star population, but they will miss the recognition of YSO with less massive or absent disks. X-ray emission, in contrast, is elevated by 1–4 orders of magnitude in YSO of all ages, irrespective of a disk presence. It thus provides a unique tool for improving the census of young star clusters.
In this article, we present the results from the HRI follow-up of the CMFA PSPC observation. The high angular resolution of these observations allows us to find counterparts to all X-ray sources without ambiguity. The comparison with the sensitive ISOCAM survey of the $`\rho `$ Oph dark cloud significantly improves the existing classification of these counterparts and allows us to do statistical studies on a well defined sample.
We first present the ROSAT HRI observations: image analysis, source detection and identification ($`\mathrm{\S }`$2). We incorporate the ISOCAM survey results from Bontemps et al. (bontemps00 (2000)) and we present the resulting IR classification for the HRI sources ($`\mathrm{\S }`$3). The next sections discuss the X-ray luminosity of the HRI detected TTS ($`\mathrm{\S }`$4), and the X-ray detectability of the embedded TTS population ($`\mathrm{\S }`$5). Next (§6), we show that the HRI census of Class III sources cannot be complete, and that numerous unknown low-luminosity Class III sources, perhaps including brown dwarfs, must exist. Summary of the main results and conclusions are presented in $`\mathrm{\S }`$7, where prospects for improvements with XMM-Newton and Chandra, are also discussed.
Appendix A gives details about the HRI X-ray source detection, and lists the X-ray detections. Optical finding charts, and identification list of the HRI X-ray sources can be found in Appendix B. Appendix C compares these HRI observations with previous PSPC ones. Appendix D discusses the status of optical/IR counterparts without IR classification.
## 2 The ROSAT HRI X-ray observations
### 2.1 The ROSAT HRI images
We have observed the dense cores A and F of the $`\rho `$ Oph dark cloud with the ROSAT HRI. The detector is sensitive to the 0.1–2.4 keV energy range, but has no spectral resolution. The two observation fields were respectively centered on approximatively the WTTS DoAr21 ($`\alpha =16^\mathrm{h}26^\mathrm{m}2\stackrel{s}{.}4,\delta =24\mathrm{°}23\mathrm{}24\mathrm{}`$ \[J2000\]) and on the Class I protostar YLW15 = IRS43 ($`\alpha =16^\mathrm{h}27^\mathrm{m}26\stackrel{s}{.}4,\delta =24\mathrm{°}40\mathrm{}48\mathrm{}`$ \[J2000\]). ROSAT HRI images have a diameter of $`40\mathrm{}`$. Fig. 1 displays the two observation fields, which include the dense DCO<sup>+</sup> cores A, B, C, E, and F (Loren et al. loren90 (1990)), most of the area studied by CMFA and ASCA, as well as the ISOCAM survey. The first observation field, centered approximatively on $`\rho `$ Oph A, will be referred to as the “core A field”; the second observation field, centered on dense core F, will be called the “core F field”. The core A field was observed between 1995 August 29 and 1995 September 12 with a total exposure of 51.3 ks. The core F field was observed at three different epochs (hereafter observations #1, #2, and #3): between 1995 March 9 and 14 (12.5 ks), between 1995 August 18 and 20 (27.5 ks), and between 1996 September 7 and 11 (37.2 ks). These three different epochs give a total exposure of 77.2 ks (see Table 1 for the log of ROSAT HRI observations details).
### 2.2 The HRI image analysis
We have analyzed separately the four data sets. Standard source detection algorithms were used to find X-ray sources, and to search optical counterparts for each set. We then selected X-ray sources with the best position accuracy (usually the brightest not too far away from the field center) and having an unambiguous counterpart, corrected all X-ray positions from the existing offsets, and used this improved astrometry to remove possible ambiguities in the identification of the X-ray sources. Details can be found in Appendix A.
We find 63 HRI X-ray sources. Fig. 2 shows the positions of these sources, superimposed on a combined IR and optical image of the $`\rho `$ Oph cloud. Coordinates, error boxes, likelihoods of existence, and count rates of these HRI X-ray sources can be found in Appendix A (Table 2).
### 2.3 HRI X-ray source identifications
Identification of the HRI X-ray sources was made by cross-correlations with published lists of confirmed or suspected cloud members (AM), IR surveys (Greene & Young 1992; BKLT; and Bontemps et al. bontemps00 (2000), see $`\mathrm{\S }`$3), $`K`$-band spectroscopic (Luhman & Rieke luhman99 (1999)), radio surveys (André, Montmerle, & Feigelson andre87 (1987); Stine et al. stine88 (1988); Leous et al. leous91 (1991)), and with previously known X-ray sources (PSPC: CMFA, Martín et al. martin98 (1998); ASCA: Kamata et al. kamata97 (1997)). For X-ray sources without published counterparts we have used SIMBAD,<sup>1</sup><sup>1</sup>1On-line version at http://simbad.u$``$strasbg.fr/Simbad . and we have also searched optical counterparts on optical red band images from the Digitized Sky Survey.<sup>2</sup><sup>2</sup>2On-line version on the ESO site: http://arch$``$http.hq.eso.org/cgi$``$bin/dss .
As shown by the finding charts in Appendix B, thanks to the good angular resolution of the ROSAT HRI ( PSF FWHM $`5\mathrm{}`$ on axis; due to the mirrors as well as the detector), the position accuracy ($`1\mathrm{}`$$`6\mathrm{}`$, see Col. 6 Table 2) allows us to find counterparts almost without ambiguity.<sup>3</sup><sup>3</sup>3When two possible counterparts are in the error box (this happens only three times, see the finding charts in Appendix B, and Table 3 notes), we take the more luminous in the $`J`$-band, since the X-ray luminosity of TTS is correlated with the $`J`$-band luminosity (see CMFA, and below $`\mathrm{\S }`$3.2, $`\mathrm{\S }`$7.2). Identification lists for the core A and core F fields are given in Appendix B, Table 3. We also discuss in Appendix B the identification of X-ray sources with a low statistical significance.
Nearly $`90\%`$ of the HRI X-ray sources are identified. We detect only $`70\%`$ of the PSPC X-ray sources (CMFA; Casanova casanova94 (1994)), but this can be explained by the difference in sensitivity between the two instruments, and the intrinsic variability of the X-ray sources (see Appendix C).
## 3 Nature and IR properties of the HRI sources
Following the results of CMFA, essentially all X-ray sources we found in the $`\rho `$ Ophiuchi dark cloud should be young stars, a number of them being still embedded in the cloud. Embedded YSO are mainly studied at IR and millimeter wavelengths, but they may all be potential X-ray emitter regardless of their IR classification. Indeed, using the results of Wilking et al. (wilking89 (1989)), AM, and GWAYL, CMFA analyzed their results in the light of the IR observations of the stars they observed in X-rays. We can then (i) use the published IR surveys to provide a list of recognized YSO members of the cluster to be compared to the observed X-ray properties, and (ii) use X-ray emission to discriminate between true cluster members and the many potential background stars seen in IR images.
### 3.1 New Class II and Class III source census after ISOCAM
Near-IR surveys of the $`\rho `$ Ophiuchi cluster, sensitive enough to detect low-luminosity embedded young stars, have recently been published (Comerón et al. comeron93 (1993); Strom et al. strom95 (1995) – hereafter SKS; BKLT). However these ground-based surveys encountered limitations in recognizing the nature of all the embedded sources, and have therefore not much increased the number of bona fide members of the $`\rho `$ Oph cluster. The mid-IR camera aboard ISO, ISOCAM, produced a map of the $`\rho `$ Oph main cloud, used by Bontemps et al. (bontemps00 (2000)) to study the young star population. This mid-IR study resulted in a significantly more complete census of the $`\rho `$ Oph cluster population. We here use this new census as a basis for discussion about the nature of the detected X-ray sources and to estimate the occurrence and properties of the X-ray emission of the different classes of YSO.
The mid-IR photometry at 6.7 and 14.3 $`\mu `$m appears invaluable to characterize sources with IR excesses, i.e., Class I sources and Class II sources (e.g., Nordh et al. nordh96 (1996); Bontemps et al. bontemps98 (1998)). In $`\rho `$ Oph, Bontemps et al. (bontemps00 (2000)) have doubled the number of Class II sources known. These authors conclude that the sample is complete down to stellar luminosities, $`L_{}`$, as low as 0.03 L, thus extending downwards the luminosity function obtained from the ground by about one order of magnitude. However, just as in the near-IR, these measurements alone cannot characterize the nature of sources without IR excess: these can be either Class III sources (diskless TTS) or background sources.
The tentative Class III classification of a number of sources coming from previous X-ray observation (CMFA), or from this article, is now confirmed by ISOCAM, which detected no IR excess. We call these sources “new” Class III sources (see Table 3, and $`\mathrm{\S }`$5).
Since we already identified one of the HRI sources as a foreground star (ROXF37 = HD148352, see Table 3), the question arises that a number of these new sources with Class III spectra may also be field stars, contaminating the genuine young star sample. Guillout et al. (guillout96 (1996)) has estimated the stellar content of flux-limited X-ray surveys, based on the age-dependent stellar population model developed by the Besançon group. For the average sensitivity of our HRI observation ($`6\times 10^4`$ cts s<sup>-1</sup>, see Fig. 6), the $`\rho `$ Oph galactic latitude ($`20\mathrm{°}`$), and our total field of view (0.5 square degrees), this model yields an estimate of 5 contaminating stars (P. Guillout, private communication). Removing HD148352, this leaves 4 possible field star candidates among the 21 “new” Class III sources and “Class II-III” sources listed in Table 3. This number is small enough not to affect significantly our discussion, and in what follows we will simply neglect the possible contamination of our various Class III samples by field stars.
The Class III source population will be studied in detail below ($`\mathrm{\S }`$4 and $`\mathrm{\S }`$6).
### 3.2 Cloud extinctions and stellar luminosities from near-IR photometry
YSO suffer from large amounts of extinction by dust – cloud dust plus circumstellar dust – (up to A$`{}_{\mathrm{V}}{}^{}=60`$ or more in $`\rho `$ Oph), which strongly affects their fluxes at all wavelengths of interest, but does not necessarily prevent their detection in soft X-rays (see, e.g., CMFA). Despite this effect, it is possible to estimate the total extinctions along the line of sight and stellar luminosities. Near-IR photometry data appear to provide the most reliable estimate for the luminosity of the embedded young stars (see discussion in Bontemps et al. 2000 and references therein): the $`J`$-band fluxes are usually almost purely photospheric and thus trace the stellar luminosities very well (e.g., GWAYL; SKS). Similarly the $`JH`$ color is a sensitive tracer of interstellar extinction. We therefore use the extinctions and stellar luminosities derived from near-IR photometry for the ISOCAM sample by Bontemps et al. (bontemps00 (2000)).
In CMFA the bolometric luminosities were approximately estimated from the $`J`$-band fluxes using the GWAYL conversion based on an observational correlation between $`J`$ fluxes and bolometric luminosities for Taurus and Chameleon TTS. The bolometric luminosity comprises in principle the total luminosity (accretion + stellar) of a YSO. However for young PMS stars like Class II sources and Class III sources, this luminosity should coincide with their stellar luminosity since their accretion luminosity (if any) has become significantly smaller than the photospheric luminosity. Finally we note that the conversion between the $`J`$-band absolute magnitude and the stellar luminosity used by Bontemps et al. (bontemps00 (2000)) is numerically similar to the CMFA conversion between the dereddened $`J`$-band and the bolometric luminosity.
Cols. 8–9 of Table 3 give for each source its more up-to-date values of interstellar extinction and stellar luminosity, determined by Bontemps et al for 140 pc. The reader will find the details of the calculations in Bontemps et al. (bontemps00 (2000)).
### 3.3 Nature of the HRI sources from IR data
The YSO evolutionary stage, inferred from IR spectral energy distributions (see $`\mathrm{\S }`$3.1), is available for most of the X-ray sources. The resulting census of the 55 X-ray sources with stellar counterparts is: one Class I protostar (YLW15=IRS43); 23 Class II sources, including 4 new ISOCAM Class II sources; 21 Class III sources, including 13 new Class III sources<sup>4</sup><sup>4</sup>44 were already PSPC Class III source candidates, and are confirmed by ISOCAM, which detected no IR excess; 3 were probably detected by the PSPC, but due to its lower angular resolution the optical or IR counterpart was uncertain (see Col. 2 of Table 3 and attached notes); 6 are genuine new X-ray detections, probably resulting from variability.; 8 whose classification is either Class II or Class III sources (see discussion in Appendix D); one Class III early-type star (S1, with spectral type B3; see André et al. andre88 (1988)), and one main sequence foreground star (HD148352 with spectral type F2V; see the Hipparcos catalogue).
In Table 3, Col. 6 lists cross-identifications with the ISOCAM survey: red (blue ) means ISOCAM sources with (without) IR excess. Col. 7 gives the IR classification, Cols. 8–9 the extinctions and stellar luminosities (from Bontemps et al. bontemps00 (2000)).
Therefore, the present analysis has revealed no new X-ray emitting Class I source, apart from the Class I protostar YLW15 which has been the subject of a specific study (Grosso et al. grosso97 (1997)). The HRI/ISOCAM sources are overwhelmingly TTS, for which improved results are described in the following sections.
## 4 X-ray luminosities of the HRI-detected T Tauri stars
### 4.1 Derivation from source count rates
We assume a fiducial TTS X-ray spectrum (see Montmerle montmerle96 (1996)) having $`kT_\mathrm{X}`$=1 keV plasma, with cosmic abundances, and Raymond-Smith line emissivities, and with interstellar absorption based on Morrison & McCammon (morrison83 (1983)) cross sections. We use the relation given by Ryter (ryter96 (1996)) to estimate the hydrogen column density, $`N_H`$, from the visual extinction, $`A_\mathrm{V}`$, determined from IR data<sup>5</sup><sup>5</sup>5Noted $`A_{\mathrm{V},\mathrm{IR}}`$ in CMFA.: $`N_H=2.23\times 10^{21}A_\mathrm{V}`$ cm<sup>-2</sup>. The intrinsic (i.e., extinction corrected) X-ray luminosities in the full ROSAT energy band (0.1–2.4 keV), $`L_\mathrm{X}`$, were calculated for classified sources from source count rates for $`d140`$ pc using EXSAS, and are given in Cols. 10–13 of Table 3. The intrinsic X-ray luminosities span the range $`L_\mathrm{X}1\times 10^{29}`$$`4\times 10^{31}`$ erg s<sup>-1</sup>. (For X-ray sources without extinction estimates we label the detection exposure with a question mark.)
### 4.2 Luminosity functions of HRI-detected Class II and Class III sources
We here compare the extinction corrected X-ray luminosities of the Class II and Class III sources detected by the HRI in the ISOCAM field<sup>6</sup><sup>6</sup>622 Class II sources and 19 Class III sources; for the core F observation we take the mean X-ray luminosity in Col. 13 of Table 3. in order to evaluate the contribution of the circumstellar disk to the X-ray absorption, or to the X-ray emission (for instance by magnetic reconnection between the star and the disk; see Montmerle et al. montmerle00 (2000)).
Fig. 3 shows the cumulative X-ray luminosity distribution functions for these two populations, estimated using the ASURV statistical software package (rev. 1.2;<sup>7</sup><sup>7</sup>7Version available at http://www.astro.psu.edu/statcodes . La Valley et al. lavalley92 (1992)), which takes upper limits into account. These distributions are mathematically identical to the maximum likelihood Kaplan-Meier estimator (Feigelson & Nelson feigelson85 (1985)). Their mean X-ray luminosities (in erg s<sup>-1</sup>) are given by $`<\mathrm{log}(L_\mathrm{X})>=30.3\pm 0.1`$ for Class II sources, and $`<\mathrm{log}(L_\mathrm{X})>=30.4\pm 0.2`$ for Class III sources. We used nonparametric two-sample tests implemented in ASURV — Gehan’s generalized Wilcoxon test, Logrank test, Peto & Peto generalized Wilcoxon test, Peto & Prentice generalized Wilcoxon test — to see whether the difference between the two luminosity functions is significant. These tests gave a high probability (46–72 $`\%`$) that they are statistically indistinguishable. This result agrees with previous deep studies of the $`\rho `$ Oph main cloud (CMFA) and Chamaeleon I (Feigelson et al. feigelson93 (1993); Lawson et al. lawson96 (1996)) YSO populations.
In contrast, Neuhäuser et al. (neuhaeuser95 (1995)) found in the Taurus-Auriga star-forming region that Class III sources are more X-ray luminous than Class II sources. Neuhäuser et al. (neuhaeuser95 (1995)) used the ROSAT All Sky Survey (RASS ) to cover a large area of the Taurus-Auriga ($`900`$ square degrees), including dense cores, but also away from this star-forming region. This shallow survey and large area must be compared to our deep pointed observations, where our field of interest covers 0.5 square degrees. In our observations we focus only on the dense cores, studying a younger population of YSO. Contamination of the T Tauri star sample by a more evolved, wide-spread, and older Class III source population, may explain the discrepancy with the result of Neuhäuser et al. (neuhaeuser95 (1995)).<sup>8</sup><sup>8</sup>8In addition one should note that the RASS includes the soft band (0.1–1 keV), which is more sensitive to the extinction than the hard band (1–2.4 keV) taken by CMFA.
Thus, we confirm that the contribution of the disk of Class II sources to their X-ray emission, or to X-ray absorption, must be small.
Next, we combine the X-ray luminosity distributions of Class II and Class III sources to obtain the cumulative X-ray luminosity distribution function of all TTS detected by the HRI. For $`\mathrm{log}(L_\mathrm{X})`$ between 29.7 ($`\mathrm{log}(L_{\mathrm{X},\mathrm{break}})`$) and 31.3 ($`\mathrm{log}(L_{\mathrm{X},\mathrm{max}})`$), the distribution is loglinear: $`𝒩(>\mathrm{log}(L_\mathrm{X}))=21.9\times (\mathrm{log}(L_\mathrm{X})31.3)`$. For $`\mathrm{log}(L_\mathrm{X})\mathrm{log}(L_{\mathrm{X},\mathrm{break}})`$, the distribution shows a downwards trend, which is due to our lower efficiency to detect weak X-ray emitting TTS. We deduce from this linear relation the total X-ray luminosity emitted by the X-ray sources with X-ray luminosity between $`L_{\mathrm{X},\mathrm{max}}`$ and $`L_{\mathrm{X},\mathrm{min}}`$ (an arbitrary value): $`L_{\mathrm{X},\mathrm{tot}}=21.9\times (L_{\mathrm{X},\mathrm{max}}L_{\mathrm{X},\mathrm{min}})`$. Thus, the total X-ray luminosity of this group of X-ray sources is dominated by the brightest sources (with $`L_{\mathrm{X},\mathrm{max}}`$), and is very weakly dependent on the loglinear fit. For $`L_{\mathrm{X},\mathrm{min}}=L_{\mathrm{X},\mathrm{break}}`$, we find $`L_{\mathrm{X},\mathrm{tot}}4.3\times 10^{32}`$ erg s<sup>-1</sup>. This value is close to the asymptotic value $`4.4\times 10^{32}`$ erg s<sup>-1</sup> obtained taking $`L_{\mathrm{X},\mathrm{min}}=0`$, thus future detections of new X-ray emitting TTS with low X-ray luminosity will not greatly affect this result.
### 4.3 HRI source variability
We present in Appendix C a study of the variability of the X-ray sources which were observed both by the HRI and the PSPC, and we show that some sources were in a high X-ray state during the HRI or the PSPC observations. Here, we study the variability of the Core F HRI sources.
Core F observations comprise three time-separated observations, which allow us to reiterate the “Christmas tree” luminosity study made with Einstein Observatory by Montmerle et al. (montmerle83 (1983)). The idea suggested by the similarity between Class II sources and Class III sources in X-rays was to assume that all the X-ray sources are basically one single type of X-ray object, seen in different states. The result was that the distribution of the flux variations could be approximated by a power-law with an index $`\beta =1.4`$.
We have estimated whenever possible for each Core F HRI source the X-ray flux variations from the observed high/low luminosity ratio, $`L_{\mathrm{X},\mathrm{high}}/L_{\mathrm{X},\mathrm{low}}`$, based on the three observations. This yields 27 values including 11 lower limits. The integral distribution for a given ratio is estimated using the maximum likelihood Kaplan-Meier estimator (see Fig. 4). We find a power-law distribution: $`𝒩(>L_{\mathrm{X},\mathrm{high}}/L_{\mathrm{X},\mathrm{low}})(L_{\mathrm{X},\mathrm{high}}/L_{\mathrm{X},\mathrm{low}})^\alpha `$, with an index $`\alpha =0.75\pm 0.03`$. This implies that the differential distribution $`d𝒩/d`$ follows a power-law distribution of slope $`\beta =\alpha 1=1.75`$.
We suggest following Montmerle et al. (montmerle83 (1983)) that this power-law behavior, may be explained in terms of variability due to stellar flares, if interpreted in terms of stochastic relaxation phenomenon (Rosner & Vaiana 1978), and dominating the X-ray activity of the underlying stars. Such a power-law behavior is seen in the solar flares in radio, optical, soft and hard X-ray emission with the power-law index $`\beta `$=1.1–3.0 (see review in Aschwanden et al. aschwanden98 (1998)). For soft X-ray emission $`\beta `$=1.7–1.9, which is consistent with our result, and supports the analogy with the solar magnetic activity.
## 5 X-ray detectability of the embedded T Tauri star population
We use the information given both by the ISOCAM survey and by our HRI deep exposure to study the X-ray detected TTS population of the $`\rho `$ Oph dense cores. We restrict the following studies to the HRI/ISOCAM overlapping area.<sup>9</sup><sup>9</sup>9We took $`19.2\mathrm{}`$ for the HRI field radius, because one of our sources (ROXRF31=SR9), is detected up to this angle from the axis in the Core F field. This area comprises 98 Class II sources (classified from ground-based and ISO observations), including 52 new ISOCAM Class II sources,<sup>10</sup><sup>10</sup>10We have excluded three new ISOCAM Class II sources for which we have only the $`K`$ magnitude, and thus no A<sub>V</sub> estimate. L for these sources must be small, and/or A<sub>V</sub> high, which implies a high upper limit on the intrinsic L<sub>X</sub>. This does not affect the statistical results. and 35 Class III sources (characterized as YSO from X-ray or radio observations, and classified as Class III sources from ground-based and ISO observations), including 21 new Class III sources (HRI or PSPC X-ray sources without IR excess observed by ISOCAM). We will call these sources the “TTS sample”.
Our HRI observation detected a large number of sources, yet these constitute only 30$`\%`$ of the “TTS sample”. In this section, we examine the reasons why the other TTS were not detected, and in particular whether the undetected TTS form a separate population of genuinely X-ray weak objects.
### 5.1 X-ray vs. stellar luminosities
First, to know more about the X-ray properties of the members of the whole TTS sample (with upper limits if they are not detected with the HRI), we examine whether a correlation exists between the X-ray luminosity and the stellar luminosity (both corrected from extinction), and if so, whether it is the same as the one found in $`\rho `$ Oph by CMFA. We chose for each Core F X-ray source its lowest X-ray luminosity (including HRI upper limits) to minimize the effects of X-ray variability. For TTS undetected by the HRI, we estimate count rate upper limits ($`3.25\sigma `$)<sup>11</sup><sup>11</sup>11When the YSO is in both Core A and Core F field, we use the longer Core F field exposure to estimate the count rate upper limit. using the EXSAS command COMPUTE/UPPER\_LIMITS, and we use the extinction estimate from Bontemps et al. (bontemps00 (2000)) to compute the corresponding limit on the intrinsic X-ray luminosity.
To establish the existence of a linear correlation between $`\mathrm{log}L_\mathrm{X}`$ and $`\mathrm{log}L_{}`$, we performed three statistical tests using ASURV: Cox’s proportional hazard model, the generalized Kendall $`\tau `$ test, and Spearman’s $`\rho `$ test. The probability of the null hypothesis (i.e. that this correlation is not present) is $`<10^4`$ for each of the three tests. Thus, a strong linear correlation between $`\mathrm{log}L_\mathrm{X}`$ and $`\mathrm{log}L_{}`$ is indeed present. We found the linear regression coefficients by using the Estimation Maximization (EM) algorithm under Gaussian assumptions and the Buckley-James method. The Buckley-James method gave results similar to those of the EM algorithm, but with a larger uncertainty on the slope. As the Buckley-James method is semi-nonparametric, this suggests that the residuals of the linear correlation may be non-Gaussian. We thus conservatively keep the slope uncertainty given by the Buckley-James method. The $`L_X`$$`L_{}`$ correlation is then given by (see Fig.5): log($`L_\mathrm{X}`$/erg s<sup>-1</sup>) = (1.0$`\pm `$0.2) $`\times `$ log($`L_{}/L_{}`$) + 30.1. We note that the censoring fraction is so high that the correlation line misses most of the data points and depends entirely on the location of the few lowest detections. The correlation dispersion may be due to X-ray variability, and also to TTS spectral type and age differences: Neuhäuser et al. (neuhaeuser95 (1995)) points out that the ratio $`L_\mathrm{X}/L_{}`$ increases with decreasing effective temperature, and shows a variation of L<sub>X</sub> with age.
This correlation spans three orders of magnitude in L<sub>X</sub> and two in L. The slope of this correlation, $`a`$, is equal to 1.0, and the TTS X-ray luminosity is then approximatively given by the simple proportionality: $`L_\mathrm{X}/L_{}10^4`$. We thus confirm that the characteristic for TTS in $`\rho `$ Oph is 10<sup>-4</sup>, with a large dispersion up to a level $`10^3`$. There is no evidence for the “saturation” effect seen at this level in late type main sequence stars, and attributed to the complete filling of the stellar surface by active regions (Fleming et al. 1989).
A similar correlation between L<sub>X</sub> and L was found for Class II and Class III sources in the previous $`\rho `$ Oph study of CMFA (the method used to estimate L was different, but very similar numerically; see $`\mathrm{\S }`$3.2), but also in other star-forming regions: Chamaeleon (Feigelson et al. feigelson93 (1993); Lawson et al. lawson96 (1996)), Taurus-Auriga (Neuhäuser et al. neuhaeuser95 (1995)), and IC 348 (Preibisch et al. preibisch96 (1996)). However, the slopes may not be identical: while Feigelson et al. (feigelson93 (1993)), Preibisch et al. (preibisch96 (1996)), and CMFA find the same $`a=1`$ slope as above. On the other hand, Lawson et al. (lawson96 (1996)) found $`a=0.55`$, on a better characterized, enlarged X-ray source sample in Chamaeleon, stessing the importance of having a sample as complete as possible. Nevertheless, the fact that we find the same slope as CMFA with an enlarged sample of X-ray sources in the same cloud is certainly a good internal consistency check between the PSPC and the HRI.
We find that the majority of the X-ray luminosity upper limits are above the $`L_\mathrm{X}`$$`L_{}`$ correlation. Only 5$`\%`$ of the X-ray undetected TTS are below the correlation, mixed with X-ray detected TTS. This is consistent with the idea that all TTS in $`\rho `$ Oph may be X-ray emitters with $`L_\mathrm{X}/L_{}10^4`$. Therefore, the TTS undetected by the HRI do not make up a separate population, but must have X-ray properties comparable to that of the detected population, verifying the same correlation.
### 5.2 The X-ray undetected T Tauri star population
Using the previous correlation between the stellar and X-ray luminosities, we can now estimate for each member of the TTS sample the intrinsic X-ray luminosity in the ROSAT energy band,<sup>12</sup><sup>12</sup>12We note that this method attributes a larger X-ray luminosity to the TTS having X-ray luminosities below the $`L_\mathrm{X}`$$`L_{}`$ correlation. By this effect X-ray undetected TTS can be put above the HRI detection threshold, but this concerns only 5 cases, see Fig. 6. and compare it with the HRI detection threshold to understand the X-ray detectability of the TTS sample with the HRI. However, the comparison is not straightforward, since the HRI detection threshold depends on both instrumental effects and extinction along the line of sight.
Fig. 6 shows the instrumental effects: the HRI count rate threshold ($`3.25\sigma `$) increases away from the pointing axis. We interpret this dependence as the consequence of the point spread function degradation and reduced sensitivity off-axis of the ROSAT mirrors (David et al. david97 (1997)). With the X-ray spectrum assumptions described in §4.1, we have determined using EXSAS the conversion factor, $`f`$, between the HRI counts and the apparent X-ray luminosity (i.e., in the absence of extinction) in the ROSAT energy band (0.1–2.4 keV), $`L_{\mathrm{X},\mathrm{app}}`$. We find: $`f=6.8\times 10^{28}`$ erg cts<sup>-1</sup> s for $`d=140`$ pc. The minimum X-ray luminosity for a $`3.25\sigma `$ HRI detection, $`L_{\mathrm{X},\mathrm{min}}`$, ranges from $`7\times 10^{27}`$ erg s<sup>-1</sup> on-axis (angle=0′) to $`7\times 10^{28}`$ erg s<sup>-1</sup> off-axis (angle=19.2′) (see Fig. 6).
In case the X-ray sources suffer some extinction equivalent to A<sub>V</sub> magnitudes, the values of $`L_{\mathrm{X},\mathrm{min}}`$ on-axis and off-axis must be corrected to obtain the corresponding intrinsic minimum X-ray luminosities as a function of A<sub>V</sub>: if a source is heavily extincted, this minimum may be up to two orders of magnitude higher or more than in the absence of extinction (see for example the high values of the upper limits of Fig. 5).
Fig. 7 plots the X-ray luminosities of the TTS sample as a function of $`A_\mathrm{V}`$ (or $`N_\mathrm{H}`$). These X-ray luminosities were estimated from the stellar luminosities using the correlation discussed in the previous section. The points are compared with the HRI threshold curves $`L_{\mathrm{X},\mathrm{min}}=f(A_\mathrm{V})`$, computed both on-axis and off-axis. The HRI detected TTS (crossed dots) are found to be rather bright ($`10^{29}10^{31}`$ erg s<sup>-1</sup>) and weakly extincted ($`A_\mathrm{V}10`$). The undetected TTS have estimated X-ray luminosities below the computed HRI detection threshold. In particular, we understand why the new ISOCAM Class II sources (Bontemps et al. bontemps00 (2000)), characterized both by low stellar luminosities ($`0.05`$ L) (and thus presumably low predicted X-ray luminosities, $`6\times 10^{28}`$ erg s<sup>-1</sup>), and relatively high extinctions ($`A_\mathrm{V}20`$), could not have been detected with our HRI observation.<sup>13</sup><sup>13</sup>13A few faint sources were however detected in spite of being below the nominal HRI detection threshold: they were probably in an X-ray flaring state at the time of the observations. In particular ROXRF32 = GY238, far below the HRI detection threshold, was detected only in the third Core F exposure, which supports this interpretation.
Now ISOCAM cannot per se recognize Class III YSO among its sources without IR excess, but X-rays can. However, a reliable census of Class III sources in $`\rho `$ Oph is de facto limited by the sensitivity of X-ray observations: Fig. 7 shows that the number of detected Class III sources decreases for low L<sub>X</sub> and high A<sub>V</sub>; roughly speaking Class III sources are mainly detected above L$`{}_{\mathrm{X},\mathrm{min}}{}^{}10^{29.6}`$ erg s<sup>-1</sup> (or equivalently L<sub>⋆,min</sub>=0.35 L), and below A$`{}_{\mathrm{V},\mathrm{min}}{}^{}30`$.
This strongly suggests that unknown Class III sources may exist. We can estimate their number at least in regions at the periphery of cloud cores, by using the fact that the WTTS/CTTS ratio (or equivalently the Class III/Class II source ratio) is $`1`$, and also that the HRI is equally sensitive to Class III and Class II sources (see §4). In the HRI/ISOCAM overlapping area, this ratio is 19/22 $`1`$; on a comparable area Martín et al. (martin98 (1998)) also found a WTTS/CTTS ratio $`1`$. Since the “TTS sample” comprises 88 Class II sources and 35 Class III sources above L$`{}_{}{}^{}0.03`$ L, we predict that $`(88\times 19/22)3540`$ Class III sources remain to be discovered in X-rays in the HRI/ISOCAM overlapping area above L$`{}_{\mathrm{X}}{}^{}3\times 10^{28}`$ erg s<sup>-1</sup>. These sources are not seen now either ($`i`$) because they are too faint in X-rays ($`L_XL_{\mathrm{X},\mathrm{min}}`$) –or equivalently from the existence of an $`L_\mathrm{X}vs.L_{}`$ proportionality, too faint in stellar luminosity ($`L_{}L_{,min}`$)– or ($`ii`$) too absorbed ($`A_VA_{\mathrm{V},\mathrm{min}}`$), or a combination of both. XMM-Newton will be an ideal tool to reveal such a large number of unknown Class III sources in the future (see $`\mathrm{\S }`$7), but, as shown in the next section, we can already figure out their nature to a large extent.
## 6 The unknown Class III source population
In this section, we seek to characterize the suspected unknown Class III source population, which is likely to exist on the basis of the Class II source findings by ISOCAM and other IR observations. We will use (i) the spatial distribution of all known Class II sources, and (ii) the extinction map derived from C<sup>18</sup>O observations.
### 6.1 Compared spatial distributions of the Class II and Class III sources
The conventional wisdom is that Class III sources are descendants of Class II sources after dispersion of their disks (e.g., Lada 1987; AM). As cloud cores have an internal velocity dispersion, stars form with an initial mean velocity distribution, implying that they drift away from their formation site (e.g., Feigelson feigelson96 (1996)). This is the usual explanation for the increase with distance from cloud cores of the Class III/Class II source ratio (or equivalently at that stage the WTTS/CTTS ratio), which is $`\stackrel{<}{}1`$ within the core region (e.g., CMFA), and reaches values $`1`$ far from the cores (e.g., Martín et al. martin98 (1998)). This implies a larger spread of the spatial distribution of the Class III source population compared to that of the Class II source population.
Let us first study the spatial distribution of the Class II source population within the HRI/ISOCAM area. We analyze the source surface density by using a 2-D Gaussian filter of a given FWHM on source position. The choice of the FWHM is optimized to enhance the contrast between regions of low and high source density, and thus reveal any clustering. Fig. 8 shows the resulting density map, in the form of dashed contours obtained with FWHM=6$`\mathrm{}`$. The Class II sources show three strong density peaks well centered on DCO<sup>+</sup> cores A, B and F, which is consistent with the idea that most of these sources were born in these cores. However, in spite of its comparable DCO<sup>+</sup> line-of-sight density, core C appears much poorer in Class II sources; the weaker star-forming activity of this core is confirmed by the presence of only one Class I source (see Bontemps et al. bontemps00 (2000)).
One can go one step further by comparing the source distribution with the matter distribution along the line-of-sight, i.e., with the extinction map. The DCO<sup>+</sup> radical is a good indicator of large densities ($`n10^5`$$`10^7`$ cm<sup>-3</sup>) in cold cores, but the relevant regions occupy a relatively small volume; by contrast, C<sup>18</sup>O, which is generally optically thin and sensitive to smaller densities, is a good column-density tracer. Using this molecule, Wilking & Lada (wilking83 (1983)) derived an extinction map of the $`\rho `$ Oph cloud center, showing that the denser regions have an equivalent visual extinction A<sub>V</sub> between $`30`$ and $`100`$.
Fig. 8 displays the C<sup>18</sup>O contours, labeled in A<sub>V</sub> by steps of A$`{}_{\mathrm{V}}{}^{}20`$, starting at A$`{}_{\mathrm{V}}{}^{}=36`$, from Wilking & Lada (wilking83 (1983)). Correspondingly, the Class II sources are represented by black dots of size decreasing with A<sub>V</sub>, from low extinctions (A$`{}_{\mathrm{V}}{}^{}<9`$) to high extinctions (A$`{}_{\mathrm{V}}{}^{}45`$), by steps of A$`{}_{\mathrm{V}}{}^{}10`$. A large majority of these sources are seen to have moderate extinctions (A$`{}_{\mathrm{V}}{}^{}<18`$), even in the areas overlapping regions of high extinctions traced by C<sup>18</sup>O. This implies that such sources are actually only moderately embedded in the cloud, in front of the densest regions traced by C<sup>18</sup>O and DCO<sup>+</sup>, rather than within them. The spatial distribution of the Class II sources can thus be more appropriately described as gaussian-like three-dimensional overlapping “haloes” around the DCO<sup>+</sup> cores A, B and F. This also implies that at least a fraction of the Class II sources with high extinctions are not necessarily really embedded in the densest regions, but may be part of these haloes behind the dense cores.
In the very same fashion, Fig. 9 displays the distribution of the known Class III sources, this time using FWHM=8′.<sup>14</sup><sup>14</sup>14A larger FWHM must be used because there are fewer Class III sources than Class II sources and because they are less clustered. This distribution is different from that of Class II sources in Fig 8. With respect to the DCO<sup>+</sup> cores, there is a strong density peak $`5\mathrm{}`$ SW of the location of core F, and no peak associated with cores A and B: in contrast with the distribution of Class II sources, there is a deficiency of Class III sources in the cloud center regions with high extinction.
The explanation for this apparent absence may be as follows. Whether they lie on the line-of-sight to regions of moderate extinction, or of high extinction, moderate-extinction Class II sources are found essentially everywhere. Therefore we also expect to have low-extinction Class III sources everywhere, in a $``$ 1:1 proportion. If they are not detected with the HRI, it can thus only mean that they are too faint in X-rays, hence have a small stellar luminosity. In addition, as for Class II sources, we must expect along the line-of-sight to the densest regions of the cloud to also have moderately embedded Class III sources at the back of the cloud. Their spatial density should be roughly comparable to that of the unseen Class III sources in the front, i.e., yield a small absolute number given the compact size of the C<sup>18</sup>O cores.
There is also the possibility of having genuinely embedded (hence very young) Class III sources in these cores: we have no information about the Class III/Class II source ratio there, so it is impossible to estimate their number. Should this number be large (such that Class III/Class II $`>1`$ for instance), this would be a problem for the earliest stages of YSO evolution; one rather expects to have Class III/Class II $`1`$ if all stars form with a disk taking at least $`\stackrel{>}{}10^5`$ yr to dissipate. However, a disk stage is perhaps not necessary for very low-mass stars, which would increase the number of very young Class III sources.
### 6.2 Constraints on the nature of the unknown Class III sources
Let us construct the H-R diagram of the 12 Class III sources in HRI/ISOCAM area for which we know the spectral types from the optical observations of Bouvier & Appenzeller (bouvier92 (1992)), and from the $`K`$-band observations of Luhman & Rieke (luhman99 (1999)), using the stellar luminosities determined by Bontemps et al. (bontemps00 (2000)).<sup>15</sup><sup>15</sup>15The relative differences between these luminosities, and the bolometric corrected ones from spectral types are lower than $`20\%`$. Fig. 10 displays the result, along with the birthline and pre-main sequence evolutionary tracks of Palla & Stahler (palla99 (1999)). According to these evolutionary tracks, the ages of the 12 Class III sources are found to be spread between $`0.2`$ and $`5`$ Myr.
It is reasonable to assume that the unknown Class III sources have the same age spread. These sources are not yet detected in X-rays either because their intrinsic X-ray luminosities are too low, and/or because they have high extinctions. In the first case they have stellar luminosities below 0.35 L, and the isochrones imply that $`M_{}0.6`$ M for the oldest ones, and $`M_{}<0.1`$ M for the youngest ones. In the second case, however, the unknown embedded Class III sources can have luminosities higher than $`L_{,\mathrm{min}}`$, i.e., so that they are not necessarily very low-mass stars.
Unless the number of Class III sources embedded in the densest regions is very high, our conclusion is that the bulk of the Class III sources which are undetected by the HRI and unrecognized by ISOCAM should be made of very low-mass stars.
## 7 Summary and conclusions
### 7.1 Main observational results
We have obtained two deep exposures of the $`\rho `$ Oph cloud core region (d=140 pc) with the ROSAT High Resolution Imager (core A: 51 ks, core F: 77 ks, in three partial exposures). The improved position accuracy (1$`\mathrm{}`$–6$`\mathrm{}`$) with respect to previous recent X-ray observations (ROSAT PSPC, Casanova et al. CMFA (1995); and ASCA, Koyama et al. koyama94 (1994) and Kamata et al. kamata97 (1997)) have allowed us to remove a number of positional ambiguities for the detected sources. We have cross-correlated the X-ray positions with IR sources found in the ISOCAM survey of the same region at 6.7 and 14.3 $`\mu `$m, in addition to sources known in the optical, IR, and radio from ground-based observations. We thus have now at our disposal the best-studied sample of X-ray emitting YSO in a star-forming region. We first summarize the main observational results of this article.
* We detect 63 HRI X-ray sources, and 55 are identified. Of the 55 identified X-ray sources 40 are PSPC sources, and 9 are ASCA sources.
* The IR classification (ground-based and ISOCAM survey) for the 55 identified X-ray sources yields: one Class I protostar (YLW15=IRS43); 23 Class II sources, including 4 new ISOCAM Class II sources; 21 Class III sources, including 13 new Class III sources; 8 new Class II or Class III source candidates; one early-type Class III source (the young magnetic B3 star S1), and one field star (the F2V star HD148352). The contamination of the sample of new X-ray sources by field stars is negligible.
* There is no statistically significant difference between the X-ray luminosity functions of HRI-detected Class II and Class III sources, i.e. T Tauri stars with and without disks, confirming that the contribution of these disks to X-ray absorption, or emission (for instance by magnetic reconnection between the star and the disk), must be small.
* X-ray variability of HRI-detected T Tauri stars has been studied by comparing the HRI data with the previously obtained PSPC data, and using HRI observations done at three different epochs. The resulting statistics show that most of the sources are variable, and that their X-ray variability is consistent with a solar-like (hence magnetic) flare origin.
* We use the information given both by the ISOCAM survey and by our HRI deep exposure to study the T Tauri star population of the $`\rho `$ Oph dense cores. We confirm that essentially all Class II and Class III YSO are X-ray emitters, and that a strong correlation ($`\mathrm{log}(L_\mathrm{X}/ergs^1)=(1.0\pm 0.2)\times \mathrm{log}(L_{}/L_{})+30.1`$) exists between the X-ray luminosity and the stellar luminosity of T Tauri stars, likely down to low luminosities (L$`{}_{}{}^{}0.1`$ L). We confirm that the characteristic $`L_\mathrm{X}/L_{}`$ for T Tauri stars is $``$10<sup>-4</sup> in the $`\rho `$ Oph cloud, albeit with a large dispersion. There is no evidence for a magnetic “saturation” seen at a level of 10<sup>-3</sup> in late-type main sequence stars.
* However, most of the new ISOCAM Class II sources are not detected by the HRI . We show that this is consistent with their intrinsic X-ray luminosities being too faint if “predicted” using the above $`L_\mathrm{X}`$$`L_{}`$ correlation.
### 7.2 What have we learned ?
* The first general conclusion we can draw from the HRI results presented above is a complete confirmation of the PSPC results obtained by CMFA. This was not a priori obvious, since the CMFA population (PSPC and near-IR) overlaps, but is different from, the HRI/near-IR/ISOCAM population presented in this paper: many PSPC sources are not detected by the HRI (see Appendix C), and some HRI sources are Class II and Class III newly classified thanks to a combined identification with ISOCAM. This shows that the $`L_\mathrm{X}`$$`L_{}`$ correlation is robust for the $`\rho `$ Oph TTS.
* The second, and perhaps most important, conclusion is the probable existence of $`40`$ unknown X-ray YSO down to a limit of L$`{}_{\mathrm{X}}{}^{}3\times 10^{28}`$ erg s<sup>-1</sup> in the HRI/ISOCAM overlapping area, which should be mainly low- to very low-mass ($`<0.1`$–0.6 M) diskless, “Class III TTS”. This prediction is based both on the use of the $`L_\mathrm{X}`$$`L_{}`$ correlation, legitimated by its robustness, and on the discovery of a large number of faint new IR sources by ISOCAM. As shown below, it may be soon verified by the next generation of X-ray satellites, namely XMM-Newton and Chandra. In this respect, the present paper can be taken as a “transition” paper between two generations of X-ray satellites.
Why is the detection of these “unknown TTS” important ? Because they are diskless, they are unlikely to be recognized as YSO by IR observations alone; and because they are likely to be as numerous as the YSO with IR excess, they have to be included in any reliable census of YSO, with an impact on such basic quantities as the initial mass function, or the star formation efficiency, especially if considered from an evolutionary point of view. For instance, from the results in this paper it is impossible to study the real connection between the distributions of the Class II and Class III sources in the densest regions, in particular to see whether the distribution of the Class III sources is also centered on the same DCO<sup>+</sup> cores as the Class II sources. The number of Class III sources embedded in the densest regions may, or may not, be comparable to that of the Class II sources, depending on the timescale for disk dispersal, especially among low-mass YSO. An X-ray improved census of Class III sources may also be crucial in determining whether a burst of star formation is presently going on in the $`\rho `$ Oph cores, as some recent indications suggest (see Martín et al. 1998). It will also allow to study the $`L_\mathrm{X}`$$`L_{}`$ correlation for Class III and Class II sources seperately, which was not possible in this paper (§) due to insufficient statistics.
### 7.3 The potential of XMM-Newton and Chandra
To quantify the prospects for improvement in the X-ray domain, we have computed the detection threshold for the X-ray camera EPIC aboard XMM-Newton, which was successfully launched in December 1999 (see Fig. 7). The improved sensitivity and enlarged energy range (0.5–12 keV) of EPIC will allow to detect the weak ISOCAM Class II sources, and also to discover numerous unknown faint or embedded Class III sources, in particular if they have high plasma temperatures (several keV) reached during flares, and extend the census of this population towards the low-mass end. In the best case, the XMM-Newton sensitivity will reach L$`{}_{\mathrm{X}}{}^{}10^{28}`$ erg s<sup>-1</sup> for A$`{}_{\mathrm{V}}{}^{}\stackrel{>}{}20`$, for long exposures ($`>`$75 ksec). This is nearly two orders of magnitude more sensitive than ROSAT. In case the faint Class III sources turn out to be so crowed that confusion problems arise, the excellent angular resolution of Chandra will be critical.
In $`\rho `$ Oph, there are already several identified bona fide and candidate brown dwarfs (see review in Neuhäuser et al. neuhaeuser99 (1999), and references therein), and four of them have been recently detected in X-ray using the ROSAT PSPC archive (Neuhäuser et al. neuhaeuser99 (1999)). Neuhäuser et al. have also shown that brown dwarfs could be X-ray emitters with the same ratio $`\mathrm{log}(L_\mathrm{X}/L_{})4`$ than for T Tauri stars. Thus Chandra and XMM-Newton should be able to detect many more of these objects with low stellar luminosity and masses, shedding a new light on their nature and early evolution.
###### Acknowledgements.
We thank Francesco Palla for fruitful discussions during the 5<sup>th</sup> French-Italian meeting in the island of Ponza, and the referee Fred Walter for his useful remarks. NG is supported by the European Union (Marie Curie Individual grant; HPMF-CT-1999-00228). EDF is partially supported by NASA contract NAS8-38252. We used SIMBAD maintained by the CDS (Strasbourg Observatory, France). We also used photographic data obtained using The UK Schmidt Telescope: original plate material is copyright $`\mathrm{\copyright }`$ of the ROE and the AAO.
## Appendix A HRI X-ray source detection
Source detection was done using EXSAS (Zimmermann et al. zimmermann97 (1997)), and the standard command DETECT/SOURCES. This command generates a local source detection by a sliding-window technique followed by a maximum likelihood test which compares the observed count distribution on the full resolution image (pixel of 0.5$`\mathrm{}`$) to a model of the point spread function (PSF) and the local background (Cruddace et al. cruddace88 (1988)). The “likelihood of existence” is defined as $`=\mathrm{ln}P_0`$, where $`P_0`$ the probability of the null hypothesis that the observed distribution of counts is only due to a statistical background fluctuation; $``$ provides a maximum likelihood measure for the presence of a source above the local background. We take $``$=6.8 as detection threshold ($`P_00.0011`$; or $`3.25\sigma `$ for Gaussian statistics) as argued in CMFA.
For source detection, the HRI report (David david97 (1997)) advises to screen out lower and higher HRI Pulse Height Analyzer (hereafter PHA) channels, which are found to have the highest background. However, source counts should always be determined using all the 1–15 channels to cancel out the uneven HRI efficiency distribution across the detector area (S. Döbereiner, private communication). We decided to search X-ray sources above a fixed detection threshold in channels 1–15 and 3–8.<sup>16</sup><sup>16</sup>16In the core F image #3, many 1–15 channel detections with high $``$ were not found in the 3–8 channel band Thus, we decided to take channels 2–8 instead of 3–8 in the three exposures. This third observation was the last of our program, one year after the Core A observation (see Table 1). According to Prestwich et al. (prestwich98 (1998)), the mean pulse height decreases by 0.5 channels year<sup>-1</sup>: this effect might explains why we must decrease the lower channel boundary from 3 to 2. The upper channel boundary seems to be less sensitive, probably due to fewer counts in upper channels (see David et al. david97 (1997)). EXSAS gave us list of X-ray detections with positions, one sigma error box ($`\sigma _X`$), $``$, and count rate. We removed sources detected in channels 1–15 but not in 3–8, considering these detections as spurious. For instance, there are two hot spots in the Core A observation in the South-East corner of the HRI. Since hot spots are usually considered as spurious detections, this criterion automatically removes them.
The astrometry must be corrected from offsets of typically a few arcseconds due to the time-dependent boresight error in the ROSAT aspect system. To do this, IR or optical counterparts in a 10$`\mathrm{}`$ radius circle around the X-ray sources were searched. We then selected a sample of X-ray sources with an unambiguous counterpart and 1 $`\sigma _X`$ (half width) error box $`1\mathrm{}`$, comparable to the IR/optical typical position error box ($`\sigma _{}1\mathrm{}`$). Then, offsets in right ascension ($`\alpha `$) and declination ($`\delta `$) were estimated by individual offset weighted mean: $`\alpha _{offset}=_{i=1}^nw_i\times (\alpha _{\mathrm{X},\mathrm{i}}\alpha _{,\mathrm{i}}),\delta _{offset}=_{i=1}^nw_i\times (\delta _{\mathrm{X},\mathrm{i}}\delta _{,\mathrm{i}})`$, with $`w_i=(1/\sigma _{\mathrm{X},\mathrm{i}}+1/\sigma _{,\mathrm{i}})/_{i=1}^n(1/\sigma _{\mathrm{X},\mathrm{i}}+1/\sigma _{,\mathrm{i}})`$. We found offsets ranging from -0.5$`\mathrm{}`$ to 2.5$`\mathrm{}`$. We subtracted these offsets, and checked the quality of our astrometry by estimating sample residuals mean before and after offsets subtraction: $`\sigma _{shift}=\frac{1}{n}\times _{i=1}^n\sqrt{(\alpha _{\mathrm{X},\mathrm{i}}\alpha _{,\mathrm{i}}\alpha _{offset})^2+(\delta _{\mathrm{X},\mathrm{i}}\delta _{,\mathrm{i}}\delta _{offset})^2}`$. We found $`\sigma _{shift}`$ ranging from 0.5$`\mathrm{}`$ to 1.2$`\mathrm{}`$. $`\sigma _{shift}`$ and $`\sigma _{}`$ ($`=1.2\mathrm{}`$) were then quadratically added to $`\sigma _X`$ to obtain an error box radius after astrometric correction ($`\sigma _{total}^2=\sigma _\mathrm{X}^2+\sigma _{shift}^2+\sigma _{}^2`$). In the case of the three different core F observations, the images were aligned and merged after astrometric correction to obtain a single deep HRI exposure of 77.2 ks. Source detection was subsequently performed as described in the article.
Table 2 lists the HRI X-ray sources, for which we adopt, in Col. 1, the same acronym as in CMFA: “ROXR” (for $`\rho `$ Oph X-ray $`ROSAT`$ source), followed by “A” or “F”.<sup>17</sup><sup>17</sup>17The “legal” designation for a new ROSAT source is RXJHHMM.$`m`$ $`\pm `$ DDMM, where J stands for coordinates in J2000 and $`m`$ is the cut (not rounded) decimal value. Our tables do not use this notation for simplicity in the discussion, but the correct designation is easy to reconstruct from the position given if required. For example, ROXRF14 = ROXs20B = RXJ162714-2430. Fig. 11 indicates the source numbering.
In order to allow easier comparisons with previous work, X-ray source positions are listed in both J2000 and B1950 equinoxes, with their $`1\sigma _{total}`$ error box, in Cols. 2–6. The likelihood of existence $``$ is in Col. 7. Count rates are indicated in Col. 8–11. For the core F field, the indicated positions ($`\alpha ,\delta `$) and $``$ values correspond to observation (#1, #2, or #3) where $``$ and the position accuracy are the best, i.e. when the count rate is highest. When an X-ray source is detected in one observation above the detection threshold, and not detected in other observations, we have estimated the corresponding count rate upper limits (3.25$`\sigma `$), using the EXSAS command COMPUTE/UPPER\_LIMITS. We have noted that the detection efficiency degrades with increasing angle to the axis, in the same way as the point spread function (this is discussed in $`\mathrm{\S }`$5.2).
## Appendix B Optical/IR counterparts of the HRI X-ray sources
We searched stellar counterparts for the 63 ROSAT HRI X-ray sources on the ESO/SERC second digitized sky survey (DSS2). Fig. 12 gives the finding charts with BKLT IR sources for each of the 63 ROSAT HRI X-ray sources.
Table 3 gives identification lists for the two fields, and cross-identification with other surveys. Col. 1 is the ROXR numbering from detection (Table 2). Cols. 2–4 are respectively cross-identification lists with the X-ray sources of CMFA (ROXR1), Casanova (casanova94 (1994); ROXR2), and Kamata et al. (kamata97 (1997)). Dots mean “X-ray source undetected”, and dash “out of observation field”. Col. 5 gives the first name attributed to this counterpart.
In the core A field, we find 26 X-ray sources, of which only one (ROXRA10) remains without optical or IR counterpart. Of the 25 identified X-ray sources, 22 were seen with the ROSAT PSPC, and 4 are new detections (ROXRA3, 10, 16, 22). In the core F field, we find 37 X-ray sources, including 7 without optical or IR counterpart. Of the 30 identified X-ray sources, 18 were seen with the ROSAT PSPC, and 12 are new detections (ROXRF3, 8, 12, 15, 18, 19, 24, 26, 28, 32, 35, 36). Altogether, 63 X-ray sources are detected, and 55 are identified. Of the 55 identified X-ray sources 40 are PSPC sources.
For sources with a low statistical significance ($`6.89.1`$, or $`10^4P_01.1\times 10^3`$; 3.25–3.9 $`\sigma `$ for Gaussian statistics) we find X-ray sources with and without optical or IR counterparts. The X-ray sources without counterparts are always weak sources and may be spurious detections (locally high background), and this may therefore also be the case for weak X-ray sources with counterparts in case of chance spatial coincidence. For instance in the Core A field (respectively Core F) there are 875 (resp. 1173) BKLT sources; this sample is dominated by background sources without detectable X-ray emission. To estimate the number of chance coincidences, we have placed in each field 10<sup>5</sup> random X-ray source positions, and searched for each whether there is a BKLT source in a circle of 10$`\mathrm{}`$ radius: we have then an estimate of the probability to find a BKLT counterpart by chance within 10$`\mathrm{}`$ from a spurious X-ray detection. This probability is 0.044 (resp. 0.049) for Core A (resp. Core F), or approximately 1/20 for both; in other words one (resp. two) spurious source identification are expected for the Core A (resp. Core F) field. As we have for Core A (resp. Core F) two (resp. 15) X-ray sources with $`9.1`$ out of 26 (resp. out of 37), this implies that one weak X-ray source in Core A (resp. 13 in Core B) is real, which is consistent with the number of identifications. We are therefore confident that the identifications of weak X-ray sources with stellar counterparts are correct.
## Appendix C Comparison between HRI and PSPC observations
Within the boundaries of our observation fields (core A and core F), there are 61 X-ray sources detected previously with the ROSAT PSPC (53 from CMFA, and 8 from Casanova casanova94 (1994)). However, 21 X-ray sources are not detected with the ROSAT HRI. This difference could be explained by lower observational sensitivity and/or source variability. To elucidate this point, we must estimate HRI count rates for PSPC sources, and compare them with the adopted HRI detection threshold (3.25 $`\sigma `$).
We first estimate the conversion factor between the PSPC count rate in the energy range of CMFA (1.0–2.4 keV) and the HRI count rate in the whole energy range (0.1–2.4 keV). We did not select X-ray sources with ambiguous PSPC detection (10 sources with notes in Cols. 2–3 of Table 3). For core F observation, the lowest detection count rate was taken to minimize variability effects (four sources with upper limits are not selected): we kept only 26 X-ray sources. Since many of these sources are variable (as shown by the core F observations), a conversion factor estimator insensitive to extreme values of the sample is needed. This is why we take the median of the PSPC/HRI count rate ratio, instead of the mean. We find PSPC $`_{12.4\mathrm{keV}}`$ count rate = 2.4 $`\times `$ HRI $`_{0.12.4\mathrm{keV}}`$ count rate.
Fig. 13 displays the HRI $`_{0.12.4\mathrm{keV}}`$ count rate vs. the PSPC $`_{12.4\mathrm{keV}}`$ count rate. It shows two classes of sources: sources near the median, and sources beyond the median (with error bars). The dispersion of points (within 1 rms) around the median value could be due to X-ray extinction effect on the conversion factor or to a variability factor $`2`$. Preibisch et al. (preibisch96 (1996)) calculated the conversion for the whole energy band of the ROSAT HRI assuming optically thin plasma emission with $`T_\mathrm{X}=10^7`$ K and different values for the X-ray extinction: for $`N_H`$ increasing from $`6.5\times 10^{19}`$ cm<sup>-2</sup> to 10<sup>22</sup> cm<sup>-2</sup>, the conversion factor decreases from 2.5 to 2.0. Our observational estimate is in agreement with these values, which also show that the dependence of the conversion factor on X-ray extinction is small compared to the dispersion of count rates and can be neglected in our plot. We conclude that the dispersion is due to variability: WL20, GSS37, VSS27, and SR9 must have been in a high state during the PSPC observation, as were ROXs4 and SR2 during the HRI observation, the other sources being essentially unchanged in both observations.
Using this conversion factor, we can estimate the HRI count rate from the PSPC count rate, and compare it with our HRI threshold computed with EXSAS. We find that 14 sources are below this threshold (ROXR1-1, 6, 7, 16, 20, 21, 27, 33, 34, 37, 40, 53, and ROXR2-16, 18), and the other 7 sources were in a high state during the PSPC observation (ROXR1-19, 30, 45, 47, and ROXR2-27, 30, 33). We conclude that the non-detection of the 21 PSPC sources by the HRI can be fully explained by the difference in sensitivity and intrinsic variability.
We can also compare our detections with ASCA (Koyama et al. koyama94 (1994); Kamata et al. kamata97 (1997); see observation field in Fig. 1) despite its lower angular resolution. Kamata et al. (kamata97 (1997)) detected 19 X-ray sources, of which 10 were previously observed by Koyama et al. (koyama94 (1994)). Compared to Einstein Observatory and ROSAT PSPC observations, 7 new X-ray sources were discovered by ASCA.<sup>18</sup><sup>18</sup>18Note that X-ray source 9B of Kamata et al. (kamata97 (1997)) is in fact ROXR1-45 of CMFA. These 7 X-ray sources are also not detected with the HRI. On the 12 sources already observed by Einstein Observatory and ROSAT PSPC, we detect 9 sources, the 3 others being below our sensitivity threshold according to the above conversion factor.
## Appendix D Optical/IR counterpart without IR classification
Nine X-ray sources have optical/IR counterpart for which the IR classification is not known. Three of these X-ray sources are found in the ISOCAM survey, but with only upper limits in LW2 and LW3 filters. We give here their spectral energy distribution (see Fig. 14), and discuss their possible IR classification. In case of doubt, the resulting Class II (or Class III) source candidates have not been included in the statistic studies of this article.
ROXRA8: The counterpart of this X-ray source is the optical star Chini8=SKS3 ($`R=16.6`$, $`K=9.5`$). The ISOCAM LW2 and LW3 upper limits exclude an IR excess. We classify this source as a Class III source.
ROXRF12: A weak star is visible in the 90$`\%`$ confidence error box on the DSS2 image, but this star is neither found in the BKLT survey, nor in the PMM USNO-A1.0 catalogue. The low optical/near IR magnitudes imply a low luminosity for this object. We propose this source as a weak Class II or Class III source candidate detected during a strong X-ray flare. This source may also be a brown dwarf.
ROXRF18: The counterpart of this X-ray source is the IR star B162720-243820 ($`K=14.6`$). This star is visible in the DSS2 (red) optical image, but it is not in the PMM USNO-A1.0 catalogue (probably because only stars appearing both in blue and red images were accepted), thus we have no estimate of its B and R magnitudes. The low near IR magnitudes imply a low luminosity for this object. This source may be a weak Class II or Class III source candidate detected during a strong X-ray flare. ISOCAM LW2 and LW3 upper limits do not exclude an IR excess for this object. This object can also be a weak Class I protostar with a strong X-ray flare.
ROXRF23: The counterpart of this X-ray source is the IR star B162728-244803 $`(K=14.1`$; see note in Table 3). The low near IR magnitudes imply a low luminosity for this object. The SED of this source peaks in the $`H`$-band. We propose this source as a weak Class II or Class III source candidate detected during a strong X-ray flare.
ROXRF24: The counterpart of this X-ray source is an optical star ($`R=16.3`$) named only in the PMM USNO-A1.0 catalogue (Monet et al. monet96 (1996)). Unfortunately, this object lies outside the BKLT survey. We propose this source as Class II or Class III source candidate.
ROXRF25: The counterpart of this source is BBRCG50 observed only in $`K`$-band ($`K=11.7`$). This source is not retrieved in BKLT. The ISOCAM LW2 and LW3 upper limits exclude a strong IR excess. We propose this source as Class II or Class III source candidate.
ROXRF34: The counterpart of this X-ray source is the optical star WSB58=B162800-244819 ($`R=17.3`$, $`K=9.3`$). Wilking et al. (wilking87 (1987)) noted a probable H<sub>α</sub> detection needing confirmation. The SED of this source peaks in the $`H`$-band. We propose this source as Class II or Class III source candidate.
ROXRF35: The counterpart of this X-ray source is the optical star B162800-245340 ($`R=16.1`$, $`K=9.6`$). The SED of this source peaks in the $`J`$-band. We propose this source as Class II or Class III source candidate.
ROXRF36: The counterpart of this X-ray source is the optical star B162812-245043 ($`R=15.7`$, $`K=9.4`$), which appears to be a close binary ($`1.5\mathrm{}`$) in the second Digitized Sky Survey (see Fig. 12). The SED of this source peaks in the $`H`$-band. We propose this source as Class II or Class III source candidate.
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# Untitled Document
hep-th/0005048
S-Duality and Noncommutative Gauge Theory
Rajesh Gopakumar, Juan Maldacena, Shiraz Minwalla and Andrew Strominger
Jefferson Physical Laboratory,
Harvard University,
Cambridge, MA 02138
Abstract
It is conjectured that strongly coupled, spatially noncommutative $`𝒩=4`$ Yang-Mills theory has a dual description as a weakly coupled open string theory in a near critical electric field, and that this dual theory is fully decoupled from closed strings. Evidence for this conjecture is given by the absence of physical closed string poles in the non-planar one-loop open string diagram. The open string theory can be viewed as living in a geometry in which space and time coordinates do not commute.
1. Introduction
Noncommutative field theories have a rich and fascinating structure. The embedding of these theories into string theory suggests that this structure may be directly relevant to understanding the inevitable breakdown of our familiar notions of space and time at short distances in quantum gravity.
Investigations to date have largely concentrated on theories with purely spatial noncommutativity (see however ). While such theories are interestingly nonlocal in space, they are local in time, admitting familiar notions like that of the Hamiltonian and a quantum state. Noncommutativity of a time-like coordinate should have even more far-reaching consequences, and it is natural to ask whether or not such theories exist.
In this paper we give one answer to this question by asking another: What is the strong coupling dual of NCYM (spatially-noncommutative $`𝒩=4`$ Yang-Mills)? This question can be addressed in the description of NCYM as a scaling limit of three-branes with a $`B`$ field in $`IIB`$ string theory . $`IIB`$ $`S`$-duality induces an $`S`$-duality on the NCYM theory, mapping the strongly coupled NCYM theory to a weakly coupled open string theory<sup>1</sup> The low energy sector of the open string theory is ordinary $`𝒩=4`$ YM, and the induced duality reduces to the standard $`S`$-duality. . This open string theory can be viewed either as living in a near critical electric field<sup>2</sup> The existence of a scaling theory at near critical electric fields, and its relevance to temporal noncommutativity was emphasized to us by N. Seiberg, L Susskind and N. Toumbas (private communications). The scaling to the critical electric field was also considered in , . <sup>3</sup> The critical value of the electric field arises when the force pulling apart the charges at either end of the string just balances the string tension, so that the string is effectively tensionless \[6--9,,10\]. Beyond this value the spectrum contains a tachyon and the vacuum is unstable. , or in a space-time with noncommuting space and time coordinates. A precise statement of the spacetime noncommutativity in this theory is that the temporal zero mode $`X^0`$ on the open string worldsheet does not commute with the spatial zero modes. The scale associated with this noncommutativity is the same as the effective open string scale. Thus the effects of the noncommutativity are inextricably tied up with the usual stringy nonlocalities.
Since the closed string sector of the $`IIB`$ theory is decoupled in the scaling limit, the dual open string theory does not have a closed string sector. The appearance of an open string theory without a closed string sector is striking. Ordinarily closed string poles appear in open string loop diagrams, and unitarity then requires the addition of asymptotic closed string states. In order to better understand this point we analyze (following \[11--16\]) the nonplanar one loop open string diagram for the bosonic case. We find that the temporally noncommutative phases lead to a precise cancellation of all the closed string poles, in accord with our expectations. This cancellation in fact occurs for branes of any dimension, indicating the existence of a family of non-commutative open string theories.
This paper is organized as follows. In section 2 we derive the $`S`$-dual of NCYM, which we refer to as NCOS (noncommutative open strings), by embedding in string theory. In section 3 we show that it is a decoupled open string theory with a near-critical electric field. In section 4 we give evidence at the one loop level for the decoupling of closed strings by computing the non-planar annulus for bosonic string theory with two incoming and two outgoing tachyons. Section 5 contains a preliminary analysis of the general higher loop diagram; no obvious closed string singularities are found. In section 6 we make some comments regarding the supergravity duals of our open string theory. We conclude with some discussion in section 7. For simplicity we concentrate on the $`U(1)`$ theories but our results generalize easily to $`U(N)`$.
Related work will appear in .
2. Inducing S-Duality
The Olive Montonen dual of ordinary $`𝒩=4`$ SYM may be deduced as a consequence of the $`S`$ duality of IIB theory in the presence of D3-branes in the zero slope limit. In this section we will determine the Olive Montonen dual of spatially noncommutative $`𝒩=4`$ SYM, using the $`S`$ duality of IIB theory in flat space in the presence of D3-branes and a background $`B_{\mu \nu }`$ field, together with the modified zero slope limit .
Consider a D3-brane, extended in the $`0,1,2,3`$ directions, in a background geometry
$$\begin{array}{cc}\hfill g_{\mu \nu }^{}& =\eta _{\mu \nu },g_{ij}^{}=\alpha _{}^{}{}_{}{}^{2}k_1\delta _{ij},g_{MN}^{}=\delta _{MN},\hfill \\ \hfill B_{ij}^{}& =Bϵ_{ij},g_{str}^{}=\alpha ^{}k_2.\hfill \end{array}$$
$`(2.1)`$
in the limit $`\alpha ^{}0`$, keeping $`k_1,k_2,B`$ fixed (we will refer to this as the NCYM limit). Here $`\mu ,\nu =0,1`$ with $`i,j=2,3`$ and $`M,N=4,\mathrm{}9`$. (We will reserve unprimed notation for the S-dual variables to be introduced in the next sub-section.) It was shown in that the decoupled theory on the brane is noncommutative $`U(1)`$ SYM propagating on a four dimensional space with (open string) metric (we use the conventions of ) $`G_{\mu \nu }^{}=\eta _{\mu \nu },G_{ij}^{}=\frac{(2\pi B)^2}{k_1}\delta _{ij}`$, noncommutativity parameter $`\theta ^{ij}=\frac{ϵ^{ij}}{B}`$, and gauge coupling $`g_{YM}^2=2\pi G_o^{}_{}{}^{}2`$, where $`G_o^2=\frac{k_2B}{k_1}`$. In order to obtain noncommutative field theory propagating on a space with unit metric we choose $`k_1=(2\pi B)^2`$. In terms of the field theory couplings $`\theta ^{^{}}`$ and $`G_o^{}`$, $`B=\frac{1}{\theta ^{^{}}}`$ and $`k_2=\frac{(2\pi )^2G_o^{}_{}{}^{}2}{\theta ^{^{}}}`$.
In order to obtain a weakly-coupled dual description of the noncommutative gauge theory at large $`G_o^{}`$ we will consider the NCYM limit described above in an $`S`$-dual picture. Before describing this in detail we note that the $`S`$-dual version has two potentially unpleasant features:
a. It seems to involve branes in the presence of an an RR 2 form potential (the $`S`$-dual of $`B_{ij}^{}`$).
b. The S-dual of the NCYM limit takes the closed string coupling $`g_{str}`$ to infinity, seeming to indicate that any description of brane dynamics obtained in this picture will be strongly rather than weakly coupled, independent of $`G_o`$.
These difficulties may both be circumvented. In order to avoid having to deal with RR fields, we gauge away the constant bulk NS-NS potential before performing the S-duality. This gauge transformation induces a magnetic field $`F_{23}^{}=B`$ on the the D3-branes, which is converted into an electric field by the $`S`$-duality; in fact an electric field that approaches its critical value in the scaling limit. This electric field may in turn be gauged into a constant background NS-NS two form potential $`B_{01}=F_{01}`$ in the bulk. But, in such a background, the open string coupling that governs the strength of interactions between brane modes is not directly related to the closed string coupling. It turns out that the open string coupling in this background is $`G_o=\frac{1}{G_o^{}}`$, i.e. it is the inverse of the original open string coupling, and therefore remains finite despite the fact that $`g_{str}\mathrm{}`$. Thus at large $`G_o^{}`$, the effective description is a weakly coupled noncommutative open string theory, with noncommutativity in the time direction!
We now consider this limit in more detail. We could consider any finite number of branes, $`N`$, but we will mostly stick to the case $`N=1`$ for simplicity.
2.1. Born-Infeld $`S`$-Duality
$`S`$-duality transforms a constant magnetic field on the three-brane to a constant electric field. Constant fields on a single D3-brane are governed by the Born-Infeld action
$$S_{BI}=\frac{1}{(2\pi )^3\alpha _{}^{}{}_{}{}^{2}g_{str}}d^4x\sqrt{det(g_{\mu \nu }2\pi \alpha ^{}F_{\mu \nu })}.$$
$`(2.2)`$
The action of $`S`$ duality on $`S_{BI}`$ will be reviewed in this subsection (See ). Consider a gauge theory on a torus. The flux of the magnetic field on any nontrivial two cycle of the torus is integrally quantized, and so must, under electromagnetic $`S`$-duality, map to a quantized electric flux. Recall why electric flux on a torus is quantized. The constant piece (zero momentum mode) of a gauge field in flat infinite space is physically unmeasurable, as it can be gauged away. This is not true, however, on a torus, as the Wilson line $`e^{i{\scriptscriptstyle A}.dx}`$ over any nontrivial cycle of the torus is a gauge invariant observable, implying that the zero momentum piece of the gauge field $`A_i`$ is a periodic physical ‘coordinate’, with period $`\frac{2\pi }{L_i}`$ ($`L_i`$ is the size of the $`i^{th}`$ spatial direction). Consequently, the momentum conjugate to the zero mode of $`A_i`$ is quantized in integral units of $`L_i`$. This quantized momentum is the electric flux that is interchanged with the quantized magnetic flux under $`S`$ duality.
In order to work out the expression for the quantized electric flux, consider the theory (2.2) on a rectangular torus, with spatial coordinate radii $`L_1,L_2,L_3`$. We are interested in background field configurations in which $`F_{01}`$ is nonzero and constant, but $`F_{ij}`$ is zero. Since $`\dot{A}_1`$ appears in the Lagrangian only through $`F_{01}`$, it is sufficient, for the purposes of computing canonical momenta in such backgrounds, to set $`F_{ij}`$ to zero in the Lagrangian. For a diagonal metric the Born Infeld action simplifies to (recall $`g^{00}`$ is negative)
$$S=\frac{1}{(2\pi )^3\alpha _{}^{}{}_{}{}^{2}g_{str}}d^4x\sqrt{g}\sqrt{1+(2\pi \alpha ^{})^2g^{11}g^{00}F_{01}^2}.$$
$`(2.3)`$
Thus, for constant $`F_{01}`$, the momentum conjugate to $`A_1`$ is
$$P^1=NL_1=\frac{1}{2\pi g_{str}}L_1L_2L_3\sqrt{g}\frac{g^{11}g^{00}F_{01}}{\sqrt{1+(2\pi \alpha ^{})^2g^{11}g^{00}F_{01}^2}}.$$
$`(2.4)`$
Thus the constant $`F_{23}^{}`$ background of the spatially noncommutative theory maps, under $`S`$ duality, to a background with constant $`F_{01}`$, whose value is given by the solutions to the equations
$$\frac{\sqrt{g}}{g_{str}}\frac{g^{11}g^{00}F_{01}}{\sqrt{1+(2\pi \alpha ^{})^2g^{11}g^{00}F_{01}^2}}=F_{23}^{}=\frac{1}{\theta ^{^{}}}$$
$`(2.5)`$
where $`g_{\mu \nu }`$ and $`F_{\mu \nu }`$ are the background metric and field strength in the $`S`$ dual description. In terms of the critical value of the electric field
$$F_{01}^{crit}=\frac{\sqrt{g_{00}g_{11}}}{2\pi \alpha ^{}}$$
$`(2.6)`$
one finds
$$F_{01}=\frac{F_{01}^{crit}}{\sqrt{1+g_{22}g_{33}(\frac{\theta ^{^{}}}{2\pi \alpha ^{}g_{str}})^2}}.$$
$`(2.7)`$
2.2. The Scaling Limits
Consider IIB theory with a D3-brane in the presence of a background NS-NS 2-form potential, $`B_{\mu \nu }`$. Prior to any scaling limit, an open string metric $`\stackrel{~}{G}^{AB}`$ (the symbol $`G^{AB}`$ will be reserved for a rescaled open string metric defined below) and a non-commutativity parameter $`\mathrm{\Theta }`$ can be deduced from disk correlators on the open string worldsheet boundaries
$$X^A(0)X^B(\tau )=\alpha ^{}\stackrel{~}{G}^{AB}\mathrm{ln}(\tau )^2+\frac{i}{2}\mathrm{\Theta }^{AB}ϵ(\tau ),A,B=0,1,2,3.$$
$`(2.8)`$
The open string coupling $`G_o`$ is similarly read off from the coefficient of the gauge theory action. These are related to closed string quantities by the formulae
$$\begin{array}{cc}\hfill 2\pi \alpha ^{}\stackrel{~}{G}^{AB}+\mathrm{\Theta }^{AB}& =(2\pi \alpha ^{})(\frac{1}{g+2\pi \alpha ^{}B})^{AB},\hfill \\ \hfill G_o^2& =g_{str}\frac{\mathrm{det}^{\frac{1}{2}}(g+2\pi \alpha ^{}B)}{\mathrm{det}^{\frac{1}{2}}(g)}.\hfill \end{array}$$
$`(2.9)`$
As discussed above, in the NCYM limit, $`\alpha ^{}0`$ while the open string metric $`G^{{}_{}{}^{}AB}`$, open string coupling $`G_o^{}`$ and the (spatial) non-commutativity matrix $`\mathrm{\Theta }^{{}_{}{}^{}AB}`$ are held fixed. We would now like to study this scaling limit in the S-dual description of Type $`IIB`$ theory. We will call this the NCOS limit. Under an S-Duality, the type $`IIB`$ closed string backgrounds transform in the usual fashion, $`g_{str}^{^{}}=\frac{1}{g_{str}}`$, $`g_{\mu \nu }^{}=\frac{g_{\mu \nu }}{g_{str}}`$ ($`\alpha ^{}`$ is unchanged). The associated open string quantities may then be read from their definitions in (2.9). The results, in the limit $`\alpha ^{}0`$, are summarized in the following table:
TABLE 1
The NCYM Limit The $`S`$-Dual NCOS Limit $`g_{\mu \nu }^{^{}}=\eta _{\mu \nu }`$ $`g_{\mu \nu }=\frac{\theta G_o^4}{2\pi \alpha ^{}}\eta _{\mu \nu }`$ $`g_{ij}^{^{}}=\frac{(2\pi \alpha ^{})^2}{\theta ^{}_{}{}^{}2}\delta _{ij}`$ $`g_{ij}=\frac{2\pi \alpha ^{}}{\theta }\delta _{ij}`$ $`B_{\mu \nu }^{}=F_{\mu \nu }^{}=0`$ $`B_{\mu \nu }=F_{\mu \nu }=F_{\mu \nu }^{crit}\left(1\frac{1}{2}\left(\frac{2\pi \alpha ^{}}{\theta G_o^2}\right)^2\right)`$ $`B_{ij}^{}=F_{ij}^{}=\frac{1}{\theta ^{^{}}}ϵ_{ij}`$ $`B_{ij}=F_{ij}=0`$ $`g_{str}^{^{}}=G_o^{}_{}{}^{}2\frac{2\pi \alpha ^{}}{\theta ^{^{}}}`$ $`g_{str}=\frac{\theta ^{^{}}}{G_o^{}_{}{}^{}22\pi \alpha ^{}}=\frac{G_o^4\theta }{2\pi \alpha ^{}}`$ $`G^{{}_{}{}^{}AB}=\eta ^{AB}`$ $`\frac{\alpha ^{}}{\alpha _{}^{}{}_{eff}{}^{}}\stackrel{~}{G}^{AB}G^{AB}=\eta ^{AB}`$ $`G^{{}_{}{}^{}MN}=g^{{}_{}{}^{}MN}=\delta ^{MN}`$ $`G^{MN}=g^{MN}=\frac{2\pi \alpha ^{}}{\theta G_o^4}\delta ^{MN}`$ $`\mathrm{\Theta }^{{}_{}{}^{}\mu \nu }=0`$ $`\mathrm{\Theta }^{\mu \nu }=\theta ^{^{}}G_o^{}_{}{}^{}2ϵ^{\mu \nu }=\theta ϵ^{\mu \nu }`$ $`\mathrm{\Theta }^{{}_{}{}^{}ij}=\theta ^{^{}}ϵ^{ij}`$ $`\mathrm{\Theta }^{ij}=0`$ $`G_o^{^{}}=G_o^{^{}}`$ $`G_o=\frac{1}{G_o^{}}`$ $`\alpha ^{}=\alpha ^{}`$ $`\alpha _{}^{}{}_{eff}{}^{}=\frac{\theta }{2\pi }`$
Here
$$\mu ,\nu =0,1,i,j=2,3,A,B=0,1,2,3,M,N=4,5,6,7,8,9.$$
In Table 1 we have expressed all open and closed string quantities as functions of $`\theta `$ and $`G_o`$, the noncommutativity parameter and open string coupling in the (S-dual) NCOS theory. We have also defined the quantities, $`\alpha _{}^{}{}_{eff}{}^{}`$ the effective open string scale and the rescaled open string metric $`G^{AB}=\frac{\alpha ^{}}{\alpha _{}^{}{}_{eff}{}^{}}\stackrel{~}{G}^{AB}`$ of the NCOS theory.
Note that
1. In the limit $`\alpha ^{}0`$, the electric field $`F_{01}`$ of the NCOS theory attains its critical value
$$F_{01}^{crit}=\frac{\theta G_o^4}{(2\pi \alpha ^{})^2}.$$
$`(2.10)`$
2. The energy per unit coordinate length of an NCOS open string stretched in the 1 direction is given by (recall that the ends of an open string are charged)
$$p_0=\frac{ϵ_{01}}{2\pi }\left(\frac{1}{\alpha ^{}}2\pi ϵ^{01}F_{01}\right)\mathrm{\Delta }x^1=\frac{1}{4\pi \alpha _{eff}^{}}\mathrm{\Delta }x^1$$
$`(2.11)`$
so these open strings have an effective tension set by $`\alpha _{}^{}{}_{eff}{}^{}`$. As a consequence, it will turn out that in the NCOS limit excited open string oscillator states are part of the decoupled theory on the brane in the NCOS limit, and that their mass scale is also set by $`\alpha _{}^{}{}_{eff}{}^{}`$ .
3. The open string coupling $`G_o`$ is the inverse of the gauge coupling $`G_o^{}`$ in the NCYM limit.
To summarize, strongly coupled spatially noncommutative Yang-Mills theory has an effective description as a weakly coupled open string theory living on D3-branes, in the presence of a near critical electric field. The parameters of this open string theory are listed in Table 1. We will explore this theory in the rest of this paper.
3. The Classical NCOS Theory
3.1. Spacetime Noncommutativity
In the NCOS limit, open strings on the brane propagate in a background electric field. This results in temporal noncommutativity, in the sense that the open string zero modes obey
$$[X^\mu ,X^\nu ]=i\theta ϵ^{\mu \nu },$$
$`(3.1)`$
as may easily be seen from (2.8).
Disk diagrams in the NCOS theory are very simple. As argued in , , open string correlation functions on the disk in the NCOS theory may be obtained from the equivalent correlation functions in the theory without the electric field, by the addition of noncommutative phases in the $`0,1`$ directions (and using the appropriate open string metric and coupling). Thus the classical action for open string modes in the NCOS limit may be obtained by turning all products in the usual open string classical action into star products. In other words if we think about the open string field theory action $`S=AQA+A_wA_wA`$ there the $`_w`$ product is Witten’s star product then the only change is that we replace Witten’s product by a modified product which just adds in the Moyal phases, and of course we replace $`\alpha ^{}\alpha _{eff}^{}`$.
Since the effective string scale $`\alpha _{}^{}{}_{eff}{}^{}`$ is the same as that of non-commutativity $`\theta `$, the noncommutative phases are non negligible only for energies of the order of those of string oscillators.
3.2. The Free Spectrum
In this subsection we will argue that the NCOS limit defines an open string theory on the 3-brane, as open string oscillators do not decouple in this limit. We will examine the spectrum in the free NCOS theory and see that the effective scale is indeed set by $`\alpha _{}^{}{}_{eff}{}^{}`$.
We first consider the scaling limit in the NCYM picture. Near the NCYM limit one has weakly coupled closed strings coupled gravitationally to open strings. Open string excitations with string frame energies obeying $`|g^{00}k_0^2|\frac{1}{\alpha ^{}}`$, or equivalently Einstein frame energies obeying $`|g_E^{00}k_0^2|m_p^2`$ , decouple from the closed strings. As $`g^{00}=1`$, open string modes with $`k_0\frac{1}{\sqrt{\alpha ^{}}}`$ decouple from closed string modes. The decoupled theory includes all brane excitations with energies that obey this inequality, namely just the $`𝒩=4`$ YM multiplet.
Now consider the same limit in the NCOS picture. The argument above ensures that open string modes with $`k_0\frac{1}{\sqrt{\alpha ^{}}}`$ decouple from closed strings. However, the open string oscillator states in this picture obey the mass shell condition set by the open string metric in the RHS of Table 1
$$\frac{\alpha _{}^{}{}_{eff}{}^{}}{\alpha ^{}}k_AG^{AB}k_B=\frac{N}{\alpha ^{}}$$
$`(3.2)`$
with $`A,B=0,1,2,3`$. This implies
$$k^2=\frac{N}{\alpha _{}^{}{}_{eff}{}^{}}\frac{1}{\alpha ^{}}$$
$`(3.3)`$
with
$$\alpha _{}^{}{}_{eff}{}^{}=\frac{\theta }{2\pi }$$
$`(3.4)`$
in the limit $`\alpha ^{}0`$. Thus the decoupled theory on the brane includes all open string oscillator states! The mass spectrum is exactly the usual free spectrum on the three brane, except with $`\alpha _{}^{}{}_{eff}{}^{}`$ replacing $`\alpha ^{}`$.
3.3. Worldsheet Correlators
Nontrivial vertex operators are functions of tangential worldsheet derivatives of $`X^A`$ and normal worldsheet derivatives of $`X^M`$. Correlation functions of such vertex operators may be computed given the two point correlators of the free fields $`X^A`$ restricted to the boundary of the world sheet, as well as the two point functions of the free fields $`X^M`$.
The boundary correlators of $`X^A`$ are finite in the limit $`\alpha ^{}0`$, and are given by
$$X^A(0)X^B(\tau )=\alpha _{}^{}{}_{eff}{}^{}G^{AB}\mathrm{ln}(\tau )^2+\frac{i}{2}\theta ^{AB}ϵ(\tau ),A,B=0,1,2,3$$
$`(3.5)`$
On the other hand, correlation functions involving the transverse directions $`X^M`$ are derived from the sigma model
$$S=\frac{1}{4\pi \alpha ^{}}G_{MN}X^MX^N=\frac{\alpha _{}^{}{}_{eff}{}^{}G_0^4}{4\pi \alpha _{}^{}{}_{}{}^{2}}X^MX^N\delta _{MN}.$$
$`(3.6)`$
In terms of the rescaled fields $`Y^M=\frac{G_0^2\alpha _{}^{}{}_{eff}{}^{}}{\alpha ^{}}X^M`$
$$S=\frac{1}{4\pi \alpha _{}^{}{}_{eff}{}^{}}Y^MY^N\delta _{MN}.$$
$`(3.7)`$
The vertex operators representing physically normalized states are functions of the normal derivatives of $`Y^M`$.
Thus all correlation functions of NCOS vertex operators on the disk will be the same as in usual open string theory except that $`\alpha ^{}\alpha _{eff}^{}`$ and we have extra non-commutative phases appearing as in . The open string coupling constant is $`G_0`$ and it is finite.
4. The One Loop Diagram
One loop open string graphs usually contain closed string poles, and unitarity then requires that closed strings be included as asymptotic states. In this section we consider the nonplanar annulus diagram in the NCOS limit, and show that it has no physical closed string poles. This demonstrates that an on shell closed string cannot be produced in collisions of open strings.
Nonplanar diagrams for spatial $`\mathrm{\Theta }`$ were computed in \[11--16\] – we will follow . The nonplanar diagram for our case can be obtained by analytic continuation. For simplicity consider the case of two initial and two final open string tachyon vertex operators $`V_T=G_oe^{ik_AX^A}`$ in the bosonic string with incoming momenta $`k_1,k_2`$ and outgoing momenta $`k_3,k_4`$. Then we get for a D-3 brane in bosonic string theory (Eq. 2.17 of )
$$\begin{array}{cc}& V_T(k_1)V_T(k_2)V_T(k_3)V_T(k_4)_{annulus}i\sqrt{G}G_o^4(4\alpha _{}^{}{}_{eff}{}^{})^2\delta ^4(k_1+k_2+k_3+k_4)\hfill \\ & \times _0^{\mathrm{}}\frac{ds}{2\pi s^{11}}\eta (\frac{is}{\pi })^{24}e^{\frac{\alpha ^{}s}{2}k_Ag^{AB}k_B}\hfill \\ & \times _0^1d\nu _1d\nu _2d\nu _3d\nu _4\mathrm{\Psi }_1\mathrm{\Psi }_2\mathrm{\Psi }_{12}e^{\frac{i}{2}\left[k_3\times k_4(2\nu _{34}ϵ(\nu _{34}))k_1\times k_2(2\nu _{12}ϵ(\nu _{12}))\right]},\hfill \end{array}$$
$`(4.1)`$
with
$$\begin{array}{cc}\hfill \mathrm{\Psi }_1& =|\frac{\theta _{11}(\nu _{12},\frac{is}{\pi })}{\theta _{11}^{}(0,\frac{is}{\pi })}|^{2\alpha _{}^{}{}_{eff}{}^{}k_1k_2},\mathrm{\Psi }_2=|\frac{\theta _{11}(\nu _{34},\frac{is}{\pi })}{\theta _{11}^{}(0,\frac{is}{\pi })}|^{2\alpha _{}^{}{}_{eff}{}^{}k_3k_4},\hfill \\ \hfill \mathrm{\Psi }_{12}& =e^{\frac{s}{4}}\underset{r=1,2s=3,4}{}|\frac{\theta _{10}(\nu _{rs},\frac{is}{\pi })}{\theta _{11}^{}(0,\frac{is}{\pi })}|^{2\alpha _{}^{}{}_{eff}{}^{}k_rk_s},(\nu _{rs}=\nu _r\nu _s),\hfill \\ \hfill k& =k_1+k_2,k_r\times k_s=k_{rA}\mathrm{\Theta }^{AB}k_{sB},k_rk_s=k_{rA}G^{AB}k_{sB}.\hfill \end{array}$$
$`(4.2)`$
The expression for the annulus amplitude in (4.1) is written in the closed string channel. (The expression for the superstring would be similar except that the factor of $`s^{11}=s^{d_t/2}s^3`$ in (4.1). This factor comes from the number of transverse dimensions $`d_t`$. ) Closed string singularities arise in the integral over $`s`$ in (4.1) as $`\eta (\frac{is}{\pi })`$ may be expanded in a series in $`e^{Ns}`$. We thus find non analyticities<sup>4</sup> These singularites are 10 dimensional poles integrated over $`d_t`$ transverse momenta. (singularities) in the amplitude when
$$\frac{\alpha ^{}}{2}k_Ag^{AB}k_B=N.$$
$`(4.3)`$
In the NCOS scaling limit, this condition may be written as
$$\frac{\pi \alpha _{}^{}{}_{}{}^{2}}{\theta G_o^4}k_\mu \eta ^{\mu \nu }k_\nu +\frac{\theta }{4\pi }k_i\delta ^{ij}k_j=N.$$
$`(4.4)`$
Singularities on the real axis occur at a squared energy
$$k_0^2=k_1^2+\left(\frac{G_0^2\theta }{2\pi \alpha ^{}}\right)^2\left(k_2^2+k_3^2+\frac{2N}{\alpha _{}^{}{}_{eff}{}^{}}\right)$$
that becomes arbitrarily large as $`\alpha ^{}`$ is made increasingly small. In the strict limit $`\alpha ^{}0`$, open string one loop amplitudes factorize on singularities of the form
$$d^{d_t}k_M\frac{1}{k_2^2+k_3^2+\frac{2N}{\alpha _{}^{}{}_{eff}{}^{}}+g^{MN}k_Mk_N}.$$
As these singularities are never in the physical region, they do not correspond to physical states.<sup>5</sup> These singularities $`(k_i^2)^{\frac{d_t2}{2}}\mathrm{ln}(k_i^2)`$ are very similar to those induced by one loop graphs in spatially noncommutative field theories, as found in , . Notice that if $`d_t2`$ ($`p`$ branes with $`p<7`$ in the supersymmetric case), this amplitude, though non analytic, is finite at $`k=0`$. For $`d_t2`$ the amplitude diverges at $`k_i^2=0`$. It is possible that stronger IR singularities appear at higher loops, specially for high dimensional branes. Recalling that $`G^{AB}`$ is fixed in the NCOS limit, it is easy to see that the amplitude (4.1) is finite (except of course for the tachyon pole which is absent in the superstring). It is also straightforward using the results of to show that there are no physical poles for any numbers of initial and final open string tachyon vertex operators. Higher mass vertex operators involve additional powers of the Green functions on the annulus. These are finite in the NCOS limit and so will not spoil the finiteness of the amplitudes. Although we have not worked out the details, we expect that the behavior of the superstring is similar.
It is instructive to contrast the behaviour of (4.1) in both the NCOS and the NCYM limits. In the latter case the $`\alpha ^{}0`$ limit is manifestly smooth when $`\frac{s}{\alpha ^{}}`$ is held fixed. This forces one into a corner of the moduli space in which the massive open string states are decoupled \[11--16\]. In the NCOS limit (4.1) receives contributions from finite $`s`$, and so from all open string oscillator states. Apart from the non-commutative phases the one loop open string diagram (4.1) has almost the same form as the corresponding diagram in a theory with $`B=0`$, with $`\alpha ^{}`$ replaced by $`\alpha _{eff}^{}`$. However, the exponential term in (4.1) coming from momentum flowing along the closed string channel has a different $`\alpha ^{}`$ dependence from standard string theory with zero $`B`$. This different dependence is responsible for the absence of physical closed string poles.
Fig. 1: Nonplanar open string diagram. In open string field theory we would build it from the cubic vertex and we would consider states carrying momentum $`q`$ and $`q+p`$ along the loop.
The absence of closed string poles in a non-commutative open string theory, whose non-commutativity parameter $`\theta `$ is $`2\pi \alpha _{eff}^{}`$ as in our NCOS theories, may be understood more directly, as we explain below. This line of reasoning also suggests that a non commutative open string theory with $`\theta <2\pi \alpha _{}^{}{}_{eff}{}^{}`$ has closed string poles, while the theory with $`\theta >2\pi \alpha _{}^{}{}_{eff}{}^{}`$ is unstable.
Consider the simple non-planar diagram represented in figure 1, in an open string field theory. Let the open string theory in question be noncommutative, with noncommutativity parameter $`\theta `$. The momentum integral for this diagram takes the form
$$d^4qe^{2ip\times q}I_{\theta =0}(q,p)d^4q_0^{\mathrm{}}𝑑te^{2ip\times q}e^{2\pi \alpha _{eff}^{}tq^2+t\beta p.q+\mathrm{}}$$
$`(4.5)`$
where $`I_{\theta =0}(p,q)`$ is the integrand at $`\theta =0`$ and $`p\times q=p_\mu q_\mu \mathrm{\Theta }^{\mu \nu }/2`$. We have exponentiated the propagators in the diagram using a Schwinger proper time representation, where $`t`$ is the total proper time along the loop and we have explicitly given the form of the leading dependence on $`q`$ ($`\beta `$ is some other Schwinger parameter, which is also integrated over; we have supressed this integral in (4.5) for simiplicity). When $`q`$ is integrated over we get the diagram as a function of $`t`$ and $`\beta `$. As in , the effect of noncommutativity on this integral is an extra term in the exponent of the form
$$e^{pop/(8\pi \alpha _{eff}^{}t).}$$
$`(4.6)`$
where $`pop=p_\mu \mathrm{\Theta }^{\mu \nu }\mathrm{\Theta }_{\nu \rho }p^\rho =\theta ^2p^2`$. This may be seen by shifting the integral over $`q`$ to one over $`q_\mu ^{}=q_\mu +i\mathrm{\Theta }_{\mu \nu }p^\nu /(4\pi \alpha _{eff}^{}t)`$. Note that terms of the form $`q.p`$ in (4.5) are unaffected by the shift due to the antisymmetry of $`\mathrm{\Theta }`$.
Thus the integrand of (4.5) is modified from its $`\theta =0`$ value only by the additional exponential factor (4.6). On shifting to the $`s=\pi /t`$ channel, the integrand has the usual terms of the form $`e^{s\frac{\alpha _{eff}^{}}{2}(p_0^2+p_1^2+\mathrm{})}`$ (terms that would produce the $`s`$-channel poles if $`\theta `$ were zero) multiplied by the additional factor $`e^{s\frac{\theta ^2}{8\pi ^2\alpha _{}^{}{}_{eff}{}^{}}(p_0^2p_1^2)}`$. When $`\theta =2\pi \alpha _{eff}^{}`$ this extra factor exactly cancels the $`p_0,p_1`$ dependence of the exponent. Here we have used the fact that we are in Lorenzian signature so that the final sign of the exponent in (4.6) is the opposite to the one in Euclidean signature. If $`\theta `$ is slightly less that its critical value, then (4.6) does not cancel the closed string poles. If $`\theta `$ is bigger than its critical value then all closed string poles turn tachyonic, a reflection of the instability of the system.
5. Higher Loop diagrams
In this section we will examine higher loop string diagrams in the NCOS limit. We will not attempt to prove that the limit is nonsingular for arbitrary diagrams, but we will observe that a simple counting of powers of $`\alpha ^{}`$ does not reveal any difficulties. Naively, a genus $`g`$ surface in the string loop expansion is weighted by $`g_{str}^{2g2}`$. As $`g_{str}`$ diverges in the NCOS limit, a perturbative expansion in genus seems impossible. However, we shall argue below that both holes and handles are really weighted by powers of $`G_0`$ and so high genus surfaces are suppressed at weak open string coupling.
5.1. Holes
The addition of a hole in the world sheet is accompanied by one power of $`g_{str}`$. It also leads to an additional integral over the zero mode momentum circulating around the loop. As shown in , these integrals have a measure factor proportional to $`\mathrm{det}^{1/2}(g+2\pi \alpha ^{}B)\mathrm{det}^{1/2}(g)`$. Hence the total weighting of a hole is
$$g_{str}\frac{\mathrm{det}^{1/2}(g+2\pi \alpha ^{}B)}{\mathrm{det}^{1/2}(g)}=G_o^2,$$
$`(5.1)`$
and is finite as $`\alpha ^{}0`$.
5.2. Handles
Consider an open string world sheet $`A`$, with open string boundary conditions corresponding to a 3-brane. The amplitude on a worldsheet $`(B)`$ with an additional handle can be factorized in the closed string channel along the handle. The resultant amplitude reads schematically as
$$S_B=S_{A_{V_a,V_a}}\lambda _{eff}^2d^6k\frac{1}{g^{IJ}k_Ik_J+m_a^2};(I,J=0\mathrm{}9).$$
Here $`S_{A_{V_a,V_a}}`$ denotes the amplitude on $`A`$ with two extra closed string insertions. The integral is over the momenta of the intermediate states in the transverse directions (momentum is not conserved in these directions).
Fig. 2: Adding a handle to a worldsheet $`A`$, we obtain a worldsheet $`B`$, which can be represented as coming from the propagation of closed string states between two points of the worldsheet. We sum over all closed string states.
The effective coupling $`\lambda _{eff}^2`$ is determined as follows: A closed string mode $`\varphi `$ with spacetime action
$$S=\frac{1}{g_{str}^2\alpha _{}^{}{}_{}{}^{4}}d^{10}x\sqrt{g}(_I\varphi _J\varphi g^{IJ}+m_a^2\varphi ^2)$$
has effective coupling
$$\lambda _{eff}=\frac{g_{str}\alpha _{}^{}{}_{}{}^{2}}{g^{\frac{1}{4}}}=\frac{\alpha _{}^{}{}_{}{}^{\frac{5}{2}}}{G_0^4\alpha _{}^{}{}_{eff}{}^{\frac{1}{2}}}$$
in the NCOS limit. The integral
$$d^6k\frac{1}{g^{IJ}k_Ik_J+m_a^2}=\alpha ^{}d^6k\frac{1}{\frac{\alpha _{}^{}{}_{}{}^{2}k_Mk_N\delta ^{MN}}{\alpha _{}^{}{}_{eff}{}^{}G_0^4}+N+\alpha _{}^{}{}_{eff}{}^{}(k_2^2+k_3^2)+\mathrm{}}$$
is of order
$$\alpha ^{}\left(\frac{\alpha _{}^{}{}_{eff}{}^{}G_0^4}{\alpha _{}^{}{}_{}{}^{2}}\right)^3.$$
Finally, in the normalization we have adopted, $`S_{A_{V_a,V_a}}`$ is of the same order as $`S_A`$.
Putting it all together, we find that
$$\frac{S_B}{S_A}G_0^4\alpha _{}^{}{}_{eff}{}^{2}.$$
$`(5.2)`$
Thus we conclude that extra handles, in the NCOS limit, are neither infinitely suppressed nor enhanced in the NCOS limit. They are instead really weighted by a factor of $`G_0^4`$, as they would have been for an ordinary weakly coupled open string theory<sup>6</sup> See also for a discussion of diagrams with many holes..
6. Supergravity duals
The considerations of the previous sections generalize to open string theories on $`N`$ coincident 3-branes. In that case since we are dealing with a deformation of $`U(N)`$ $`𝒩=4`$ SYM we expect that it should have a supergravity dual for large $`N`$. The relevant supergravity solutions were written in \[24,,25\]. We start from the Lorentzian version of the solution (2.3) in , with $`B_{23}=0`$. Then we do the following scaling of parameters
$$\begin{array}{cc}\hfill r=& \sqrt{\alpha ^{}}u\hfill \\ \hfill \mathrm{cosh}\theta ^{}=& \frac{\stackrel{~}{b}^{}}{\alpha ^{}}\hfill \\ \hfill g=& \frac{\stackrel{~}{g}\stackrel{~}{b}^{}}{\alpha ^{}}\hfill \\ \hfill x_{0,1}=& \frac{\stackrel{~}{b}^{}}{\sqrt{\alpha ^{}}}\stackrel{~}{x}_{0,1}\hfill \\ \hfill x_{2,3}=& \sqrt{\alpha ^{}}\stackrel{~}{x}_{2,3}\hfill \\ \hfill R^4=& fixed=4\pi \stackrel{~}{g}N\hfill \end{array}$$
$`(6.1)`$
We obtain<sup>7</sup> Here we normalize the $`B`$ field as in the previous sections, in it was normalized differently by a factor of $`2\pi \alpha ^{}`$, $`B_{MR}=2\pi \alpha ^{}B_{here}`$.
$$\begin{array}{cc}\hfill ds_{str}^2& =\alpha ^{}f^{1/2}\left[\frac{u^4}{R^4}(d\stackrel{~}{x}_{0}^{}{}_{}{}^{2}+d\stackrel{~}{x}_{1}^{}{}_{}{}^{2})+f^1(d\stackrel{~}{x}_{2}^{}{}_{}{}^{2}+d\stackrel{~}{x}_{3}^{}{}_{}{}^{2})+du^2+u^2d\mathrm{\Omega }_5^2\right]\hfill \\ \hfill 2\pi \alpha ^{}B_{01}& =\alpha ^{}\frac{u^4}{R^4},\hfill \\ \hfill e^{2\varphi }& =\stackrel{~}{g}^2f\frac{u^4}{R^4}\hfill \\ \hfill A_{23}=& \alpha ^{}\frac{1}{\stackrel{~}{g}}f^1,\hfill \\ \hfill F_{0123u}& =\alpha _{}^{}{}_{}{}^{2}\frac{1}{\stackrel{~}{g}}\frac{4f^1}{u}\hfill \\ \hfill f& =1+\frac{R^4}{u^4}\hfill \end{array}$$
$`(6.2)`$
The particular scalings that we have to do to reproduce this solution are, up to constants, the same as those in section 2.2. The only scaling that is not so obvious is the scaling of the radial coordinate. Notice that in the $`𝒩=4`$ SYM case we rescale the radial coordinate as $`r\alpha ^{}u`$. The fact that we have $`r\sqrt{\alpha ^{}}u`$ in this case is related to the fact that the closed string metric has a factor of $`1/\alpha ^{}`$ in section 2.2. We see from 5.2 that for small $`u`$ we recover the usual $`AdS_5\times S^5`$ solution as we expect, since the open string theory reduces to $`𝒩=4`$ SYM at low energies. In particular we see that we should identify $`\stackrel{~}{g}=G_0^2`$. As we increase $`u`$ the metric becomes different than the metric of $`AdS`$ and we also see that the dilaton becomes large. This suggests that for large $`u`$ we should do an S-duality to analyze the solution. After we do the S-duality we obtain a solution which is the same as the supergravity solution which corresponds to a D3 brane with spatial non-commutativity in the directions 23, see , eqn. (2.7). This suggests that at very high energies the open string theory we are studying would have a dual description in terms of the theory with spatial non-commutativity.
7. Discussion
7.1. Open String Dipoles and UV/IR
Free open string states in the NCOS limit behave quite differently from ordinary open strings propagating in the same metric, despite having the same spectrum. In the presence of background fields, (as discussed for example in \[7,,8,,9\] and especially in for the magnetic case) the mode expansion reads
$$X^\mu (\sigma ,\tau )=x_0^\mu +2i\alpha _{}^{}{}_{eff}{}^{}p^\mu \tau +\frac{1}{\pi }\mathrm{\Theta }^{\mu \nu }p_\nu \sigma +(\mathrm{oscillators}).$$
$`(7.1)`$
For strings in the NCOS limit this implies that the distance along the direction of the field between the two ends of the string , as measured in the metric $`G^{AB}`$, is
$$\mathrm{\Delta }X^1=2\pi \alpha _{}^{}{}_{eff}{}^{}k_0,$$
$`(7.2)`$
plus oscillator contributions which time average to zero. (Note that $`\mathrm{\Delta }X^1`$ is the distance between the endpoints of the string worldsheet along a line of constant worldsheet time rather than along a line of constant $`X^0`$. As we argue below, the proper length of the string is given by a formula analogeous to (7.2) with $`k_0`$ replaced by the centre of mass energy of the state.)
The invariant energy and proper length of an oscillator state may be estimated as follows. The tension of an open string aligned with a near critical electric field is almost canceled by a negative contribution from its dipole interaction with the field. In the NCOS limit, the effective tension $`T_{eff}=\frac{1}{4\pi \alpha _{}^{}{}_{eff}{}^{}}`$ (see (2.11)) . The energy of such a string, with an oscillation number $`N`$, is $`E=T_{eff}L+\frac{\pi N}{L}`$. This is minimized for $`L=2\pi \sqrt{\alpha _{eff}^{}N}=2\pi \alpha _{eff}^{}E`$, as in (7.2).
7.2. Thermodynamics
At low energies, compared to $`\frac{1}{\sqrt{\alpha _{}^{}{}_{eff}{}^{}}}`$ the NCOS theory reduces to ordinary $`𝒩=4`$ SYM, and its free energy scales like $`T^4`$. At intermediate energies, the thermodynamics of a weakly coupled NCOS theory ($`\lambda =G_o^2N1`$), may be expected to reflect its Hagedorn density of states.
However, as argued in this paper, the weakly coupled NCOS theory has a dual description as a strongly coupled NCYM theory. In a spatially noncommutative field theory, at weak coupling, planar diagrams dominate over nonplanar diagrams for energies $`k_0\frac{1}{\sqrt{\theta }}`$. It is plausible that this result to continues to hold at strong coupling<sup>8</sup> This statement is true at least in the ‘supergravity’ limit $`\lambda 1`$, $`G_o^21`$; in that limit , , supergravity suggests that planar diagrams dominate for $`k_0\frac{1}{(\lambda \theta ^2)^{\frac{1}{4}}}`$., with a crossover scale renormalized by a function of the coupling. If true, this assertion implies that, at high temperatures, the free energy of spatially noncommutative SYM is proportional to the free energy of ordinary large $`N`$ SYM, and so scales with temperature like $`T^4`$, even at large $`G_o^{}_{}{}^{}2`$.
It would be interesting to investigate this issue further.
7.3. Generalizations to other Dimensions
In this paper we have ‘derived’ the existence of a decoupled four dimensional open string theory, NCOS, by $`S`$ dualizing spatially noncommutative SYM. We presented evidence that, independent of this derivation, the resultant theory is well defined, and weakly coupled over a range of parameters.
It is easy to extend our construction of the NCOS to other dimensions, even though we do not have an independent ($`S`$ duality) argument for the decoupling of closed strings. The NCOS scaling limit for a $`p`$ brane is, once again, defined by table 1, where the indices $`i,j`$ run from $`2\mathrm{}p`$ and $`A,B`$ from $`0\mathrm{}p`$. In other words, this limit still describes a near critical electric field turned on in the $`1`$ direction. The open string coupling defined in (2.9) and the effective low energy Yang Mills coupling constant $`g_{YM}^2G_0^2\alpha _{}^{}{}_{eff}{}^{\frac{p3}{2}}`$ are finite. In the NCOS limit, open strings appear to decouple from closed strings for all $`p`$. The annulus amplitude is finite in arbitrary dimension, and always factorizes on unphysical closed string poles. As in the 3-brane, string diagrams with handles and holes are suppressed by powers of the open string coupling, and may be neglected at weak coupling.
In fact, these open string theories appear to be non-gravitational UV finite completions of low energy (supersymmetric) Yang-Mills. This statement appears to be true even in high dimensions where the gauge theory is non-renormalizable.
Acknowledgements
We are grateful to C. Bachas, M. Berkooz, M. Gutperle, J. Harvey, S. Kachru, H. Liu, J. McGreevy, P. Kraus, N. Seiberg, A. Sen, S. Shenker, E. Silverstein, L. Susskind and N. Toumbas for useful discussions. R.G. and S.M. are grateful to the high energy theory group at Stanford University for their hospitality.
This work was supported in part by DOE grant DE-FG02-91ER40654.
References
relax Alain Connes, Michael R. Douglas and Albert Schwarz, Noncommutative Geometry and Matrix Theory: Compactification on Tori, hep-th/9711162, JHEP 9802 (1998) 003. relax N. Seiberg, L. Susskind and N. Toumbas, The Teleological Behavior of Rigid Regge Rods, hep-th/0005015. relax N. Seiberg and E. Witten String Theory and Noncommutative Geometry, hep-th/9912072. relax S. Gukov, I. R. Klebanov and A. M. Polyakov, Dynamics of (n,1) strings, Phys. Lett. B423, 64 (1998) \[hep-th/9711112\]. relax H. Verlinde, A matrix string interpretation of the large N loop equation, hep-th/9705029. relax C.G. Callan, C. Lovelace, C.R. Nappi, S.A. Yost, Open Strings in Background gauge Fields, Nucl.Phys. B288 525,1985. relax C.P.Burgess, Open String Instability in Background electric Fields, Nuc. Phys. B294427-444, 1987. relax C. Bachas and M. Porrati, Pair Creation of Open Strings in an Electric Field, Phys.Lett.B296:77-84,1992. relax V.V.Nesterenko, The Dynamics of Open Strings in a Background Electromagnetic Field, Int. J. Mod. Phys. A4 (1989) 2627-2652. relax E.S. Fradkin, A.A. Tseytlin, Nonlinear Electrodynamics from Quantized Strings Phys.Lett.B163 123,1985. relax Oleg Andreev and Harald Dorn, Diagrams of Noncommutative Phi-Three Theory from String Theory , hep-th/0003113. relax Youngjai Kiem and Sangmin Lee, UV/IR Mixing in Noncommutative Field Theory via Open String Loops, hep-th/0003145. relax Jaume Gomis, Matthew Kleban, Thomas Mehen, Mukund Rangamani and Stephen Shenker, Noncommutative Gauge Dynamics From The String Worldsheet,hep-th/0003215. relax Hong Liu and Jeremy Michelson, Stretched Strings in Noncommutative Field Theory, hep-th/0004013. relax Adel Bilal, Chong-Sun Chu, Rodolfo Russo, String Theory and Noncommutative Field Theories at One Loop, hep-th/0003180. relax Chong-Sun Chu, Rodolfo Russo and Stefano Sciuto, Multiloop String Amplitudes with B-Field and Noncommutative QFT, hep-th/0004183. relax N. Seiberg, L. Susskind and N. Toumbas, Strings in background electric field, space / time noncommutativity and a new noncritical string theory, hep-th/0005040. relax A.A. Tseytlin, Born-Infeld action, supersymmetry and string theory hep-th/9908105, and references therein. relax V. Schomerus, D-branes and Deformation Quantization, hep-th/9903205, JHEP 9906 (1999) 030. relax E. Witten, Noncommutative Geometry And String Field Theory, Nucl. Phys. B268, 253 (1986). relax Shiraz Minwalla, Mark Van Raamsdonk, Nathan Seiberg, Noncommutative Perturbative Dynamics , hep-th/9912072. relax M. van Raamsdonk and N. Seiberg, Comments on Noncommutative Perturbative Dynamics,JHEP 0003 (2000) 035 relax O. Andreev, A note on open strings in the presence of constant B-field, Phys. Lett. B481, 125 (2000) hep-th/0001118. relax A. Hashimoto and N. Itzhaki, Non-Commutative Yang-Mills and the AdS/CFT Correspondence, Phys.Lett. B465 (1999) 142 relax J. Maldacena and J. Russo, The Large N limit of non-commutative gauge theories, hep-th/9908134, JHEP 9909:025,1999. relax T. Filk, Divergences in a Field Theory on Quantum Space, Phys.Lett. B376:53-58,1996.
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# Interlayer c-axis transport in the normal state of cuprates
## Abstract
A theoretical model of c-axis transport properties in cuprates is proposed. Inter-plane and in-plane charge fluctuations make hopping between planes incoherent and diffusive (the in-plane momentum is not conserved after tunneling). The non-Drude optical conductivity $`\sigma _c(\omega )`$ and the power-law temperature dependence of the dc conductivity are generically explained by the strong fluctuations excited in the process of tunneling. Several microscopic models of the charge fluctuation spectrum are considered.
Despite the strongly two-dimensional layered structure of the high-temperature cuprate superconductors, features associated with the third dimension, perpendicular to the $`CuO_2`$ planes, may be an important ingredient in their superconductivity. In fact, it is well accepted that a certain degree of Josephson-type coupling between different planes is necessary to suppress the two-dimensional fluctuations, which will otherwise destroy the superconducting long-range order. However, the systematic dependence of the critical temperature $`T_c`$ on the number of layers in the unit cell (together with the absence of evidence for strong fluctuations effects above $`T_c`$ at optimal doping, which suggests that these fluctuations are not the major reason for this systematic dependence) points almost unambiguously to the conclusion that theories formulated for a single plane cannot be the whole story. Either hopping between planes, or Coulomb interaction between them, or both, is an important factor in raising the critical temperature (and, perhaps, in some cases also for lowering it, see Ref. ). In the light of this, the study of the c-axis optical and transport properties is more than just a minor diversion from the main issue.
These c-axis optical and transport properties are very puzzling and anomalous. Most remarkable is the fact that the temperature dependence of the dc c-axis resistivity $`\rho _c(T)`$, in sharp contrast to the well-known linear $`T`$-dependence of the in-plane resistivity $`\rho _{ab}(T)`$, is non-universal, being described in most cases by a power law $`\rho _c(T)T^\gamma `$, where however the exponent $`\gamma `$ can be anything in between approximately $`+1`$ and $`1`$. The optical conductivity $`\sigma _c(\omega )`$ is roughly frequency independent from low frequencies up to mid-infrared frequencies (except in the case of some overdoped cuprates ($`YBa_2Cu_3O_7`$ and $`La_{2x}Sr_xCuO_4`$)); the numerical value is below the Mott-Ioffe-Regel minimum metallic conductivity. This behavior, which is dramatically different from the behavior of the in-plane resistivity, has been christened “confinement”. Thus, despite the dramatic differences in the raw data between the different cuprate families, one can isolate at least two elements which may be legitimately called “universal”: a non-Drude optical conductivity of a magnitude below the Mott limit, and a power-law temperature dependence of the dc resistivity $`\rho _c(T)`$ (albeit with a material-specific exponent). In this paper we develop a framework for the explanation of these universalities, which will hopefully shed light on the role of the inter-plane and in-plane Coulomb interaction as well as on the more obvious one of the inter-layer hopping.
On the basis of the above experimental observations, we can make the assumption of incoherent transport (diffusive tunneling) in the c-direction. The inter-plane (or rather inter-unit cell) hopping time $`\tau _{hop}`$ can be estimated from the dc c-axis resistivity. Relating the diffusion constant $`D`$ to the hopping time $`\tau _{hop}`$ by $`D=d^2/\tau _{hop}`$, we can derive a model-independent relation between the conductivity $`\sigma _c`$ and the hopping time: $`\sigma _c=e^2\nu _{2d}\frac{d}{\tau _{hop}}`$, where $`\nu _{2d}`$ is a two-dimensional density of states. Using this formula and the experimental values for the c-axis conductivity we can estimate the c-axis hopping time $`\tau _{hop}`$. The strong two-dimensionality of the electron motion becomes obvious if we compare the hopping time $`\tau _{hop}`$ with the in-plane scattering time $`\tau _{ab}`$. As is well known, the in-plane scattering time $`\tau _{ab}`$ is of order of $`\mathrm{}/kT`$. Direct comparison shows that for most materials the c-axis hopping time is much longer than the scattering time in the plane. These two times are comparable only for overdoped $`YBa_2Cu_3O_7`$ and $`La_{2x}Sr_xCuO_4`$ suggesting a crossover to a different regime of c-axis transport; this is confirmed by the experimental observation of the Drude-like frequency dependence of the conductivity $`\sigma _c(\omega )`$ for these compounds.
Many approaches have been suggested to describe the c-axis transport properties. Most of them stem from phenomenologically assumed in-plane Green functions. One remarkable example is a non-Fermi “Luttinger” liquid theory which explains naturally the “confinement” for c-axis motion in the normal state . Many other theories are essentially based on the Fermi liquid theory modified by strong correlations.
We take a quite different approach to the problem. Although in our approach the in-plane motion (expressed by in-plane Green functions) is undeniably important, we show that most of the c-axis properties can be qualitatively understood on the basis of knowledge of the spectrum of in-plane and inter-plane charge density fluctuations excited during the process of the inter-plane tunneling. The spectrum of charge fluctuations can be directly measured experimentally by optical reflectivity and electron-energy-loss spectroscopy (EELS).
The essential physical picture of our approach is that the c-axis tunneling is strongly suppressed by charge fluctuations excited in the process of tunneling. This anomaly (the so called Coulomb blockade) is widely observed in many strongly correlated and mesoscopic systems. Other examples of this class of phenomena are orthogonality catastrophes and a zero-bias anomaly in diffusive systems. The ubiquity of the Coulomb blockade phenomena (static or dynamic) in correlated systems indicates that the anomalous c-axis transport properties may be merely a consequence of strong correlations in the cuprates. The necessary condition for the appearance of the Coulomb blockade phenomenon is the strong effective coupling of the tunneling electron to the collective excitations of the liquid. In fact, we can think about the Coulomb blockade as a “high-energy” phenomenon of order of Coulomb energy per electron and independent of the low-energy quasiparticle spectral properties, be they fermi- or non-fermi liquid. In other words, the tunneling electron couples to the excitations in the broad range of frequencies from low up to high frequencies. The non-fermi liquid property (property not present for a three-dimensional Fermi liquid or for a electron gas in the RPA approximation) which is responsible for the anomalous c-axis properties is simply the large density of “detuning” charge fluctuations (as observed by optical and Raman spectroscopy) over a broad frequency range up to mid-infrared frequencies.
Empirically, an overview of the in-plane conductivity and the c-axis conductivity in various families of cuprates does not reveal any obvious correlation between their temperature dependencies. On the other hand, the perturbative diagrammatic expression (assuming the equivalent Green’s functions in each plane and uncorrelated impurities) for the in-plane conductivity and out-of-plane conductivity are equivalent up to the vertex functions. Therefore the difference in the temperature and frequency dependence of the in- and out- plane conductivities must come exclusively from the interplane tunneling probability. For instance, the notion of the “two-dimensional Luttinger” liquid is not sufficient by itself to explain the difference between the in-plane and out-of-plane resistivities .
One particular approach to explaining the difference between the resistivities along the different directions is based on the highly anisotropic form of the tunneling matrix element $`t_{}(k_x,k_y)`$ as a function of the in-plane momentum. Several authors developed a phenomenological approach (assuming as well a strong anisotropy of quasiparticle lifetimes and the density of states around the fermi surface) which seem to fit successfully the experimental data. Two remarks are in order. First, this approach assumes that the tunneling conserves the in-plane momentum. This may be the case in certain situations (possibly, in the superconducting state and in the overdoped regime), but in general this assumption deserves a close scrutiny by experiment and theory. In fact, we will argue that the in-plane momentum is not conserved when $`\tau _{hop}\tau _{ab}`$. Second, the anisotropy of the tunneling matrix element $`t_{}(k_x,k_y)`$ is different for some variations of cuprates, and , in general, it can be doping dependent. Thus in some cases, $`t_{}(k_x,k_y)`$ may not vanish along the diagonals of the Brillouin zone making the contribution of the diagonal quasiparticles ($`k_x=\pm k_y`$) to the c-axis conductivity non-vanishing contrary to the assumptions of references.
In the rest of the paper, we begin by generalizing the standard tunneling formalism, introducing a non-trivial tunneling probability which accounts for the inelastic and elastic (momentum scattering) processes. Then, we calculate this tunneling probability from assumed spectra of “the detuning fluctuations”. After that, we calculate the experimentally measured optical conductivity $`\sigma _c(\omega ,T)`$ and the tunneling conductance $`\sigma _c(V)`$. Finally, we discuss the complex experimental situation and possible extensions of the proposed theory. In the appendix we discuss several microscopic models giving the spectrum of the charge fluctuations.
Theory. The tunneling formalism. The tunneling of an electron from one plane to another plane can be considered by using the time-dependent tunneling hamiltonian formalism. The “blockade effect” due to the excitation of the electromagnetic modes is accounted by the modulation of the tunneling matrix element by the Coulomb interaction. Thus the part of the hamiltonian responsible for the transfer of electrons between planes is:
$`H_c={\displaystyle \underset{r_1,r_2}{}}t_{}(r_1,r_2;t)(a_1^+(r_1)a_2(r_2)+a_2^+(r_2)a_1(r_1)),`$ (1)
where the quasi-classical tunneling matrix element $`t_{}(r_1,r_2;t)`$ is equal to $`t_{}(r_1,r_2)exp\left(\frac{ie}{c}_{r_1}^{r_2}A(z,t)𝑑z\right).`$ Gauge invariance dictates the presence of the phase factor $`\phi (r_1,t)=(e/c)_{r_1}^{r_2}A(z,t)𝑑z`$, where the integral is taken over a path connecting two points $`r_1`$ and $`r_2`$ on the different planes. Since the optimal tunneling trajectory is perpendicular to the planes (along the $`z`$ axis), the tunneling matrix element can be written as $`t_{}(r_1,r_2)=t_{}\delta (r_1r_2)exp\left(\frac{ie}{c}_{z_1}^{z_2}A(z,t)𝑑z\right)`$. It conserves the in-plane momentum. Using the tunneling hamiltonian formalism , we get the following expression for the tunneling current $`I(t)`$ between two planes:
$`I(t)={\displaystyle \frac{2e}{\mathrm{}^2}}Re{\displaystyle 𝑑r𝑑r^{}_{\mathrm{}}^+\mathrm{}𝑑t^{}e^{i\frac{eVt}{\mathrm{}}}|t_{}|^2P(rr^{},tt^{})S(rr^{},tt^{})},P(rr^{},tt^{})e^{i\phi (r,t)}e^{i\phi (r^{},t^{})},`$ (2)
$`S(rr^{},tt^{})\mathrm{\Theta }(tt^{})<G_R^<(rr^{},tt^{})G_L^>(r^{}r,t^{}t)G_L^<(rr^{},tt^{})G_R^>(r^{}r,t^{}t)>`$ (3)
where $`P(rr^{},tt^{})`$ is a phase-phase correlation function between two planes averaged over the equilibrium fluctuations, and $`V`$ is an applied voltage. The definitions of Green’s functions and essential details of the derivation can be found in the Ref.. The fact that the hopping time $`\tau _{hop}`$ is much longer than the in-plane scattering time $`\tau _{ab}`$ allows us to separate the in-plane propagation $`S(rr^{},tt^{})`$ and the tunneling probability $`P(rr^{},tt^{})`$. It is important to remark at this stage that the long-wavelength fluctuation modes (with the wavelength much longer than the in-plane mean free path) can suppress the tunneling probability without effecting the in-plane motion. The scattering by the short-wavelength fluctuations is accounted by the spectral properties of the in-plane propagation $`S(rr^{},tt^{})`$. The understanding of the properties of the tunneling probability $`P(rr^{},tt^{})`$, describing the effect of the “detuning fluctuations”, is imperative for any particular problem of the tunneling. The importance of this correlation function was first described in Ref.. The novel element here is the discussion of the spatial dependence of the tunneling probability function $`P(rr^{},tt^{})`$. The spatial dependence appears to be very important for many questions of the c-axis transport properties. As mentioned above, in previous studies of c-axis transport in cuprates, specific properties of the tunneling probability (e.g. the tunneling with or without the conservation of the in-plane momentum $`k_{}`$) were assumed. Here we analyze and calculate the tunneling probability from the fluctuation spectrum of the electromagnetic field. We can call the tunneling “diffusive” if the in-plane momentum is not conserved (if the momentum is conserved, it can be called specular). In other words, the tunneling is diffusive, if the tunneling probability $`P(rr^{},tt^{})`$ is significant only if $`|rr^{}|/l\stackrel{<}{}1`$ (where $`l`$ is a short length scale of order of a lattice constant). It should be noted that, generally speaking, the question of the conservation of in-plane momentum during tunneling is another aspect of tunneling not equivalent to the question of coherence or incoherence of tunneling (that is the question of the dephasing of an electron). Equation 2 can be rewritten in the following form:
$`I(V)={\displaystyle \frac{2eSt_{}^2}{\mathrm{}}}{\displaystyle }dEdE^{}dkdk^{}A_1(k,E)A_2(k^{},E^{})(f(E)(1f(E^{}))P(E+eVE^{},kk^{})`$ (4)
$`f(E^{})(1f(E))P(E^{}eVE,kk^{}))`$ , (5)
where $`A_{1,2}(k,E)`$ are the spectral functions, and $`f(E)`$ is a Fermi function.
The tunneling probability. We need to calculate the correlation function: $`e^{i\widehat{\phi }(r,t)}e^{i\widehat{\phi }(r^{},t^{})}`$. The averaging can be done if we assume the field $`\phi (r,\tau )`$ is Gaussian correlated. Using gauge invariance, the phase $`\phi (r,\tau )`$ can be rewritten as $`\phi (r,\tau )=_{\mathrm{}}^\tau \delta V(r,t)𝑑t`$, where $`\delta V(r,t)`$ is the local voltage difference between two planes. This way we get an expression for the tunneling probability (the calculation is a generalization of the exact calculation from Ref.(, p.273-277)):
$`P(\delta r=rr^{},\delta t=tt^{})exp(R(\delta r,\delta t)),`$ (6)
$`R(\delta r,\delta t){\displaystyle \frac{d\omega }{\omega ^2}d^2q<\delta V_{q,\omega }^2>(1cos(\omega \delta t+\stackrel{}{q}\stackrel{}{\delta r}))coth\frac{\omega }{2T}}.`$ (7)
In this paper, we assume that the most effective “detuning” fluctuations are the voltage fluctuations (or related charge fluctuations). It is important to point out that the same method can be used to calculate the “blocking” of the tunneling due to any mechanism of in-plane scattering. Since the nature of the ground state of cuprates and therefore the spectrum of the fluctuations is not known, later we examine several general forms of the spectrum. The problem of incoherent tunneling between a couple of two-dimensional planes is a natural generalization of the spin-boson model of quantum dissipation.
In the view of the importance of the spatial dependence of the tunneling probability, we give several different arguments proving the diffusive nature of the tunneling (if $`\tau _{hop}\tau _{ab}`$) in the normal state. First of all a qualitative argument: if the hopping time is much longer the in-plane scattering time, an electron experiences many inelastic and elastic scattering processes (both not conserving the direction of the in-plane momentum) before the hopping between planes. Thus it is intuitively natural to think that the momentum is not conserved after the hopping. A straightforward quantitative argument is given by the analysis of the function $`R(\delta r,\delta t)`$ in the exponent of the expression for the tunneling probability (Eqn. 7). To separate the spatial and time dependence of the function $`R(\delta r,\delta t)`$, we rewrite the multiplier in the integral as $`(1cos(\omega \delta t+q_x\delta r))=(1cos(\omega \delta t))+cos(\omega \delta t)(1cos(q_x\delta r))sin(q_x\delta r)sin(\omega \delta t)`$. The integral with the last term vanishes, because this term makes the integral expression antisymmetric with respect to integration over $`q_x`$. Thus, we can write
$`R(\delta r,\delta t)=R_0(\delta r,\delta t)+R_1(\delta r,\delta t),`$ (8)
$`R_0(\delta t)={\displaystyle \frac{d\omega }{\omega ^2}𝑑q_xq_y<\delta V_{q,\omega }^2>coth\frac{\omega }{2T}(1cos(\omega \delta t))},`$ (9)
$`R_1(\delta r,\delta t)=2{\displaystyle \frac{d\omega }{\omega ^2}𝑑q_xq_y<\delta V_{q,\omega }^2>coth\frac{\omega }{2T}cos(\omega \delta t)sin^2((q_x\delta r)/2)}.`$ (10)
The space-independent part $`R_0(\delta t)`$ is calculated later in the paper (see Eqn. 21). Below we calculate the function $`R_1(\delta r,\delta t)`$ which describes the spatial dependence of the tunneling probability. For this calculation we assume that the fluctuations are uncorrelated in two planes (see below for a more general discussion), and the interplane noise is just twice the in-plane Johnson-Nyquist (JN) noise (voltage fluctuations). The Johnson-Nyquist noise can be calculated from the spectral density of Coulomb noise in the two-dimensional plane valid in the hydrodynamic approximation:
$`<\delta V_{q,\omega }^2>4\pi \sigma _Q{\displaystyle \frac{\sigma _2\omega }{\omega ^2+4\pi ^2\sigma _2^2q^2}},`$ (11)
where $`\sigma _2`$ is a two-dimensional conductance, and $`\sigma _Qe^2/\mathrm{}`$. We substitute Eqn.11 into the expression (10) for $`R_1(\delta r,\delta t)`$ and impose an upper cutoff $`q_c`$ on the $`q`$-integration to take into account of the fact that the expression (10) is strictly valid only in the long-wavelength limit; thus we take $`qq_F`$ (but still of the general order of $`q_F`$). We also impose on the $`\omega `$-integration a lower cutoff $`\omega _l`$, the choice of which will be discussed below. Then, taking into account the fact that we are interested in values of $`\delta t`$ which are of the general order of magnitude $`\tau _{hop}`$ and thus several orders of magnitude larger than $`(\sigma _2q_c)^1`$, we see that to a good approximation $`R_1(\delta r,\delta t)`$ factorizes into a product of a function of $`\delta r`$ and a function of $`\delta t`$:
$`R_1(\delta r,\delta t)F(\delta r)G(\delta t),`$ (12)
$`F(\delta r){\displaystyle _0^{2\pi }}𝑑\theta {\displaystyle _0^{q_c\delta r}}𝑑x{\displaystyle \frac{1cos(xcos\theta )}{x}},`$ (13)
$`G(\delta t){\displaystyle \frac{2\sigma _Q}{\pi \sigma _2}}{\displaystyle _{\omega _l}^{\mathrm{}}}{\displaystyle \frac{d\omega }{\omega }}coth({\displaystyle \frac{\omega }{2kT}})cos(\omega \delta t).`$ (14)
The expression $`F(\delta r)`$ is approximately $`\frac{\pi }{8}(q_c\delta r)^2`$ for $`q_c\delta r1`$ and $`2\pi ln(q_c\delta r)`$ for $`q_c\delta r1`$. As for the function $`G(\delta t)`$, we will see below that the “interesting” values of $`\delta t`$ (i.e. those for which the function $`R_0(\delta t)`$ does not suppress $`P(\delta r,\delta t)`$ too badly) are less or of order of $`\mathrm{}/2\pi \alpha kT`$, where the dimensionless quantity $`\alpha `$ is typically less or of order of $`1`$. Under these conditions, provided $`\mathrm{}\omega _lkT`$ which will be satisfied by our choice of $`\omega _l`$ (cf. below), the integral defining $`G(\delta t)`$ is dominated by its lower limit and approximately given by the $`\delta t`$-independent expression
$`G(\delta t){\displaystyle \frac{2\sigma _Q}{\pi \sigma _2}}{\displaystyle \frac{kT}{\mathrm{}\omega _l}}.`$ (15)
We will make the choice $`\omega _l1/\tau _{hop}`$, on the grounds that once we need to allow for appreciable interplane hopping, Eqn.11 for the noise is no longer applicable (and we expect the expression for $`<\delta V_{q,\omega }^2>`$ to decrease as a higher power of $`\omega `$ for $`\omega 0`$, thereby effectively cutting off the integral (14)). We thus have for $`R_1(\delta r,\delta t)`$ the approximate expression
$`R_1(\delta r,\delta t)(kT\tau _{hop}/\mathrm{})(\sigma _Q/\sigma _2)(2/\pi )F(q_c\delta r).`$ (16)
The salient point, now, is that the quantity $`kT\tau _{hop}/\mathrm{}`$, which is essentially the ratio of the ab-plane and c-axis conductivities, is of order $`10^210^4`$ for most of the cuprates, while the ratio $`\sigma _2/\sigma _Q`$ is never greater than about $`10`$. Thus, the quantity $`R_1`$ has a value large compared to unity for values of $`\delta r`$ small compared to $`1/q_c`$, and we can approximate the expression $`F(\delta r)`$ by its limiting form $`\frac{\pi }{8}(q_c\delta r)^2`$. Thus, the “effective area” $`S_{eff}(\delta r_{eff})^2`$ for which $`R_1`$ is appreciable is defined by
$`S_eff{\displaystyle \frac{4(\sigma _2/\sigma _Q)}{q_c^2(kT\tau _{hop}/\mathrm{})}}`$ (17)
and by the above argument this is much smaller than $`1/q_c^2`$ and thus at most of the order of $`1/q_F^2`$. At distances of this order formulas such as (11) should no longer be taken seriously, but the crucial upshot of the argument is that coherence between tunneling events separated in space by more than $`1/q_F`$ can be simply neglected. To put it differently, the tunneling is effectively local (diffusive); the effective rms change in momentum in the course of a tunneling event is of order of $`q_F`$ (cf. below). It is noteworthy that this is so even if the momentum cutoff $`q_c`$ on the voltage fluctuations is only a small fraction of $`q_F`$; consideration of shorter-wavelength fluctuations can only strengthen this conclusion.
Another argument estimates directly the change of momentum in the process of tunneling. The change of the momentum due to a fluctuation of the electromagnetic potential is $`\delta p(e\delta A)/c`$, therefore $`<\delta p^2>=(e/c)^2<\delta A^2>`$. We can relate the correlation function of the vector potential with the correlation function of the scalar potential by gauge transformation (assuming only longitudinal fluctuations): $`<A_{q,\omega }^2>=\frac{c^2q^2}{\omega ^2}<\delta V_{q,\omega }^2>`$. Using Eqn. 11 at low frequencies $`\omega <\sigma _2q`$, the correlation function can be approximated as $`<A_{q,\omega }^2>=\frac{4\pi \sigma _2c^2}{\omega ^2}\frac{q^2}{\omega ^2+\sigma _2^2q^2}\frac{4\pi c^2}{\omega ^2\sigma _2}`$. Thus the variance of momentum $`<\delta p^2>`$ is $`<\delta p^2>=\frac{e^2}{\mathrm{}c^2}𝑑qq𝑑\omega \frac{c^2}{\omega ^2}\frac{\omega coth\frac{\omega }{2T}}{\sigma _2}.`$ We see that the integral over frequency is diverging for $`\omega <T`$ as $`\frac{d\omega }{\omega ^2}`$. This integral can be cut off again on $`1/\tau _{hop}`$, therefore $`<\delta p^2>q_c^2\frac{\sigma _Q}{\sigma _2}(\tau _{hop}kT)`$ (at any finite temperature). This estimate gives a result equivalent to the earlier calculation. This indicates again that the in-plane momentum is completely randomized after the process of tunneling. All these arguments validate theories of c-axis transport in the normal state assuming non-conservation of in-plane momentum during tunneling. The fact of the non-conservation of momentum $`k_{}`$ in the normal state (if $`\tau _{hop}\tau _{ab}`$) is quite general, a sufficient condition as can be seen from the above discussion is the ohmic density of the noise $`<\delta V_r(\omega )^2>\omega `$ for $`\omega 0`$. It is important to stress that the “diffusivity” of the tunneling is due to specific form of the spectrum of the voltage fluctuations (and not due to short links or impurities!), for instance, it may not be true in the superconducting state. It is important to point out that the question of the in-plane momentum conservation for c-axis tunneling can be examined experimentally supporting or disproving the above arguments.
When the tunneling is diffusive, the “detuning fluctuations” are simply the local voltage fluctuations
$`<\delta V_\omega ^2>={\displaystyle d^2q<\delta V_{q,\omega }^2>}\alpha \omega ,`$ (18)
where $`\alpha `$ is the microscopic parameter describing the ohmic density of the noise. As can be seen below this microscopic parameter $`\alpha `$ is sufficient to describe all dc and ac dependencies of the c-axis conductivity. When the two planes are widely separated and isolated (the situation possibly realized in $`Bi2201`$), the noise spectra in each plane are uncorrelated, so that the interplane noise spectrum is just the sum of the noise spectra in each plane. In such a case, the coefficient $`\alpha `$ should be determined only by the properties of the copper-oxygen plane. If the planes are moved closer together, so that the inter-plane Coulomb interaction become relevant, the intensity of the interplane noise (assuming no interplane hopping) increases due to the presence of the acoustic (out-of-phase) plasmon in this bi-layer structure. Unfortunately, realistically the noise between planes can be suppressed and correlated at low frequencies because of interplane hopping and become dependent on the interlayer structure thus implying different values of $`\alpha `$ for different cuprate materials. It is known experimentally that the detailed temperature dependence of the c-axis resistivity is very sensitive to several factors (sample preparation, interlayer structure and doping). The role of the interlayer structure (different intercalating atoms, chains and additional layers present in some compounds) and the structure of the tunneling matrix element $`t_{}(k_x,k_y)`$ (e.g. in-plane anisotropy) is not clear, it makes the question of the interlayer noise and the c-axis transport properties even more complex.
We have examined several models describing the voltage noise to estimate the microscopic parameter $`\alpha `$ (see the appendix for the discussion of this question, also ). The goal of such exercise is to verify a crude consistency between the estimate of $`\alpha `$ from microscopic noise and the parameter $`\alpha `$ required to describe the c-axis transport dependencies. The difficulty (not surprising since these calculations assume Fermi-liquid or diffusive spectra of density fluctuations) common to all of the calculations (see the appendix) is that the “microscopic” value of $`\alpha `$ is significantly smaller than the value of order of $`1`$ necessary to explain the c-axis transport properties. For our approach to be valid a large ($`\alpha 1`$) density of “detuning fluctuations” (not present in RPA or Fermi liquid pictures) is vital. An alternative approach is to extract the charge fluctuation noise directly from the optical reflectivity measurements. The most dramatic difference between good metals and cuprates seen in Raman and optical measurements of the in-plane dielectric constant $`ϵ_{ab}(q,\omega )`$ is that the low-frequency noise for cuprates is ohmic (linear) even for $`q0`$ . If we write $`Im(1/ϵ_{ab}(q,\omega ))=\gamma \omega `$ for $`q0`$, then from the experimental data $`\gamma 0.2eV^1`$. Since $`<\delta V_\omega ^2>=\frac{dqq}{(2\pi )^2}V_2(q)Im(1/ϵ_{k,\omega })\frac{e^2q_c}{2\pi }\gamma \omega `$, we get $`\alpha \frac{e^2q_c}{2\pi }\gamma `$. If we take the upper cut-off wave vector $`q_c\frac{2\pi }{a}`$ ($`a`$ is a in-plane lattice constant), we get $`\alpha \frac{e^2}{a}\gamma 1.2`$. It shows that $`\alpha `$ may be of order of $`1`$, exactly what is required to explain the c-axis transport properties. We now use the “local” approximation justified above to calculate the dc and ac c-axis conductivity.
First of all, we calculate the local tunneling probability $`P(r=0,t)`$ (or rather its Fourier transform $`P(ϵ,kk^{})`$). The quantity $`P(ϵ,kk^{})`$ can be interpreted as a probability to exchange an energy $`ϵ`$ and momentum $`(kk^{})`$ with fluctuating fields. Since, as argued above, the tunneling probability is strongly peaked at $`r=0`$ (if $`\tau _{hop}\tau _{ab}`$), the Fourier transform in momentum space $`P(ϵ,kk^{})`$ is essentially independent of $`(kk^{})`$, therefore we omit the index $`(kk^{})`$ below. For the small values of the coefficient $`\alpha `$ ($`\alpha 1`$) describing the spectral density of the local voltage noise $`<\delta V_\omega ^2>=\alpha \omega `$, the tunneling probability can be calculated analytically. The function $`R_0(\delta t)`$ in the exponent of the $`P(\delta r=0,\delta t)`$ is
$`R_0(\delta t)={\displaystyle \frac{d\omega }{\omega ^2}<\delta V_\omega ^2>(1cos(\omega \delta t))coth(\frac{\omega }{2kT})}=`$ (19)
$`={\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{d\omega }{\omega ^2}}\alpha \omega (1cos(\omega \delta t))coth({\displaystyle \frac{\omega }{2kT}})`$ (20)
$`{\displaystyle \frac{\delta t}{\tau _\varphi }}+2\alpha ln(\omega _c/kT),`$ (21)
where $`\tau _\varphi =1/2\pi \alpha kT`$ and $`\omega _c`$ is a high-frequency cutoff which we associate with the inverse of the transversal time $`\mathrm{}/\tau _{tr}1eV`$. After a Fourier transform, it gives the probability of tunneling $`P(ϵ)=S_{eff}(\frac{kT}{\omega _c})^{2\alpha }\frac{2\pi \alpha kT}{ϵ^2+(2\pi \alpha kT)^2}`$, where $`S_{eff}`$ is the effective area of tunneling discussed above. For $`\alpha `$ of order of 1, the tunneling probability is strongly suppressed and weakly depends on $`ϵ`$ (for $`ϵkT`$). It cannot be calculated explicitly analytically, but at small $`ϵ`$ can be approximated as a function independent of $`ϵ`$: $`P(ϵ)S_{eff}(\frac{kT}{\omega _c})^{2\alpha }\frac{1}{2\pi \alpha kT}`$.
The optical conductivity $`\sigma _c(\omega ,kT)`$. The tunneling conductance $`\sigma _c(V)`$. Equation 5 can be further transformed assuming the diffusive tunneling probability. In this case, we can integrate over the in-plane momenta ($`k,k^{}`$) separately. The next transformation is due to the detailed balance condition (see Ref.), eventually we can write the expression for the dc conductivity $`\sigma _c(V)`$ as a function of the applied voltage $`V`$:
$`\sigma _c(V)={\displaystyle \frac{et_{}^2d\nu _{2D}^2}{\mathrm{}}}{\displaystyle \frac{1e^{\beta eV}}{V}}{\displaystyle _{\mathrm{}}^+\mathrm{}}𝑑ϵ{\displaystyle \frac{ϵP(eVϵ)}{1e^{\beta ϵ}}}.`$ (22)
From Eqn. 22 in the limit of small voltage ($`VkT`$) (for all $`\alpha `$ as long as the tunneling is diffusive) appropriate for the dc measurements, we calculate the temperature dependence of the c-axis conductivity
$`\sigma _c(T)={\displaystyle \frac{e^2}{\mathrm{}}}t_{}^2d\nu _{2D}^2S_{eff}\left({\displaystyle \frac{kT}{\omega _c}}\right)^{2\alpha }{\displaystyle \frac{1}{2\pi \alpha }}.`$ (23)
This result describes the c-axis resistivity $`\rho _c(T)`$ either constant or diverging at low temperatures found experimentally in several compounds ($`Bi`$-family, $`Hg1201`$, $`Tl2212`$,$`Tl1212`$ and slightly underdoped $`La214`$).
In order to calculate the optical conductivity $`\sigma _c(\omega )`$, we make the following observation. If we apply the external dc voltage, the tunneling probability acquires an additional phase factor $`P(tt^{})exp(\frac{i}{\mathrm{}}_t^{}^tV𝑑\tau )`$ (in Eqn.5, it gives a shift in the energy difference (after a Fourier transform) $`P(E+eVE^{})`$). Thus schematically, the conductance is $`\sigma _c(V)\frac{1}{V}𝑑t^{}e^{\frac{i}{\mathrm{}}V(tt^{})}(\mathrm{})`$. For the ac voltage,
$`P(tt^{})exp({\displaystyle \frac{i}{\mathrm{}}}{\displaystyle _t^{}^t}Ve^{i\omega \tau }𝑑\tau )=`$ (24)
$`=exp({\displaystyle \frac{i}{\mathrm{}}}V{\displaystyle \frac{e^{i\omega t}e^{i\omega t^{}}}{i\omega }})`$ (25)
$`1+{\displaystyle \frac{V}{\mathrm{}}}{\displaystyle \frac{e^{i\omega t}}{\omega }}(1e^{i\omega (tt^{})}).`$ (26)
In the linear response (the optical conductivity is a linear response) and separating a corresponding harmonic of the current proportional to $`e^{i\omega t}`$, we conclude that the dependence of the conductivity $`\sigma _c(\omega )`$ on the frequency $`\omega `$ is equivalent to the dependence on the voltage $`V`$, such that $`\sigma _c(\omega )=\sigma _c(V\omega )`$ (if the tunneling is incoherent and diffusive). Namely,
$`\sigma _c(\omega )={\displaystyle \frac{et_{}^2d\nu _{2D}^2}{\mathrm{}}}{\displaystyle \frac{1e^{\beta \omega }}{\omega }}{\displaystyle _{\mathrm{}}^+\mathrm{}}𝑑ϵ{\displaystyle \frac{ϵP(\omega ϵ)}{1e^{\beta ϵ}}}.`$ (27)
A ubiquitous nearly flat optical response for $`\sigma _c(\omega )`$ is observed experimentally (if measurements exist) in these compounds. Indeed, a qualitative and numerical analysis of the frequency dependence $`\sigma _c(\omega )`$ of Eqn.27 (or equivalently dependence on the voltage) indicates a very weak dependence on frequency.
It is interesting that under the conditions of incoherent tunneling ohmic $`IV`$ curves (constant conductance $`\sigma _c(V)`$ as a function of voltage) correspond to a flat optical conductivity $`\sigma _c(\omega )`$. The correspondence $`\sigma _c(\omega )=\sigma _c(V\omega )`$ can be directly checked experimentally by comparing the tunneling conductance $`\sigma _c(V)=I(V)/V`$ and the frequency dependence of $`\sigma _c(\omega )`$ from the optical reflectivity measurements. The graphs of $`\sigma _c(V)`$ and $`\sigma _c(\omega )`$ can be taken from experimental papers . It appears to be roughly true for optimally doped and underdoped compounds in the normal state.
Discussion of experiments. The above results are insufficient to describe all experimental data for different dopings. In optimally doped $`La214`$ and $`YBaCuO`$ and other overdoped cuprates the c-axis resistivity has a linear temperature dependence with a large “residual” value (intercept at $`T=0K`$) . In these compounds the hopping time becomes comparable to the in-plane scattering time. Thus we can assume that specular tunneling becomes possible. In this crossover situation between two limiting pictures of the tunneling between planes and the anisotropic band along c-axis, we can think about two channels of conduction. One channel is diffusive, while another one is specular. The tunneling probability is the sum of probabilities to tunnel without and with the conservation of the in-plane momentum. In this case, the total c-axis conductivity is the sum of conductivities in each channel. If the fermi surface has very anisotropic properties, electrons from one part of the fermi surface can tunnel specularly, while electrons from other parts tunnel diffusively. In this scenario, the second channel of conduction (conserving $`k_{}`$) can be due to the diagonal parts of the fermi-surface in the normal state. It appears from photoemission experiments that the quasiparticles along diagonals of the Brillouin zone ($`k_x=\pm k_y`$) have longer life-times $`\tau _{ab,diag}`$ (which should be compared with the hopping time). These quasiparticles can tunnel then with conserved momentum. It is natural to suggest that for overdoped cuprates the c-axis transport is dominated by incoherent, but specular channel, while for underdoped cuprates the diffusive channel is only present.
We can calculate the conductivity of a specular channel, if we assume that the time dependence of the specular tunneling probability is $`P(\delta t)exp(\frac{\delta t}{\tau _\varphi })`$ with $`\tau _\varphi `$ calculated for a weak detuning, that is $`\tau _\varphi =1/2\pi \alpha kT`$. Thus we substitute the tunneling probability of the form $`P_2(ϵϵ^{},kk^{})=\delta (kk^{})\frac{2\pi \alpha kT}{(ϵϵ^{})^2+(2\pi \alpha kT)^2}`$ (or any form $`P_2(ϵϵ^{},kk^{})\delta (kk^{})\frac{1}{kT}f(\frac{ϵϵ^{}}{kT})`$) to the Eqn. 5, we get the contribution to the c-axis conductivity
$`\sigma _{2,c}(T)={\displaystyle \frac{e^2}{\mathrm{}}}t_{}^2d\nu _{2D}{\displaystyle \frac{A}{kT}},`$ (28)
where $`A`$ is a numerical coefficient. This result is a well-known result for incoherent tunneling with conservation of the in-plane momentum . It is suggested to explain the linear temperature dependence of c-axis resistivity observed in some compounds . Another important consideration is that the band calculations predict a significant angular dependence of $`t_{}`$ in some families of cuprates. Due to this reason, the contribution from specular channel (from diagonal parts of the fermi-surface) can be reduced.
Yet another complexity of cuprates with multiple planes per unit cell is the question of intra-unit cell and inter-cell conduction. In this case, the resistance associated with hopping between planes of the unit cell (intra-cell resistance) and the resistance associated with hopping between different unit cells (inter-cell resistance) should be discussed. The total resistance is certainly a sum of intra- and inter-unit cell resistances. It may not be correct to assume (as is frequently done in the literature) that the intra-cell resistance is negligible; a systematic experimental investigation of this question is necessary. We hope to discuss this question elsewhere.
In conclusion, a picture of c-axis interlayer (and inter-cell) tunneling strongly suppressed by voltage fluctuations is proposed. This approach can provide a consistent understanding of observed temperature and frequency dependencies of c-axis conductivity in the normal state.
This work was supported by the National Science Foundation through the Science and Technology Center for Superconductivity (grant no. DMR-91-20000) and through grant no. DMR-99-86199. We thank A. Shnirman and L.-Y. Shieh for useful discussions.
Appendix.
In this appendix, we demonstrate an example of a calculation of the parameter $`\alpha `$ describing the low frequency voltage fluctuations based on a particular microscopic model of the density fluctuation spectrum. At the end of the appendix, we list several results calculated for various other microscopic models.
If the two planes are widely separated, then the voltage fluctuations can be treated independently in each plane. In this case, the interplane noise is just twice the in-plane local Johnson-Nyquist noise in each plane. For a high-density electron gas in a hydrodynamic approximation ($`\omega \tau 1`$ and $`ql1`$), the charge density-density susceptibility of the two-dimensional electron gas can be written as
$`\chi (k,\omega )={\displaystyle \frac{s^2k_{TF}k^2}{2\pi e^2}}{\displaystyle \frac{1}{\omega (\omega +i/\tau )s^2k^2s^2k_{TF}k}},`$ (29)
where $`s^2=v_F^2/2`$, $`k_{TF}=\frac{2\pi n_0e^2}{ms^2}`$ is the Thomas-Fermi wave number, and $`\tau `$ is a phenomenological relaxation time. The spectral density of the voltage fluctuations is
$`<V_{k,\omega }^2>=V_k^2Im\chi (k,\omega ),`$ (30)
where $`V_k=\frac{2\pi e^2}{k}`$ is the two-dimensional Coulomb interaction. Eventually, we need to calculate the partial frequency-dependent spectral density of the voltage noise $`<V_\omega ^2>=\frac{dkk}{2\pi }<V_{k,\omega }^2>`$. The calculation gives
$`<V_\omega ^2>={\displaystyle \frac{dkk}{2\pi }\frac{(2\pi e^2)^2}{k^2}\frac{s^2k_{TF}k^2}{2\pi e^2}}`$ (31)
$`{\displaystyle \frac{\omega (1/\tau )}{(\omega /\tau )^2+(s^2k^2+s^2k_{TF}k)^2}}=`$ (32)
$`{\displaystyle \frac{e^2}{s^2k_{TF}\tau }}\omega ln({\displaystyle \frac{s^2k_{TF}^2\tau }{\omega _c}}),`$ (33)
where $`\omega _c`$ is an infra-red cut-off frequency. It implies to the accuracy of the value of the logarithm that the parameter $`\alpha `$ is
$`\alpha {\displaystyle \frac{e^2}{s^2k_{TF}\tau }}{\displaystyle \frac{1}{\mathrm{}}}={\displaystyle \frac{\mathrm{}}{2ϵ_F\tau }}.`$ (34)
We can rewrite this expression in the following form:
$`\alpha ={\displaystyle \frac{1}{\pi }}{\displaystyle \frac{\sigma _Q}{\sigma _{2D}}},`$ (35)
where $`\sigma _Q=e^2/\mathrm{}`$, $`\sigma _{2D}=e^2\nu _{2D}D=\frac{2}{\pi }\frac{e^2(ϵ_F\tau )}{\mathrm{}^2}`$ is the two-dimensional conductivity (it can be shown to have a Drude frequency dependence). The above results may be used to estimate the value of the parameter $`\alpha `$ in the cuprates. Two difficulties can be seen from such literal application of the above model. First of all, the value of $`\alpha `$ is significantly smaller than one. Second, and more importantly, the parameter $`\alpha `$ appears to be temperature dependent. It should be realized that the spectrum of charge fluctuations in the cuprates is much more complex and not represented correctly by the above simple model. We investigated several other simple microscopic models in order to get further insight into this question. At this moment, it seems more reasonable to extract the charge fluctuation spectrum directly from experiment as shown in the text of this paper.
The results for several other microscopic models are summarized below. If the spectrum of charge fluctuations is given by the weakly damped acoustic two-dimensional plasmon, then the parameter $`\alpha `$ is $`\alpha \frac{\sigma _Q^2}{s^2(1+k_{TF}d)}`$, where $`d`$ is the inter-plane distance. Another calculation taking the voltage noise due to electron-hole pairs of the two-dimensional Fermi liquids gives $`\alpha 1/4\pi `$. If the interplane Coulomb interaction is taken into account in the RPA approximation for the same calculation (el-hole pairs), $`\alpha =1/(2\pi k_Fd)`$. We hope to present the details of calculations and expanded arguments elsewhere.
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# Quantum nucleation in ferromagnets with tetragonal and hexagonal symmetries
## I. Introduction
The tunneling of macroscopic object, known as Macroscopic Quantum Tunneling (MQT), is one of the most fascinating phenomena in condensed matter physics. In the last decade, the problem of quantum tunneling of magnetization in nanometer-scale magnets has attracted a great deal of theoretical and experimental interest. MQT in magnetic systems are interesting from a fundamental point of view as they can extend our understanding of the transition from quantum to classical behavior. On the other hand, these phenomena are important to the reliability of small magnetic units in memory devices and the designing of quantum computers in the future. And the measurement of magnetic MQT quantities such as the tunneling rates could provide independent information about microscopic parameters such as the magnetocrystalline anisotropies and the exchange constants. All this makes magnetic quantum tunneling an exciting area for theoretical research and a challenging experimental problem.
The problem of quantum nucleation of a stable phase from a metastable one in ferromagnetic films is an interesting fundamental problem which allows direct comparison between theory and experiment. Consider a ferromagnetic film with its plane perpendicular to the easy axis determined by the magnetocrystalline anisotropy energy depending on the crystal symmetry. A magnetic field H is applied in a direction between perpendicular and opposite to the initial easy direction of the magnetization M, which favors the reversal of the magnetization. The reversal occurs via the nucleation of a critical bubble, which then the nucleus does not collapse, but grows unrestrictedly in volume. If the temperature is sufficiently high, the nucleation of a bubble is a thermal overbarrier process, and the rate of thermal nucleation follows the Arrhenius law $`\mathrm{\Gamma }_T\mathrm{exp}\left(U/k_BT\right)`$, with $`k_B`$ being the Boltzmann constant and $`U`$ being the height of energy barrier. In the limit of $`T0`$, the nucleation is purely quantum-mechanical and the rate goes as $`\mathrm{\Gamma }_Q\mathrm{exp}\left(𝒮_{cl}/\mathrm{}\right)`$, with $`𝒮_{cl}`$ being the classical action or the WKB exponent which is independent of temperature. Because of the exponential dependence of the thermal rate on $`T`$, the temperature $`T_c`$ characterizing the crossover from quantum to thermal regime can be estimated as $`k_BT_c=\mathrm{}U/𝒮_{cl}`$.
The problem of quantum nucleation was studied by Privorotskii who estimated the exponent in the rate of quantum nucleation based on the dimensional analysis. Chudnovsky and Gunther studied the quantum nucleation of a thin ferromagnetic film in the presence of an external magnetic field along the opposite direction to the easy axis at zero temperature by applying the instanton method in the spin-coherent-state path-integral representation. Ferrera and Chudnovsky extended the quantum nucleation to a finite temperature. Kim studied the quantum nucleation in a thin ferromagnetic film placed in a magnetic field at an arbitrary angle.
It is noted that the previous results of quantum nucleation were obtained for ferromagnetic sample with the simplest form of the magnetocrystalline anisotropy energy such as the biaxial symmetry, and the model considered in Refs. 4 and 5 was confined to the condition that the magnetic field be applied along the opposite direction to the easy axis. The purpose of this paper is to extend the previous results of quantum nucleation in ferromagnetic system with simple biaxial symmetry to that of system with a more general symmetry, such as tetragonal and hexagonal symmetry. The generic quantum nucleation problem, however, and the easiest to implement in practice, is that of ferromagnets with a general structure of magnetocrystalline anisotropy in a magnetic field applied at a some angle $`\theta _H`$ to the anisotropy axis. This problem does not possess any symmetry and for that reason is more difficult mathematically. It is worth pursuing, however, because of its significance for experiments. In this paper the magnetic field is applied in an arbitrary direction between perpendicular and opposite to the initial easy axis ($`\widehat{z}`$ axis). Our interest in studying quantum nucleation of magnetic bubbles with a more general structure of magnetocrystalline anisotropy in an arbitrarily directed magnetic field is stimulated by the fact that the corresponding experiment would be most easy to perform and to interpret. Within the instanton approach, we present the numerical results for the WKB exponent in quantum nucleation of a thin ferromagnetic film with the magnetic field applied in a range of angles $`\pi /2<\theta _H<\pi `$, where $`\theta _H`$ is the angle between the initial easy axis ($`\widehat{z}`$ axis) and the field. We also discuss the $`\theta _H`$ dependence of the crossover temperature $`T_c`$ from purely quantum nucleation to thermally assisted processes. Our results show that the distinct angular dependence, together with the dependence of the WKB exponent on the strength of the external magnetic field, may provide an independent experimental test for quantum nucleation in a ferromagnetic film. Quantum nucleation (the description involves space-time instantons), being a field theory problem, is more difficult than tunneling of magnetization in single-domain particles, both at the conceptual and at the technical level. Therefore, this paper provides a nontrivial generalization of uniform rotation of magnetization vector (homogeneous spin tunneling) in single-domain magnets to a nonuniform rotation of magnetization in bulk magnets with a more general structure of magnetocrystalline anisotropy in the presence of a magnetic field at an arbitrary angle. Compared with the tunneling in single-domain particles, a local tunneling event in a bulk magnet can trigger instability on a much greater scale, which leads to really macroscopic consequences. In experiments, it may be easier to monitor single nucleation events in a thin film than to detect the magnetization reversal in a nanometer-scale particle. Therefore, our theoretical results for a general structure of magnetocrystalline anisotropy in an arbitrarily directed field will be more applicable for experimental tests of quantum nucleation. Besides the importance from the fundamental point of view, processes of quantum nucleation and collapse of magnetic bubbles are potentially important for quantum limitations on the density and long-term reliability of the data storage in magnetic memory devices and designing of quantum computer.
This paper is structured in the following way. In Sec. II, we review briefly some basic ideas of quantum nucleation of magnetization in ferromagnets. In Secs. III and IV, we study quantum nucleation of magnetization in ferromagnets with tetragonal and hexagonal symmetry in an external magnetic field applied in the $`ZX`$ plane with a range of angles $`\pi /2\theta _H<\pi `$. The conclusions and discussions are presented in Sec. V.
## II. The physical model
For a spin tunneling problem, the rate of magnetization reversal by quantum tunneling is determined by the imaginary-time transition amplitude from an initial state $`|i`$ to a final state $`|f`$ as
$$U_{fi}=f\left|e^{HT}\right|i=𝒟\left\{𝐌(𝐫,\tau )\right\}\mathrm{exp}\left(𝒮_E/\mathrm{}\right),$$
(1)
where $`𝒮_E`$ is the Euclidean action which includes the Euclidean Lagrangian density $`_E`$ as
$$𝒮_E=𝑑\tau d^3𝐫_E.$$
(2)
For ferromagnets at sufficiently low temperature, all the spins are locked together by the strong exchange interaction, and therefore only the orientation of magnetization $`𝐌(𝐫,\tau )`$ can change but not its absolute value. For that reason the field $`𝐌(𝐫,\tau )`$ is equivalent to the fields $`\theta (𝐫,\tau )`$ and $`\varphi (𝐫,\tau )`$, which are spherical coordinates of $`𝐌`$. In this case the measure of the path integral $`𝒟\left\{𝐌(𝐫,\tau )\right\}`$ in Eq. (1) is equivalent to
$$𝒟\left\{\theta (𝐫,\tau )\right\}𝒟\left\{\varphi (𝐫,\tau )\right\}=\underset{\epsilon 0}{lim}\underset{k=1}{\overset{N}{}}\left(\frac{2S+1}{4\pi }\right)\mathrm{sin}\theta _kd\theta _kd\varphi _k,$$
(3)
where $`\epsilon =\mathrm{max}\left(\tau _{k+1}\tau _k\right)`$ and $`S=M_0/\mathrm{}\gamma `$ is the total spin of ferromagnet. Here $`\gamma `$ is the gyromagnetic ratio and $`M_0`$ is the magnitude of magnetization.
In the spin-coherent-state representation, the magnetic Lagrangian is given by
$$_E=i\frac{M_0}{\gamma }\left(\frac{d\varphi (𝐫,\tau )}{d\tau }\right)\left[1\mathrm{cos}\theta (𝐫,\tau )\right]+E(\theta ,\varphi ).$$
(4)
The first term in Eq. (4) is a total imaginary-time derivative, which has no effect on the classical equations of motion, but it is crucial for the spin-parity effects. However, for the closed instanton trajectory described in this paper (as shown in the following), this time derivative gives a zero contribution to the path integral, and therefore can be omitted.
The energy density in Eq. (4) is
$$E(\theta ,\varphi )=E_a(\theta ,\varphi )+E_{ex}(\theta ,\varphi ),$$
(5)
where $`E_a`$ includes the magnetocrystalline anisotropy energy and the energy due to the external magnetic field, and $`E_{ex}`$ is the exchange energy
$$E_{ex}=\frac{\alpha }{2}\left(_iM_j\right)^2=\frac{\alpha }{2}M_0^2\left[\left(\theta \right)^2+\mathrm{sin}^2\theta \left(\varphi \right)^2\right],$$
(6)
where $`\alpha `$ is the exchange stiffness.. The magnetocrystalline anisotropy energy for tetragonal and hexagonal symmetry is shown in Sec. III and IV, respectively. In the semiclassical limit, the rate of quantum nucleation, with an exponential accuracy, is given by
$$\mathrm{\Gamma }_Q\mathrm{exp}\left[𝒮_E^{\mathrm{min}}/\mathrm{}\right],$$
(7)
where $`𝒮_E^{\mathrm{min}}`$ is obtained along the trajectory that minimizes the Euclidean action $`𝒮_E`$.
## III. Ferromagnets with tetragonal symmetry
In this section, we study the quantum nucleation of magnetization in ferromagnets with tetragonal symmetry in the presence of a magnetic field at arbitrary angles in the $`ZX`$ plane, which has the following magnetocrystalline anisotropy energy
$$E_a(\theta ,\varphi )=K_1\mathrm{sin}^2\theta +K_2\mathrm{sin}^4\theta K_2^{}\mathrm{sin}^4\theta \mathrm{cos}\left(4\varphi \right)M_0H_x\mathrm{sin}\theta \mathrm{cos}\varphi M_0H_z\mathrm{cos}\theta ,$$
(8)
where $`K_1`$, $`K_2`$ and $`K_2^{}`$ are the magnetic anisotropy coefficients, and $`K_1>0`$. In the absence of the magnetic field, the easy axes of this system are $`\pm \widehat{z}`$ for $`K_1>0`$. And the field is applied in the $`ZX`$ plane at $`\pi /2<\theta _H<\pi `$. Then the total energy is given by
$`E[\theta (𝐫,\tau ),\varphi (𝐫,\tau )]`$ $`=`$ $`K_1\mathrm{sin}^2\theta +K_2\mathrm{sin}^4\theta K_2^{}\mathrm{sin}^4\theta \mathrm{cos}\left(4\varphi \right)+{\displaystyle \frac{\alpha }{2}}M_0^2\left[\left(\theta \right)^2+\mathrm{sin}^2\theta \left(\varphi \right)^2\right]`$ ()
$`M_0H_x\mathrm{sin}\theta \mathrm{cos}\varphi M_0H_z\mathrm{cos}\theta +E_0,`$
where $`E_0`$ is a constant which makes $`E(\theta ,\varphi )`$ zero at the initial state. By applying the similar method in Ref. 15, we can perform a Gaussian integration over the variable $`\varphi `$ in the path integral and reduce the system to that with only one variable $`\delta `$ (as shown in the following). Then it is possible to perform the rest of the calculation by using the instanton method. This method simplifies the problem tremendously, compared to the problem where the action depended on $`\theta \left(\tau \right)`$ and $`\varphi \left(\tau \right)`$, though a complete mathematical equivalence to the initial problem is preserved.
By introducing the dimensionless parameters as
$$\overline{K}_2=K_2/2K_1,\overline{K}_2^{}=K_2^{}/2K_1,\overline{H}_x=H_x/H_0,\overline{H}_z=H_z/H_0,$$
(10)
Eq. (9) can be rewritten as
$`\overline{E}(\theta ,\varphi )`$ $`=`$ $`{\displaystyle \frac{1}{2}}\mathrm{sin}^2\theta +\overline{K}_2\mathrm{sin}^4\theta \overline{K}_2^{}\mathrm{sin}^4\theta \mathrm{cos}\left(4\varphi \right)+{\displaystyle \frac{\alpha M_0^2}{4K_1}}\left[\left(\theta \right)^2+\mathrm{sin}^2\theta \left(\varphi \right)^2\right]`$ ()
$`\overline{H}_x\mathrm{sin}\theta \mathrm{cos}\varphi \overline{H}_z\mathrm{cos}\theta +\overline{E}_0,`$
where $`E(\theta ,\varphi )=2K_1\overline{E}(\theta ,\varphi )`$, and $`H_0=2K_1/M_0`$. At finite magnetic field, the plane given by $`\varphi =0`$ is the easy plane, on which $`\overline{E}_a(\theta ,\varphi )`$ reduces to
$$\overline{E}_a\left(\theta ,\varphi =0\right)=\frac{1}{2}\mathrm{sin}^2\theta +\left(\overline{K}_2\overline{K}_2^{}\right)\mathrm{sin}^4\theta \overline{H}\mathrm{cos}\left(\theta \theta _H\right)+\overline{E}_0.$$
(12)
The initial angle $`\theta _0`$ is determined by $`\left[d\overline{E}_a(\theta ,0)/d\theta \right]_{\theta =\theta _0}=0`$, and the critical angle $`\theta _c`$ and the dimensionless critical field $`\overline{H}_c`$ are determined by both $`\left[d\overline{E}_a(\theta ,0)/d\theta \right]_{\theta =\theta _c,\overline{H}=\overline{H}_c}=0`$ and $`\left[d^2\overline{E}_a(\theta ,0)/d\theta ^2\right]_{\theta =\theta _c,\overline{H}=\overline{H}_c}=0`$, which leads to
$`{\displaystyle \frac{1}{2}}\mathrm{sin}\left(2\theta _0\right)+\overline{H}\mathrm{sin}\left(\theta _0\theta _H\right)+4\left(\overline{K}_2\overline{K}_2^{}\right)\mathrm{sin}^3\theta _0\mathrm{cos}\theta _0`$ $`=`$ $`0,`$ ()
$`{\displaystyle \frac{1}{2}}\mathrm{sin}\left(2\theta _c\right)+\overline{H}_c\mathrm{sin}\left(\theta _c\theta _H\right)+4\left(\overline{K}_2\overline{K}_2^{}\right)\mathrm{sin}^3\theta _c\mathrm{cos}\theta _c`$ $`=`$ $`0,`$ ()
$`\mathrm{cos}\left(2\theta _c\right)+\overline{H}_c\mathrm{cos}\left(\theta _c\theta _H\right)+4\left(\overline{K}_2\overline{K}_2^{}\right)\left(3\mathrm{sin}^2\theta _c\mathrm{cos}^2\theta _c\mathrm{sin}^4\theta _c\right)`$ $`=`$ $`0.`$ ()
Assuming that $`\left|\overline{K}_2\overline{K}_2^{}\right|1`$, we obtain the critical magnetic field and the critical angle as
$`\overline{H}_c`$ $`=`$ $`{\displaystyle \frac{1}{\left[\left(\mathrm{sin}\theta _H\right)^{2/3}+\left|\mathrm{cos}\theta _H\right|^{2/3}\right]^{3/2}}}\left[1+{\displaystyle \frac{4\left(\overline{K}_2\overline{K}_2^{}\right)}{1+\left|\mathrm{cot}\theta _H\right|^{2/3}}}\right],`$ ()
$`\mathrm{sin}\theta _c`$ $`=`$ $`{\displaystyle \frac{1}{\left(1+\left|\mathrm{cot}\theta _H\right|^{2/3}\right)^{1/2}}}\left[1+{\displaystyle \frac{8}{3}}\left(\overline{K}_2\overline{K}_2^{}\right){\displaystyle \frac{\left|\mathrm{cot}\theta _H\right|^{2/3}}{1+\left|\mathrm{cot}\theta _H\right|^{2/3}}}\right].`$ ()
In the low barrier limit, i.e., $`ϵ=1\overline{H}/\overline{H}_c0`$, by using Eqs. (13b) and (13c) we obtain the approximate equation for $`\eta \left(\theta _c\theta _0\right)`$ in the order of $`ϵ^{3/2}`$,
$`ϵ\overline{H}_c\mathrm{sin}\left(\theta _c\theta _H\right)+\eta ^2\left[{\displaystyle \frac{3}{2}}\overline{H}_c\mathrm{sin}\left(\theta _c\theta _H\right)+3\left(\overline{K}_2\overline{K}_2^{}\right)\mathrm{sin}\left(4\theta _c\right)\right]`$ (13)
$`+\eta \left\{ϵ\overline{H}_c\mathrm{cos}\left(\theta _c\theta _H\right)\eta ^2\left[{\displaystyle \frac{1}{2}}\overline{H}_c\mathrm{cos}\left(\theta _c\theta _H\right)+4\left(\overline{K}_2\overline{K}_2^{}\right)\mathrm{cos}\left(4\theta _c\right)\right]\right\}=0.`$ ()
Introducing $`\delta \theta \theta _0`$ ($`\left|\delta \right|1`$ in the small $`ϵ`$ limit), we derive the energy $`\overline{E}(\theta ,\varphi )`$ as
$`\overline{E}(\delta ,\varphi )`$ $`=`$ $`\overline{K}_2^{}\left[1\mathrm{cos}\left(4\varphi \right)\right]\mathrm{sin}^4\left(\theta _0+\delta \right)+\overline{H}_x\left(1\mathrm{cos}\varphi \right)\mathrm{sin}\left(\theta _0+\delta \right)`$ ()
$`+{\displaystyle \frac{\alpha M_0^2}{4K_1}}\left[\left(\theta \right)^2+\mathrm{sin}^2\theta \left(\varphi \right)^2\right]+\overline{E}_1\left(\delta \right),`$
where $`\overline{E}_1\left(\delta \right)`$ is a function of only $`\delta `$ given by
$`\overline{E}_1\left(\delta \right)`$ $`=`$ $`\left[{\displaystyle \frac{1}{2}}\overline{H}_c\mathrm{sin}\left(\theta _c\theta _H\right)+\left(\overline{K}_2\overline{K}_2^{}\right)\mathrm{sin}\left(4\theta _c\right)\right]\left(\delta ^33\delta ^2\eta \right)`$ ()
$`+\left[{\displaystyle \frac{1}{8}}\overline{H}_c\mathrm{cos}\left(\theta _c\theta _H\right)+\left(\overline{K}_2\overline{K}_2^{}\right)\mathrm{cos}\left(4\theta _c\right)\right]\left(\delta ^44\delta ^3\eta +6\delta ^2\eta ^24\delta ^2ϵ\right)`$
$`+4\left(\overline{K}_2\overline{K}_2^{}\right)ϵ\delta ^2\mathrm{cos}\left(4\theta _c\right).`$
It can be shown that in the region of $`\pi /2<\theta _H<\pi `$, $`0<\theta _c<\pi /2`$, $`\eta `$ and $`\delta `$ are of the order of $`\sqrt{ϵ}`$, the second or third term in Eq. (17) is smaller than the first term in the small $`ϵ`$ limit. It is convenient to use dimensionless variables
$$𝐫^{}=ϵ^{1/4}𝐫/r_0,\tau ^{}=ϵ^{1/4}\omega _0\tau ,\overline{\delta }=\delta /\sqrt{ϵ},$$
(17)
where $`r_0=\sqrt{\frac{\alpha M_0^2}{2K_1}}`$, and $`\omega _0=2\gamma K_1/M_0`$. Then the Euclidean action Eq. (2) for $`\pi /2<\theta _H<\pi `$ becomes
$`𝒮_E[\overline{\delta }(𝐫^{},\tau ^{}),\varphi (𝐫^{},\tau ^{})]`$ $`=`$ $`{\displaystyle \frac{\mathrm{}Sr_0^3}{ϵ}}{\displaystyle }d\tau ^{}d^3𝐫^{}\{iϵ^{1/4}\mathrm{sin}(\theta _0+\sqrt{ϵ}\overline{\delta })\varphi \left({\displaystyle \frac{\overline{\delta }}{\tau ^{}}}\right)`$ ()
$`+2\overline{K}_2^{}\mathrm{sin}^2\left(2\varphi \right)\mathrm{sin}^4\left(\theta _0+\sqrt{ϵ}\overline{\delta }\right)+2\overline{H}_x\mathrm{sin}^2\left({\displaystyle \frac{\varphi }{2}}\right)\mathrm{sin}\left(\theta _0+\sqrt{ϵ}\overline{\delta }\right)`$
$`+{\displaystyle \frac{1}{2}}ϵ^{3/2}\left(^{}\overline{\delta }\right)^2+{\displaystyle \frac{1}{2}}ϵ^{1/2}\mathrm{sin}^2\left(\theta _0+\sqrt{ϵ}\overline{\delta }\right)\left(^{}\varphi \right)^2`$
$`+{\displaystyle \frac{A}{4}}ϵ^{3/2}(\sqrt{6}\overline{\delta }^2\overline{\delta }^3)\},`$
where
$$A=2\frac{\left|\mathrm{cot}\theta _H\right|^{1/3}}{1+\left|\mathrm{cot}\theta _H\right|^{2/3}}\left[1+\frac{4}{3}\left(\overline{K}_2\overline{K}_2^{}\right)\frac{74\left|\mathrm{cot}\theta _H\right|^{2/3}}{1+\left|\mathrm{cot}\theta _H\right|^{2/3}}\right].$$
(21)
In Eq. (19) we have performed the integration by part for the first term and have neglected the total imaginary-time derivative. In can be showed that for $`\pi /2<\theta _H<\pi `$, only small values of $`\varphi `$ contribute to the path integral, so that one can replace $`\mathrm{sin}^2\varphi `$ in Eq. (19) by $`\varphi ^2`$ and neglect the term including $`\left(^{}\varphi \right)^2`$ which is of the order $`ϵ^2`$ while the other terms are of the order $`ϵ^{3/2}`$. Then the Gaussian integration over $`\varphi `$ leads to
$$𝒟\left\{\delta (𝐫^{},\tau ^{})\right\}\mathrm{exp}\left(\frac{1}{\mathrm{}}𝒮_E^{eff}\right),$$
(22)
where the effective action is
$$𝒮_E^{eff}\left[\overline{\delta }(𝐫^{},\tau ^{})\right]=\mathrm{}Sϵ^{1/2}r_0^3𝑑\tau ^{}d^3𝐫^{}\left[\frac{1}{2}M\left(\frac{\overline{\delta }}{\tau ^{}}\right)^2+\frac{1}{2}\left(^{}\overline{\delta }\right)^2+\frac{A}{4}\left(\sqrt{6}\overline{\delta }^2\overline{\delta }^3\right)\right].$$
(23)
The effect mass in Eq. (22) is found to be
$$M=\frac{\left(1+\left|\mathrm{cot}\theta _H\right|^{2/3}\right)\left[1+\frac{8}{3}\left(\overline{K}_2\overline{K}_2^{}\right)\frac{\left|\mathrm{cot}\theta _H\right|^{2/3}}{1+\left|\mathrm{cot}\theta _H\right|^{2/3}}\right]}{1ϵ+16\overline{K}_2^{}+4\left(\overline{K}_2\overline{K}_2^{}\right)\frac{1}{1+\left|\mathrm{cot}\theta _H\right|^{2/3}}+128\overline{K}_2^{}\left(\overline{K}_2\overline{K}_2^{}\right)\frac{\left|\mathrm{cot}\theta _H\right|^{2/3}}{1+\left|\mathrm{cot}\theta _H\right|^{2/3}}}.$$
(24)
Introducing the variables $`\overline{\tau }=\tau ^{}\sqrt{A/M}`$ and $`\overline{𝐫}=𝐫^{}\sqrt{A}`$, the effective action Eq. (22) is simplified as
$$𝒮_E^{eff}\left[\overline{\delta }(\overline{𝐫},\overline{\tau })\right]=\mathrm{}Sϵ^{1/2}r_0^3\frac{\sqrt{M}}{A}𝑑\overline{\tau }d^3\overline{𝐫}\left[\frac{1}{2}\left(\frac{\overline{\delta }}{\overline{\tau }}\right)^2+\frac{1}{2}\left(\overline{}\overline{\delta }\right)^2+\frac{1}{4}\left(\sqrt{6}\overline{\delta }^2\overline{\delta }^3\right)\right].$$
(25)
For the quantum reversal of magnetization $`𝐌`$ in a small particle of volume $`Vr_0^3`$, $`𝐌`$ is uniform within the particle and $`\overline{\delta }`$ does not depend on the space $`\overline{𝐫}`$, Eq. (24) reduces to
$$𝒮_E^{eff}\left[\overline{\delta }(\overline{𝐫},\overline{\tau })\right]=\mathrm{}Sϵ^{5/4}\sqrt{MA}V𝑑\overline{\tau }\left[\frac{1}{2}\left(\frac{d\overline{\delta }}{d\overline{\tau }}\right)^2+\frac{1}{4}\left(\sqrt{6}\overline{\delta }^2\overline{\delta }^3\right)\right].$$
(26)
The corresponding classical trajectory satisfies the equation of motion
$$\frac{d^2\overline{\delta }}{d\overline{\tau }^2}=\frac{1}{2}\sqrt{6}\overline{\delta }\frac{3}{4}\overline{\delta }^2.$$
(27)
Eq. (26) has the instanton solution
$$\overline{\delta }\left(\overline{\tau }\right)=\frac{\sqrt{6}}{\mathrm{cosh}^2\left(3^{1/4}\times 2^{5/4}\overline{\tau }\right)},$$
(28)
corresponding to the variation of $`\overline{\delta }`$ from $`\overline{\delta }=0`$ at $`\overline{\tau }=\mathrm{}`$, to $`\overline{\delta }=\sqrt{6}`$ at $`\overline{\tau }=0`$, and then back to $`\overline{\delta }=0`$ at $`\overline{\tau }=\mathrm{}`$. Eq. (27) agrees well with the result in Refs. 13 and 15. The associated classical action is found to be
$`𝒮_{cl}`$ $`=`$ $`{\displaystyle \frac{2^{17/4}\times 3^{1/4}}{5}}\mathrm{}Sϵ^{5/4}`$ ()
$`\times {\displaystyle \frac{\left|\mathrm{cot}\theta _H\right|^{1/6}\left[1+\frac{2}{3}\left(\overline{K}_2\overline{K}_2^{}\right)\frac{72\left|\mathrm{cot}\theta _H\right|^{2/3}}{1+\left|\mathrm{cot}\theta _H\right|^{2/3}}\right]}{\sqrt{1ϵ+16\overline{K}_2^{}+4\left(\overline{K}_2\overline{K}_2^{}\right)\frac{1}{1+\left|\mathrm{cot}\theta _H\right|^{2/3}}+128\overline{K}_2^{}\left(\overline{K}_2\overline{K}_2^{}\right)\frac{\left|\mathrm{cot}\theta _H\right|^{2/3}}{1+\left|\mathrm{cot}\theta _H\right|^{2/3}}}}}.`$
Now we turn to the nonuniform problem. In case of a thin film of thickness $`h`$ less than the size $`r_0/ϵ^{1/4}`$ of the critical nucleus and its plane is perpendicular to the initial easy axis, we obtain the action Eq. (24) after performing the integration over the $`\overline{z}`$ variable,
$$𝒮_E^{eff}\left[\overline{\delta }(\overline{𝐫},\overline{\tau })\right]=\mathrm{}Sϵ^{3/4}r_0^2h\sqrt{\frac{M}{A}}𝑑\overline{\tau }d^2\overline{𝐫}\left[\frac{1}{2}\left(\frac{\overline{\delta }}{\overline{\tau }}\right)^2+\frac{1}{2}\left(\overline{}\overline{\delta }\right)^2+\frac{1}{4}\left(\sqrt{6}\overline{\delta }^2\overline{\delta }^3\right)\right].$$
(30)
At zero temperature the classical solution of the effective action Eq. (29) has $`O\left(3\right)`$ symmetry in two spatial plus one imaginary time dimensions. Therefore, the solution $`\overline{\delta }`$ is a function of $`u`$, where $`u=\left(\overline{\rho }^2+\overline{\tau }^2\right)^{1/2}`$, and $`\overline{\rho }=\left(\overline{x}^2+\overline{y}^2\right)^{1/2}`$ is the normalized distance from the $`𝐳`$ axis. Now the effective action Eq. (29) becomes
$$𝒮_E^{eff}\left[\overline{\delta }(\overline{𝐫},\overline{\tau })\right]=4\pi \mathrm{}Sϵ^{3/4}r_0^2h\sqrt{\frac{M}{A}}𝑑uu^2\left[\frac{1}{2}\left(\frac{d\overline{\delta }}{du}\right)^2+\frac{1}{4}\left(\sqrt{6}\overline{\delta }^2\overline{\delta }^3\right)\right].$$
(31)
The corresponding classical trajectory satisfies the following equation of motion
$$\frac{d^2\overline{\delta }}{du^2}+\frac{2}{u}\frac{d\overline{\delta }}{du}=\frac{\sqrt{6}}{2}\overline{\delta }\frac{3}{4}\overline{\delta }^2.$$
(32)
By applying the similar method, the instanton solution of Eq. (31) can be found numerically and is illustrated in Fig. 1. The maximal rotation of $`𝐌`$ is $`\overline{\delta }_{\mathrm{max}}6.8499`$ at $`\overline{\tau }=0`$ and $`\overline{\rho }=0`$. Numerical integration in Eq. (30), using this solution, gives the rate of quantum nucleation for a thin ferromagnetic film as
$`\mathrm{\Gamma }_Q`$ $``$ $`\mathrm{exp}\left(𝒮_E/\mathrm{}\right)`$ (33)
$`=`$ $`\mathrm{exp}\{74.39Sϵ^{3/4}r_0^2h{\displaystyle \frac{1+\left|\mathrm{cot}\theta _H\right|^{2/3}}{\left|\mathrm{cot}\theta _H\right|^{1/6}}}[1{\displaystyle \frac{2}{3}}(\overline{K}_2\overline{K}_2^{}){\displaystyle \frac{76\left|\mathrm{cot}\theta _H\right|^{2/3}}{1+\left|\mathrm{cot}\theta _H\right|^{2/3}}}]`$ ()
$`\times {\displaystyle \frac{1}{\sqrt{1ϵ+16\overline{K}_2^{}+4\left(\overline{K}_2\overline{K}_2^{}\right)\frac{1}{1+\left|\mathrm{cot}\theta _H\right|^{2/3}}+128\overline{K}_2^{}\left(\overline{K}_2\overline{K}_2^{}\right)\frac{\left|\mathrm{cot}\theta _H\right|^{2/3}}{1+\left|\mathrm{cot}\theta _H\right|^{2/3}}}}}\}.`$
At high temperature, the nucleation of $`𝐌`$ is due to thermal activation, and the rate of nucleation follows $`\mathrm{\Gamma }_T\mathrm{exp}\left(W_{\mathrm{min}}/k_BT\right)`$, where $`W_{\mathrm{min}}`$ is the minimal work necessary to produce a nucleus capable of growing. In this case the instanton solution becomes independent of the imaginary-time variable $`\overline{\tau }`$. In order to obtain $`W_{\mathrm{min}}`$, we consider the effective potential of the system
$$U_{eff}=d^3𝐫E=d^3𝐫\left[\frac{\alpha }{2}M_0^2\left(\left(\theta \right)^2+\mathrm{sin}^2\theta \left(\varphi \right)^2\right)+E_a(\theta ,\varphi )\right].$$
(35)
For a cylindrical bubble Eq. (33) reduces to
$$U_{eff}=4\pi K_1hϵr_0^2_0^{\mathrm{}}𝑑\overline{\rho }\overline{\rho }\left[\frac{1}{2}\left(\frac{d\overline{\delta }}{d\overline{\rho }}\right)^2+\frac{1}{4}\left(\sqrt{6}\overline{\delta }^2\overline{\delta }^3\right)\right].$$
(36)
From the saddle point of the functional the shape of the critical nucleus satisfies
$$\frac{d^2\overline{\delta }}{d\overline{\rho }^2}+\frac{1}{\overline{\rho }}\frac{d\overline{\delta }}{d\overline{\rho }}=\frac{\sqrt{6}}{2}\overline{\delta }\frac{3}{4}\overline{\delta }^2.$$
(37)
The solution can be found by numerical method similar to the one in Refs. 4 and 6. Fig. 2 shows the shape of the critical bubble in thermal nucleation, and the maximal size is $`3.906`$ at $`\overline{\rho }=0`$. Using this result, the minimal work corresponding the thermal nucleation is
$$W_{\mathrm{min}}=41.3376K_1hϵr_0^2.$$
(38)
Comparing this with Eq. (32), we obtain the approximate formula for the temperature characterizing the crossover from thermal to quantum nucleation as
$`k_BT_c`$ $``$ $`0.55{\displaystyle \frac{K_1ϵ^{1/4}}{S}}{\displaystyle \frac{\left|\mathrm{cot}\theta _H\right|^{1/6}}{1+\left|\mathrm{cot}\theta _H\right|^{2/3}}}\left[1+{\displaystyle \frac{2}{3}}\left(\overline{K}_2\overline{K}_2^{}\right){\displaystyle \frac{76\left|\mathrm{cot}\theta _H\right|^{2/3}}{1+\left|\mathrm{cot}\theta _H\right|^{2/3}}}\right]`$ ()
$`\times [1ϵ+16\overline{K}_2^{}+4(\overline{K}_2\overline{K}_2^{}){\displaystyle \frac{1}{1+\left|\mathrm{cot}\theta _H\right|^{2/3}}}`$
$`+128\overline{K}_2^{}(\overline{K}_2\overline{K}_2^{}){\displaystyle \frac{\left|\mathrm{cot}\theta _H\right|^{2/3}}{1+\left|\mathrm{cot}\theta _H\right|^{2/3}}}]^{1/2}.`$
To observe the quantum nucleation one needs a large crossover temperature and not too small a nucleation rate. Eq. (37) shows that ferromagnets with large anisotropy, i.e., small ration of $`K_2^{}`$ to $`K_1`$, and small saturated magnetization are preferable for experimental study. In Fig. 3, we plot the $`\theta _H`$ dependence of the crossover temperature $`T_c`$ for typical values of parameters for nanometer-scale ferromagnets: $`K_1=10^7`$ erg/cm<sup>3</sup>, $`\overline{K}_2^{}=0.1`$, $`\overline{K}_2\overline{K}_2^{}=0.01`$, $`M_0=500`$ emu/cm<sup>3</sup>, $`ϵ=0.01`$ in a wide range of angles $`\pi /2<\theta _H<\pi `$. Fig. 3 shows that the maximal value of $`T_c`$ is about 225 mK at $`\theta _H=1.743`$. The maximal value of $`T_c`$ as well as $`\mathrm{\Gamma }_Q`$ is expected to be observed in experiment. The similar $`\theta _H`$ dependence of the crossover temperature $`T_c`$ was first observed in Ref. 15, while the problem considered in Ref. 15 was homogeneous spin tunneling in single-domain particles with uniaxial symmetry.
## IV. Ferromagnets with hexagonal symmetry
In this section, we study the quantum nucleation of magnetization in nanometer-scale ferromagnets with hexagonal symmetry in an external magnetic field at an arbitrary angle in the $`ZX`$ plane. Now the magnetocrystalline anisotropy energy $`E_a(\theta ,\varphi )`$ can be written as
$`E_a(\theta ,\varphi )`$ $`=`$ $`K_1\mathrm{sin}^2\theta +K_2\mathrm{sin}^4\theta +K_3\mathrm{sin}^6\theta K_3^{}\mathrm{sin}^6\theta \mathrm{cos}\left(6\varphi \right)`$ ()
$`M_0H_x\mathrm{sin}\theta \mathrm{cos}\varphi M_0H_z\mathrm{cos}\theta ,`$
where $`K_1`$, $`K_2`$, $`K_3`$, and $`K_3^{}`$ are the magnetic anisotropic coefficients. The easy axes are $`\pm \widehat{z}`$ for $`K_1>0`$. By choosing $`K_3^{}>0`$, we take $`\varphi =0`$ to be the easy plane, at which the anisotropy energy can be written in terms of the dimensionless parameters as
$$\overline{E}_a\left(\theta ,\varphi =0\right)=\frac{1}{2}\mathrm{sin}^2\theta +\overline{K}_2\mathrm{sin}^4\theta +\left(\overline{K}_3\overline{K}_3^{}\right)\mathrm{sin}^6\theta \overline{H}\mathrm{cos}\left(\theta \theta _H\right)+\overline{E}_0,$$
(42)
where $`\overline{K}_3=K_3/2K_1`$ and $`\overline{K}_3^{}=K_3^{}/2K_1`$.
Then the initial angle $`\theta _0`$ is determined by $`\left[d\overline{E}_a(\theta ,0)/d\theta \right]_{\theta =\theta _0}=0`$, and the critical angle $`\theta _c`$ and the dimensionless critical field $`\overline{H}_c`$ by both $`\left[d\overline{E}_a(\theta ,0)/d\theta \right]_{\theta =\theta _c,\overline{H}=\overline{H}_c}=0`$ and $`\left[d^2\overline{E}_a(\theta ,0)/d\theta ^2\right]_{\theta =\theta _c,\overline{H}=\overline{H}_c}=0`$, which leads to
$`{\displaystyle \frac{1}{2}}\mathrm{sin}\left(2\theta _0\right)+\overline{H}\mathrm{sin}\left(\theta _0\theta _H\right)+4\overline{K}_2\mathrm{sin}^3\theta _0\mathrm{cos}\theta _0+6\left(\overline{K}_3\overline{K}_3^{}\right)\mathrm{sin}^5\theta _0\mathrm{cos}\theta _0=0,`$ ()
$`{\displaystyle \frac{1}{2}}\mathrm{sin}\left(2\theta _c\right)+\overline{H}_c\mathrm{sin}\left(\theta _c\theta _H\right)+4\overline{K}_2\mathrm{sin}^3\theta _c\mathrm{cos}\theta _c+6\left(\overline{K}_3\overline{K}_3^{}\right)\mathrm{sin}^5\theta _c\mathrm{cos}\theta _c=0,`$ ()
$`\mathrm{cos}\left(2\theta _c\right)+\overline{H}_c\mathrm{cos}\left(\theta _c\theta _H\right)+4\overline{K}_2\left(3\mathrm{sin}^2\theta _c\mathrm{cos}^2\theta _c\mathrm{sin}^4\theta _c\right)`$ (43)
$`+6\left(\overline{K}_3\overline{K}_3^{}\right)\left(5\mathrm{sin}^4\theta _c\mathrm{cos}^2\theta _c\mathrm{sin}^6\theta _c\right)=0,`$ ()
Under the assumption that $`\left|\overline{K}_2\right|`$, $`\left|\overline{K}_3\overline{K}_3^{}\right|1`$, we obtain the dimensionless critical field $`\overline{H}_c`$ and the critical angle as
$`\overline{H}_c`$ $`=`$ $`{\displaystyle \frac{1}{\left[\left(\mathrm{sin}\theta _H\right)^{2/3}+\left|\mathrm{cos}\theta _H\right|^{2/3}\right]^{3/2}}}\left[1+{\displaystyle \frac{4\overline{K}_2}{1+\left|\mathrm{cot}\theta _H\right|^{2/3}}}+{\displaystyle \frac{6\left(\overline{K}_3\overline{K}_3^{}\right)}{\left(1+\left|\mathrm{cot}\theta _H\right|^{2/3}\right)^2}}\right],`$ ()
$`\mathrm{sin}\theta _c`$ $`=`$ $`{\displaystyle \frac{1}{\left(1+\left|\mathrm{cot}\theta _H\right|^{2/3}\right)^{1/2}}}[1+{\displaystyle \frac{8}{3}}\overline{K}_2{\displaystyle \frac{\left|\mathrm{cot}\theta _H\right|^{2/3}}{1+\left|\mathrm{cot}\theta _H\right|^{2/3}}}`$ ()
$`+8(\overline{K}_3\overline{K}_3^{}){\displaystyle \frac{\left|\mathrm{cot}\theta _H\right|^{2/3}}{\left(1+\left|\mathrm{cot}\theta _H\right|^{2/3}\right)^2}}].`$
In the limit of small $`ϵ=1\overline{H}/\overline{H}_c`$, Eq. (40a) becomes
$`ϵ\overline{H}_c\mathrm{sin}(\theta _c\theta _H)+\eta ^2[(3/2)\overline{H}_c\mathrm{sin}(\theta _c\theta _H)+3\overline{K}_2\mathrm{sin}\left(4\theta _c\right)`$ (45)
$`+12(\overline{K}_3\overline{K}_3^{})\mathrm{sin}^3\theta _c\mathrm{cos}\theta _c(58\mathrm{sin}^2\theta _c)]+\eta \{ϵ\overline{H}_c\mathrm{cos}(\theta _c\theta _H)`$ (46)
$`\eta ^2[(1/2)\overline{H}_c\mathrm{cos}(\theta _c\theta _H)+4\overline{K}_2\mathrm{cos}\left(4\theta _c\right)`$ (47)
$`+12(\overline{K}_3\overline{K}_3^{})\mathrm{sin}^2\theta _c(520\mathrm{sin}^2\theta _c+16\mathrm{sin}^4\theta _c)]\}=0,`$ ()
where $`\eta \theta _c\theta _0`$ which is small for $`ϵ1`$. By introducing a small variable $`\delta \theta \theta _0`$ $`\left(\left|\delta \right|1\text{ in the limit of }ϵ1\right)`$, the anisotropy energy becomes
$$\overline{E}_a(\delta ,\varphi )=\overline{K}_3^{}\left[1\mathrm{cos}\left(6\varphi \right)\right]\mathrm{sin}^6\left(\theta _0+\delta \right)+\overline{H}_x\left(1\mathrm{cos}\varphi \right)\mathrm{sin}\left(\theta _0+\delta \right)+\overline{E}_1\left(\delta \right),$$
(48)
where $`\overline{E}_1\left(\delta \right)`$ is a function of only $`\delta `$ given by
$`\overline{E}_1\left(\delta \right)`$ $`=`$ $`\left[{\displaystyle \frac{1}{2}}\overline{H}_c\mathrm{sin}\left(\theta _c\theta _H\right)+\overline{K}_2\mathrm{sin}\left(4\theta _c\right)+4\left(\overline{K}_3\overline{K}_3^{}\right)\left(5\mathrm{sin}^3\theta _c\mathrm{cos}^3\theta _c3\mathrm{sin}^5\theta _c\mathrm{cos}\theta _c\right)\right]`$ ()
$`\times (\delta ^33\delta ^2\eta )+[{\displaystyle \frac{1}{8}}\overline{H}_c\mathrm{cos}(\theta _c\theta _H)+\overline{K}_2\mathrm{cos}\left(4\theta _c\right)+3(\overline{K}_3\overline{K}_3^{})\mathrm{sin}^2\theta _c(\mathrm{sin}^4\theta _c`$
$`10\mathrm{sin}^2\theta _c\mathrm{cos}^2\theta _c+5\mathrm{cos}^4\theta _c)\left]\right(\delta ^44\delta ^3\eta +6\delta ^2\eta ^24\delta ^2ϵ)+ϵ\delta ^2[4\overline{K}_2\mathrm{cos}\left(4\theta _c\right)`$
$`+12(\overline{K}_3\overline{K}_3^{})\mathrm{sin}^2\theta _c(\mathrm{sin}^4\theta _c10\mathrm{sin}^2\theta _c\mathrm{cos}^2\theta _c+5\mathrm{cos}^4\theta _c)].`$
By applying the similar procedure in Sec. III, we obtain the transition amplitude Eqs. (21) and (22) by integrating out $`\varphi `$. For this case the effective mass is
$`M`$ $`=`$ $`\left(1+\left|\mathrm{cot}\theta _H\right|^{2/3}\right)\left[1+{\displaystyle \frac{8}{3}}\overline{K}_2{\displaystyle \frac{\left|\mathrm{cot}\theta _H\right|^{2/3}}{1+\left|\mathrm{cot}\theta _H\right|^{2/3}}}+8\left(\overline{K}_3\overline{K}_3^{}\right){\displaystyle \frac{\left|\mathrm{cot}\theta _H\right|^{2/3}}{\left(1+\left|\mathrm{cot}\theta _H\right|^{2/3}\right)^2}}\right]`$ ()
$`\times [1ϵ+36\overline{K}_3^{}{\displaystyle \frac{1}{1+\left|\mathrm{cot}\theta _H\right|^{2/3}}}+4\overline{K}_2(1+120\overline{K}_3^{}{\displaystyle \frac{\left|\mathrm{cot}\theta _H\right|^{2/3}}{1+\left|\mathrm{cot}\theta _H\right|^{2/3}}}){\displaystyle \frac{1}{1+\left|\mathrm{cot}\theta _H\right|^{2/3}}}`$
$`+6(\overline{K}_3\overline{K}_3^{})(1+240\overline{K}_3^{}{\displaystyle \frac{\left|\mathrm{cot}\theta _H\right|^{2/3}}{1+\left|\mathrm{cot}\theta _H\right|^{2/3}}}){\displaystyle \frac{1}{\left(1+\left|\mathrm{cot}\theta _H\right|^{2/3}\right)^2}}]^1,`$
and the prefactor $`A`$ is
$$A=2\frac{\left|\mathrm{cot}\theta _H\right|^{1/3}}{1+\left|\mathrm{cot}\theta _H\right|^{2/3}}\left[1+\frac{4}{3}\overline{K}_2\frac{74\left|\mathrm{cot}\theta _H\right|^{2/3}}{1+\left|\mathrm{cot}\theta _H\right|^{2/3}}+2\left(\overline{K}_3\overline{K}_3^{}\right)\frac{1116\left|\mathrm{cot}\theta _H\right|^{2/3}}{\left(1+\left|\mathrm{cot}\theta _H\right|^{2/3}\right)^2}\right].$$
(54)
In case of a small ferromagnet of volume $`Vr_0^3`$, the result of quantum nucleation is $`\mathrm{\Gamma }_Q\mathrm{exp}\left(𝒮_{cl}/\mathrm{}\right)`$, where the classical action for hexagonal symmetry is found to be
$`𝒮_{cl}`$ $`=`$ $`{\displaystyle \frac{2^{17/4}\times 3^{1/4}}{5}}\mathrm{}Sϵ^{5/4}\left|\mathrm{cot}\theta _H\right|^{1/6}`$ ()
$`\times \left[1+{\displaystyle \frac{2}{3}}\overline{K}_2{\displaystyle \frac{72\left|\mathrm{cot}\theta _H\right|^{2/3}}{1+\left|\mathrm{cot}\theta _H\right|^{2/3}}}+\left(\overline{K}_3\overline{K}_3^{}\right){\displaystyle \frac{1112\left|\mathrm{cot}\theta _H\right|^{2/3}}{\left(1+\left|\mathrm{cot}\theta _H\right|^{2/3}\right)^2}}\right]`$
$`\times [1ϵ+36\overline{K}_3^{}{\displaystyle \frac{1}{1+\left|\mathrm{cot}\theta _H\right|^{2/3}}}+4\overline{K}_2(1+120\overline{K}_3^{}{\displaystyle \frac{\left|\mathrm{cot}\theta _H\right|^{2/3}}{1+\left|\mathrm{cot}\theta _H\right|^{2/3}}}){\displaystyle \frac{1}{1+\left|\mathrm{cot}\theta _H\right|^{2/3}}}`$
$`+6(\overline{K}_3\overline{K}_3^{})(1+240\overline{K}_3^{}{\displaystyle \frac{\left|\mathrm{cot}\theta _H\right|^{2/3}}{1+\left|\mathrm{cot}\theta _H\right|^{2/3}}}){\displaystyle \frac{1}{\left(1+\left|\mathrm{cot}\theta _H\right|^{2/3}\right)^2}}]^{1/2},`$
In case of a thin film of thickness $`h`$ less than the size $`r_0/ϵ^{1/4}`$ of the critical nucleus, we obtain the quantum nucleation as $`\mathrm{\Gamma }_Q\mathrm{exp}\left(𝒮_E/\mathrm{}\right)`$, with the classical action
$`𝒮_E`$ $`=`$ $`74.39Sϵ^{3/4}r_0^2h{\displaystyle \frac{1+\left|\mathrm{cot}\theta _H\right|^{2/3}}{\left|\mathrm{cot}\theta _H\right|^{1/6}}}`$ ()
$`\times \left[1{\displaystyle \frac{2}{3}}\overline{K}_2{\displaystyle \frac{72\left|\mathrm{cot}\theta _H\right|^{2/3}}{1+\left|\mathrm{cot}\theta _H\right|^{2/3}}}\left(\overline{K}_3\overline{K}_3^{}\right){\displaystyle \frac{1112\left|\mathrm{cot}\theta _H\right|^{2/3}}{\left(1+\left|\mathrm{cot}\theta _H\right|^{2/3}\right)^2}}\right]`$
$`\times [1ϵ+36\overline{K}_3^{}{\displaystyle \frac{1}{1+\left|\mathrm{cot}\theta _H\right|^{2/3}}}+4\overline{K}_2(1+120\overline{K}_3^{}{\displaystyle \frac{\left|\mathrm{cot}\theta _H\right|^{2/3}}{1+\left|\mathrm{cot}\theta _H\right|^{2/3}}}){\displaystyle \frac{1}{1+\left|\mathrm{cot}\theta _H\right|^{2/3}}}`$
$`+6(\overline{K}_3\overline{K}_3^{})(1+240\overline{K}_3^{}{\displaystyle \frac{\left|\mathrm{cot}\theta _H\right|^{2/3}}{1+\left|\mathrm{cot}\theta _H\right|^{2/3}}}){\displaystyle \frac{1}{\left(1+\left|\mathrm{cot}\theta _H\right|^{2/3}\right)^2}}]^{1/2}.`$
And the crossover temperature for hexagonal symmetry is found to be
$`k_BT_c`$ $``$ $`0.55{\displaystyle \frac{K_1ϵ^{1/4}}{S}}{\displaystyle \frac{\left|\mathrm{cot}\theta _H\right|^{1/6}}{1+\left|\mathrm{cot}\theta _H\right|^{2/3}}}`$ ()
$`\times \left[1+{\displaystyle \frac{2}{3}}\overline{K}_2{\displaystyle \frac{72\left|\mathrm{cot}\theta _H\right|^{2/3}}{1+\left|\mathrm{cot}\theta _H\right|^{2/3}}}+\left(\overline{K}_3\overline{K}_3^{}\right){\displaystyle \frac{1112\left|\mathrm{cot}\theta _H\right|^{2/3}}{\left(1+\left|\mathrm{cot}\theta _H\right|^{2/3}\right)^2}}\right]`$
$`\times [1ϵ+36\overline{K}_3^{}{\displaystyle \frac{1}{1+\left|\mathrm{cot}\theta _H\right|^{2/3}}}+4\overline{K}_2(1+120\overline{K}_3^{}{\displaystyle \frac{\left|\mathrm{cot}\theta _H\right|^{2/3}}{1+\left|\mathrm{cot}\theta _H\right|^{2/3}}}){\displaystyle \frac{1}{1+\left|\mathrm{cot}\theta _H\right|^{2/3}}}`$
$`+6(\overline{K}_3\overline{K}_3^{})(1+240\overline{K}_3^{}{\displaystyle \frac{\left|\mathrm{cot}\theta _H\right|^{2/3}}{1+\left|\mathrm{cot}\theta _H\right|^{2/3}}}){\displaystyle \frac{1}{\left(1+\left|\mathrm{cot}\theta _H\right|^{2/3}\right)^2}}]^{1/2}.`$
## V. Conclusions and discussions
In summary we have investigated the quantum nucleation of magnetization in nanometer-scale ferromagnets in the presence of an external magnetic field at arbitrary angle. We consider the magnetocrystalline anisotropy with tetragonal symmetry and that with hexagonal symmetry, respectively. By applying the instanton method in the spin-coherent-state path-integral representation, we obtain the analytical formulas for quantum reversal of magnetization in small magnets and the numerical formulas for quantum nucleation in thin ferromagnetic film in a wide range of angles $`\pi /2<\theta _H<\pi `$. The temperature characterizing the crossover from the quantum to thermal nucleation is clearly shown for each case. Our results show that the rate of quantum nucleation and the crossover temperature depend on the orientation of the external magnetic field distinctly. Therefore, both the orientation and the strength of the external magnetic field are the controllable parameters for the experimental test of quantum nucleation of magnetization in nanometer-scale ferromagnets. If the experiment is to be performed, there are three control parameters for comparison with theory: the angle of the external magnetic field $`\theta _H`$, the strength of the field in terms of $`ϵ`$, and the temperature $`T`$. Our results show that ferromagnetic samples with large anisotropy and small saturated magnetization are suitable for experimental study the phenomenon of quantum nucleation.
Recently, Wernsdorfer and co-workers have performed the switching field measurements on individual ferrimagnetic and insulating BaFeCoTiO nanoparticles containing about $`10^5`$-$`10^6`$ spins at very low temperatures (0.1-6K). They found that above 0.4K, the magnetization reversal of these particles is unambiguously described by the Néel-Brown theory of thermal activated rotation of the particle’s moment over a well defined anisotropy energy barrier. Below 0.4K, strong deviations from this model are evidenced which are quantitatively in agreement with the predictions of the MQT theory without dissipation. It is noted that the observation of quantum nucleation in ferromagnets would be extremely interesting as the next example, after single-domain nanoparticles, of macroscopic quantum tunneling. The theoretical results presented here may be useful for checking the general theory in a wide range of systems, with more general symmetries. The experimental procedures on single-domain ferromagnetic nanoparticles of Barium ferrite with uniaxial symmetry may be applied to the systems with more general symmetries. Note that the inverse of the WKB exponent $`B^1`$ is the magnetic viscosity $`S`$ at the quantum-tunneling-dominated regime $`TT_c`$ studied by magnetic relaxation measurements. Therefore, the quantum nucleation of magnetization should be checked at any $`\theta _H`$ by magnetic relaxation measurements. Over the past years a lot of experimental and theoretical works were performed on the spin tunneling in molecular Mn<sub>12</sub>-Ac and Fe<sub>8</sub> clusters having a collective spin state $`S=10`$ (in this paper $`S=10^6`$). Further experiments should focus on the level quantization of collective spin states of $`S=10^2`$-$`10^4`$.
The ferromagnet is typically an insulating particle with as many as $`10^310^6`$ magnetic moments. For the dynamical process, it is important to include the effect of the environment on quantum tunneling caused by phonons, nucleation spins, and Stoner excitations and eddy currents in metallic magnets. However, many studies showed that even though these couplings might be crucial in macroscopic quantum coherence, they are not strong enough to make the quantum tunneling unobservable. The theoretical calculations performed in this paper can be extended to the AFM bubbles, where the relevant quantity is the excess spin due to the small noncompensation of two sublattices. Work along this line is still in progress. We hope that the theoretical results presented in this paper may stimulate more experiments whose aim is observing quantum nucleation in nanometer-scale ferromagnets.
## Acknowledgments
Y. Zhou and R.L. would like to acknowledge Dr. Su-Peng Kou, Professor Zhan Xu, Professor Mo-Lin Ge, Professor Jiu-Qing Liang and Professor Fu-Cho Pu for stimulating discussions. R. L. would like to thank Professor W. Wernsdorfer and Professor R. Sessoli for providing their paper (Ref. 11), and Professor Kim for providing his paper (Ref. 13).
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# References
On Teleportation of a Completely Unknown State
of Relativistic Photon
R.Laiho, S.N.Molotkov, and S.S.Nazin
Wihuri Physical Laboratory,
Department of Physics, University of Turku, 20014 Turku, Finland
Institute of Solid State Physics of Russian Academy of Sciences,
Chernogolovka, Moscow District, 142432 Russia
## Abstract
The process of teleportation of a completely unknown single-photon relativistic state is considered. Analysis of the relativistic case reveals that the teleportation as it is understood in the non-relativistic quantum mechanics is impossible if no a priori information on the state to be teleported is available. It is only possible to speak of the amplitude of the propagation of the field (taking into account the measurement procedure) since the existence of a common vacuum state together with the microcausality principle (the field operators commutation relations) make the concept of the propagation amplitude for the individual subsystems physically meaningless. When partial a priori information is available (for example, only the polarization state of the photon is unknown while its spatial state is specified beforehand), the teleportation does become possible in the relativistic case. In that case the a priori information can be used to “label” the identical particles to make them effectively distinguishable.
PACS numbers: 03.67.-a, 03.65.Bz, 42.50Dv
One of the basic results of the non-relativistic quantum information theory consists in the possibility of the ideal teleportation of an unknown quantum state by means of a classical and distributed quantum communication channel . The latter is realized by a non-local entangled state (an EPR-pair ). The teleportation of an unknown quantum state from the user A to user B is performed in the following way . The user A has a quantum system 1 in an unknown state $`\rho _1`$ (for example, a spin-1/2 particle) whose state is to be teleported to user B. To do this, the user A prepares a composite system consisting of the two particles 2 and 3 in a maximally entangled state $`\rho _{23}`$ (EPR-pair) so that the user B can only access the particle 3 while the particles 1 and 2 remain available to user A. Then user A performs an appropriate joint measurement over the particles 1 (whose state is to be teleported) and 2. The measurement brings the composite system consisting of the particles 1, 2, and 3 from the initial state $`\rho _1\rho _{23}`$ to the new state $`\rho _{123}^{^{}}`$ which depends on the measurement outcome $`i`$ (in the teleportation of the system described by a finite-dimensional state space the set of all possible outcomes is discrete). The crucial point here is the existence of the measurements possessing the special property that the state of the third particle after the measurement $`\rho _3^{^{}}`$ obtained by taking a partial trace with respect to the states of particles 1 and 2, $`\rho _3^{^{}}=\text{Tr}_{12}\{\rho _{123}^{^{}}\}`$, coincides with the unknown initial state $`\rho _1`$ to within a unitary transformation which is completely determined by the measurement outcome $`i`$ obtained by user A only:
$$\rho _1=U_i\rho _3^{^{}}U_i^1.$$
(1)
Here we identify the isomorphic state space of the subsystems 1 and 3. The classical communication channel in the outlined teleportation scheme is required to convey the measurement outcome $`i`$ from user A to the distant user B. Preforming the measurement, the user A acquires no information on the teleported state since all the measurement outcomes are equiprobable.
Later, a number of different non-relativistic teleportation schemes were proposed. All the schemes based on the non-relativistic quantum mechanics employ the fact that the state space of a composite system can be represented as the tensor product of the state spaces of the constituent subsystems. In addition, it is important that the measurement performed over a composite system can be represented as a tensor product of the appropriate measurements over the constituent subsystems.
The possibility of representation of both the state space and the measurement as a tensor product stems from the following argument. If one assumes that the composite system consists of physically distinct (and hence distinguishable) particles, its state space can be written as a tensor product of the corresponding states. Then the quantum teleportation can be described using the standard mathematical apparatus of quantum mechanics.
However, the non-relativistic quantum mechanics provides only an approximate description of real world. A more complete description is provided by the quantum field theory. In the quantum field theory the superselection rules prohibit the superposition of pure states belonging to different coherent sectors . The analysis of teleportation of identical particles belonging to the same coherent sector in the quantum field theory proves to be substantially different from the non-relativistic theory. First, the state space of a composite system can no longer be considered as the tensor product of the state spaces of constituent subsystems because of the existence of the common cyclic vacuum vector . Second, the microcausality principle (commutation or anticommutation relations) requires that the field operators commute (anticommute) if their supports are separated by a space-like interval .
Using the photon field as an example, we shall show below that the teleportation of the state of one of the identical particles as it is understood in the non-relativistic quantum mechanics is impossible in the quantum field theory because of the existence of a common vacuum vector and impossibility of writing a measurement over a composite system as a tensor product of appropriate measurements over the constituent systems. It is only possible to speak of the propagation amplitude of the field as a whole.
It should be emphasized that our arguments do not formally prohibit the possibility of teleportation of a completely unknown state in a system of physically distinguishable particles (for example, a photon state can be teleported to the electron degrees of freedom and vice versa). However, a correct analysis of this question within the framework of the quantum field theory for interacting fields is much more difficult.
Consider now the teleportation of a completely unknown single-photon state of the electromagnetic field.
The electromagnetic field operators can be written as
$$A_\mu ^\pm (\widehat{x})=\frac{1}{(2\pi )^{3/2}}\frac{d𝐤}{\sqrt{2k^0}}\text{e}^{\pm i\widehat{k}\widehat{x}}e_\mu ^m(𝐤)a_m^\pm (𝐤),$$
(2)
$$\mu ,m=1\mathrm{}4,\widehat{x}=(x^0,𝐤),\widehat{k}=(k^0,𝐤),k^0=|𝐤|,$$
$$(𝐞^m𝐞^n)=\delta _{m,n},(m,n=1÷3),e_0^m=0,𝐞^3=𝐤/|𝐤|.$$
The four-dimensional vector potential operators satisfy the commutation relations
$$[A_\mu ^{}(\widehat{x}),A_\nu ^+(\widehat{x}^{})]=ig_{\mu \nu }D_0^{}(\widehat{x}\widehat{x}^{}),$$
(3)
where $`g_{\mu \nu }`$ ($`g_00=g_ii`$, $`i=1÷3`$), and $`D_0^{}(\widehat{x})`$ is the commutator function for the massless field
$$D_0^\pm (\widehat{x})=\pm \frac{1}{i(2\pi )^{3/2}}\text{e}^{\pm i\widehat{k}\widehat{x}}\theta (k^0)\delta (\widehat{k}^2)𝑑\widehat{k}=\pm \frac{1}{i(2\pi )^{3/2}}\frac{d𝐤}{2k^0}\text{e}^{\pm i\widehat{k}\widehat{x}}=$$
(4)
$$\frac{1}{4\pi }\frac{\delta (x^0|𝐱|)\delta (x^0+|𝐱|)}{2|𝐱|}.$$
There are for types of the creation operators ($`a_m^\pm (𝐤)`$) in Eq. (2) describing two transverse, one temporal, and longitudinal photons The last two photon types are unphysical and can be eliminated by introducing an indefinite metrics . For our purposes the shortest way to the required result involves employment of a particular gauge. We shall further work in the physical subspace of the two types of transverse photons in the Coulomb gauge ($`A_\mu =(𝐀,\phi =0)`$).
In that case the operator-valued distribution is a three-dimensional vector
$$𝝍(\widehat{x})=\frac{1}{(2\pi )^{3/2}}\underset{s=\pm }{}\frac{d𝐤}{\sqrt{2k^0}}𝐰(𝐤,s)\{a^{}(𝐤,s)\text{e}^{i\widehat{k}\widehat{x}}+a^+(𝐤,s)\text{e}^{i\widehat{k}\widehat{x}}\},$$
(5)
where $`𝐰(𝐤,s)`$ is a three-dimensional vector describing the state helicity ($`s=\pm `$)
$$𝐰(𝐤,s)=\frac{1}{\sqrt{2}}(𝐞_1(𝐤)\pm i𝐞_2(𝐤)),𝐞_2(𝐤)𝐞_2(𝐤)=0,𝐞_1(𝐤)𝐞_2(𝐤)𝐤,$$
(6)
where $`𝐞_{1,2}(𝐤)`$ are the linear polarization vectors. The creation operators in Eq.(5) satisfy the commutation relations
$$[a^{}(𝐤s),a^+(𝐤,s^{})]=\delta (𝐤𝐤^{})\delta _{s,s^{}}.$$
(7)
For convenience we introduce the new operators
$$𝐀^\pm (𝐤,s)=𝐰(𝐤,s)a^\pm (𝐤,s).$$
(8)
The operator-valued distribution satisfies the Maxwell equation for the free electromagnetic field
$$\times 𝝍(\widehat{x})=i\frac{}{x^0}𝝍(\widehat{x}),$$
(9)
$$𝝍(\widehat{x})=0.$$
Any state of the free electromagnetic field can be obtained by acting with the field operators on the cyclic vacuum vector
$$|𝚿=(f_0+\underset{n=1}{\overset{\mathrm{}}{}}\underset{s_1=\pm }{}\underset{s_2=\pm }{}\mathrm{}\underset{s_n=\pm }{}\mathrm{}\frac{d𝐤_1}{\sqrt{2k_1^0}}\frac{d𝐤_2}{\sqrt{2k_2^0}}\mathrm{}\frac{d𝐤_n}{\sqrt{2k_n^0}}$$
(10)
$$f(𝐤_1,s_1,𝐤_2,s_2,\mathrm{}𝐤_n,s_n,)𝐀^+(𝐤_1,s_1)𝐀^+(𝐤_2,s_2)\mathrm{}𝐀^+(𝐤_n,s_n))|0,$$
where $`f(𝐤_1,s_1,𝐤_2,s_2,\mathrm{}𝐤_n,s_n,)`$ is the amplitude in the $`𝐤`$-representation.
The creation operator of the unknown single-particle field state to be teleported is written as
$$𝝍^+(𝐟)=\frac{d𝐤}{\sqrt{2k^0}}\left(f(𝐱,+)\text{e}^{ik^0x^0}𝐀^+(𝐤,+)+f(𝐱,)\text{e}^{ik^0x^0}𝐀^+(𝐤,)\right).$$
(11)
It is convenient to write the temporal factor in Eq.(11) explicitly although it can be incorporated in the definition of the single-particle amplitude. In the position representation Eq.(11) becomes
$$𝝍^+(𝐟)=𝑑𝐱\left(f(𝐱,+)𝝍^+(\widehat{x},+)+f(𝐱,)𝝍^+(\widehat{x},)\right).$$
(12)
We have included the temporal factor in the argument of the field operator in Eq.(12). Here $`f(𝐱,\pm )`$ is the Fourier transform of the amplitude in the $`𝐤`$-representation. The quantity $`f(𝐱,\pm )`$ is interpreted as the packet shape in the $`𝐱`$-representation (for helicities $`\pm `$) at time $`x^0`$. Note that because of the transverse nature of the electromagnetic field it is generally impossible to factorize the polarization and spatial degrees of freedom (as it is frequently done without any justifications in the description of experiments when only the two-dimensional polarization states space is considered).
Let us now construct the entangled EPR-state. The relativistic counterpart of a maximally entangled EPR-state with respect to the polarization (and, because of the transverse nature of the field, automatically with respect to the momentum) is
$$|𝚿_{00}=\left(\frac{d𝐤}{2k^0}\text{e}^{2ik^0x^0}𝐀^+(𝐤,+)𝐀^+(𝐤,)\right)|0=\left(𝑑𝐱𝝍^+(\widehat{x},+)𝝍^+(\widehat{x},)\right)|0,$$
(13)
where for convenience the temporal factor is again introduced which is later included into field operator argument in Eq.(13). The EPR-state formally correspond to the two-particle amplitude chosen in the form $`(\widehat{x}_1,s_1,\widehat{x}_2,s_2)=\delta _{s_1,s_2}\delta (x_1^0x^0)\delta (x_2^0x^0)\delta (𝐱_1𝐱_2)const(𝐱_1+𝐱_2)`$. The latter should be understood as a limit of smooth (test) functions. Formally, the EPR state (13) can be interpreted as describing the creation of two photons with $`s_1=s_2`$ at time $`x^0`$ at the point $`𝐱_1=𝐱_2`$ in a correlated non-local way simultaneously in the entire space as suggested by the factor $`const(𝐱_1+𝐱_2)`$.
The EPR-state (13) corresponds to a pair of photons with zero (in the chosen reference frame) total momentum. Since the polarization degrees of freedoms generally do not factor out from the spatial ones, there exist a continuum of maximally entangled EPR-states with different total momentum. Although any of these states can equally well be used for our purposes, we shall use the pair with zero total momentum to make the analogy with the non-relativistic Bell basis more clear.
The states analogous to the non-relativistic Bell basis are
$$|𝚿_{\mathrm{𝐘𝐏}}^{(\pm )}=𝚿_{\mathrm{𝐘𝐏}}^{(\pm )+}|0=\frac{1}{\sqrt{2}}\left(𝑑𝝃\text{e}^{i𝝃𝐏}(𝝍^+(\widehat{\xi },+)𝝍^+(\widehat{\xi }𝐘,)\pm 𝝍^+(\widehat{\xi },)𝝍^+(\widehat{\xi }𝐘,+))\right)|0,$$
(14)
$$|𝚽_{\mathrm{𝐘𝐏}}^{(\pm )}=𝚽_{\mathrm{𝐘𝐏}}^{(\pm )+}|0=\frac{1}{\sqrt{2}}\left(𝑑𝝃\text{e}^{i𝝃𝐏}(𝝍^+(\widehat{\xi },+)𝝍^+(\widehat{\xi }𝐘,+)\pm 𝝍^+(\widehat{\xi },)𝝍^+(\widehat{\xi }𝐘,))\right)|0.$$
It is easily checked that these states constitute an orthogonal identity resolution in the subspace of two-particle states. Indeed,
$$I_2=\frac{d𝐘d𝐏}{(2\pi )^3}$$
(15)
$$\left(|𝚿_{\mathrm{𝐘𝐏}}^{(+)}𝚿_{\mathrm{𝐘𝐏}}^{(+)}|+|𝚿_{\mathrm{𝐘𝐏}}^{()}𝚿_{\mathrm{𝐘𝐏}}^{()}|+|𝚽_{\mathrm{𝐘𝐏}}^{(+)}𝚽_{\mathrm{𝐘𝐏}}^{(+)}|+|𝚽_{\mathrm{𝐘𝐏}}^{()}𝚽_{\mathrm{𝐘𝐏}}^{()}|\right)=$$
$$\underset{s_1,s_2=\pm }{}\frac{d𝐤_1}{2k_2^0}\frac{d𝐤_2}{2k_2^0}\left(𝐀^+(𝐤_1,s_1)𝐀^+(𝐤_2,s_2)|0\right)\left(0|𝐀^{}(𝐤_1,s_1)𝐀^{}(𝐤_2,s_2)\right)$$
It should be emphasized once again that because of the transverse nature of the electromagnetic field it is impossible to construct an identity resolution (as well as an EPR-pair) with factorized polarization and spatial degrees of freedom.
The initial common (because of the common vacuum vector) state of the field describing the completely unknown state to be teleported and the EPR-pair is written as
$$|𝝍(𝐟)𝚿_{\mathrm{𝐘𝐏}}^{()}=𝝍^+(𝐟)𝚿_{\mathrm{𝐘𝐏}}^{()+}|0=𝑑𝐱𝑑𝝃\text{e}^{i𝝃𝐏}$$
(16)
$$\{𝚿^{()+}(\widehat{\xi },\widehat{x})(f(𝐱,+)𝝍^+(\widehat{\xi }𝐘,+)+f(𝐱,)𝝍^+(\widehat{\xi }𝐘,))+$$
$$𝚿^{(+)+}(\widehat{\xi },\widehat{x})(f(𝐱,+)𝝍^+(\widehat{\xi }𝐘,+)+f(𝐱,)𝝍^+(\widehat{\xi }𝐘,))+$$
$$𝚽^{()+}(\widehat{\xi },\widehat{x})(f(𝐱,+)𝝍^+(\widehat{\xi }𝐘,)+f(𝐱,)𝝍^+(\widehat{\xi }𝐘,+))+$$
$$𝚽^{(+)+}(\widehat{\xi },\widehat{x})(f(𝐱,+)𝝍^+(\widehat{\xi }𝐘,)f(𝐱,)𝝍^+(\widehat{\xi }𝐘,+))\}|0.$$
Here the following notation is introduced:
$$𝚿^{(\pm )+}(\widehat{\xi },\widehat{x})=\frac{1}{\sqrt{2}}\left(𝝍^+(\widehat{x},+)𝝍^+(\widehat{\xi },)\pm 𝝍^+(\widehat{x},)𝝍^+(\widehat{\xi },+)\right),$$
(17)
$$𝚽^{(\pm )+}(\widehat{\xi },\widehat{x})=\frac{1}{\sqrt{2}}\left(𝝍^+(\widehat{x},+)𝝍^+(\widehat{\xi },+)\pm 𝝍^+(\widehat{x},)𝝍^+(\widehat{\xi },)\right),$$
The temporal factors are again included in the arguments $`\widehat{x}`$ and $`\widehat{\xi }`$. The identical transformations in Eq. (16) are performed to make the analogy with the non-relativistic case more graphical.
Let us now discuss the construction of the appropriate measurement. We shall go back for a moment to the non-relativistic case. The change of the state of a quantum system after a measurement act is completely described by the corresponding instrument (superoperator) . In the context of the teleportation problem, when the measurement affects only the particle in the unknown state to be teleported (particle 1) and one of the particles in the EPR-pair (particle 2) while the particle 3 is not directly affected by the measurement, such an instrument with the space of all possible outcomes $`\mathrm{\Theta }`$ is written as $`\text{T}_{123}(d\theta )=\text{T}_{12}(d\theta )I_3`$, where $`I_3`$ is the identity operator in the state space of the third particle. The system state just after the measurement which gave the outcome in the interval $`(\theta ,\theta +d\theta )`$ is
$$\rho _{123}^{^{}}=\frac{\text{T}_{123}(d\theta )\rho _{1233}}{\text{Tr}_{123}\{\text{T}_{123}(d\theta )\rho _{123}\}}$$
(18)
The density matrix of particle 3 after the measurement is obtained by taking the partial trace over the states of particles 1 and 2:
$$\rho _3^{^{}}=\frac{\text{ Tr}_{12}\{\text{T}_{123}(d\theta )\rho _{123}\}}{\text{Tr}_{123}\{\text{T}_{123}(d\theta )\rho _{123}\}}.$$
(19)
If the composite system S consists of two subsystems A and B with the density matrix $`\rho _{AB}`$ (for the teleportation problem, the subsystem A consists of particles 1 and 2 while subsystem B consists of particle 3), the state of subsystem B just after the measurement is (to within the normalization constant)
$$\rho _B^{^{}}=\text{ Tr}_A\{\text{T}_{AB}(d\theta )\rho _{AB}\}=\text{ Tr}_A\{(_A(d\theta )I_B)\rho _{AB}\};$$
(20)
here $`_A(d\theta )=[\text{T}_{AB}(d\theta )]^{}I_{AB}`$ is a positive operator-valued measure on the space of possible outcomes $`\mathrm{\Theta }`$ generated by the instrument. Therefore, in the teleportation problem it is sufficient to know only the measurement $`_A(d\theta )`$ rather than the instrument itself (generally, to determine the system state just after the measurement one should know the corresponding instrument itself). However, this statement is only valid if the instrument (and measurement) can be represented as a tensor product of the instruments (measurements) acting in the state spaces of the subsystems.
In the quantum field theory the measurement cannot be in principle represented as a tensor product of the measurements related to the individual subsystems because of the following two reasons. The first one is the existence of the common vacuum vector. The second reason is the microcausality principle (commutation relation given by Eq. (2)). The possibility of the representation of the measurement over identical particles in the form of a tensor product implies that the field operators related to two different factors (particles) always commute irrespective of the relative position of their supports in the Minkowski space. Hence in the quantum field theory the measurement over a composite system consisting of three particles can only be written as a general identity resolution in the entire space of three-particle states. However, then the teleportation itself as it is understood in non-relativistic quantum mechanics becomes meaningless and one can only speak of the amplitude of propagation of the field as a whole rather than the individual subsystems.
Nevertheless, it is still interesting to find out at the level of concrete formulas where the two indicated circumstances come into play.
We shall further need the following identity resolution in the one-particle state space:
$$I_1=\underset{s=\pm }{}\frac{d𝐤}{2k^0}\left(𝐀^+(𝐤,s)|0\right)\left(0|𝐀^{}(𝐤,s)\right)=$$
(21)
$$\underset{s=\pm }{}𝑑𝐱\left(𝝍^+(𝐱,s)|0\right)\left(0|𝝍^{}(𝐱,s)\right).$$
The set of possible measurement outcomes $`\mathrm{\Theta }`$ for the two-particle state space is $`\mathrm{\Theta }=`$ $`\{i,k,s,𝐗,𝐐:i\times k\times (\pm )\times 𝐑_𝐗\times 𝐑_𝐐\}`$, where $`i=1,2`$ and $`k=\pm `$ label the states of the Bell basis. The total identity resolution in the three-particle state space is
$$I_3=\underset{s=\pm ,k=\pm ,i=1,2}{}_{is}^k(d\theta ),$$
(22)
where
$$_{1s}^\pm (d\theta )=\left(𝝍^+(\widehat{x},s)|𝚿_{\mathrm{𝐗𝐐}}^{(\pm )}\right)\left(𝚿_{\mathrm{𝐗𝐐}}^{(\pm )}|𝝍^{}(\widehat{x},s)\right)\frac{d𝐱d𝐗d𝐐}{(2\pi )^3},$$
$$_{2s}^\pm (d\theta )=\left(𝝍^+(\widehat{x},s)|𝚽_{\mathrm{𝐗𝐐}}^{(\pm )}\right)\left(𝚽_{\mathrm{𝐗𝐐}}^{(\pm )}|𝝍^{}(\widehat{x},s)\right)\frac{d𝐱d𝐗d𝐐}{(2\pi )^3}.$$
The resolution (22) is an analogue of the identity resolution of the form $`_{12}I_1`$ arising in the non-relativistic analysis although Eq.(22) does not reduce to the latter because of the two indicated reasons. It should be noted that the same time $`x^0`$ appears in all arguments in Eq. (22). Such a measurement can be interpreted as a spatially non-local measurement performed at time $`x^0`$.
Note also that unlike the orthogonal identity resolution (15) in the two-particle state space the resolution in the three-particle space (in contrast to the non-relativistic case) becomes non-orthogonal. Of course, this measurement can be made orthogonal by choosing the symmetry-adapted basis functions realizing the irreducible representations of the permutation group in the three-particle space, but it cannot be made orthogonal when restricted to the two-particle space.
The probabilities of various measurement outcomes from the space $`\mathrm{\Theta }`$ are given by the standard formula
$$\text{Pr}\{s,i,k,d\theta \}=\text{Tr}\{|𝝍(𝐟)𝚿_{\mathrm{𝐘𝐏}}^{()}𝚿_{\mathrm{𝐘𝐏}}^{()}𝝍(𝐟)|_{is}^k(d\theta )\}=|𝒜_{is}^k(d\theta )|^2,$$
(23)
where the amplitude $`𝒜_{is}^k(d\theta )`$ is defined, for example, in the channel $`_{s1}^{}(d\theta )`$ as
$$𝒜_{1s}^{}(d\theta )=0|𝚿_{\mathrm{𝐗𝐐}}^{()}𝝍^{}(\widehat{x},s)𝝍^+(𝐟)𝚿_{\mathrm{𝐘𝐏}}^{()+}|0d\theta =$$
(24)
$$\frac{1}{2}𝑑𝐱^{}𝑑𝝃^{}𝑑𝝃\text{e}^{i(𝝃𝐏𝝃^{}𝐐)}$$
$$\{0|𝚿_{\mathrm{𝐗𝐐}}^{()}𝝍^{}(\widehat{x},s)𝚿^{()+}(\widehat{\xi },\widehat{x}^{})(f(𝐱^{},+)𝝍^+(\widehat{\xi }𝐘,+)+f(𝐱^{},)𝝍^+(\widehat{\xi }𝐘,))|0+$$
$$0|𝚿_{\mathrm{𝐗𝐐}}^{()}𝝍^{}(\widehat{x},s)𝚿^{()+}(\widehat{\xi },\widehat{x}^{})\left(f(𝐱^{},+)𝝍^+(\widehat{\xi }𝐘,+)+f(𝐱^{},)𝝍^+(\widehat{\xi }𝐘,)\right)|0+$$
$$0|𝚿_{\mathrm{𝐗𝐐}}^{()}𝝍^{}(\widehat{x},s)𝚽^{()+}(\widehat{\xi },\widehat{x}^{})\left(f(𝐱^{},+)𝝍^+(\widehat{\xi }𝐘,)+f(𝐱^{},)𝝍^+(\widehat{\xi }𝐘,+)\right)|0+$$
$$0|𝚿_{\mathrm{𝐗𝐐}}^{()}𝝍^{}(\widehat{x},s)𝚽^{(+)+}(\widehat{\xi },\widehat{x}^{})(f(𝐱^{},+)𝝍^+(\widehat{\xi }𝐘,)f(𝐱^{},)𝝍^+(\widehat{\xi }𝐘,+))|0\}d\theta .$$
It is impossible to identify the individual contribution of different subsystems to the field propagation amplitude. Nevertheless, it is interesting to consider the transition to the non-relativistic limit. In that case, if one neglects the commutation relations and assumes that the operators $`𝚿_{\mathrm{𝐗𝐐}}^{(\pm )\pm }`$ ($`𝚽_{\mathrm{𝐗𝐐}}^{(\pm )\pm }`$), $`𝝍^{}(\widehat{x},s)`$ and $`\left(f(𝐱^{},+)𝝍^+(\widehat{\xi }𝐘,+)+f(𝐱^{},)𝝍^+(\widehat{\xi }𝐘,)\right)`$, etc. commute then, because of the orthogonality of different $`𝚿_{\mathrm{𝐗𝐐}}^{(\pm )\pm }`$ and $`𝚽_{\mathrm{𝐗𝐐}}^{(\pm )\pm }`$, only one term is left in Eq. (24), for example, only the outcome related to the projection on $`𝚿_{\mathrm{𝐗𝐐}}^{()}`$ in which case we obtain
$$𝒜_{1s}^{}(d\theta )\frac{1}{2}𝑑𝐱^{}𝑑𝝃^{}𝑑𝝃\text{e}^{i(𝝃𝐏𝝃^{}𝐐)}$$
(25)
$$0|𝚿_{\mathrm{𝐗𝐐}}^{()}𝝍^{}(\widehat{x},s)𝚿^{()+}(\widehat{\xi },\widehat{x}^{})\left(f(𝐱^{},+)𝝍^+(\widehat{\xi }𝐘,+)+f(𝐱^{},)𝝍^+(\widehat{\xi }𝐘,)\right)|0d\theta .$$
For the free field the vacuum average in Eq. (25) factorizes into the pair-wise averages. Assuming again that $`𝚿_{\mathrm{𝐗𝐐}}^{()}`$ and $`𝚿^{()+}(\widehat{\xi },\widehat{x}^{})`$ are related to the particles 1 and 2 (particle in the unknown state and one of the EPR-pair particles) and act in the appropriate subsystem state space while the operators $`𝝍^{}(\widehat{x},s)`$ and $`\left(f(𝐱^{},+)𝝍^+(\widehat{\xi }𝐘,+)+f(𝐱^{},)𝝍^+(\widehat{\xi }𝐘,)\right)`$ act in the third particle state space (the second photon in the EPR-pair), one obtains
$$𝒜_{1s}^{}(d\theta )\frac{1}{2}𝑑𝐱^{}𝑑𝝃^{}𝑑𝝃\text{e}^{i(𝝃𝐏𝝃^{}𝐐)}0|𝚿_{\mathrm{𝐗𝐐}}^{()}𝚿^{()+}(\widehat{\xi },\widehat{x}^{})|0$$
(26)
$$0|𝝍^{}(\widehat{x},s)\left(f(𝐱^{},+)𝝍^+(\widehat{\xi }𝐘,+)+f(𝐱^{},)𝝍^+(\widehat{\xi }𝐘,)\right)|0.$$
The amplitude in Eq. (26) has a rather transparent interpretation. The first factor describes the propagation of a “new” EPR-pair after the measurement affecting the particle 1 (which was initially in a completely unknown state) and one of the particles of the original EPR-pair. This is described by the operator $`𝚿^{()+}(\widehat{\xi },\widehat{x}^{})`$. Corresponding to the measurement in the Bell basis is the action of the operator of annihilation of the new EPR-pair $`𝚿_{\mathrm{𝐗𝐐}}^{()}`$.
The second factor describes the propagation of the second particle from the original EPR-pair which due to the initial correlations in the EPR-pair is brought after the performed Bell measurement into the state ($`\left(f(𝐱^{},+)𝝍^+(\widehat{\xi }𝐘,+)+f(𝐱^{},)𝝍^+(\widehat{\xi }𝐘,)\right)`$ which coincides to within a unitary transformation completely determined by the measurement outcome with the unknown state of the particle 1 which had to be teleported. The fact that the measurement did not affect the particle 3 corresponds to the annihilation operator $`𝝍^{}(\widehat{x},s)`$ formally describing the projection onto all one-particle states (because of the integration over $`𝐱`$ in the expression (23) for the probability).
One can advance further in this way using the explicit expression for the first factor in Eq. (26) which gives
$$0|𝚿_{\mathrm{𝐗𝐐}}^{()}𝚿^{()+}(\widehat{\xi },\widehat{x}^{})|0=2\left(𝒟_0^+(\widehat{\xi }^{}\widehat{x}^{})𝒟_0^+(\widehat{\xi }^{}\widehat{\xi }𝐗)𝒟_0^+(\widehat{\xi }^{}\widehat{\xi })𝒟_0^+(\widehat{\xi }^{}\widehat{x}^{}𝐗)\right).$$
(27)
The commutator function $`𝒟_0^+(\widehat{x}\widehat{y})`$ describes the process of particle creation at point $`\widehat{x}`$, its propagation, and annihilation at point $`\widehat{y}`$ (for $`y^0>x^0`$) :
$$0|𝝍^{}(\widehat{y},s),𝝍^+(\widehat{x},s)|0=i𝒟_0^+(\widehat{x}\widehat{y}).$$
(28)
The second term in Eq. (27) arises because of the particles indistinguishability (propagation with the exchange of the particles) If one again assumes that the creation and annihilation operators related to the particles 1 and 2 act in the different state spaces (which would correspond to the distinguishability of the particles), then the second term should be discarded.
Finally, to obtain the non-relativistic limit, the commutator functions $`𝒟_0^+(\widehat{x})`$ in Eq. (27) should be replaced by the ordinary $`\delta `$-functions of the three-dimensional argument ($`\delta (𝐱)`$). This replacement of $`𝒟_0^+(\widehat{x})`$-functions by ordinary $`\delta (𝐱)`$-functions should be done because in the non-relativistic case the integration is performed with a Galilei-invariant measure $`d\mu (𝐤)=d𝐤`$ while in the relativistic case one employs the Lorentz-invariant measure $`d\mu (𝐤)=\theta (k^0)\delta (\widehat{k}^2)d\widehat{k}=d𝐤/2k^0`$ which finally yields
$$𝒟_0^+(\widehat{x})=\frac{1}{i(2\pi )^{3/2}}\frac{d𝐤}{2k_0}\text{e}^{i\widehat{k}\widehat{x}}\frac{1}{(2\pi )^{3/2}}𝑑𝐤\text{e}^{i\mathrm{𝐤𝐱}}=\delta (𝐱).$$
(29)
Then we obtain the final expression for the amplitude of the teleported state (in the non-relativistic case the amplitude actually coincides with the wave function)
$$𝒜_{1s}^{}(d\theta )\frac{1}{2}𝑑𝐱^{}𝑑𝝃^{}𝑑𝝃\text{e}^{i(𝝃𝐏𝝃^{}𝐐)}\delta (𝝃𝐱^{})\delta (𝝃^{}𝝃𝐗)$$
(30)
$$\delta (𝐱𝝃𝐘)\left(\delta _{s,+}f(𝐱^{^{}},+)+\delta _{s,}f(𝐱^{^{}},)\right)=$$
$$\left(\delta _{s,+}f(𝐱𝐘,+)+\delta _{s,}f(𝐱𝐘,)\right)\text{e}^{i(𝐱𝐘)(𝐏𝐐)+i\mathrm{𝐗𝐐}},$$
which coincides (to within a unitary transformation) with the original completely unknown state to be teleported.
Thus, if the photon state to be teleported is completely unknown, the teleportation as it is understood in the non-relativistic quantum mechanics of distinguishable particles becomes impossible and one can only speak of the amplitude of the propagation (taking into account the measurement procedure) as whole.
However, if the state to be teleported is only partly unknown, and one has only to teleport several degrees of freedom (for example, only the polarization state), the rest degrees of freedom can be used to “label” the individual particles to make them effectively distinguishable in the teleportation.
To be more precise, we mean the following. Suppose that it is known in advance that the single-photon state to be teleported has a specified momentum $`𝐤_1`$ and only the polarization state is unknown. Such a state can be written as
$$𝝍^+(𝐟)=\frac{1}{\sqrt{2k_1^0}}\left(f(+)𝐀^+(𝐤_1,+)+f()𝐀^+(𝐤_1,)\right),$$
(31)
$$|f(+)|^2+|f()|^2=1,$$
where $`f(\pm )`$ are the amplitudes for different polarizations $`\pm `$. The phase factors are omitted as unimportant. A plane wave is infinitely extended in space and it would be more correct to consider a ray-like state. However, for our purposes the difference between these states is insignificant and we shall use the simplest plane wave state.
The EPR-state entangled with respect to the polarization degrees of freedom is described by the creation operator
$$𝚿^{()+}(𝐤_2,𝐤_3)=\frac{1}{\sqrt{2}\sqrt{2k_2^0}\sqrt{2k_3^0}}\left(𝐀^+(𝐤_2,+)𝐀^+(𝐤_3,)𝐀^+(𝐤_2,)𝐀^+(𝐤_3,+)\right),$$
(32)
where the wave vectors $`𝐤_2`$ and $`𝐤_3`$ are also known beforehand. In the experiments this state is associated with a pair of photons in an entangled state leaving a non-linear crystal and propagating along the vectors $`𝐤_2`$ and $`𝐤_3`$ .
The required measurement is described by the identity resolution (22). It is more convenient here to rewrite it in the momentum representation choosing $`\mathrm{\Theta }=`$ $`\{s,i,𝐤_1^{^{}},`$ $`𝐤_2^{^{}},𝐤_3^{^{}}:s\times i\times 𝐤_1^{^{}}\times 𝐤_2^{^{}}\times 𝐤_3^{^{}}\}`$ as the space of possible outcomes (here $`i`$ labels different states in the Bell basis ($`i=1÷4`$),
$$𝚿^{(\pm )+}(𝐤_2^{^{}},𝐤_3^{^{}})=\frac{1}{\sqrt{2}\sqrt{2k_2^0}\sqrt{2k_3^0}}\left(𝐀^+(𝐤_2^{^{}},+)𝐀^+(𝐤_3^{^{}},)\pm 𝐀^+(𝐤_2^{^{}},)𝐀^+(𝐤_3^{^{}},+)\right),$$
(33)
$$𝚽^{(\pm )+}(𝐤_2^{^{}},𝐤_3^{^{}})=\frac{1}{\sqrt{2}\sqrt{2k_2^0}\sqrt{2k_3^0}}\left(𝐀^+(𝐤_2^{^{}},+)𝐀^+(𝐤_3^{^{}},+)\pm 𝐀^+(𝐤_2^{^{}},)𝐀^+(𝐤_3^{^{}},)\right).$$
The identity resolution itself is
$$I_3=\underset{s=\pm }{}\frac{d𝐤_1^{^{}}}{\sqrt{2k_1^0^{}}}\frac{d𝐤_2^{^{}}}{\sqrt{2k_2^0^{}}}\frac{d𝐤_3^{^{}}}{\sqrt{2k_3^0^{}}}$$
(34)
$$\{𝐀^+(𝐤_1^{^{}},s)(𝚿^{(+)+}(𝐤_2^{^{}},𝐤_3^{^{}})|00|𝚿^{(+)}(𝐤_2^{^{}},𝐤_3^{^{}})+𝚿^{()+}(𝐤_2^{^{}},𝐤_3^{^{}})|00|𝚿^{()}(𝐤_2^{^{}},𝐤_3^{^{}})+$$
$$𝚽^{(+)+}(𝐤_2^{^{}},𝐤_3^{^{}})|00|𝚽^{(+)}(𝐤_2^{^{}},𝐤_3^{^{}})+𝚽^{()+}(𝐤_2^{^{}},𝐤_3^{^{}})|00|𝚽^{()}(𝐤_2^{^{}},𝐤_3^{^{}}))𝐀^{}(𝐤_1^{^{}},s)\}$$
The initial state (an EPR-pair and a photon in the unknown polarization state) is written as
$$|𝝍(𝐟)𝚿^{()}(𝐤_2,𝐤_3)=𝝍^+(𝐟)𝚿^{()+}(𝐤_2,𝐤_3)|0=$$
(35)
$$\{𝚿^{()+}(𝐤_1,𝐤_2)(f(+)𝝍^+(𝐤_3,+)+f()𝝍^+(𝐤_3,))+$$
$$𝚿^{(+)+}(𝐤_1,𝐤_2)(f(+)𝝍^+(𝐤_3,+)+f()𝝍^+(𝐤_3,))+$$
$$𝚽^{()+}(𝐤_1,𝐤_2)(f(+)𝝍^+(𝐤_3,)+f()𝝍^+(𝐤_3,+))+$$
$$𝚽^{(+)+}(𝐤_1,𝐤_2)(f(+)𝝍^+(𝐤_3,)f()𝝍^+(𝐤_3,+))\}|0.$$
The fact that the momenta of all the particles are known allows to preform the post-selection after the measurement, i.e. to discard all the outcomes except for those in which the total momentum of the measured particles $`𝐤_1^{^{}}+𝐤_2^{^{}}=𝐤_2+𝐤_2`$, and the teleported state is selected by the condition $`𝐤_3^{^{}}=𝐤_3`$. In other words, kept in the two-particle channel of the measurement are only the outcomes where the total momentum is equal to the sum of the momentum of one of the particles of the original EPR-pair and the momentum of the particle in the unknown polarization state.
The probabilities of different outcomes after the post-selection in the outcome space in the vicinity of $`𝐤_1^{^{}}+𝐤_2^{^{}}=𝐤_1+𝐤_2`$ and $`𝐤_3^{^{}}=𝐤_3`$) are determined by a formula similar to Eqs. (23–24); for example, in the channel associated with the projection on the state $`𝚿^{()}(𝐤_1,𝐤_2)`$, we obtain for the amplitude
$$𝒜_s^{}(𝐤_1,𝐤_2,𝐤_3)=0|𝚿^{()}𝐤_1,𝐤_2)𝝍^{}(𝐤_3,s)𝝍^+(𝐟)𝚿^{()+}𝐤_2,𝐤_3)|0=$$
(36)
$$\frac{1}{2(2k_1^02k_2^02k_3^0)}$$
$$\{0|𝚿^{()}(𝐤_1,𝐤_2)𝝍^{}(𝐤_3,s)𝚿^{()+}(𝐤_1,𝐤_2)(f(+)𝝍^+(𝐤_3,+)+f()𝝍^+(𝐤_3,))|0+$$
$$0|𝚿^{()}(𝐤_1,𝐤_2)𝝍^{}(𝐤_3,s)𝚿^{()+}(𝐤_1,𝐤_2)\left(f(+)𝝍^+(𝐤_3,+)+f()𝝍^+(𝐤_3,)\right)|0+$$
$$0|𝚿^{()}(𝐤_1,𝐤_2)𝝍^{}(𝐤_3,s)𝚽^{()+}(𝐤_1,𝐤_2)\left(f(+)𝝍^+(𝐤_3,)+f()𝝍^+(𝐤_3,+)\right)|0+$$
$$0|𝚿^{()}(𝐤_1,𝐤_2)𝝍^{}(𝐤_3,s)𝚽^{(+)+}(𝐤_1,𝐤_2)(f(+)𝝍^+(𝐤_3,)f()𝝍^+(𝐤_3,+))|0\}.$$
Finally for the propagation amplitude one obtains
$$𝒜_s^{}(𝐤_1,𝐤_2,𝐤_3)0|𝐀^{}(s,𝐤_3)\left(f(+)𝐀^+(𝐤_3,+)+f()𝐀^+(𝐤_3,)\right)|0,$$
(37)
which actually coincides with the amplitude of propagation of the state with the initial unknown polarization along the direction defined by $`𝐤_3`$ instead of $`𝐤_1`$. A similar situation takes place for different outcomes in the other channels where the teleported state coincides to within a unitary transformation completely determined by the channel number (Bell basis vector number) with the unknown state which had to be teleported.
Thus, if the state to be teleported is completely unknown, one can only speak of the total amplitude of the propagation of the field as a whole. On the other hand, if the state is only partly unknown, the available a priori information can be used to “label” the particles to make the identical particles effectively distinguishable.
This work was supported by the Russian Foundation for Basic Research (project No 99-02-18127), the project “Physical Principles of the Quantum Computer”, and the program “Advanced Devices and Technologies in Micro- and Nanoelectronics”.
This work was also supported by the Wihuri Foundation, Finland.
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# A superspace gauge-invariant formulation of a massive tridimensional 2-form field
## I Introduction
Antisymmetric tensor fields appear in many field theories. In particular, the Kalb-Ramond gauge field plays an important role in strong-weak coupling dualities among string theories and in axionic cosmic strings . On the other hand, a first order formulation of the non-Abelian Yang-Mills gauge theory ( BF-YM model) makes use of a two form gauge potential $`B`$ to contribute to a discussion of the problem of quark confinement in continuum QCD . Another interesting aspect of the (3+1) dimensional $`BF`$ term ($`F=dA`$ is the field strength of a one form gauge potential $`A`$) is its ability to give rise to gauge invariant mass to the gauge field . This property has been used to obtain an axion field topologically massive and an axionic charge on a black hole as well . In addition, the existence of the Higgs mechanism to the Kalb-Ramond gauge fields was demonstrated by S.-J. Key in the context of closed strings. On the other hand, if coupled to open strings, the KB field becomes a massive vector field through the Stückelberg mechanism. Also, we can mention a topologically massive Kalb-Ramond field in a $`D=3`$ context that was introduced in ref. .
It is known that massless string excitations may be described by a low-energy supergravity theory and that a massless gravity supermultiplet of graviton, dilaton and Kalb-Ramond fields appears in all known string theories. However, the spectrum of the $`D=4`$ and $`D=3`$ compactified theory from $`D=10`$ supergravity, contains the massive antisymmetric tensor fields. Thus, since supersymmetry places severe constraints on the ground state and the mass spectrum of the excitations, supersymmetric mechanisms of mass generation are of considerable importance.
The purpose of this letter is twofold. First we construct an $`N=1`$ $`D=4`$ superspace version of the $`U(1)`$ BF model. By means of a dimensional reduction procedure, we obtain a massive antisymmetric tensor field into a $`N=2D=3`$ supersymmetric topological massive gauge invariant theory. In contrast to several works on $`D=3`$ BF models, we have considered here a topological term which involves a KB and a pseudoscalar field with derivative coupling. Secondly, we have addressed a $`N=1`$ superspace mechanism to generate mass for Kalb-Ramond field without loss of gauge invariance. Actually, this mechanism is a superspace version of the topological massive formulation of Deser, Jackiw and Templeton . On the other hand, an alternative model with an explicit mass breaking term is constructed in $`N=1`$ superspace and a supersymmetric version of the Stückelberg transformation is used to restore the gauge invariance of the model.
## II The $`N=1D=4`$ Extended BF Model
Let us begin by introducing the $`N=1D=4`$ supersymmetric BF extended model. For extended we mean that we include mass terms for the Kalb-Ramond field. This mass term will be introduced here for later comparison to the tridimensional case. Actually, this construction can be seen as a superspace and Abelian version of the so called BF-Yang-Mills models .
As our basic superfield action we takeOur spinorial notations and other conventions follow ref. .
$$S_{BF}^{SS}=\frac{1}{8}d^4x\{i\kappa [d^2\theta B^\alpha W_\alpha d^2\overline{\theta }\overline{B}_{\dot{\alpha }}\overline{W}^{\dot{\alpha }}]+\frac{g^2}{2}[d^2\theta B^\alpha B_\alpha +d^2\overline{\theta }\overline{B}_{\dot{\alpha }}\overline{B}^{\dot{\alpha }}]\}\text{ }.$$
(1)
where $`W_\alpha `$ is a spinor superfield-strenght, $`B_\alpha `$ is a chiral spinor superfield, $`\overline{D}_{\dot{\alpha }}B_\alpha =0`$, $`\kappa `$ and $`g`$ are massive parameters. Their corresponding $`\theta `$-expansions are:
$`W_\alpha (x,\theta ,\overline{\theta })`$ $`=`$ $`4i\lambda _\alpha (x)[4\delta _\alpha ^\beta D(x)+2i(\sigma ^\mu \overline{\sigma }^\nu )_\alpha ^\beta F_{\mu \nu }(x)]\theta _\beta `$ (3)
$`+4\theta ^2\sigma _{\alpha \dot{\alpha }}^\mu _\mu \overline{\lambda }^{\dot{\alpha }}`$
$$B_\alpha (x,\theta ,\overline{\theta })=e^{i\theta \sigma ^\mu \overline{\theta }_\mu }[i\psi _\alpha (x)+\theta ^\beta T_{\alpha \beta }(x)+\theta \theta \xi _\alpha (x)]\text{ ,}$$
(4)
where
$$T_{\alpha \beta }=T_{(\alpha \beta )}+T_{[\alpha \beta ]}=4i(\sigma ^{\mu \nu })_{\alpha \beta }B_{\mu \nu }+2\epsilon _{\alpha \beta }(M+iN)\text{ }.$$
(5)
Our conventions for supersymmetric covariant derivatives are
$`D_\alpha `$ $``$ $`{\displaystyle \frac{}{\theta ^\alpha }}+i\sigma _{\alpha \dot{\alpha }}^\mu \overline{\theta }^{\dot{\alpha }}_\mu `$ (6)
$`\overline{D}_{\dot{\alpha }}`$ $``$ $`{\displaystyle \frac{}{\overline{\theta }^{\dot{\alpha }}}}i\theta ^\alpha \sigma _{\alpha \dot{\alpha }}^\mu _\mu \text{ .}`$ (7)
We call attention for the electromagnetic field-strenght and the antisymmetric gauge field which are contained in $`W_\alpha `$ and $`B_\alpha `$, respectively. In terms of the components fields, the action (1) can be read as
$`S`$ $`=`$ $`{\displaystyle }d^4x\{[{\displaystyle \frac{i\kappa }{2}}(\xi \lambda \overline{\xi }\overline{\lambda })+{\displaystyle \frac{\kappa }{2}}(\psi ^\alpha \sigma _{\alpha \dot{\alpha }}^\mu _\mu \overline{\lambda }^{\dot{\alpha }}+\overline{\psi }_{\dot{\alpha }}\left(\overline{\sigma }^\mu \right)^{\dot{\alpha }\alpha }_\mu \lambda _\alpha )+{\displaystyle \frac{\kappa }{2}}B^{\mu \nu }\stackrel{~}{F}_{\mu \nu }`$ (9)
$`\kappa DN]+g^2[{\displaystyle \frac{1}{8}}(\psi \xi +\overline{\psi }\overline{\xi })+{\displaystyle \frac{1}{2}}B^{\mu \nu }B_{\mu \nu }{\displaystyle \frac{1}{2}}(M^2+N^2)]\}`$
$`=`$ $`{\displaystyle }d^4x[({\displaystyle \frac{i\kappa }{2}}\overline{\mathrm{\Xi }}\gamma ^5\mathrm{\Lambda }+{\displaystyle \frac{\kappa }{2}}\overline{\mathrm{\Psi }}\gamma ^\mu _\mu \mathrm{\Lambda }+{\displaystyle \frac{\kappa }{2}}B^{\mu \nu }\stackrel{~}{F}_{\mu \nu }\kappa DN)`$ (11)
$`+g^2({\displaystyle \frac{1}{8}}\overline{\mathrm{\Psi }}\mathrm{\Xi }+{\displaystyle \frac{1}{2}}B^{\mu \nu }B_{\mu \nu }{\displaystyle \frac{1}{2}}(M^2+N^2))]\text{ }.`$
In the last equality above, the fermionic fields have been organized as four-component Majorana spinors as follows
$$\mathrm{\Xi }=\left(\begin{array}{c}\xi _\alpha \\ \overline{\xi }^{\dot{\alpha }}\end{array}\right)\text{ };\text{ }\mathrm{\Lambda }=\left(\begin{array}{c}\lambda _\alpha \\ \overline{\lambda }^{\dot{\alpha }}\end{array}\right)\text{ };\text{ }\mathrm{\Psi }=\left(\begin{array}{c}\psi _\alpha \\ \overline{\psi }^{\dot{\alpha }}\end{array}\right)\text{ },$$
(12)
and we denote the dual field-strenght defining $`\stackrel{~}{F}_{\mu \nu }\frac{1}{2}\epsilon _{\mu \nu \alpha \beta }F^{\alpha \beta }`$. Furthermore, we use the following identities
$`\overline{\mathrm{\Psi }}\mathrm{\Lambda }`$ $`=`$ $`\overline{\psi }\overline{\lambda }+\psi \lambda `$ (13)
$`\overline{\mathrm{\Psi }}\gamma ^5\mathrm{\Lambda }`$ $`=`$ $`\overline{\psi }\overline{\lambda }\psi \lambda `$ (14)
$`\overline{\mathrm{\Psi }}\gamma ^\mu \mathrm{\Lambda }`$ $`=`$ $`\psi \sigma ^\mu \overline{\lambda }+\overline{\psi }\overline{\sigma }^\mu \lambda \text{ .}`$ (15)
The superfield action (1) is a particular case of the action proposed in ref. . However, a point of difference must be noted. In contrast with , we have not considered coupling with matter fields and a propagation term for the gauge fields. On the other hand, our supespace BF term was constructed in a distinct and simpler way. A quite similar construction was introduced by Clark et al. .
The off-diagonal mass term $`\xi \lambda `$ (or $`\overline{\mathrm{\Xi }}\gamma ^5\mathrm{\Lambda }`$) has been shown by Brooks and Gates, Jr. in the context of super-Yang-Mills theory. Note that the identity
$$\gamma _5\sigma ^{\mu \nu }=\frac{i}{2}\epsilon _{\mu \nu \alpha \beta }\sigma ^{\alpha \beta }$$
(16)
reveals a connection between the topological behaviour denoted by the Levi-Civita tensor $`\epsilon _{\mu \nu \alpha \beta },`$ and the pseudo-escalar $`\gamma _5.`$
So, it is worthwhile to mention that this term has topological origin and it can be seen as a fermionic counterpart of the BF term. In our opinion, this fermionic mass term deserves more attention and will be investigated elsewhere.
## III The $`N=2D=3`$ Topological Model
As it is well known, the BF model in $`D=3`$ consists in a one form field ($`\mathrm{"}B\mathrm{"}`$ field) and one form gauge field $`A`$. So, the Chern-Simons term is simply the identification of $`B`$ and $`A`$. However, as has been shown in ref. , after dimensional reduction of the four dimensional BF model, an interesting additional term arises, namely, a topological term which involves a 2-form and a 0-form. We will call it a $`B\phi `$ term. A quite similar model was presented in a Yang-Mills version by Del Cima et. al. , and its finiteness was proved in the framework of algebraic renormalization.
Following the procedure of ref. , we will carry out a dimensional reduction in the bosonic sector of (11). Dimensional reduction is usually done by expanding the fields in normal modes corresponding to the compactified extra dimensions, and integrating out the extra dimensions. This approach is very useful in dual models and superstrings . Here, however, we only consider the fields in higher dimensions to be independent of the extra dimensions. In this case, we assume that our fields are independent of the extra coordinate $`x_3.`$
Therefore, after dimensional reduction, the bosonic sector of (11) can be written as
$`S_{bos.}`$ $`=`$ $`{\displaystyle }d^3x\{[\kappa \epsilon _{\mu \alpha \beta }V^\mu F^{\alpha \beta }+\kappa \epsilon _{\mu \nu \alpha }B^{\mu \nu }^\alpha \phi \kappa DN]`$ (18)
$`+g^2[{\displaystyle \frac{1}{2}}B^{\mu \nu }B_{\mu \nu }V^\mu V_\mu {\displaystyle \frac{1}{2}}(M^2+N^2)]\}\text{ }.`$
Notice that the first term in r.h.s. of (18) can be transformed in the Chern-Simons term if we identify $`V^\mu A^\mu `$. The second one is the so called $`B\phi `$ term.
Now let us proceed to the dimensional reduction of the fermionic sector of the model. First, note that the Lorentz group in three dimensions is $`SL(2,R)`$ rather than $`SL(2,C)`$ in $`D=4`$. Therefore, Weyl spinors with four degrees of freedom will be mapped into Dirac spinorsFor details about spinorial dimensional reduction, we suggest refs. and .. So the correct associations keeping the degrees of freedom are sketched as
$`\mathrm{\Xi }`$ $`=`$ $`\left(\begin{array}{c}\xi _\alpha \\ \overline{\xi }^{\dot{\alpha }}\end{array}\right)\mathrm{\Xi }_\pm =\xi _\alpha \pm i\tau _\alpha `$ (21)
$`\mathrm{\Lambda }`$ $`=`$ $`\left(\begin{array}{c}\lambda _\alpha \\ \overline{\lambda }^{\dot{\alpha }}\end{array}\right)\mathrm{\Lambda }_\pm =\lambda _\alpha \pm i\rho _\alpha `$ (24)
$`\mathrm{\Psi }`$ $`=`$ $`\left(\begin{array}{c}\psi _\alpha \\ \overline{\psi }^{\dot{\alpha }}\end{array}\right)\mathrm{\Psi }_\pm =\psi _\alpha \pm i\chi _\alpha \text{ }.`$ (27)
From (27), we find that
$`\mathrm{\Psi }\overline{\mathrm{\Xi }}`$ $``$ $`{\displaystyle \frac{1}{2}}\left(\mathrm{\Psi }_+\mathrm{\Xi }_{}+\mathrm{\Psi }_{}\mathrm{\Xi }_+\right)`$ (28)
$`\overline{\mathrm{\Psi }}\gamma ^\mu _\mu \mathrm{\Lambda }`$ $``$ $`{\displaystyle \frac{1}{2}}(\mathrm{\Psi }_+\gamma ^{\widehat{\mu }}_{\widehat{\mu }}\mathrm{\Lambda }_{}+\mathrm{\Psi }_{}\gamma ^{\widehat{\mu }}_{\widehat{\mu }}\mathrm{\Lambda }_+)`$ (29)
$`\mathrm{\Xi }\gamma ^5\mathrm{\Lambda }`$ $``$ $`{\displaystyle \frac{1}{2}}(\mathrm{\Xi }_+\mathrm{\Lambda }_++\mathrm{\Xi }_{}\mathrm{\Lambda }_{})\text{ }.`$ (30)
where $`hatted`$ index means three-dimensional space-time.
Thus, the dimensionally reduced fermionic sector of (11) may be written
$`S_{ferm.}`$ $`=`$ $`{\displaystyle }d^3x\{{\displaystyle \frac{i\kappa }{4}}(\mathrm{\Xi }_+\mathrm{\Lambda }_++\mathrm{\Xi }_{}\mathrm{\Lambda }_{})+{\displaystyle \frac{\kappa }{4}}(\mathrm{\Psi }_+\gamma ^{\widehat{\mu }}_{\widehat{\mu }}\mathrm{\Lambda }_{}+\mathrm{\Psi }_{}\gamma ^{\widehat{\mu }}_{\widehat{\mu }}\mathrm{\Lambda }_+)`$ (32)
$`+{\displaystyle \frac{g^2}{16}}(\mathrm{\Psi }_+\mathrm{\Xi }_{}+\mathrm{\Psi }_{}\mathrm{\Xi }_+)\}\text{ .}`$
The action $`S=S_{bos.}+S_{ferm.}`$ is invariant under the following supersymmetry transformations
$`\delta \lambda _\alpha `$ $`=`$ $`iD\eta _\alpha \left(\sigma ^\mu \sigma ^\nu \right)_\alpha ^\beta \eta _\beta F_{\mu \nu }`$ (33)
$`\delta \rho _\alpha `$ $`=`$ $`iD\zeta _\alpha \left(\sigma ^\mu \sigma ^\nu \right)_\alpha ^\beta \zeta _\beta F_{\mu \nu }`$ (34)
$`\delta F^{\mu \nu }`$ $`=`$ $`i^\mu \left(\eta \sigma ^\nu \rho \lambda \sigma ^\nu \zeta \right)i^\nu \left(\eta \sigma ^\mu \rho \lambda \sigma ^\mu \zeta \right)`$ (35)
$`\delta D`$ $`=`$ $`_\mu \left(\eta \sigma ^\mu \rho +\lambda \sigma ^\mu \zeta \right)`$ (36)
$`\delta \left(\psi _\alpha \pm i\chi _\alpha \right)`$ $`=`$ $`\delta \mathrm{\Psi }_\pm =i\eta ^\beta \stackrel{~}{T}_{\beta \alpha }\pm \zeta ^\beta \stackrel{~}{T}_{\beta \alpha }`$ (37)
$`\delta \stackrel{~}{T}_{\beta \alpha }`$ $`=`$ $`\eta _\beta \xi _\alpha +\zeta ^\lambda \sigma _{\beta \lambda }^\mu _\mu \psi _\alpha `$ (38)
$`\delta \left(\xi _\alpha \pm i\tau _\alpha \right)`$ $`=`$ $`\delta \mathrm{\Xi }_\pm =i\zeta _\lambda \left(\sigma ^\mu \right)^{\lambda \beta }T_{\beta \alpha }\eta _\lambda \left(\overline{\sigma }^\mu \right)^{\beta \lambda }T_{\beta \alpha }\text{ ,}`$ (39)
where $`\eta `$ and $`\zeta `$ are supersymmetric parameters, which indicates that we have two supersymmetries in the aforementioned action.
## IV Remarks on Some 3D Supersymmetric Models and Stückelberg Formulation
From the two topological terms introduced in (18) we can set up two supersymmetric models. The first one, which involves a two and a zero form, can be expressed as
$$S=d^3xd^2\theta (D^\alpha \mathrm{\Phi }B_\alpha +\frac{1}{2}g^2B^\alpha B_\alpha )\text{ },$$
(40)
where $`B_\alpha `$ and $`\mathrm{\Phi }`$ are spinor and real scalar superfields, which are defined by projection as
$`B_\alpha `$ $``$ $`=\chi _\alpha `$ (41)
$`D_{(\beta }B_{\alpha )}`$ $``$ $`=2iM_{\beta \alpha }=M_{\alpha \beta }=B^{\mu \nu }\left(\sigma _{\mu \nu }\right)_{\alpha \beta }`$ (42)
$`D^\alpha B_\alpha `$ $``$ $`=2N`$ (43)
$`D^\beta D_\alpha B_\beta `$ $``$ $`=2\omega _\alpha `$ (44)
and
$`\mathrm{\Phi }`$ $``$ $`=\phi `$ (45)
$`D_\alpha \mathrm{\Phi }`$ $``$ $`=\psi _\alpha `$ (46)
$`D^2\mathrm{\Phi }`$ $``$ $`=F\text{ }.`$ (47)
Here the supersymmetry covariant derivative is given by $`D_\alpha =_\alpha +i\theta ^\beta _{\alpha \beta }`$ . So, in terms of components fields, the action (40) becomes
$`S`$ $`=`$ $`{\displaystyle }d^3x[(\kappa ^{\alpha \beta }\phi M_{\beta \alpha }+2\kappa \psi ^\alpha \omega _\alpha 2\kappa FN)`$ (49)
$`+{\displaystyle \frac{1}{2}}g^2(4\omega ^\alpha \chi _\alpha +2i\chi _\alpha ^{\beta \alpha }\chi _\beta +M^{\beta \alpha }M_{\alpha \beta }+2N^2)]\text{ }.`$
Starting from the definitions of two spinor superfields given by
$`\mathrm{\Lambda }_\alpha `$ $``$ $`=\xi _\alpha `$ (50)
$`D_{(\beta }\mathrm{\Lambda }_{\alpha )}`$ $``$ $`=2iV_{\beta \alpha }`$ (51)
$`D^\alpha \mathrm{\Lambda }_\alpha `$ $``$ $`=2G`$ (52)
$`D^\beta D_\alpha \mathrm{\Lambda }_\beta `$ $``$ $`=2\rho _\alpha `$ (53)
and
$`W_\alpha `$ $``$ $`=\lambda _\alpha `$ (54)
$`D_\alpha W_\beta `$ $``$ $`=f_{\alpha \beta }\text{ },`$ (55)
where
$$V_{\beta \alpha }V^\mu \left(\stackrel{~}{\sigma }_\mu \right)_{\beta \alpha }\text{ };\text{ }f_{\alpha \beta }\left(\stackrel{~}{\sigma }_\mu \right)_{\alpha \beta }f^\mu \text{ };\text{ }f^\mu =\frac{i}{2}\epsilon ^{\mu \nu \rho }F_{\nu \rho }\text{ },$$
(56)
we can propose another supersymmetric action, now involving two 1-forms, namely
$`S`$ $`=`$ $`{\displaystyle d^3xd^2\theta (\mathrm{\Lambda }^\alpha W_\alpha g^2\mathrm{\Lambda }^\alpha \mathrm{\Lambda }_\alpha )}`$ (57)
$`=`$ $`{\displaystyle }d^3x[(2\rho ^\alpha \lambda _\alpha iV^{\alpha \beta }f_{\beta \alpha })`$ (59)
$`g^2(4\rho ^\alpha \omega _\alpha +2i\xi _\alpha ^{\beta \alpha }\xi _\beta +V^{\beta \alpha }V_{\beta \alpha }+2G^2)]\text{ .}`$
It is easy to see that the superspace actions (40) and (59) are not invariant under the following gauge transformations
$`\delta B^\alpha `$ $`=`$ $`D^\beta D^\alpha \mathrm{\Pi }_\beta `$ (60)
$`\delta \mathrm{\Phi }`$ $`=`$ $`0`$ (61)
$`\delta \mathrm{\Lambda }^\alpha `$ $`=`$ $`D^\alpha \mathrm{\Omega }`$ (62)
$`\delta W^\alpha `$ $`=`$ $`0\text{ }.`$ (63)
However, if we reparameterize $`\mathrm{\Lambda }^\alpha `$ and $`B^\alpha `$ through introduction of the Stückelberg superfields<sup>§</sup><sup>§</sup>§For historical reasons, it is important to cite here the first work, to the best of our knowledge, in the framework of supersymmetric Stückelberg formalism, namely ref. . $`\mathrm{\Theta }`$ and $`\mathrm{\Sigma }_\alpha `$ such that
$`\mathrm{\Lambda }^\alpha `$ $``$ $`\left(\mathrm{\Lambda }^\alpha \right)^{}=\mathrm{\Lambda }^\alpha +{\displaystyle \frac{1}{g}}D^\alpha \mathrm{\Theta }`$ (64)
$`B^\alpha `$ $``$ $`\left(B^\alpha \right)^{}=B^\alpha +D^\beta D^\alpha \mathrm{\Pi }_\beta \text{ },`$ (65)
and imposing that $`\mathrm{\Theta }`$ and $`\mathrm{\Sigma }_\alpha `$ transform like
$`\delta \mathrm{\Theta }`$ $`=`$ $`g\mathrm{\Omega }`$ (66)
$`\delta \mathrm{\Sigma }^\beta `$ $`=`$ $`\mathrm{\Pi }^\beta \text{ ,}`$ (67)
we ensure gauge invariance for that superactions.
We remark that integrating out the superfield $`B_\alpha `$ in the equation (40) we arrive at a supersymmetric Klein-Gordon action and, if we do the same for $`\mathrm{\Lambda }_\alpha `$ in (59), we obtain a Maxwell superaction. Observe that both these relations may be understood as two duality tranformations. We recall here that an analogous connection in $`4D`$ pure bosonic BF-theory was viewed as a perturbative expansion in the coupling $`g`$ around the topological pure BF theory . Thereupon, it may be interesting to perform a similar investigation in the framework of action (40).
## V $`N=1`$ Superspace Topological Mass Generation
In order to show the topological mass generation for the Kalb-Ramond two form field, we will construct a variation from the model (40), by introducing the propagation term for it. Before that, for ilustration purpose, we quote the bosonic action introduced in ref. :
$$S=d^3x\left[\frac{1}{6}H_{\mu \nu \rho }H^{\mu \nu \rho }+kϵ_{\mu \nu \rho }B^{\mu \nu }^\rho \varphi +\frac{1}{2}_\mu \varphi ^\mu \varphi \right]\text{ },$$
(68)
where $`H_{\mu \nu \rho }`$, a three form field-strength of the $`B^{\mu \nu }`$ field, is defined as
$$H_{\mu \nu \rho }=_{[\mu }B_{\mu \rho ]}=_\mu B_{\nu \rho }+_\nu B_{\rho \mu }+_\rho B_{\mu \nu }\text{ }.$$
(69)
The $`N=1`$ superspace construction of the supersymmetric version of (68) proceeds as follows. First, we introduce a scalar superfield $`G`$ defined by
$$G=D^\alpha B_\alpha \text{ },$$
(70)
where $`B_\alpha `$ is the super-Kalb-Ramond field defined in (44). Then, after looking the expression (40), we find the action
$$S=d^3xd^2\theta [\frac{1}{2}\left(D^\alpha G^2\right)+kB^\alpha D_\alpha \mathrm{\Phi }\frac{1}{2}D^\alpha \mathrm{\Phi }D_\alpha \mathrm{\Phi }]\text{ }.$$
(71)
Now it is straightforward to show that the topological term $`kB^\alpha D_\alpha \mathrm{\Phi }`$ gives rise to a mass term for the super-Kalb-Ramond field. The equation of motion associated with $`\mathrm{\Phi }`$ is,
$$D^\alpha \left(kB_\alpha D_\alpha \mathrm{\Phi }\right)=0\text{ }.$$
(72)
Consequently,
$$kB_\alpha D_\alpha \mathrm{\Phi }=𝒞\text{ .}$$
(73)
Since that the constant $`𝒞`$ can be absorbed by $`B_\alpha `$, we conclude that
$$kB_\alpha D_\alpha \mathrm{\Phi }=0\text{ .}$$
(74)
Therefore the original action (71) can be rewritten as
$$S=d^3xd^2\theta [\left(D^\alpha G^2\right)+\frac{1}{2}k^2B^\alpha B_\alpha ]\text{ }.$$
(75)
This exhibits a topological mechanism of mass generation for the Kalb-Ramond field. Naturally, the topological mass terms arise due to the coupling of the $`B_\alpha `$ and $`\mathrm{\Phi }`$ superfields. In other words, this mass term results of the breakdown of the gauge invariance (61).
Incidentally let us mention a possible equivalence similar to that between massive topologically and self-dual theories in $`D=3`$ . Indeed, starting from (40), we can construct an action by introduction of a mass term for the superfield $`\mathrm{\Phi },`$ namely
$$S=d^3xd^2\theta (D^\alpha \mathrm{\Phi }B_\alpha +\frac{1}{2}g^2B^\alpha B_\alpha +m\mathrm{\Phi }^2)\text{ }.$$
(76)
It is easy to see that the equations of motion of (76) and (71) are equivalent. So, the action (76) can be considered locally equivalent to action (71). On the other hand, it would be interesting to investigate if this equivalence is preserved at quantum level.
## VI Conclusions
In this letter, we have constructed an $`N=1`$ $`D=3`$ superspace action for a model involving an antisymmetric gauge field. Our main point is a topological term that consists in a coupling of this 2-form field and a scalar field. To the best of our knowledge, in the form presented here, this model is completely new in the literature. A similar approach, but involving a 3-form and a scalar fields in $`N=1D=4`$, was introduced in ref. .
Starting from the so called $`BF`$ model in $`N=1D=4`$ superspace, we carried out a dimensional reduction to the three-dimensional space-time, in order to obtain our basic model. The superspace construction for the $`BF`$ is known, but we point out the appearance of a fermionic counterpart of the $`BF`$ term.
We have introduced two massive gauge invariant models for an antisymmetric tensor field into a $`N=1D=3`$ superspace. In the first, we resort to the Stückelberg formalism and in the other, we construct an abelian topologically massive theory, and a topologically generated mass for the Kalb-Ramon superfield is exhibited. An equivalence of both massive models is suggested. Furthermore, a component field analysis is performed, showing a second supersymmetry in the model.
ACKNOWLEDGMENTS
We wish to thank M. S. Cunha for helpful discussions. Prof. J. E. Moreira is acknowledged by a critical reading of the manuscript. This work was supported in part by Conselho Nacional de Desenvolvimento Científico e Tecnológico-CNPq and Fundação Cearense de Amparo à Pesquisa-FUNCAP.
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# Restoration of Macroscopic Isotropy on (𝑑+1)-Simplex Fractal Conductor Networks
## 1 INTRODUCTION
Restoration of isotropy in an anisotropic system is of great interest in a variety of disciplines where much attention has been focused on it, particularly on the problem of diffusion in inhomogeneous materials . In general, diffusion on lattices can be formulated in terms of an AC electric problem and DC electric response in a percolating structures can be viewed as a very special case of diffusion in disordered medium . The purpose of this paper is to investigate the restoration of macroscopic isotropy in $`(𝐝+\mathrm{𝟏})`$-simplex fractal conductor networks with microscopic anisotropy. In general, deterministic fractal lattices, as proposed by Kirkpatrick, mimic some properties of percolation clusters in random media and disordered systems, and among fractal objects, the $`(𝐝+\mathrm{𝟏})`$-simplex fractal is the simplest one to study various physical problems from random walk to electrical problem on it. Using the exact renormalization technique based on the minimization of total dissipative power (TDP) in these networks, we present a rigorous proof that the conductivity becomes isotropic for large scales, and anisotropy vanishes with a scaling exponent $`\overline{\lambda }`$, as $`𝐋^{\overline{\lambda }}`$. We exactly compute $`\overline{\lambda }`$ for arbitrary values of $`𝐝`$ and decimation numbers $`𝐛=\mathrm{𝟐},\mathrm{𝟑},\mathrm{𝟒}`$ and $`\mathrm{𝟓}`$. The contents of this paper is as follows:
Section II presents a brief description of $`(𝐝+\mathrm{𝟏})`$-simplex fractals with decimation number $`𝐛`$ together with an explanation of labelling their subfractal and vertices with the partitions of positive integers, where this coding plays a very important role throughout the article. In section III we consider the most general network that can be built in a deterministic way, by putting circuit elements on the bonds of $`(𝐝+\mathrm{𝟏})`$-simplex fractal of a given generation $`𝐧`$ with decimation number $`𝐛`$. In order for the self-similarity of the structure to be preserved in the presence of anisotropy at microscopic level, the nature of the circuit elements, namely its resistances, must depend on the orientation of the bonds. It is clear that in $`(𝐝+\mathrm{𝟏})`$-simplex there are $`\frac{𝐝(𝐝+\mathrm{𝟏})}{\mathrm{𝟐}}`$ different orientations. Then we try to establish recursion equations for the connection resistances which represent the conductivity of these networks, on two successive length scales $`𝐋`$ and $`𝐋^{}=\mathrm{𝐛𝐋}`$. In general, these recursion relations are very involved. Fortunately, we do not need to have the explicit form of these recurrence equations, for the investigation of the restoration of isotropy. All we need here is the general properties of these maps, which can be obtained through some physical requirements and assumptions. It should be stressed that these circuits are not fictitious, since $`(𝐝+\mathrm{𝟏})`$-simplex fractals are embedible in Euclidean $`\mathrm{𝟐}`$-dimensions, hence they can be considered as two-dimensional networks, see Fig. 1. Section IV is devoted to a very rigorous proof of the uniqueness of the fixed point of the real space renormalization group transformation of the ratios of the connection resistances. Here in this section we show that all flows of the real space renormalization group transformation of the connection resistances, stemming from the finite physical region of connection resistance space, diverges to a direction which makes equal angle with all coordinates axes. The proof is based on some theorems and definitions of fixed point theory of the maps on complete metric spaces with the Hilbert metric. We have quoted the required theorems without presenting their proofs, since this section would be otherwise more mathematical in style. We refer the readers to reference for proofs of all theorems and for more details. Those readers who are only interested in the results of this section can skip it. In section V by minimizing the TDP in isotropic state, we get linear equations for the inner inward flowing currents in terms of input currents with Lagrange multipliers as their coefficients. Then using $`𝐒_{(𝐝+\mathrm{𝟏})}`$ symmetry group of the $`(𝐝+\mathrm{𝟏})`$-simplex, we suggest an ansatz for the Lagrange multipliers which leads to determination of the inner flowing currents in terms of the input one for any values of $`𝐝`$ and decimation number $`𝐛=\mathrm{𝟐},\mathrm{𝟑},\mathrm{𝟒}`$ and $`\mathrm{𝟓}`$. Section VI contains the main results of the article. Here in this section, by linearising the recurrence relation of the connection resistances near the isotropy state, we calculate power scaling exponent and the scaling exponent of the suppression of the anisotropy, for arbitrary values of $`𝐝`$ and decimation numbers $`𝐛=\mathrm{𝟐},\mathrm{𝟑},\mathrm{𝟒}`$ and $`𝐛=\mathrm{𝟓}`$, which are in agreement with the results of references in special cases. Also these results hold true for $`(𝐝+\mathrm{𝟏})`$-honeycomb fractal conductor network with decimation number $`𝐛=\mathrm{𝟐}`$, see Fig. 2, which is in agreemet with reference for $`𝐝=\mathrm{𝟐}`$ and $`𝐛=\mathrm{𝟐}`$ case. The paper ends with a brief conclusion.
## 2 (d+1)-Simplex Fractals
$`(𝐝+\mathrm{𝟏})`$-simplex fractal is a generalization of a two dimensional Sierpinski gasket to $`𝐝`$-dimensions such that its subfractals are $`(𝐝+\mathrm{𝟏})`$-simplices or $`𝐝`$-dimensional polyhedra with $`𝐒_{(𝐝+\mathrm{𝟏})}`$-symmetry. In order to obtain a fractal with decimation number $`𝐛`$, we choose a $`(𝐝+\mathrm{𝟏})`$-simplex and divide all the links (that is the lines connecting sites ) into $`𝐛`$ parts and then draw all possible $`𝐝`$-dimensional hyperplanes through the links parallel to the transverse $`𝐝`$-simplices. Next, having omitted every other innerpolyhedra, we repeat this process for the remaining simplices or for the subfractals of next higher generation. This way through $`(𝐝+\mathrm{𝟏})`$-simplex fractals are constructed. In order to calculate the fractal dimension, also to determine the current distribution, it is convenient to label subfractals of generation ($`n`$+1) in terms of partition of $`(𝐛\mathrm{𝟏})`$ into $`(𝐝+\mathrm{𝟏})`$ positive integers $`\lambda _\mathrm{𝟏},\lambda _\mathrm{𝟐},\mathrm{},\lambda _{𝐝+\mathrm{𝟏}}`$. Each partition represents a subfractal of generation $`𝐧`$, and $`\lambda `$ shows the distance of the corresponding subfractal from $`𝐝`$-dimensional hyper-planes which construct the $`(𝐝+\mathrm{𝟏})`$ simplex. On the other hand, each vertex denoted by partition of $`𝐛`$ into $`(𝐝+\mathrm{𝟏})`$ non-negative integers $`\eta _\mathrm{𝟏},\eta _\mathrm{𝟐},\mathrm{},\eta _{𝐝+\mathrm{𝟏}}`$ and obviously the $`𝐢`$-th vertex of subfractal $`(\lambda _\mathrm{𝟏},\lambda _\mathrm{𝟐},\mathrm{},\lambda _{𝐝+\mathrm{𝟏}})`$ is denoted by $`\eta _𝐣=\lambda _𝐣+\delta _{𝐢,𝐣}`$, where $`𝐣=\mathrm{𝟏},\mathrm{𝟐},\mathrm{},𝐝+\mathrm{𝟏}`$. As an illustrating example we show in Fig. 3 the method of labelling a Sierpenski gasket with decimation number $`𝐛=\mathrm{𝟑}`$.
Obviously the number of all possible partitions is equal to the distribution of $`(𝐛\mathrm{𝟏})`$ objects amongst $`(𝐝+\mathrm{𝟏})`$ boxes, which is the same as the Bose-Einstein distribution of $`(𝐛\mathrm{𝟏})`$ identical bosons in $`(𝐝+\mathrm{𝟏})`$ quantum states. This is equal to
$$C=\frac{(b+d1)!}{(b1)!.d!}.$$
(2-1)
As is well known, the fractal dimension $`𝐃_𝐟`$ of a self similar object is defined according to
$$NL_f^D=1$$
where $`𝐍`$ is the number of similar objects, up to translation and rotation, here being equal to the number of subfractals of generation $`𝐧`$, and $`𝐋`$ is the scale of subfractal of generation $`𝐧`$.
Hence
$$N=C^r,L=b^{}r$$
Therefore,
$$D_f=\frac{\mathrm{ln}c}{\mathrm{ln}b},$$
or
$$D_f=\frac{\mathrm{ln}(\frac{(b+d1)!}{(b1)!)}}{\mathrm{ln}b}.$$
(2-2)
## 3 Fractal Connection Resistances and their Exact <br>Renormalization Group Transformations
A two-dimensional anisotropic $`(𝐝+\mathrm{𝟏})`$-simplex resistor network consists of $`(𝐝+\mathrm{𝟏})`$ nodes, with $`𝐈_𝐢`$ denoting the amount of current injected into the network through the node $`𝐢`$ and $`\frac{𝐝(𝐝+\mathrm{𝟏})}{\mathrm{𝟐}}`$ different resistors (coated with insulator) mutually connecting all the nodes of the network (see Fig (2)).
As usual, total dissipative power TDP in these networks can be written in terms of the resistances and the currents flowing in them. But it is more convenient and also advantageous throughout this article to express TDP in terms of the inward flowing currents $`𝐈_𝐢,𝐢=\mathrm{𝟏},\mathrm{𝟐},\mathrm{},𝐝+\mathrm{𝟏}`$. In that case, it is clear that TDP is a bilinear function of the input currents with the coefficients which have the dimensions of the resistance.
Hence we call these coefficients, connection resistances denoted by $`R_{jk},j,k=1,2,\mathrm{},d+1`$. Therefore, TDP of the network assumes the following form
$$\mathrm{𝐓𝐃𝐏}(network)=\underset{j,k=1}{\overset{(d+1)}{}}R_{jk}I_jI_k.$$
(3-1)
It is clear from equation (3-1) that $`𝐑_{\mathrm{𝐣𝐤}}`$ is symmetric with respect to the interchange of indices $`𝐢`$ and $`𝐣`$. Also the diagonal elements $`𝐑_{\mathrm{𝐣𝐣}},𝐣=\mathrm{𝟏},\mathrm{𝟐},\mathrm{},𝐝+\mathrm{𝟏}`$ can be eliminated from the expression (3-1), if we use Kirchhoff’s current law for the input currents
$$\underset{j=1}{\overset{(d+1)}{}}I_j=0.$$
(3-2)
Thus, The expression (3-1) takes the following form
$$\mathrm{𝐓𝐃𝐏}(network)=\underset{jk=1}{\overset{(d+1)}{}}R_{jk}I_jI_k=2\underset{k>j=1}{\overset{(d+1)}{}}R_{jk}I_jI_k.$$
(3-3)
From positive definiteness of TDP for all arbitrary values of input inward flowing currents consistent with Kichhoff’s current law, it follows that all connection resistances are positive, that is we have:
$$R_{jk}>0forallk>j=1,2,\mathrm{},d+1.$$
From the form of the TDP given in (3-3), it also follows that there is a bijective map between these sets of the independent connection resistances $`\{𝐑_{\mathrm{𝐣𝐤}},𝐤>𝐣=\mathrm{𝟏},\mathrm{𝟐},\mathrm{},𝐝+\mathrm{𝟏}\}`$ and $`\frac{𝐝(𝐝+\mathrm{𝟏})}{\mathrm{𝟐}}`$ mutual resistors of $`(𝐝+\mathrm{𝟏})`$-simplex network. Accordingly, these independent connection resistances can represent the mutual resistors of the network and in the case of an anisotropic network the connection resistances will be different . Consequently, for the investigation of the restoration of macroscopic isotropy in $`(𝐝+\mathrm{𝟏})`$-simplex fractal resistor lattices, by real space renomalization group method , we need to know the recursion relations between the connection resistances of a given generation and the connection resistances of one generation below it.
These recursion relations can easily be obtained if we compare the total dissipative power TDP of generation $`𝐧`$ given in (3-3) with the same quantity, calculated as sum of power of its $`(𝐧\mathrm{𝟏})`$th generated subfractals which can be expressed as a function of connection resitances of generation $`𝐧\mathrm{𝟏}`$, provided that in calculating the power of its subfractals, the inner inward flowing currents are stated in terms of input currents. To determine these currents it is convenient to denote the $`𝐣`$-th inward flowing current of subfractal corresponding to the partition $`\lambda _\mathrm{𝟏},\lambda _\mathrm{𝟐},\mathrm{},\lambda _{𝐝+\mathrm{𝟏}}`$ by $`𝐈_{\lambda _\mathrm{𝟏},\lambda _\mathrm{𝟐},\mathrm{},\lambda _{𝐝+\mathrm{𝟏}}(\lambda _\mathrm{𝟏},\mathrm{},\lambda _{𝐣\mathrm{𝟏}},\lambda _𝐣+\mathrm{𝟏},\lambda _{𝐣+\mathrm{𝟏}},\mathrm{},\lambda _{𝐝+\mathrm{𝟏}})}`$. Thence $`𝐈_𝐣`$, the $`𝐣`$-th inward flowing current of $`(𝐝+\mathrm{𝟏})`$-simplex fractal, is given by
$`I_{0,0,,0,\underset{jth}{\underset{}{1}},0,,0}(0,0,,0,\underset{jth}{\underset{}{2}},0,,0)=I_j.`$
To determine the inner inward flowing currents, besides applying Kirchhoff’s current law at each node and subfractal, we have to minimize the total dissipative power of $`(𝐝+\mathrm{𝟏})`$-simplex fractal of generation $`𝐧`$, calculated as the sum of the TDP of its subfractals as:
$`{\displaystyle \underset{_{_{_{_{_{_{sumoverpartitionof(b1)}}}}}}}{}}{\displaystyle \underset{j,k=1}{\overset{d+1}{}}}R_{jk}(n1)I_{\lambda _1,\mathrm{},\lambda _{d+1}}(\lambda _1,,\lambda _j+1,,\lambda _{d+1})I_{\lambda _1,,\lambda _{d+1}}(\lambda _1,,\lambda _k+1,,\lambda _{d+1})`$
$`{\displaystyle \underset{sumoverpartitionof(b1)}{}}2\mu _{\lambda _1,,\lambda _{d+1}}I_{\lambda _1,,\lambda _{d+1}}(\lambda _1,,\lambda _k+1,,\lambda _{d+1})`$
$`{\displaystyle \underset{sumoverpartitionofb}{}}2\nu _{\eta _1,,\eta _{d+1}}I_{\eta _,,\eta _{k1},,\eta _{d+1}}(\eta _1,,\eta _{d+1}),`$
where $`\mu _{\lambda _\mathrm{𝟏},\mathrm{},\lambda _{𝐝+\mathrm{𝟏}}}`$ and $`\nu _{\eta _\mathrm{𝟏},\mathrm{},\eta _{𝐝+\mathrm{𝟏}}}`$ are lagrange multipliers due to Kirchhoff’s law on each subfractal, and also on each node,respectively. Minimizing the expression (3-4), we get linear equations between inner input flowing currents and lagrange multipliers together with the Kirchhoff’s law for each subfractal and each vertex, respectively. Solving the equations thus obtained we can write all inner inward flowing currents as a linear function in terms of input ones. Substituting the expressions thus obtained for the inner currents in Eq. (3-4), we determine TDP of generation $`𝐧`$ which is obviously a bilinear function of input currents with coefficients which are in general very involved functions of the connection resistances of the generation $`𝐧\mathrm{𝟏}`$. Comparing the final result with the expression (3-4), connection resistance of generation $`𝐧`$ as its coefficient, we get the required transformation between connection resistances of generations $`𝐧`$ and $`𝐧\mathrm{𝟏}`$, respectively:
$$R_{jk}(r+1))=f_{jk}(R_{lm}(r)_{m>l}),k>j=1,2,,d+1.$$
(3-5)
Here in this article, we show that the power and the anisotropy suppression exponents can be calculated, without having any knowledge of the explicit form of the functions $`𝐟_{\mathrm{𝐣𝐤}}`$. All we need to know is some general properties of these functions which can be obtained rather easily from some physical requirements and also from dimensional analysis: these functions are homogeneous functions of degree one mapping positive connection resistances of generation $`𝐧\mathrm{𝟏}`$ into positive connection resistances of generation $`𝐧`$, that is they form positive homogeneous map of degree one.
All connection resistances are positive; none of them can be negative or zero. The physical reason behind it is that if, for example, the connection resistance $`𝐑_{\mathrm{𝐣𝐤}}`$ becomes negative or if it vanishes, then for inward flowing currents $`𝐈_𝐣=𝐈_𝐤=𝐈`$, and $`𝐈_𝐥=\mathrm{𝟎}`$ if $`𝐥𝐣𝐤`$, we obviously get negative or zero power which is not physical in either cases. Analogously, we can rather easily deduce that none of them can be infinite, since all resistors of the network are finite, otherwise we will have infinite total dissipative power which is not again physical.
Definitely the transformation (3-5) is monotonically increasing, since by increasing the connection resistances at a given generation $`𝐧\mathrm{𝟏}`$, without changing the input currents, the total dissipative power of generation $`𝐧`$ will increase, that is the connection resistances of generation $`𝐧`$ will increase. Naturally, under the action of the point group $`𝐒_{(𝐝+\mathrm{𝟏})}`$, the connection resistances simply permute among themselves. For example, the exchange of the vertices $`𝐣`$ and $`𝐤`$ in $`(𝐝+\mathrm{𝟏})`$-simplex induces the following transformation among the connection resistances:
$`R_{jk}`$ $`R_{kj}`$ $`=R_{jk}`$
$`R_{jl}`$ $`R_{lk}`$ $`forljk`$
$`R_{kl}`$ $`R_{lj}`$ $`forljk`$
$`R_{lm}`$ $`R_{lm}`$ $`forlmjk.`$ (3-6)
As an example, we give the explicit form of the transformation for the special case of $`𝐝=\mathrm{𝟐}`$ and $`𝐛=\mathrm{𝟐}`$
$`R_{12}(n)={\displaystyle \frac{R_{12}[R_{12}(n1)+2R_{13}(n1)+2R_{23}(n1)]}{R_{12}(n1)+R_{13}(n1)+R_{23}(n1)}},`$
$`R_{13}(n)={\displaystyle \frac{R_{13}[R_{13}(n1)+2R_{12}(n1)+2R_{23}(n1)]}{R_{12}(n1)+R_{13}(n1)+R_{23}(n1)}},`$
$`R_{23}(n)={\displaystyle \frac{R_{12}[R_{23}(n1)+2R_{12}(n1)+2R_{13}(n1)]}{R_{12}(n1)+R_{13}(n1)+R_{23}(n1)}}.`$
## 4 Fixed Point of Recurrence Equation of Connection <br>Resistances
In this section we present a rigorous proof that the renormalization group transformation of the connection resistances has a unique fixed direction in the space of connection resistances. That is, all of the flows stemming from the finite physical region of connection resistance space converge to infinity at a direction which has the same angle with all coordinates.
For simplicity we denote the connection resitances of generation $`𝐧\mathrm{𝟏}`$ $`𝐑_{\mathrm{𝐣𝐤}}(𝐧\mathrm{𝟏})`$ $`(𝐤,𝐣=\mathrm{𝟏},\mathrm{𝟐},\mathrm{},𝐝+\mathrm{𝟏})`$, by $`𝐗_\alpha `$ $`(\alpha =\mathrm{𝟏},\mathrm{𝟐},\mathrm{},\frac{𝐝(𝐝+\mathrm{𝟏})}{\mathrm{𝟐}})`$ and the connection resistances of generation $`𝐧`$ $`𝐑_{\mathrm{𝐣𝐤}}(𝐧)`$ $`(𝐤,𝐣=\mathrm{𝟏},\mathrm{𝟐},\mathrm{},𝐝+\mathrm{𝟏})`$, by $`𝐗_\alpha ^{}`$ $`(\alpha =\mathrm{𝟏},\mathrm{𝟐},\mathrm{},\frac{𝐝(𝐝+\mathrm{𝟏})}{\mathrm{𝟐}})`$, respectively. Then the transformations (3-5) can be written as
$$X_\alpha ^{}=f_\alpha (X_\beta )for\alpha =1,2,\mathrm{},\frac{d(d+1)}{2}.$$
(4-1)
Now we consider $`𝐗_\alpha >\mathrm{𝟎}`$ $`(\alpha =\mathrm{𝟏},\mathrm{𝟐},\mathrm{},\frac{𝐝(𝐝+\mathrm{𝟏})}{\mathrm{𝟐}})`$ as coordinates of the interior points of a cone in $`\frac{𝐝(𝐝+\mathrm{𝟏})}{\mathrm{𝟐}}`$ dimensional Euclidean space denoted by $`\stackrel{˘}{𝒞}`$. Denoting the cone itself by $`𝒞`$, the transformation (4-1) can be considered as the map of this cone into itself:
$$𝒞\stackrel{F}{}𝒞,$$
where we have denoted the extension of the map (4-1) over the cone itself by $`𝐅`$.
From the action of the permutation group $`𝐒_{(𝐝+\mathrm{𝟏})}`$ on the space of connection resistances (the cone $`𝒞`$) given in (3-6), it follows that the transformation (4-1) is equivariant with respect to the action of $`𝐒_{(𝐝+\mathrm{𝟏})}`$, that is we have the following commutative diagram :
$`𝒞`$ $`\stackrel{g}{}𝒞`$
$`F`$ $``$ $`FforeverygS_{(d+1)},`$
$`𝒞`$ $`\stackrel{g}{}𝒞`$
or
$`g(F(𝒳))=F(g(𝒳))foreverygS_{(d+1)}and𝒳𝒞.`$ (4-2)
A cone in Euclidean space has the following properties:
1. $`𝒞`$+$`𝒞𝒞`$
2. $`\lambda 𝒞𝒞`$ for all $`\lambda >\mathrm{𝟎}`$
3. $`𝒞𝒞=\mathrm{𝟎}`$.
We denote interior of this cone by $`\stackrel{˘}{𝒞}`$. One can define an ”order relation” as follows:
$`𝒳𝒴`$ if $`𝒳𝒴𝒞`$ and $`𝒳>𝒴`$ if $`𝒳𝒴\stackrel{˘}{𝒞}`$. If one defines the numbers $`𝐦(𝒳,𝒴)`$ and $`𝐌(𝒳,𝒴)`$ such that
$`m(𝒳,𝒴)=max_i{\displaystyle \frac{x_i}{y_i}}`$
$`M(𝒳,𝒴)=min_i{\displaystyle \frac{x_i}{y_i}},`$ (4-3)
then for any $`𝒳,𝒴\stackrel{˘}{𝒞}`$ the following relation holds
$`m(𝒳,𝒴)𝒴𝒳M(𝒳,𝒴)𝒴.`$ (4-4)
The above relation and also all the other theorems of this section have been proved in reference . Using the numbers defined in (4-4), one can define the Hilbert metric $`𝐝(𝒳,𝒴)`$ for any $`𝒳,𝒴\stackrel{˘}{𝒞}`$ as
$`d(𝒳,𝒴)=\mathrm{log}[{\displaystyle \frac{M(𝒳,𝒴)}{m(𝒳,𝒴)}}]=\mathrm{log}[max_{i,j}{\displaystyle \frac{x_iy_j}{y_ix_j}}],`$ (4-5)
with the usual property of a pseudometric on $`\stackrel{˘}{𝒞}`$ and a metric on $`\stackrel{˘}{𝒞}𝐒(\mathrm{𝟎},\mathrm{𝟏})`$, where $`𝐒(\mathrm{𝟎},\mathrm{𝟏})`$ denotes the set of points of sphere of radius one with the origin as its center. The metric space $`\stackrel{˘}{𝒞}𝐒(\mathrm{𝟎},\mathrm{𝟏})`$ is complete under Hilbert metric(4-6). Also, it is straightforward to see that the following assertions about this metric are valid :
1. For any $`𝒳,𝒴\stackrel{˘}{𝒞}`$ and $`𝐚,𝐛𝐑`$ we have
$$d(a𝒳,b𝒳)=d(𝒳,𝒳).$$
2. $`𝐝(𝒳,𝒳)`$ =0 if and only if $`𝒳=\lambda 𝒳`$.
3. For any $`𝒳,𝒴\stackrel{˘}{𝒞}`$ the metric $`𝐝(𝒳,𝒳)`$ is finite. A given map of the cone into itself:
$$F:𝒞𝒞$$
is called positive homogeneous of degree $`𝐧`$ and monotonically increasing if, for all $`𝒳𝒞`$ and $`𝐚>\mathrm{𝟎}`$, we have
$$F(a𝒳)=a^nF(𝒳),$$
and
$$𝒳𝒴F(𝒳)F(𝒴).$$
According to the arguments given in section III, the transformation (3-5) or (4-1) is monotonically increasing, positive and homogeneous map of degree one.
Using the relation (4-5) one can prove that for a positive homogeneous map of degree one and monotonically increasing map like the transformation (4-1), the following inequality holds:
$$m(𝒳,𝒴)T(𝒳)T(𝒳).M(𝒳),(𝒴)T(𝒴).$$
(4-6)
It is straightforward to get the following inequality from the inequality (4-7)
$$m(𝒳,𝒴)m(𝒯𝒳,𝒯𝒴)M(𝒳,𝒴)M(𝒯𝒳,𝒯𝒴).$$
(4-7)
From the above inequality and also from the definition of Hilbert metric (4-6), it follows that for all $`𝒳,𝒴𝒞𝐒(\mathrm{𝟎},\mathrm{𝟏})`$ we have
$$d(𝒯𝒳,𝒯𝒴)d(𝒳,𝒴).$$
(4-8)
Therefore, the transformation (4-1) satisfies the Lipschitz condition and is of contractive type. Thus, according to the Principle of Contraction Mapping Theorem, the contracting mapping (4-1) has a unique fixed point $`𝒳_\mathrm{𝟎}`$ in the complete metric space $`(𝒞𝐒(\mathrm{𝟎},\mathrm{𝟏}),𝐝)`$ (d is the Hilbert metric given in (4-4) and
$$\underset{n\mathrm{}}{lim}\stackrel{n}{\stackrel{}{F(F(\mathrm{}F(F}}(𝒳))\mathrm{}))=𝒳_0,forevery𝒳𝒞.$$
(4-9)
But, because of the equivariant property (4-3)of the transformation (4-1), any fixed point of the point group $`𝐒_{(𝐝+\mathrm{𝟏})}`$(or the stability point of the point group ) will definitely be the fixed point of the transformation (4-1) acting on the space $`(𝒞𝐒(\mathrm{𝟎},\mathrm{𝟏}))`$. Obviously, the direction $`𝐗_\mathrm{𝟏}=𝐗_\mathrm{𝟐}==𝐗_{\frac{𝐝(𝐝+\mathrm{𝟏})}{\mathrm{𝟐}}}`$ is the only fixed point of the permutation group $`𝐒_{(𝐝+\mathrm{𝟏})}`$ acting on the space $`(𝒞𝐒(\mathrm{𝟎},\mathrm{𝟏}))`$. Hence, because of the uniqueness of the fixed point of the transformation (4-1) on the space $`(𝒞𝐒(\mathrm{𝟎},\mathrm{𝟏}))`$, the direction $`𝐗_\mathrm{𝟏}=𝐗_\mathrm{𝟐}==𝐗_{\frac{𝐝(𝐝+\mathrm{𝟏})}{\mathrm{𝟐}}}`$ is the only fixed direction of the connection resistances renormalization group transformation. This direction corresponds to the isotropic $`(𝐝+\mathrm{𝟏})`$-simplex, which indicates that the macroscopic conductivity becomes isotropic on large scales.
## 5 Determination of Inner Inward Flowing Currents <br>of Subfractals in Isotropic State
In order to determine the inner inward flowing currents in terms of the input currents $`𝐈_𝐣`$ $`(𝐣=\mathrm{𝟏},\mathrm{𝟐},,𝐝+\mathrm{𝟏})`$ in isotropic state, we have to minimize the TDP given in (3-4). But here in isotropic state all connection resistances are the same, hence they can be put equal to one in (3-4), simply by rescaling the lagrange multipliers of current conservations of vertices and subfractals. Now, by minimizing TDP, we get the following equation for $`𝐈`$
$$I_{\lambda _1,\mathrm{},\lambda _{d+1}}(\lambda _1,\mathrm{},\lambda _j+1,\mathrm{},\lambda _{d+1})\mu _{\lambda _1,\mathrm{},\lambda _j,\lambda _{d+1}}nu_{\lambda _1,\mathrm{},\lambda _j+1,\mathrm{},\lambda _{j+1}}=0$$
(5-1)
together with the Kirchhoff’s law for each subfractal and each vertex, respectively
$$\underset{j=1}{\overset{d+1}{}}I_{\lambda _1,,\lambda _j,,\lambda _{d+1}}(\lambda _1,,\lambda _j+1,,\lambda _{d+1})=0.$$
(5-2a)
$$\underset{j=1}{\overset{d+1}{}}I_{\eta _1,,\eta _j1,,\eta _{d+1}}(\eta _1,,\eta _j,,\eta _{d+1})=0.$$
(5-2b)
We assume the following ansatz for the Lagrange multipliers:
$$\mu _{\lambda _1,\lambda _2,,\lambda _{d+1}}=\underset{k=1}{\overset{d+1}{}}a_{\lambda _1,\lambda _2,,\lambda _{d+1}}(\lambda _k)I_k$$
(5-3a)
$$\nu _{\eta _1,\eta _2,,\eta _{d+1}}\underset{k=1}{\overset{d+1}{}}b_{\eta _1,\eta _2,,\eta _{d+1}}(\eta _k)I_k$$
(5-3b)
with $`a_{\lambda _1,\lambda _2,,\lambda _{d+1}}(0)`$ and $`b_{\eta _1,\eta _2,,\eta _{d+1}}(0)`$ taken to be zero.
Using the ansatz (5-3a) and (5-3b) in equation (5-1), the inflowing currents can be given in terms of $`𝐚`$ and $`𝐛`$ respectively, that is
$$I_{\lambda _1,,\lambda _{d+1}}(\eta _1,,\eta _{d+1})=\underset{k=1}{\overset{d+1}{}}a_{\lambda _1,\lambda _2,,\lambda _{d+1}}(\lambda _k)I_k+\underset{k=1}{\overset{d+1}{}}b_{\eta _1,\eta _2,,\eta _{d+1}}(\eta _k)I_k.$$
(5-4)
Due to the $`𝐒_{(𝐝+\mathrm{𝟏})}`$ permutation symmetry of $`(𝐝+\mathrm{𝟏})`$-simplex fractal, the parameters $`𝐚_{\lambda _\mathrm{𝟏},\lambda _\mathrm{𝟐},,\lambda _{𝐝+\mathrm{𝟏}}}(\lambda _𝐤)`$ and $`𝐛_{\eta _\mathrm{𝟏},\eta _\mathrm{𝟐},,\eta _{𝐝+\mathrm{𝟏}}}(\eta _𝐤)`$ depend only on the corresponding partition $`\{\lambda _\mathrm{𝟏},\lambda _\mathrm{𝟐},,\lambda _{𝐝+\mathrm{𝟏}}\}`$ and $`\{\eta _\mathrm{𝟏},\eta _\mathrm{𝟐},,\eta _{𝐝+\mathrm{𝟏}}\}`$, respectively. They do not change under the permutation of $`\lambda _𝐢`$ or $`\eta _i`$ within a given partition. From now on, as far as $`a`$ and $`b`$ are concerned, only nonzero values are going to be quoted in their partition.
Actually one could write the currents in terms of input ones as in (5-4) by simply using the symmetry of simplex fractal, and the minimization of power is not required. Finally $`𝐚`$ and $`𝐛`$ can be determined through the equations (5-2a) and (5-2b). Obviously the number of equations are the same as the number of unknowns, hence the unknowns $`𝐚`$ and $`𝐛`$ can be determined uniquely. Here we determine the currents only for $`𝐛=\mathrm{𝟐},\mathrm{𝟑},\mathrm{𝟒}`$ and $`\mathrm{𝟓}`$, respectively.
Let us first consider the case where $`𝐛=\mathrm{𝟐}`$
$$I_{0,0,,0,\underset{jth}{\underset{}{1}},0,,0}(0,0,,0,\underset{jth}{\underset{}{2}},0,,0)=I_j$$
$$I_{0,0,,0,\underset{jth}{\underset{}{1}},0,,0}(0,0,,0,\underset{jth}{\underset{}{1}},0,,0,\underset{kth}{\underset{}{1}},0,,0)=a_1(1)I_J+b_1(1)I_j+b_(1)1I_k$$
Using equation (5-2b) we have
$$a_1(1)+2b_1(1)=0$$
and from equation (5-2a) we get
$$1+da_1(1)+(d1)b_1(1)=0.$$
Solving the above equations we get the following result
$$I_{0,,0,\underset{jth}{\underset{}{1}},0,,0}(0,,0,\underset{jth}{\underset{}{1}},0,,0,\underset{kth}{\underset{}{1}},0,,0)=\frac{(I_kI_j)}{(d+1)}.$$
Via the the procedure explained above, we can similarly calculate the inner inward flowing currents corresponding to $`𝐛=\mathrm{𝟑},\mathrm{𝟒}`$ and $`𝐛=\mathrm{𝟓}`$, where the details of calculation appear in Appendices I, II and III, respectively and below we quote only the results:
I: Inner inward flowing currents corresponding to decimation number $`𝐛=\mathrm{𝟑}`$
$$I_{0,,0,\underset{jth}{\underset{}{2}},0,,0}(0,,0,\underset{jth}{\underset{}{2}},0,,0,\underset{kth}{\underset{}{1}},0,,0)=$$
$$\frac{2d+5}{(2d+3)(d+1)}I_j+\frac{3}{(2d+3)(d+1)}I_k$$
$$I_{0,,0,\underset{jth}{\underset{}{1}},0,,0,\underset{kth}{\underset{}{1}},0,,0}(0,,0,\underset{jth}{\underset{}{2}},0,,0,\underset{kth}{\underset{}{1}},0,,0)=$$
$$\frac{2d+5}{(2d+3)(d+1)}I_j\frac{3}{(2d+3)(d+1)}I_k$$
$$I_{0,,0,\underset{jth}{\underset{}{1}},0,,0,\underset{kth}{\underset{}{1}},0,,0}(0,,0,\underset{jth}{\underset{}{1}},0,,0,\underset{kth}{\underset{}{1}},0,,0)=$$
$$\frac{2}{(2d+3)(d+1)}(2I_lI_jI_k).$$
II: Inner inward flowing currents corresponding to decimation number $`𝐛=\mathrm{𝟒}`$
$$I_{0,,0,\underset{jth}{\underset{}{3}},0,,0}(0,,0,\underset{jth}{\underset{}{3}},0,,0,\underset{kth}{\underset{}{1}},0,,0)=$$
$$\frac{8d^3+52d^2+5(25d+21)}{P}I_j+\frac{25d+49}{P}I_k$$
$$I_{0,,0,\underset{jth}{\underset{}{2}},0,,0,\underset{kth}{\underset{}{1}},0,,0}(0,,0,\underset{jth}{\underset{}{3}},0,,0,\underset{kth}{\underset{}{1}},0,,0)=$$
$$\frac{8d^3+52d^2+125d+105}{P}I_j\frac{25d+49}{P}I_k$$
$$I_{0,,0,\underset{jth}{\underset{}{2}},0,,0,\underset{kth}{\underset{}{1}},0,,0}(0,,0,\underset{jth}{\underset{}{2}},0,,0,\underset{kth}{\underset{}{2}},0,,0)=$$
$$\frac{12d^2+79d+91}{P}(I_jI_k)$$
$$I_{0,,0,\underset{jth}{\underset{}{2}},0,,0,\underset{kth}{\underset{}{1}},0,,0}(0,,0,\underset{jth}{\underset{}{2}},0,,0,\underset{kth}{\underset{}{1}},0,,0,\underset{lth}{\underset{}{1}},0,,0)=$$
$$8\frac{d^2+6d+7}{P}I_j4\frac{3d+7}{P}I_k+2\frac{19d+35}{P}I_l$$
$$I_{0,,0,\underset{jth}{\underset{}{1}},0,,0,\underset{kth}{\underset{}{1}},0,,0,\underset{lth}{\underset{}{1}},0,,0}(0,\mathrm{\hspace{0.17em}0},\underset{jth}{\underset{}{2}},0,\mathrm{\hspace{0.17em}0},\underset{kth}{\underset{}{1}},0,\mathrm{\hspace{0.17em}0},\underset{lth}{\underset{}{1}},0,\mathrm{\hspace{0.17em}0})=$$
$$16\frac{d^2+6d+7}{P}I_j2\frac{13d+21}{P}(I_k+I_l)$$
$$I_{0,.,0,\underset{jth}{\underset{}{1}},0,.,0,\underset{kth}{\underset{}{1}},0,.,0,\underset{lth}{\underset{}{1}},0,.,0}(0,.,0,\underset{jth}{\underset{}{1}},0,.,0,\underset{kth}{\underset{}{1}},0,.,0,\underset{lth}{\underset{}{1}},0,,\underset{mth}{\underset{}{1}},0,.,0)=$$
$$4\frac{4d+7}{P}(I_j+I_k+I_l3I_m)$$
where $`P`$ is defined as
$$P=(8d^3+44d^2+81d+49)(d+1).$$
III: Inner inward flowing currents corresponding to decimation number $`𝐛=\mathrm{𝟓}`$
$$I_{0,.,0,\underset{jth}{\underset{}{4}},0,,0}(0,,0,\underset{jth}{\underset{}{4}},0,,0,\underset{kth}{\underset{}{1}},0,,0)=$$
$$\frac{192d^6+2720d^5+16332d^4+53648d^3+102215d^2+106746d+47255}{Q}I_j$$
$$3\frac{542d^3+3803d^2+8576d+6275}{Q}I_k$$
$$I_{0,,0,\underset{jth}{\underset{}{3}},0,,0,\underset{kth}{\underset{}{1}},0,,0}(0,,0,\underset{jth}{\underset{}{4}},0,,0,\underset{kth}{\underset{}{1}},0,,0)=$$
$$\frac{192d^6+2720d^5+16332d^4+53648d^3+102215d^2+106746d+47255}{Q}$$
$$+3\frac{542d^3+3803d^2+8576d+6275}{Q}I_k$$
$$I_{0,,0,\underset{jth}{\underset{}{3}},0,,0,\underset{kth}{\underset{}{1}},0,,0}(0,,0,\underset{jth}{\underset{}{3}},0,,0,\underset{kth}{\underset{}{2}},0,,0)=$$
$$\frac{288d^5+4104d^4+24466d^3+70139d^2+94822d+48213}{Q}I_j$$
$$\frac{600d^4+8602d^3+36869d^2+62290d+36303}{Q}I_k$$
$$I_{0,,0,\underset{jth}{\underset{}{3}},0,,0,\underset{kth}{\underset{}{1}},0,,0}(0,,0,\underset{jth}{\underset{}{3}},0,,0,\underset{kth}{\underset{}{1}},0,,0,\underset{lth}{\underset{}{1}},0,,0)=$$
$$2\frac{96d^5+1312d^4+7426d^3+20691d^2+27782d+14261}{Q}I_j$$
$$+2\frac{300d^3+2462d^2+6245d+5043}{Q}I_k$$
$$2\frac{1326d^3+8947d^2+19483d+13782}{Q}I_l$$
$$I_{0,,0,\underset{jth}{\underset{}{2}},0,,0,\underset{kth}{\underset{}{2}},0,,0}(0,,0,\underset{jth}{\underset{}{3}},0,,0,\underset{kth}{\underset{}{2}},0,,0)=$$
$$\frac{288d^5+4104d^4+24466d^3+70139d^2+94822d+48213}{Q}I_j$$
$$\frac{600d^4+8602d^3+36869d^2+62290d+36303}{Q}I_k$$
$$I_{0,,0,\underset{jth}{\underset{}{2}},0,,0,\underset{kth}{\underset{}{2}},0,,0}(0,,0,\underset{jth}{\underset{}{2}},0,,0,\underset{kth}{\underset{}{2}},0,,0,\underset{lth}{\underset{}{1}},0,,0)=$$
$$2\frac{144d^4+1896d^3+8260d^2+14607d+9059}{Q}I_j$$
$$+\frac{144d^4+1896d^3+8260d^2+14607d+9059}{Q}I_k$$
$$2\frac{784d^3+5144d^2+10907d+7507}{Q}I_l$$
$$I_{0,,0,\underset{jth}{\underset{}{2}},0,,0,\underset{kth}{\underset{}{1}},0,,0,\underset{lth}{\underset{}{1}},0,,0}(0,\mathrm{\hspace{0.17em}0},\underset{jth}{\underset{}{3}},0,\mathrm{\hspace{0.17em}0},\underset{kth}{\underset{}{1}},0,\mathrm{\hspace{0.17em}0},\underset{lth}{\underset{}{1}},0,\mathrm{\hspace{0.17em}0})=$$
$$2\frac{1026d^3+6485d^2+13238d+8739}{Q}(I_k+I_l)$$
$$4\frac{96d^5+1312d^4+7426d^3+20691d^2+27782d+14261}{Q}I_j$$
$$I_{0,,0,\underset{jth}{\underset{}{2}},0,,0,\underset{kth}{\underset{}{1}},0,,0,\underset{lth}{\underset{}{1}},0,,0}(0,\mathrm{\hspace{0.17em}0},\underset{jth}{\underset{}{2}},0,\mathrm{\hspace{0.17em}0},\underset{kth}{\underset{}{2}},0,\mathrm{\hspace{0.17em}0},\underset{lth}{\underset{}{1}},0,\mathrm{\hspace{0.17em}0})=$$
$$2\frac{312d^4+4084d^3+16326d^2+26083d+14489}{Q}I_j$$
$$2\frac{228d^4+2990d^3+12293d^2+20345d+11774}{Q}I_k$$
$$+\frac{784d^3+5144d^2+10907d+7507}{Q}I_l$$
$$I_{0,,0,\underset{jth}{\underset{}{2}},0,,0,\underset{kth}{\underset{}{1}},0,,0,\underset{lth}{\underset{}{1}},0,,0}(0,,0,\underset{jth}{\underset{}{2}},0,,0,\underset{kth}{\underset{}{1}},0,,0,\underset{lth}{\underset{}{1}},0,,\underset{mth}{\underset{}{1}},0,,0)=$$
$$8\frac{48d^4+596d^3+2389d^2+3899d+2238}{Q}I_j$$
$$+6\frac{152d^3+1062d^2+2361d+1691}{Q}I_k$$
$$+\frac{152d^3+1062d^2+2361d+1691}{Q}I_l$$
$$2\frac{316d^3+2041d^2+4273d+2908}{Q}I_m$$
$$I_{0,,0,\underset{jth}{\underset{}{1}},0,,0,\underset{kth}{\underset{}{1}},0,,0,,0,\underset{lth}{\underset{}{1}},0,,0,,0,\underset{mth}{\underset{}{1}},0,,0}(0,,0,\underset{jth}{\underset{}{2}},0,,0,\underset{kth}{\underset{}{1}},0,,0,\underset{lth}{\underset{}{1}},0,,0,,0,\underset{mth}{\underset{}{1}},0,,0)=$$
$$24\frac{48d^4+596d^3+2389d^2+3899d+2238}{Q}I_j$$
$$\frac{164d^3+979d^2+1912d+1217}{Q}(I_k+I_l+I_m)$$
$$I_{0,,0,\underset{jth}{\underset{}{1}},0,,0,\underset{kth}{\underset{}{1}},0,,0,\underset{lth}{\underset{}{1}},0,,0,,0,\underset{mth}{\underset{}{1}},0,,0}(0,,0,\underset{jth}{\underset{}{1}},0,,0,\underset{kth}{\underset{}{1}},0,,0,\underset{lth}{\underset{}{1}},0,,0,\underset{mth}{\underset{}{1}},0,,0,\underset{nth}{\underset{}{1}},0,,0)=$$
$$12\frac{96d^3+604d^2+1237d+825}{Q}(I_j+I_k+I_l+I_m4I_n)$$
where $`Q`$ is defined as
$$Q=192d^7+2720d^6+16332d^5+53648d^4+103841d^3+118155d^2+72983d+18825.$$
## 6 Scaling Exponent of Anisotropy Suppression of <br>$`(𝐝+\mathrm{𝟏})`$-Simplex Fractal Conductor Network
To investigate the abolition of anisotropy and also in order to calculate the scaling exponent of its suppression, we linearize the recursion map (4-1) near the fixed direction (isotropy state) of this map:
$$\underset{n\mathrm{}}{lim}R_{jk}(n)=Rfork>j=1,2,,d+1.$$
This leads us to write
$$R_{jk}(n)=R+\epsilon _{jk}(n);fork>j=1,2,,d+1$$
(6-1)
with $`\epsilon _{\mathrm{𝐣𝐤}}(𝐧)`$ as an infinitesimal deviation of the connection resistances of generation $`𝐧`$ from the isotropic state, for large values of $`𝐧`$. Now, all we need to know is the recursion relations between the deviation of connection resistances of the generation $`𝐧`$ and the infinitesimal deviation of connection resistances of the generation $`𝐧\mathrm{𝟏}`$, for large values of $`𝐧`$.
These recursion relations can easily be obtained, if we compare the deviation of TDP from the isotropic state of generation $`𝐧`$ with the deviation of the same quantity, calculated as the sum of deviation of TDP of its subfractals of generation $`𝐧\mathrm{𝟏}`$. Clearly the deviation of TDP of generation $`𝐧`$ can be obtained from the expression (3-3), provided that in (3-3) we replace the connection resistances with the deviation of the connection resistances of generation $`𝐧`$. Also TDP of generation $`𝐧`$ is the sum of TDP of subfractals generations $`𝐧\mathrm{𝟏}`$, where again the latter can be obtained from (3-3), if we replace in (3-3) the input currents with the inner inflowing currents (which have been expressed in terms of input currents in section V) and the connection resistances with the deviation of the generation $`𝐧\mathrm{𝟏}`$, respectively. Proceeding as above we obtain the recursion relations of the following form for the deviation of connection resistances of generations $`𝐧\mathrm{𝟏}`$ and $`𝐧`$ in a $`(𝐝+\mathrm{𝟏})`$-simplex fractal conductor network, for large values of $`𝐧`$:
$`{}_{jk}{}^{}\epsilon _{jk}^{}(n)=f(d,b)_{jk}\epsilon _{jk}(n1)+g(d,b)({\displaystyle \underset{l=1jk}{\overset{d+1}{}}}\epsilon _{jl}(n1)+{\displaystyle \underset{l=1jk}{\overset{d+1}{}}}\epsilon _{lk}(n1))`$
$`+g(d,b)({\displaystyle \underset{lm=1jk}{\overset{d+1}{}}}\epsilon _{lm}(n1)).`$ (6-2)
For a given value of $`𝐣𝐤`$ we denote $`\epsilon _{\mathrm{𝐣𝐤}}(𝐧)`$($`\epsilon _{jk}(n1)`$) by $`X(X^{})`$. Next we assume that $`\epsilon _{\mathrm{𝐣𝐥}}(𝐧)`$($`\epsilon _{\mathrm{𝐣𝐥}}(𝐧\mathrm{𝟏})`$ and $`\epsilon _{\mathrm{𝐥𝐤}}(𝐧)`$ $`(\epsilon _{\mathrm{𝐥𝐤}}(𝐧\mathrm{𝟏}))`$ with ($`𝐥𝐣𝐤`$) are all equal which are denoted by $`𝐘`$($`𝐘^{}`$). Finally we assume that the remaining deviation of connection resistances, that is, $`\epsilon _{\mathrm{𝐥𝐦}}(𝐧)`$($`\epsilon _{\mathrm{𝐥𝐦}}(𝐧\mathrm{𝟏})`$) with ($`𝐦𝐥𝐣𝐤`$) are all equal which are denoted by $`𝐙(𝐙^)`$. Then the recursion relations (6-2) take the following form:
$`\left(\begin{array}{c}X^{}\\ Y^{}\\ Z^{}\end{array}\right)=`$ (6-6)
$`\left(\begin{array}{ccc}f(d,b)& 2(d1)g(d,b)& (d2)(d1)h(d,b)\\ g(d,b)& f+g(d1)+2h(d2)& g(d2)+h(d3)(d2)\\ 2h(d,b)& 4g(d,b)+4(d3)h(d,b)& f(d,b)+2(d3)g(d,b)+(d4)(d3)h(d,b)\end{array}\right)\left(\begin{array}{c}X\\ Y\\ Z\end{array}\right).`$ (6-11)
$`(63)`$
The following eigen-values are obtained by diagonalizing the $`\mathrm{𝟑}\times \mathrm{𝟑}`$ matrix (6-3) . The eigen-values are quoted in decreasing order as:
$`\lambda _{max}`$ $`=`$ $`{\displaystyle \frac{d^2h(d,b)+d(4g(d,b)3h(d,b))+2f(d,b)4g(d,b)+2h(d,b)}{2}}`$
$`\lambda _{midle}`$ $`=`$ $`d(g(d,b)h(d,b))+f(d,b)3g(d,b)+2h(d,b)`$
$`\lambda _{min}`$ $`=`$ $`f(d,b)2g(d,b)+h(d,b),`$ (6-12)
where the corresponding eigen-vectors are given as the rows of the following matrix:
$`\left(\begin{array}{ccc}1& 1& 1\\ 2(d1)& d3& 2\\ (d1)(d2)& (d2)& 2\end{array}\right).`$ (6-16)
We see that the eigen-directions do not depend on the decimation number $`𝐛`$ but rather only on the dimension $`𝐝`$. This is again due to equivariancy of the map (4-1) with respect to the action of the point group $`𝐒_{(𝐝+\mathrm{𝟏})}`$ on the space of the connection resistances given in (4-3).
As expected, the maximum eigen-value corresponds to isotropy state, which gives the power scaling exponent of $`(𝐝+\mathrm{𝟏})`$-simplex fractal conductor networks . Therefore, anisotropy vanishes with a scaling exponent which can be obtained in terms of the eigenvalues in the following way: from the recurence relation (3-5) and its linearized form (6-2), it follows that, for large values of $`𝐧`$ we have
$$\underset{n\mathrm{}}{lim}R_{jk}=L_n^{D_2},foreveryjk=1,2,,d+1$$
(6-17)
where $`𝐋_𝐧=𝐛^𝐧`$, and the power scaling exponent $`𝐃_\mathrm{𝟐}`$ is defined as:
$$D_2(d,b)=\frac{\mathrm{log}\lambda _{max}}{\mathrm{log}b}$$
(6-18)
and
$`\underset{n\mathrm{}}{lim}{\displaystyle \frac{R_{jk}}{R_{jl}}}1=L_n^{\overline{\lambda }}foreveryjkl=1,2,,d+1`$
$`\underset{n\mathrm{}}{lim}{\displaystyle \frac{R_{jk}}{R_{lm}}}1=L_n^{\overline{\lambda }}foreveryjklm=1,2,,d+1`$
where the scaling exponent of suppression of the anisotropy, $`\overline{\lambda }`$ , is defined as
$$\overline{\lambda }(d,b)=\frac{\mathrm{log}\frac{maximumeigenvalue}{thenextgreatesteigenvalue}}{\mathrm{log}b}=\frac{\mathrm{log}\frac{\lambda _{max}}{\lambda _{midle}}}{\mathrm{log}b}.$$
(6-20)
In the remaining part of this section we quote the results for $`𝐛=\mathrm{𝟐},\mathrm{𝟑},\mathrm{𝟒}`$ and $`\mathrm{𝟓}`$, respectively. $`𝐈:𝐛=\mathrm{𝟐}`$
$`f(d,2)={\displaystyle \frac{5d+3}{(d+1)^2}}`$
$`g(d,2)={\displaystyle \frac{d1}{(d+1)^2}}`$
$`h(d,2)={\displaystyle \frac{2}{(d+1)^2}}`$
$`D_2(d,2)={\displaystyle \frac{\mathrm{log}\frac{d+3}{d+1}}{\mathrm{log}2}}`$
$`\overline{\lambda }(d,2)={\displaystyle \frac{\mathrm{log}\frac{d+3}{d+2}}{\mathrm{log}2}}`$
$`\mathrm{𝐈𝐈}:𝐛=\mathrm{𝟑}`$
$`f(d,3)`$ $`=`$ $`{\displaystyle \frac{(8d^3+88d^2+145d+59)}{2d+3)^2(d+1)^2}}`$
$`g(d,3)`$ $`=`$ $`{\displaystyle \frac{4d^3+20d^2+d24}{(2d+3)^2(d+1)^2}}`$
$`h(d,3)`$ $`=`$ $`{\displaystyle \frac{(4d^2+24d+25)}{(2d+3)^2(d+1)^2}}`$
$`D_2(d,3)`$ $`=`$ $`{\displaystyle \frac{\mathrm{log}\frac{2d^2+9d+19}{(2d+3)(d+1)}}{\mathrm{log}3}}`$
$`\overline{\lambda }(d,3)`$ $`=`$ $`{\displaystyle \frac{log\frac{(2d+3)(2d^2+9d+19)}{4d^3+20d^2+41d+31}}{\mathrm{log}3}}`$
III:b=4
$`f(d,4)`$ $`=`$ $`{\displaystyle \frac{128d^7+1792d^6+13936d^5+59116d^4+137757d^3+175421d^2+113267d+28567}{(8d^4+52d^3+125d^2+130d+49)^2}}`$
$`g(d,4)`$ $`=`$ $`{\displaystyle \frac{64d^7+832d^6+4768d^5+13448d^4+16313d^3701d^218501d11319}{(8d^4+52d^3+125d^2+130d+49)^2}}`$
$`h(d,4)`$ $`=`$ $`{\displaystyle \frac{(64d^6+896d^5+5664d^4+19096d^3+35061d^2+32970d+12397)}{(8d^4+52d^3+125d^2+130d+49)^2}}`$
$`D_2(d,4)`$ $`=`$ $`{\displaystyle \frac{\mathrm{log}\frac{8d^4+68d^3+253d^2+588d+539}{8d^4+52d^3+125d^2+130d+49}}{\mathrm{log}4}}`$
$`\overline{\lambda }(d,4)`$ $`=`$ $`{\displaystyle \frac{\mathrm{log}\frac{8d^4+68d^3+253d^2+588d+539}{8d^4+52d^3+125d^2+130d+49}}{\mathrm{log}4}}`$
IV:b=5
$`f(d,5)={\displaystyle \frac{1}{Q^2}}(73728d^{13}+2162688d^{12}+29576192d^{11}+258155264d^{10}+1614743456d^9`$
$`+7530179904d^8+26333589428d^7+68567523880d^6+131200269465d^5`$
$`+180815964435d^4+173716650934d^3+109891587638d^2+40940417277d`$
$`+6768087791)`$
$`g(d,5)={\displaystyle \frac{1}{Q^2}}(36864d^{13}+1044480d^{12}+13706752d^{11}+110574336d^{10}+607189008d^9`$
$`+2357625920d^8+6492213656d^7+12314821608d^6+14686625629d^5`$
$`+7431447086d^46360742466d^313852397352d^29571181547d`$
$`2499415382)`$
$`h(d,5)={\displaystyle \frac{1}{Q^2}}(36864d^{12}+1081344d^{11}+14788096d^{10}+125353216d^9+732200464d^8`$
$`+3083843024d^7+9521524696d^6+21544010108d^5+35237710633d^4`$
$`+40459001364d^3+30870831766d^2+14032614408d+2871299265)`$
$`D_2(d,5)={\displaystyle \frac{log\frac{192d^7+3104d^6+22348d^5+95720d^4+280525d^3+559419d^2+652155d+320017}{Q}}{\mathrm{log}5}}`$
$`\overline{\lambda }(d,5)={\displaystyle \frac{\mathrm{log}\frac{(192d^7+3104d^6+22348d^5+95720d^4+280525d^3+559419d^2+652155d+320017)Q}{(d+1)P}}{\mathrm{log}5}}`$
with $`𝐏`$ and $`𝐐`$ defined as
$`P=(36864d^{13}+1044480d^{12}+13706752d^{11}+110574336d^{10}+614600976d^9`$
$`+2499189440d^8+7689210552d^7+18190236812d^6+33121369305d^5`$
$`+45749851193d^4+46378189714d^3+32473949342d^2+13985550557d`$
$`+2781136877)`$
and
$$Q=(192d^6+2528d^5+13804d^4+39844d^3+63997d^2+54158d+18825)(d+1).$$
It is straightforward to see that these results will also hold true for $`(𝐝+\mathrm{𝟏})`$-honeycomb fractal conductor network with decimation number $`𝐛`$, which can be constructed from a given $`(𝐝+\mathrm{𝟏})`$-simplex fractal conductor network, simply by replacing the resistors in the links with the resistors which connect the center of a subfractal to its vertices, (see Fig. 2) where, this has also been shown for $`𝐝=\mathrm{𝟐}`$ and $`𝐛=\mathrm{𝟐}`$ case in reference .
Appendix I: Calculation of currents of $`𝐛=\mathrm{𝟑}`$.
Here in this Appendix we give the detail of calculation of inner inward flowing currents corresponding to decimation number $`𝐛=\mathrm{𝟑}`$
Following the procedure of section IV, for $`𝐛=\mathrm{𝟑}`$ we have
$$I_{0,,0,\underset{jth}{\underset{}{2}},0,,0}(0,,0,\underset{jth}{\underset{}{3}},0,,0)=I_j$$
$$I_{0,,0,\underset{jth}{\underset{}{2}},0,,0}(0,,0,\underset{jth}{\underset{}{2}},0,,0,\underset{kth}{\underset{}{1}},0,,0)=$$
$$a_2(2)I_j+b_{21}(2)I_j+b_{21}(1)I_k$$
$$I_{0,,0,\underset{jth}{\underset{}{1}},0,,0,\underset{kth}{\underset{}{1}},0,,0}(0,,0,\underset{jth}{\underset{}{2}},0,,0,\underset{kth}{\underset{}{1}},0,,0)=$$
$$a_{11}(1)(I_j+I_k)+b_{21}(2)I_j+b_{21}(1)I_k$$
$$I_{0,,0,\underset{jth}{\underset{}{1}},0,,0,\underset{kth}{\underset{}{1}},0,,0}(0,,0,\underset{jth}{\underset{}{1}},0,,0,\underset{kth}{\underset{}{1}},0,,0,\underset{lth}{\underset{}{1}},0,,0)=$$
$$a_{11}(1)(I_j+I_k)+b_{111}(1)(I_j+I_k+I_l).$$
Using equation (5-2a) in subfractal $`(\mathrm{𝟎},,\mathrm{𝟎},\underset{𝐣\mathrm{𝐭𝐡}}{\underset{}{\mathrm{𝟐}}},\mathrm{𝟎},,\mathrm{𝟎})`$, we get
$$1+d(a_2(2)+b_{21}(2))b_{21}(1)=0,$$
also using equation (5-2a) in subfractal $`(\mathrm{𝟎},\mathrm{},\mathrm{𝟎},\mathrm{𝟏}_𝐣,\mathrm{𝟎},\mathrm{},\mathrm{𝟏}_𝐤,\mathrm{𝟎},\mathrm{},\mathrm{𝟎})`$ we get
$$(d+1)a_{11}(1)+b_{21}(1)+b_{21}(2)+(d2)b_{111}(1)=0,$$
also, for vertices equation (5-2b) gives
$`a_2(2)+2b_{21}(2)+a_{11}(1)=0`$
$`a_{11}(1)+2b_{21}(1)=0`$
$`2a_{11}(1)+3b_{111}(1)=0.`$
By solving the above equations we can determine inner inward flowing currens corresponding to decimation number $`𝐛=\mathrm{𝟑}`$ which is given in section V.
Appendix II: Calculation of currents of $`𝐛=\mathrm{𝟒}`$.
Here in this Appendix we give the detail of calculation of inner inward flowing currents corresponding to decimation number $`𝐛=\mathrm{𝟒}`$
Similarly, following the procedure of section IV, for $`𝐛=\mathrm{𝟒}`$ we have
$$I_{0,,0,\underset{jth}{\underset{}{3}},0,,0}(0,,0,\underset{jth}{\underset{}{4}},0,,0)=I_j$$
$$I_{0,,0,\underset{jth}{\underset{}{3}},0,,0}(0,,0,\underset{jth}{\underset{}{3}},0,,0,\underset{kth}{\underset{}{1}},0,,0)=$$
$$a_3(3)I_j+b_{31}(3)I_j+b_{31}(1)I_k$$
$$I_{0,,0,\underset{jth}{\underset{}{2}},0,,0,\underset{kth}{\underset{}{1}},0,,0}(0,,0,\underset{jth}{\underset{}{3}},0,,0,\underset{kth}{\underset{}{1}},0,,0)=$$
$$a_{21}(2)I_j+a_{21}(1)I_k+b_{31}(3)I_j+b_{31}(1)I_k$$
$$I_{0,,0,\underset{jth}{\underset{}{2}},0,,0,\underset{kth}{\underset{}{1}},0,,0}(0,,0,\underset{jth}{\underset{}{2}},0,,0,\underset{kth}{\underset{}{2}},0,,0)=$$
$$a_{21}(2)I_j+a_{21}(1)I_k+b_{22}(2)(I_j+I_k)$$
$$I_{0,,0,\underset{jth}{\underset{}{2}},0,,0,\underset{kth}{\underset{}{1}},0,,0}(0,,0,\underset{jth}{\underset{}{2}},0,,0,\underset{kth}{\underset{}{1}},0,,0,\underset{lth}{\underset{}{1}},0,,0)=$$
$$a_{21}(2)I_j+a_{21}(1)+b_{211}(2)I_j+b_{211}(1)(I_k+I_l)$$
$$I_{0,,0,\underset{jth}{\underset{}{1}},0,,0,\underset{kth}{\underset{}{1}},0,,0,\underset{lth}{\underset{}{1}},0,,0}(0,\mathrm{\hspace{0.17em}0},\underset{jth}{\underset{}{2}},0,\mathrm{\hspace{0.17em}0},\underset{kth}{\underset{}{1}},0,\mathrm{\hspace{0.17em}0},\underset{lth}{\underset{}{1}},0,\mathrm{\hspace{0.17em}0})=$$
$$a_{111}(1)(I_j+I_k+I_l)+b_{211}(2)I_j+b_{211}(1)(I_k+I_l)$$
$$I_{0,.,0,\underset{jth}{\underset{}{1}},0,.,0,\underset{kth}{\underset{}{1}},0,.,0,\underset{lth}{\underset{}{1}},0,.,0}(0,.,0,\underset{jth}{\underset{}{1}},0,.,0,\underset{kth}{\underset{}{1}},0,.,0,\underset{lth}{\underset{}{1}},0,,\underset{mth}{\underset{}{1}},0,.,0)=$$
$$a_{111}(1)(I_j+I_k+I_l)+b_{1111}(1)(I_j+I_k+I_l+I_m).$$
Now, imposing Kirchhoff’s law on subfractals and vertices, we get the following equations for $`𝐚`$ and $`𝐛`$
$`1+da_3(3)+b_{31}(3)b_{31}(1)=0`$
$`(d+1)a_{21}(2)+b_{31}(3)+b_{22}(2)+(d1)b_{211}(2)b_{211}(1)=0`$
$`(d+1)a_{21}(2)+b_{31}(1)+b_{22}(2)+(d2)b_{211}(1)=0`$
$`(d+1)a_{111}(2)+b_{211}(2)+(d3)b_{1111}(1)=0`$
$`a_{21}(1)+2b_{31}(1)=0`$
$`a_3(3)+b_{21}(2)+2b_{31}(3)=0`$
$`a_{21}(2)+a_{21}(1)+2b_{22}(2)=0`$
$`2a_{21}(2)+a_{111}(1)+3b_{211}(2)=0`$
$`a_{21}(1)+a_{111}(1)+3b_{211}(1)=0`$
$`3a_{111}(1)+4b_{1111}(1)=0.`$
By solving the above equations we can determine inner inward flowing currents corresponding to decimation number $`𝐛=\mathrm{𝟒}`$ which appear in section V.
Appendix III: Calculation of currents of $`𝐛=\mathrm{𝟓}`$.
Here in this Appendix we give the detail of calculation of inner inward flowing currents corresponding to decimation number $`𝐛=\mathrm{𝟓}`$
Finally following the procedure of section IV, for $`𝐛=\mathrm{𝟓}`$ we have
$$I_{0,,0,\underset{jth}{\underset{}{4}},0,,0}(0,,0,\underset{jth}{\underset{}{5}},0,,0)=I_j$$
$$I_{0,,0,\underset{jth}{\underset{}{4}},0,,0}(0,,0,\underset{jth}{\underset{}{4}},0,,0,\underset{kth}{\underset{}{1}},0,,0)=$$
$$a_4(4)I_j+b_{41}(4)I_j+b_{41}(1)I_k$$
$$I_{0,,0,\underset{jth}{\underset{}{3}},0,,0,\underset{kth}{\underset{}{1}},0,,0}(0,,0,\underset{jth}{\underset{}{4}},0,,0,\underset{kth}{\underset{}{1}},0,,0)=$$
$$a_{31}(3)I_j+a_{31}(1)I_k+b_{41}(4)I_j+b_{41}(1)I_k$$
$$I_{0,,0,\underset{jth}{\underset{}{3}},0,,0,\underset{kth}{\underset{}{1}},0,,0}(0,,0,\underset{jth}{\underset{}{3}},0,,0,\underset{kth}{\underset{}{2}},0,,0)=$$
$$a_{31}(3)I_j+a_{31}(1)I_k+b_{32}(3)I_j+b_{32}(3)I_k$$
$$I_{0,,0,\underset{jth}{\underset{}{3}},0,,0,\underset{kth}{\underset{}{1}},0,,0}(0,,0,\underset{jth}{\underset{}{3}},0,,0,\underset{kth}{\underset{}{1}},0,,0,\underset{lth}{\underset{}{1}},0,,0)=$$
$$a_{31}(3)I_j+a_{31}(1)+b_{311}(3)I_j+b_{311}(1)(I_k+I_l)$$
$$I_{0,,0,\underset{jth}{\underset{}{2}},0,,0,\underset{kth}{\underset{}{2}},0,,0}(0,,0,\underset{jth}{\underset{}{3}},0,,0,\underset{kth}{\underset{}{2}},0,,0)=$$
$$a_{22}(2)(I_j+I_k)+b_{32}(3)I_j+b_{32}(2)I_k$$
$$I_{0,.,0,\underset{jth}{\underset{}{2}},0,.,0,\underset{kth}{\underset{}{2}},0,.,0}(0,.,0,\underset{jth}{\underset{}{2}},0,.,0,\underset{kth}{\underset{}{2}},0,.,0,\underset{lth}{\underset{}{1}},0,.,0)=$$
$$a_{22}(2)(I_j+I_k)+b_{221}(2)(I_j+I_k)+b_{221}(1)I_l$$
$$I_{0,.,0,\underset{jth}{\underset{}{2}},0,.,0,\underset{kth}{\underset{}{1}},0,.,0,\underset{lth}{\underset{}{1}},0,.,0}(0,.,0,\underset{jth}{\underset{}{3}},0,.,0,\underset{kth}{\underset{}{1}},0,.,0,\underset{lth}{\underset{}{1}},0,.,0)=$$
$$a_{211}(2)I_j+a_{211}(1)(I_k+I_l)+b_{311}(3)I_j+b_{311}(1)(I_k+I_l)$$
$$I_{0,.,0,\underset{jth}{\underset{}{2}},0,.,0,\underset{kth}{\underset{}{1}},0,.,0,\underset{lth}{\underset{}{1}},0,.,0}(0,.,0,\underset{jth}{\underset{}{2}},0,.,0,\underset{kth}{\underset{}{2}},0,.,0,\underset{lth}{\underset{}{1}},0,.,0)=$$
$$a_{211}(2)I_j+a_{211}(1)(I_k+I_l)+b_{221}(2)(I_j+I_k)+b_{221}(1)I_l$$
$$I_{0,.,0,\underset{jth}{\underset{}{2}},0,.,0,\underset{kth}{\underset{}{1}},0,.,0,\underset{lth}{\underset{}{1}},0,.,0}(0,.,0,\underset{jth}{\underset{}{2}},0,.,0,\underset{kth}{\underset{}{1}},0,.,0,\underset{lth}{\underset{}{1}},0,,\underset{mth}{\underset{}{1}},0,.,0)=$$
$$a_{211}(2)I_j+a_{211}(1)(I_k+I_l)+b_{2111}(2)I_j+b_{2111}(1)(I_k+I_l+I_m)$$
$$I_{0,.,0,\underset{jth}{\underset{}{1}},0,.,0,\underset{kth}{\underset{}{1}},0,.,0,\underset{lth}{\underset{}{1}},0,.,0,.,0,\underset{mth}{\underset{}{1}},0,.,0}(0,.,0,\underset{jth}{\underset{}{2}},0,.,0,\underset{kth}{\underset{}{1}},0,.,0,\underset{lth}{\underset{}{1}},0,.,0,.,0,\underset{mth}{\underset{}{1}},0,.,0)=$$
$$a_{111}(1)(I_j+I_k+I_l+I_m)+b_{2111}(2)I_j+b_{2111}(1)(I_k+I_l+I_m)$$
$$I_{0,.,0,\underset{jth}{\underset{}{1}},0,.,0,\underset{kth}{\underset{}{1}},0,.,0,\underset{lth}{\underset{}{1}},0,.,0,.,0,\underset{mth}{\underset{}{1}},0,.,0}(0,.,0,\underset{jth}{\underset{}{1}},0,.,0,\underset{kth}{\underset{}{1}},0,.,0,\underset{lth}{\underset{}{1}},0,.,0,.,0,\underset{mth}{\underset{}{1}},0,.,0\underset{nth}{\underset{}{1}},0,.,0)=$$
$$a_{111}(1)(I_j+I_k+I_l+I_m)+b_{11111}(1)(I_j+I_k+I_l+I_m+I_n).$$
Again imposing Kirchhoff’s law on subfractals and vertices, we get the following equations for $`𝐚`$ and $`𝐛`$
$`1+da_4(4)+b_{41}(3)b_{41}(1)=0`$
$`(d+1)a_{31}(3)+b_{41}(4)+b_{32}(3)+(d1)b_{311}(3)b_{311}(1)=0`$
$`(d+1)a_{31}(1)+b_{41}(1)+b_{32}(2)+(d2)b_{311}(1)=0`$
$`(d+1)a_{22}(1)+b_{32}(3)+b_{32}(2)+(d1)b_{221}(2)b_{221}(1)=0`$
$`(d+1)a_{211}(2)+b_{311}(3)+2b_{221}(2)+(d2)b_{2111}(2)b_{2111}(1)=0`$
$`(d+1)a_{211}(2)+b_{311}(1)+b_{221}(2)+b_{221}(1)+(d3)b_{2111}(1)=0`$
$`(d+1)a_{1111}(2)+b_{2111}(2)+3b_{2111}(1)+(d4)b_{11111}(1)=0`$
$`a_4(4)+a_{31}(3)+2b_{41}(4)=0`$
$`a_{31}(1)+2b_{41}(1)=0`$
$`a_{31}(3)+a_{22}(2)+2b_{32}(3)=0`$
$`a_{31}(1)+a_{22}(2)+2b_{32}(2)=0`$
$`2a_{31}(3)+a_{211}(2)+3b_{311}(3)=0`$
$`a_{31}(1)+a_{211}(1)+3b_{311}(1)=0`$
$`a_{22}(2)+a_{211}(2)+a_{211}(1)+3b_{221}(2)=0`$
$`2a_{211}(1)+3b_{221}(1)=0`$
$`3a_{211}(2)+a_{1111}(1)+4b_{2111}(2)=0`$
$`2a_{211}(1)+a_{1111}(1)+4b_{2111}(1)=0`$
$`4a_{1111}(1)+5b_{11111}(1)=0.`$
By solving the above equations we can determine inner inward flowing currens corresponding to decimation number $`𝐛=\mathrm{𝟓}`$ which appear in section V.
.
Conclusion
Here in this work it has been rigorously shown that the macroscopic isotropy will be restored if the corresponding renormalization map between two different scales has properties such as: positivity, homogeneity of first order and most of all monotonically increasing property. Obviously, homogeneity is enough and the order of homogeneity does not play a very important role. It is clear that this can be true in many physical phenomena, where we quote only very few of them here: Diffusion in inhomogeneous media , elasticity property of rubber or the network of polymer chains, conductivity in random resitor network , flux distribution in josephson junction networks. It would be rather interesting to see whether there exist the restoration isotropy which is not due to positive, homogeneous, and monotonically increasing renormalization map of two different scales.
ACKNOWLEDGEMENT
We wish to thank Dr. S. K. A. Seyed Yagoobi for his careful reading the article and for his constructive comments.
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# Localization length and impurity dielectric susceptibility in the critical regime of the metal-insulator transition in homogeneously doped 𝑝-type Ge
## Abstract
We have determined the localization length $`\xi `$ and the impurity dielectric susceptibility $`\chi _{\mathrm{imp}}`$ as a function of Ga acceptor concentrations ($`N`$) in nominally uncompensated <sup>70</sup>Ge:Ga just below the critical concentration ($`N_c`$) for the metal-insulator transition. Both $`\xi `$ and $`\chi _{\mathrm{imp}}`$ diverge at $`N_c`$ according to the functions $`\xi (1N/N_c)^\nu `$ and $`\chi _{\mathrm{imp}}(N_c/N1)^\zeta `$, respectively, with $`\nu =1.2\pm 0.3`$ and $`\zeta =2.3\pm 0.6`$ for $`0.99N_c<N<N_c`$. Outside of this region ($`N<0.99N_c`$), the values of the exponents drop to $`\nu =0.33\pm 0.03`$ and $`\zeta =0.62\pm 0.05`$. The effect of the small amount of compensating dopants that are present in our nominally uncompensated samples, may be responsible for the change of the critical exponents at $`N0.99N_c`$.
The metal-insulator transition (MIT) in doped semiconductors is a unique quantum phase transition in the sense that both disorder and electron-electron interaction play a key role. Important information about the MIT is provided by the values of the critical exponents for the zero-temperature conductivity, correlation length, localization length, and impurity dielectric susceptibility. From a theoretical point of view, the correlation length in the metallic phase and the localization length in the insulating phase diverge at the critical point with the same exponent $`\nu `$, i.e., they are proportional to $`|N/N_c1|^\nu `$ in the critical regime of the MIT. ($`N`$ is the dopant concentration and $`N_c`$ is the critical concentration for the MIT.) Since direct experimental determination of $`\nu `$ is extremely difficult, researchers have usually determined, instead of $`\nu `$, the value of $`\mu `$ defined by $`\sigma (0)(N/N_c1)^\mu `$ where $`\sigma (0)`$ is the conductivity extrapolated to $`T=0`$ It is also possible to evaluate $`\mu `$ from finite-temperature scaling of the form $`\sigma (N,T)T^xf(|N/N_c1|/T^y)`$ where $`x/y`$ is equivalent to $`\mu `$ Values of $`\nu `$ are then obtained assuming $`\nu =\mu `$ for three-dimensional systems.
In this work we have determined directly the localization length $`\xi `$ and the impurity dielectric susceptibility $`\chi _{\mathrm{imp}}`$ in neutron-transmutation-doped (NTD), nominally uncompensated <sup>70</sup>Ge:Ga just below $`N_c`$. The application of NTD to isotopically enriched <sup>70</sup>Ge leads to unsurpassed doping homogeneity and precisely controlled doping concentration. As a result, we have been able to approach the transition as close as $`0.999N_c`$ from the insulating side and $`1.0004N_c`$ from the metallic side. In zero magnetic field, the low-temperature resistivity of the samples is described by variable-range hopping (VRH) conduction within the Coulomb gap. Magnetic field and temperature dependence of the resistivity are subsequently measured in order to determine directly $`\xi `$ and $`\chi _{\mathrm{imp}}`$ in the context of the VRH theory.
This kind of determination of $`\xi `$ and $`\chi _{\mathrm{imp}}`$ was performed for compensated Ge:As by Ionov et al. They found $`\xi (1N/N_c)^\nu `$ and $`\chi _{\mathrm{imp}}(N_c/N1)^\zeta `$ with $`\nu =0.60\pm 0.04`$ and $`\zeta =1.38\pm 0.07`$, respectively, for $`N<0.96N_c`$. The significance of their result is the experimental verification of the relation $`2\nu \zeta `$ that had been predicted by scaling theories. However, the critical exponents of compensated samples are known to be different from those of nominally uncompensated samples. Therefore, the present work which probes $`\xi `$ and $`\chi _{\mathrm{imp}}`$ in nominally uncompensated samples is relevant for the fundamental understanding of the MIT. The previous effort to measure $`\chi _{\mathrm{imp}}`$ has also contributed. Hess et al. found $`\zeta =1.15\pm 0.15`$ in nominally uncompensated Si:P. Since $`\mu 0.5`$ was determined for the same series of Si:P samples, $`2\mu \zeta `$ was again valid. Katsumoto has found $`\zeta 2`$ and $`\mu 1`$ for compensated Al<sub>x</sub>Ga<sub>1-x</sub>As:Si, i.e., again, $`2\mu \zeta `$ applies. Thus, in these cases the conclusion $`2\nu \zeta `$ was reached indirectly, by assuming $`\mu =\nu `$. The work reported here, on the other hand, determines $`\nu `$ directly, i.e., we do not have to rely on the assumption $`\mu =\nu `$ in order to study the behavior of $`\xi `$ near $`N_c`$.
All of the <sup>70</sup>Ge:Ga samples used in this study were prepared by NTD of isotopically enriched <sup>70</sup>Ge single crystals. We use the NTD process since it is known to produce the most homogeneous dopant distribution. Details of the sample preparation and characterization are described elsewhere. In this study, we determined the low-temperature ($`0.050.5`$ K) resistivity of nine samples in weak magnetic fields ($`<0.4`$ T) applied in the direction perpendicular to the current flow.
The electrical conduction of doped semiconductors on the insulating side of the MIT is often dominated by VRH at low temperatures. The temperature dependence of the resistivity $`\rho (T)`$ for VRH is written in the form of
$$\rho (T)=\rho _0(T)\mathrm{exp}[(T_0/T)^p],$$
(1)
where $`p=1/2`$ for the excitation within a parabolic-shaped energy gap (the Coulomb gap), and $`p=1/4`$ for a constant density of states around the Fermi level. In our earlier work, we reported that $`p=1/2`$ for $`N<0.991N_c`$ ($`N_c=1.860\times 10^{17}`$ cm<sup>-3</sup>) and that $`p`$ decreases rapidly as $`N`$ approaches $`N_c`$ from $`0.991N_c`$ and becomes even smaller than 1/4 when we neglect the temperature variation of $`\rho _0(T)`$. However, the variation contributes greatly to the temperature dependence of $`\rho (T)`$ near $`N_c`$ because the factor $`T_0/T`$ in the exponential terms become very small, i.e., the temperature dependencies of $`\rho _0(T)`$ and that of the exponential term become comparable. Theoretically, $`\rho _0(T)`$ is expected to vary as $`\rho _0T^r`$ but the value of $`r`$ including the sign has not been derived yet for VRH with both $`p=1/2`$ (Ref. REFERENCES) and $`p=1/4`$ (Ref. REFERENCES).
Recently, we have shown that the temperature dependence of the conductivity of the same series of <sup>70</sup>Ge:Ga samples within $`\pm `$0.3% of $`N_c`$ is proportional to $`T^{1/3}`$ at $`0.021`$ K. Since both the $`T^{1/3}`$ dependence of the conductivity and the Efros-Shklovskii VRH are results of the electron-electron interaction in disordered systems, they can be expressed, in principle, in a unified form. Moreover, the electronic transport in barely metallic samples and that in barely insulating samples should be essentially the same at high temperatures so long as the inelastic scattering length and the thermal diffusion length are smaller than, or at most comparable to the correlation length or the localization length. So, the temperature dependence of conductivity at high temperatures should be the same on both sides of the transition. Such behavior is confirmed experimentally in the present system, i.e., the conductivity of samples very close to $`N_c`$ shows a $`T^{1/3}`$ dependence at $`T0.5`$ K, irrespective of the phase (metal or insulator) to which they belong at $`T=0`$. Based on this consideration we fix $`r=1/3`$. Figure 1 shows $`\rho T^r`$ with $`r=1/3`$ for four samples ($`N/N_c=0.993`$, 0.994, 0.996, and 0.998) as a function of (a) $`T^{1/2}`$ and (b) $`T^{1/4}`$. All the data points lie on straight lines with $`p=1/2`$ in Fig. 1(a) while they curve upward with $`p=1/4`$ in Fig. 1(b). This dependence is maintained even when we change the values of $`r`$ between 1/2 and 1/4. Thus we conclude that the resistivity of all samples for $`N`$ up to 0.998$`N_c`$ is described by the VRH theory where the excitation occurs within the Coulomb gap, i.e., Eq. (1) with $`p=1/2`$.
Based on these findings, we evaluate $`T_0`$ in Eq. (1) with $`p=1/2`$ and $`r=1/3`$, and show it as a function of $`1N/N_c`$ in Fig. 2. The vertical and horizontal error bars have been estimated based on the values of $`T_0`$ obtained with $`r=1/2`$ and $`r=1/4`$, and the values of $`1N/(1.858\times 10^{17}`$ cm$`{}_{}{}^{3})`$ and $`1N/(1.861\times 10^{17}`$ cm$`{}_{}{}^{3})`$, where $`1.858\times 10^{17}`$ cm<sup>-3</sup> is the highest concentration in the insulating phase and $`1.861\times 10^{17}`$ cm<sup>-3</sup> is the lowest in the metallic phase, respectively.
According to theory, $`T_0`$ in Eq. (1) is given by
$$k_BT_02.8e^2/4\pi ϵ_0ϵ(N)\xi (N)$$
(2)
in SI units, where $`ϵ(N)`$ is the dielectric constant. Here, we should note that the condition $`T<T_0`$ is needed for the theory to be valid, i.e., $`T_0`$ has to be evaluated only from the data obtained at temperatures low enough to satisfy the condition. This requirement is fulfilled in Fig. 2 for all the samples except for the one with $`N=0.998N_c`$. Concerning this latter sample, we will include it for the determination of $`\xi `$ and $`\chi _{\mathrm{imp}}`$ (Fig. 5) but not for the calculation of the critical exponents.
Our next step is to separate $`T_0`$ into $`ϵ`$ and $`\xi `$. For $`\xi /\lambda 1`$, the magnetoresistance is expressed as
$$\mathrm{ln}[\rho (B,T)/\rho (0,T)]0.0015(\xi /\lambda )^4(T_0/T)^{3/2},$$
(3)
where $`\lambda \sqrt{\mathrm{}/eB}`$ is the magnetic length in SI units. According to Eq. (3), the magnetic-field variation of $`\mathrm{ln}\rho `$ at $`T=\mathrm{const}.`$ is proportional to $`B^2`$, i.e., $`\mathrm{ln}\rho (B,T)=\mathrm{ln}\rho (0,T)+C(T)B^2`$, and the slope $`C(T)`$ is proportional to $`T^{3/2}`$. In order to demonstrate that these relations hold for our samples, we show for the $`N=0.989N_c`$ sample $`\mathrm{ln}\rho (B,T)`$ vs $`B^2`$ in Fig. 3 and $`C(T)`$ determined by least-square fitting of $`\mathrm{ln}\rho /B^2`$ vs $`T^{3/2}`$ in Fig. 4. Since Eq. (3) is equivalent to
$$\gamma C(T)/T^{3/2}0.0015(e/\mathrm{})^2\xi ^{\mathrm{\hspace{0.17em}4}}T_0^{\mathrm{\hspace{0.17em}3}/2},$$
(4)
$`\xi `$ is given by
$$\xi 5.1(\mathrm{}/e)^{1/2}\gamma ^{1/4}T_0^{3/8}.$$
(5)
In this way we have determined $`\gamma `$ for nine samples. The inset of Fig. 5 shows $`\gamma `$ as a function of $`T_0`$. The value of $`\gamma `$ is almost independent of $`T_0`$, and if one assumes $`\gamma T_0^\delta `$, one obtains a small value of $`\delta =0.094\pm 0.005`$ from least-square fitting. Figure 5 shows $`\xi `$ and $`\chi _{\mathrm{imp}}=ϵϵ_h`$ determined from Eqs. (2) and (5). Here, $`ϵ_h`$ is the dielectric constant of the host Ge, and hence, $`\chi _{\mathrm{imp}}`$ is the dielectric susceptibility of the Ga acceptors. We should note that both $`\xi `$ and $`\chi _{\mathrm{imp}}`$ are sufficiently larger than the Bohr radius (8 nm for Ge) and $`ϵ_h=15.4`$ (Ref. REFERENCES), respectively. According to the theories of the MIT, both $`\xi `$ and $`\chi _{\mathrm{imp}}`$ diverge at $`N_c`$ as $`\xi (N)(1N/N_c)^\nu `$ and $`\chi _{\mathrm{imp}}(N)(N_c/N1)^\zeta `$, respectively. We find, however, both $`\xi `$ and $`\chi _{\mathrm{imp}}`$ do not show such simple dependencies on $`N`$ in the range shown in Fig. 5, and that there is a sharp change of both dependencies at $`N0.99N_c`$. On both sides of the change in slope, the concentration dependence of $`\xi `$ and $`\chi _{\mathrm{imp}}`$ are expressed well by the scaling formula as shown in Fig. 5. Theoretically, the quantities should show the critical behavior when $`N`$ is very close to $`N_c`$. So $`\nu =1.2\pm 0.3`$ and $`\zeta =2.3\pm 0.6`$ may be concluded from the data in $`0.99<N/N_c`$. However, the other region ($`0.9<N/N_c<0.99`$), where we obtain $`\nu =0.33\pm 0.03`$ and $`\zeta =0.62\pm 0.05`$, is also very close to $`N_c`$ in a conventional experimental sense.
As a possible origin for the change in slope, we refer to the effect of compensation. Although our samples are nominally uncompensated, doping compensation of less than 0.1% may present due to residual isotopes that become $`n`$-type impurities after NTD. In addition to the doping compensation, the effect known as “self compensation” may play an important role near $`N_c`$ It is empirically known that the doping compensation affects the value of the critical exponents. Rentzsch et al. studied VRH conduction of n-type Ge in the concentration range of $`0.2<N/N_c<0.91`$, and showed that $`T_0`$ vanishes as $`T_0(1N/N_c)^\alpha `$ with $`\alpha 3`$ for $`K=38`$% and 54%, where $`K`$ is the compensation ratio. Since $`\alpha \nu +\zeta `$ \[Eq. (2)\], we find for our <sup>70</sup>Ge:Ga samples $`\alpha =3.5\pm 0.8`$ for $`0.99<N/N_c<1`$ and $`\alpha =0.95\pm 0.08`$ for $`0.9<N/N_c<0.99`$. Interestingly, $`\alpha =3.5\pm 0.8`$ agrees with $`\alpha 3`$ found for compensated samples. Moreover, we have recently proposed the possibility that the conductivity critical exponent $`\mu 1`$ in the same <sup>70</sup>Ge:Ga only within the very vicinity of $`N_c`$ (up to about +0.1% of $`N_c`$). An exponent of $`\mu =0.50\pm 0.04`$, on the other hand, holds for a wide region of $`N`$ up to $`1.4N_c`$ Again, $`\mu 1`$ near $`N_c`$ may be viewed as the effect of compensation. Therefore, it may be possible that the region of $`N`$ around $`N_c`$ where $`\nu 1`$ and $`\mu 1`$ changes its width as a function of the doping compensation. In the limit of zero compensation, the part which is characterized by $`\nu 1`$ and $`\mu 1`$ vanishes, i.e., we propose $`\nu =0.33\pm 0.03`$, $`\zeta =0.62\pm 0.05`$, and $`\mu =0.50\pm 0.04`$ for truly uncompensated systems and that Wegner’s scaling law of $`\nu =\mu `$ is not satisfied. In compensated systems, on the other hand, Wegner’s law may hold as it does in the very vicinity of $`N_c`$. The experiment on compensated Al<sub>x</sub>Ga<sub>1-x</sub>As:Si that showed $`\zeta 2`$ and $`\mu 1`$ (Ref. REFERENCES) is also consistent with the law. However, our preceding discussion needs to be proven in the future by experiments in samples with precisely and systematically controlled compensation ratios.
In summary, we have determined directly the localization length and the dielectric susceptibility arising from the impurities in nominally uncompensated NTD <sup>70</sup>Ge:Ga samples near the critical point for the MIT. While the relation $`2\nu \zeta `$ predicted by scaling theory holds for $`0.9<N/N_c<1`$, the critical exponents for localization length and impurity susceptibility change at $`N/N_c0.99`$. The small amount of doping compensation that is unavoidably present in our samples may be responsible for such a change in the exponents.
We are thankful to T. Ohtsuki, B. I. Shklovskii, and M. P. Sarachik for valuable comments, and V. I. Ozhogin for the supply of the Ge isotope. Most of the low-temperature measurements were carried out at the Cryogenic Center, the University of Tokyo. M. W. would like to thank the Japan Society for the Promotion of Science (JSPS) for financial support. The work at Keio was supported by a Grant-in-Aid for Scientific Research from the Ministry of Education, Science, Sports, and Culture, Japan. The work at Berkeley was supported by the Director, Office of Energy Research, Office of Basic Energy Science, Materials Sciences Division of the U. S. Department of Energy under Contract No. DE-AC03-76SF00098 and U. S. NSF Grant No. DMR-97 32707.
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# Constant Crunch Coordinates for Black Hole Simulations
## I Constant Crunch Surfaces
In this paper, we address a single question: Is there a numerically-viable coordinatization of a Schwarzschild black hole spacetime foliated by hypersurfaces of constant (not necessarily zero) mean extrinsic curvature? In other words, can we coordinatize the Schwarzschild spacetime with constant mean extrinsic curvature ($`Tr(K)=\text{constant}`$) hypersurfaces so as to bound the growth of metric components and their gradients? We demonstrate here that the single shift freedom yields a spacetime metric that is static, and therefore bounds the growth in time of such gradients. A more complete analysis of the stability of our coordinatization, and a more thorough canvassing of the parameter space, will appear elsewhere. Shoemaker2000 Our foliation is consistent with that of Iriondo et al. Iriondo1996 , who provided a generic constant mean curvature (CMC) foliation of the Reissner-Nordström spacetime for the purpose of finding trapped surfaces. In this paper we focus on the utility of CMC slicings for the numerical simulation of black holes, in support of the emerging field of gravity-wave astrophysics.
The trace of the extrinsic curvature tensor ($`Tr(𝐊)=K^a{}_{a}{}^{}=K`$) at a point on a spacelike hypersurface measures the fractional rate of contraction of 3-volume along a unit normal to the surface. It represents the amount of “crunch” the 3-surface is experiencing at the point, at a given time. If all the observers throughout a spacelike hypersurface moving in time orthogonal to the surface experience the same amount of contraction per unit proper time, we say that the surface is a $`K`$-surface or a “constant crunch” surface. In this paper we examine foliations of a single spherically-symmetric, static black hole where each spacelike hypersurface has the same constant value of the extrinsic curvature, $`K`$.
Generic $`K`$-surface foliations have found great utility in the numerical simulation of cosmological spacetimes. Centrella1984 In addition to decoupling the three momentum constraint equations from the Hamiltonian constraint, these surfaces (in the case of compact or W-model universes) provide a convenient cosmological time parameter ($`K`$, or York, time). Wheeler1988 Furthermore, for such cosmological spacetimes one has powerful existence and uniqueness theorems. Tipler1980 ; Tipler1985 Extensive work into the characteristics of these surfaces for Schwarzschild spacetimes has been done by Brill et al.Brill1978 and foundational work into their use in numerical relativity was done by D. Eardley et al.Eardley1978 More recently Pervez et al. Pervez1995 provided a foliation partially covering the Schwarzschild spacetime with $`K`$-surfaces, with $`K`$ ranging from $`\mathrm{}`$ to $`\mathrm{}`$, and Iriondo et al. Iriondo1996 provided a generic constant mean curvature foliation of the Reissner-Nordström spacetime for the purpose of finding trapped surfaces. In this paper we build upon the work of these investigators by examining the utility of these surfaces for numerical relativity in support of gravity-wave detectors.
Although surfaces of constant $`K`$ were thoroughly investigated decades ago, their use in current numerical simulations of black holes is conspicuously absent (apart from the use of maximal ($`K=0`$) surfaces). Alcubierre2000 One reason why these slicing methods have not been more fully developed is that they lag in time close in, to avoid crashing into the singularity, while they simultaneously evolve forward normally at the outer edge of the grid to allow for wave extraction. This tension, many fear, will unavoidingly lead to unbounded growth in the metric and extrinsic curvature components in the intermediate region, as is indeed found in maximal slicing. This computational concern has been referred to by the numerical relativity community as “grid stretching” or “grid sucking.” We show in this paper that a proper choice of radial shift can yield a constant crunch foliation of a spherically-symmetric black hole without such pathologies. In fact, we foliate a Schwarzschild black hole such that the 3-metric and extrinsic curvature are both bounded and static (i.e. unchanging in time).
To numerically evolve a black hole 3-space in time it is desirable to have a foliation, and its coordinatization, which satisfy the following four properties:
1. Avoids black hole singularities or facilitates their excision.
2. Possesses asymptotically null hypersurfaces to aid in radiation extraction.
3. Minimizes steep gradients in the lapse, shift, 3-metric and extrinsic curvature tensor.
4. Maximizes the future development of the initial data for the purpose of gravity-wave extraction.
As a first step towards achieving these goals for systems containing multiple black holes, we explore the families of $`K`$-surfaces in the Schwarzschild spacetime, and find a CMC foliation satisfying the above properties.
In the next section we construct the $`K`$-surfaces outside and inside the horizon. In section III we explore the properties of the $`K`$-surfaces, dwelling in particular on their approaches to the singularity. We also examine and illustrate the three families of $`K`$-surfaces. In section IV we derive a metric for Schwarzschild whose constant time slices are $`K`$-surfaces. We restrict our attention to a subfamily of $`K`$-surfaces – surfaces which, when generalized to the colliding black hole spacetimes, support the gravity wave detection problem. We also present some preliminary numerical simulations using constant-crunch coordinates. We conclude with general comments on the applicability of $`K`$-surfaces to numerical calculations of more general black hole spacetimes.
## II Construction of Constant Crunch Surfaces
Over three decades ago Eardley and Smarr Eardley1978 carried out a generic classification of the spacetimes that could be simulated numerically, and investigated the limitations that the presence of singularities would impose. In their paper they argued that CMC slicings are particularly useful for numerical purposes. In particular, they demonstrated this explicitly by constructing numerical solutions to a wide array of dust collapse models. In a similar vein, Brill et al. Brill1978 explored the nature of CMC slices of the Schwarzschild spacetime, and they also presented some numerical examples. In the present work, we explore the numerical utility of CMC slicings in the case of single black hole spacetimes. For the sake of clarity, we will commence with a re-derivation of the equations governing the CMC surfaces starting from Schwarzschild coordinates. From there, we will explore specific properties of the surfaces, paying particular attention to implications for numerical relativity.
We wish to find a spacelike hypersurface in the Schwarzschild spacetime such that every point on the surface has the same constant value of the trace of the extrinsic curvature tensor. We have at our disposal the specification of the initial-value data, as well as the freedom to choose the lapse and shift throughout the evolution. To begin, let us take the standard coordinate system of a single black hole spacetime of mass $`M`$ in Schwarzschild coordinates:
$$ds^2=B(r)dt^2+C(r)dr^2+r^2\left(d\theta ^2+\mathrm{sin}^2\theta d\phi ^2\right),$$
(1)
with $`B(r)=(12M/r)`$ and $`C(r)=1/B(r)`$. It will be convenient to treat separately the regions inside and outside of the horizon. We will find that the two are related by an isometry.
### II.1 Outside the Horizon ($`r>2M`$)
Outside the horizon, CMC surfaces will be labeled by $`t=T(r)`$ (Fig. 1). The requirement that the trace of the extrinsic curvature be constant throughout this surface yields a first order differential equation for $`T(r)`$, determined by examining the behavior of the normals to the surface. The normal $`𝐧`$ to the spacelike hypersurface $`T`$ is given by,
$$𝐧=N_0(tT(r))=n_tdt+n_rdr=N_0(dtT^{}dr),$$
(2)
where $`N_0`$ is a normalization constant and primes denote differentiation with respect to $`r`$. The normalization is fixed by demanding that
$`𝐧𝐧`$ $`=`$ $`1`$ (3)
$`=`$ $`g^{tt}n_tn_t+g^{rr}n_rn_r`$ (4)
$`=`$ $`N_0^2\left({\displaystyle \frac{1}{B}}+{\displaystyle \frac{1}{C}}T^2\right).`$ (5)
Therefore,
$$N_0=\frac{1}{\sqrt{CBT^2}},$$
(6)
and
$`n_r`$ $`=`$ $`{\displaystyle \frac{T^{}}{\sqrt{CBT^2}}},`$ (7)
$`n_t`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{CBT^2}}}.`$ (8)
The contravariant components of the normal are given by
$`n^r`$ $`=`$ $`g^{rr}n_r={\displaystyle \frac{T^{}}{C\sqrt{CBT^2}}},`$ (9)
$`n^t`$ $`=`$ $`g^{tt}n_t={\displaystyle \frac{1}{B\sqrt{CBT^2}}}.`$ (10)
The trace of the extrinsic curvature is the fractional rate of contraction of 3-volume per unit proper time along the normal, namely
$$K=n_{;\alpha }^\alpha =\frac{1}{r^2}\frac{d}{dr}\left(r^2n^r\right).$$
(11)
Substitution of $`n^r`$ into Eq. (11) yields a second-order ordinary differential equation for $`T`$. Integrating this equation, we find
$$H=\left(\frac{Br^2T^{}}{\sqrt{CBT^2}}\right)+J,$$
(12)
with $`H`$ an integration constant and $`J`$ an indefinite integral given by
$$J=^rK\sqrt{BC}r^2𝑑r=\frac{1}{3}Kr^3.$$
(13)
Along the surface the rate of change of proper time, $`d\tau `$, with proper distance, $`ds`$, is related to the slope of the surface $`T`$,
$$\frac{d\tau }{ds}=\sqrt{\frac{B}{C}}T^{}.$$
(14)
From Eq. (12) we find
$$\left(\frac{d\tau }{ds}\right)^2=\frac{(HJ)^2}{(HJ)^2+Br^4}.$$
(15)
### II.2 Inside the Horizon ($`r<2M`$)
Finding the $`K`$-constant slices of Eq. (1) within the horizon is similar to the calculation done in the previous section; however, as the roles of time and space coordinates reverse within the horizon, we will find it useful to parameterize our spacelike surface as a function of coordinate $`t`$, and look for $`K`$-constant surfaces of the form (Fig. 1)
$$r=R(t).$$
(16)
The normal $`𝐧`$ to the spacelike hypersurface $`R`$ is given by
$$𝐧=N_0(R(t)r)=n_tdt+n_rdr=N_0(\dot{R}dtdr),$$
(17)
with differentiation with respect to $`t`$ denoted by dots and the $`N_0`$ a normalization constant fixed by:
$`𝐧𝐧`$ $`=`$ $`1`$ (18)
$`=`$ $`g^{tt}n_tn_t+g^{rr}n_rn_r`$ (19)
$`=`$ $`N_0^2\left({\displaystyle \frac{1}{C}}{\displaystyle \frac{1}{B}}\dot{R}^2\right).`$ (20)
We have, therefore,
$$N_0=\frac{1}{\sqrt{C\dot{R}^2B}},$$
(21)
and
$`n_r`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{C\dot{R}^2B}}},`$ (22)
$`n_t`$ $`=`$ $`{\displaystyle \frac{\dot{R}}{\sqrt{C\dot{R}^2B}}}.`$ (23)
The contravariant components of the normal are given by
$`n^r`$ $`=`$ $`g^{rr}n_r={\displaystyle \frac{1}{C\sqrt{C\dot{R}^2B}}},`$ (24)
$`n^t`$ $`=`$ $`g^{tt}n_t={\displaystyle \frac{\dot{R}}{B\sqrt{C\dot{R}^2B}}}.`$ (25)
From Eq. (11), we once again find that fixing the trace of the extrinsic curvature gives us a second-order differential equation for $`R(t)`$, namely,
$$K=\frac{2}{CR\sqrt{C\dot{R}^2B}}\frac{2CB^{}\dot{R}^2B\left(B^{}+C^{}\dot{R}^2+2C\ddot{R}\right)}{2\left(C\dot{R}^2B\right)^{3/2}},$$
(26)
which can be simplified to
$$KR^2\dot{R}=\frac{d}{dt}\left(\frac{BR^2}{\sqrt{C\dot{R}^2B}}\right).$$
(27)
Paralleling the approach from the last section, we introduce an integration constant, $`H`$, and an indefinite integral $`J`$ (given by Eq. (13)), to obtain the first integral:
$$H=\left(\frac{BR^2}{\sqrt{C\dot{R}^2B}}\right)+J.$$
(28)
From Eq. (27) we find that the “proper velocity” along the surface, $`ds/d\tau =\sqrt{C/B}\dot{R}`$, results in the same equation both inside and outside of the horizon (Eq. (15)). This can be rewritten as
$$\left(\frac{ds}{d\tau }\right)^2=1+\frac{BR^4}{(HJ)^2}.$$
(29)
The spacelike $`K`$-surfaces obtained from the first integrals, Eqs. (15) and (29), differ only by an isometry,
$$T^{}\frac{1}{\dot{R}}.$$
(30)
## III Properties of the $`K`$-surfaces
The spatial metric of a $`K`$-surface outside of the horizon is given by
$`ds^2`$ $`=`$ $`d\mathrm{}^2+r^2d\mathrm{\Omega }^2`$ (31)
$`=`$ $`\left(CBT^2\right)dr^2+r^2d\mathrm{\Omega }^2.`$ (32)
Within the horizon it becomes
$`ds^2`$ $`=`$ $`d\mathrm{}^2+r^2d\mathrm{\Omega }^2`$ (33)
$`=`$ $`\left(C\dot{R}^2B\right)dt^2+r^2d\mathrm{\Omega }^2.`$ (34)
These two expressions differ by the isometry of Eq. (30). Using Eq. (15) we can rewrite them in terms of $`H`$ and $`K`$:
$$ds^2=\frac{r^4}{\left(HJ\right)^2+Br^4}dr^2+r^2d\mathrm{\Omega }^2.$$
(35)
From this we arrive at the scalar curvature of the $`K`$-surface:
$${}_{}{}^{(3)}R=\frac{2}{3}K^2+\frac{6H^2}{r^6}.$$
(36)
Similarly, by using Eqs. (15) and (29), the extrinsic curvature associated with observers moving on world lines orthogonal to the $`K`$-slices are also expressible in terms of $`K`$ and $`H`$:
$`K_{\widehat{r}}^{\widehat{r}}`$ $`=`$ $`{\displaystyle \frac{1}{3}}K+{\displaystyle \frac{2H}{r^3}},`$ (37)
$`K_{\widehat{\theta }}^{\widehat{\theta }}`$ $`=`$ $`K_{\widehat{\varphi }}^{\widehat{\varphi }}={\displaystyle \frac{1}{3}}K{\displaystyle \frac{H}{r^3}}.`$ (38)
The $`K`$-surfaces are therefore parametrized by two constants: the trace of the extrinsic curvature tensor, $`K`$, and the constant of integration, $`H`$. In addition one must fix a single point on the surface, $`t_o=T(r_o)`$, which amounts to setting a time translation parameter. As can be seen in Eqs. (36)–(38), the constant $`H`$ controls the variation of the intrinsic and extrinsic curvatures over the $`K`$-surface.
To elucidate the nature of the $`K`$-surfaces, we numerically integrate Eqs. (15) and (29). We find that within the horizon there are 3 classes of $`K`$-constant surfaces, differentiated by their asymptotic behavior. The singularity-singularity surfaces (SS ) begin at the singularity aligned with the null surface, reach up towards the horizon, and then fall back, reaching the singularity along the null cone. The horizon-horizon surfaces (HH ), which we have also dubbed “horizon-hugging” surfaces, asymptote to the horizon ($`r2M`$ for $`|t|\mathrm{}`$ in Schwarzschild), dipping down towards the singularity in between. This feature was previously remarked upon by Brill et al. Brill1978 . The asymptotes converge toward a null surface at the horizon. Finally, the horizon-singularity (HS ) surfaces begin at the horizon, and asymptote in to the singularity. Representative surfaces for the value $`K=0.1`$ are shown in Fig. 2. We integrate the HH and HS surfaces across the horizon into the region $`r>2M`$ by imposing continuity of the surface and its first derivative at the horizon. Because of the isometry, Eq. (30), the surfaces outside of the horizon are characteristically similar to those on the inside; in particular, both sets are null at their asymptotes.
We have chosen to use the acronym HH for the horizon-to-horizon hypersurfaces, in lieu of referring to them as “regular” hypersurfaces Brill1978 , as each of the three types of $`K`$-surfaces are, in a strict sense, regular. In particular, each surface asymptotes to a null surface, be it at the singularity or the horizon. Observers on such a surface, or more precisely, observers that are time synchronized throughout the surface, are never seen crossing the horizon, nor do they ever reach the singularity! Outside the horizon, every HH and HS K-surface ($`K0`$) asymptotes for large $`r`$ to null infinity. $`K`$-surfaces corresponding to positive values of $`K`$ asymptote to past null infinity, and asymptote to future null infinity for $`K<0`$.
To gain a qualitative understanding of the $`K`$-constant foliation, it is useful to analyze Eq. (29) as an energy conservation equation for a particle of unit total energy ($`E=1`$) moving in the potential
$$V(r)=\frac{\left(1\frac{2M}{r}\right)r^4}{\left(H\frac{1}{3}Kr^3\right)^2}.$$
(39)
By using this energy equation we can determine the closest approach to the singularity, $`R_{min}`$, of a given HH $`K`$-constant surface. Two conditions must be satisfied to determine $`R_{min}`$. First, the closest approach occurs when $`\dot{R}=(ds/d\tau )=0`$, which is equivalent to demanding
$$V(R_{min})=1.$$
(40)
This condition leads to a sixth-order polynomial in $`R_{min}`$:
$$\frac{K^2}{9}R_{min}^6+R_{min}^42\left(M+\frac{HK}{3}\right)R_{min}^3+H^2=0.$$
(41)
Second, the solution for this surface at $`R(t_{min})=R_{min}`$ must be concave, so as to rule out the SS surfaces (which bend towards the singularity rather than the horizon). This is enforced by demanding:
$$\ddot{R}(t_{min})|_{\dot{R}(t_{min})=0}=\frac{(R_{min}2)\left[3+2R_{min}+KR_{min}^2\sqrt{\frac{2}{R_{min}}1}\right]}{R_{min}^3}0.$$
(42)
The solution of these two conditions, as shown in Fig. 3, gives rise to the emergence of two critical values for $`H`$, namely $`H_\pm `$, for a given value of $`K`$. In addition, an HH surface can be made to approach arbitrarily closely to the singularity at $`r=0`$ by choosing an appropriately large positive value of $`K`$. Negative values of $`K`$ tend to “hug the horizon.”
For fixed $`K`$ and $`H`$ we know how to compute how closely a CMC surface comes to the singularity. But, for a given value of $`K`$, what value of $`H`$ gives the closest overall approach to the singularity? We can determine the critical values for $`R`$ and $`H`$, given by $`R_c`$ and $`H_\pm `$ respectively, by looking at the point where the first and second time derivatives of $`R(t)`$ vanish. $`H_{}`$ occurs along the lower boundary of the contour plot in Fig. 3, and corresponds to the $`K`$-surface that reaches down the furthest towards the singularity for a given value of $`K`$. The vanishing of $`\ddot{R}|_{\dot{R}=0}`$ leads to the following equation for $`R_c`$:
$$\left(R_c2\right)\left(3+2R_c+KR_c^2\sqrt{\frac{2}{R_c}1}\right)=0,$$
(43)
which can be rewritten as a 4th order polynomial in $`R_c`$
$$K^2R_c^42K^2R_c^3+4R_c^212R_c+9=0.$$
(44)
One must take care in examining the roots of this equation, as there are more solutions to Eq. (44) than there are for Eq. (43). Nevertheless this equation gives two distinct real roots, depending on the sign of $`K`$:
$`R_c|_{K>0}`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \frac{1}{2}}\sqrt{3{\displaystyle \frac{8}{K^2}}\chi +{\displaystyle \frac{2}{\sqrt{\chi }}}+{\displaystyle \frac{16}{K^2\sqrt{\chi }}}}+{\displaystyle \frac{\sqrt{\chi }}{2}},`$ (45)
$`R_c|_{K<0}`$ $`=`$ $`{\displaystyle \frac{1}{2}}+{\displaystyle \frac{1}{2}}\sqrt{3{\displaystyle \frac{8}{K^2}}\chi +{\displaystyle \frac{2}{\sqrt{\chi }}}+{\displaystyle \frac{16}{K^2\sqrt{\chi }}}}+{\displaystyle \frac{\sqrt{\chi }}{2}},`$ (46)
where
$`\xi `$ $``$ $`32+108K^2+243K^4+27K^2\sqrt{16+56K^2+81K^4},`$ (47)
$`\chi `$ $``$ $`1{\displaystyle \frac{8}{3K^2}}+{\displaystyle \frac{16+36K^2}{\mathrm{3\; 2}^{1/3}K^2\xi ^{1/3}}}+{\displaystyle \frac{2^{1/3}\xi ^{1/3}}{3K^2}}.`$ (48)
When $`K=0`$ we see from Eq. (43) that $`R_c=3/2`$. The regular $`K`$-surfaces are thus bounded between the horizon at $`R_+=2M`$ and $`R_{}=R_cM`$. Using Eq. (28), and setting $`\dot{R}=0`$, we obtain, for the case of a black hole in Schwarzschild coordinates
$$H_\pm =\frac{B_\pm R_\pm ^2}{\sqrt{B_\pm }}+\frac{KR_\pm ^3}{3},$$
(49)
with
$$B_\pm =\left(1\frac{2M}{R_\pm }\right).$$
(50)
We therefore have
$$H_+=\frac{8}{3}M^3K.$$
(51)
For large values of $`|K|`$ one can show that the critical value of $`R`$, $`R_{}`$, depends upon the sign of $`K`$. In particular,
$$R_c\{\begin{array}{cc}2\frac{1}{8}K^2\hfill & \text{for }K1\hfill \\ \left(\frac{9}{4}\right)^{1/3}K^{2/3}\hfill & \text{for }K1\hfill \end{array}.$$
(52)
This in turn gives the following asymptotic values for $`H_m`$:
$$H_{}\{\begin{array}{cc}\frac{8}{3}K\frac{1}{2K}\hfill & \text{for }K1\hfill \\ \frac{9}{2K}\hfill & \text{for }K1\hfill \end{array}.$$
(53)
## IV Constant Crunch Coordinates: A Spacetime Metric for a $`K`$-surface Foliation of the Schwarzschild Black Hole.
A number of features of the $`K`$-constant surfaces presented in the previous sections seem particularly well suited to the numerical analysis of generic black hole spacetimes. First, the surfaces asymptote to a null surface, making them effective for gravity wave extraction. Second, the $`K`$-surfaces naturally avoid the crushing singularity. Finally, for large negative values of $`K`$ the surfaces “hug the horizon.” This last feature, illustrated in Fig. 4, allows one to focus attention on the region relevant for gravity wave generation—the region outside the horizon.
In this section we generate a $`K`$-constant foliation for the Schwarzschild black hole that, in addition to the properties just mentioned, also has regular and static metric and extrinsic curvature components. To generate this $`K`$-constant slicing we use the coordinate transformation
$`\overline{t}`$ $`=`$ $`tT(r),`$ (54)
$`\rho `$ $`=`$ $`r.`$ (55)
Under this transformation, the metric from Eq. (1) becomes
$$ds^2=(1\frac{2M}{\rho })d\overline{t}^2+2\frac{(JH)}{\sqrt{(HJ)^2+(1\frac{2M}{\rho })\rho ^4}}d\overline{t}d\rho +\frac{\rho ^4}{(HJ)^2+(1\frac{2M}{\rho })\rho ^4}d\rho ^2+\rho ^2d\mathrm{\Omega }^2.$$
(56)
The constant $`\overline{t}`$ slices of this metric are $`K`$-constant surfaces. It is to be noted that Eq. (56) agrees with Eq. (53) of Iriondo et al. (for the case of constant $`K`$ and vacuum). Iriondo1996 However, in order to regularize the $`g_{\rho \rho }`$ metric component at the throat, we add the isotropic-like radial transformation Matzner2000 :
$$\overline{r}=\frac{1}{2}\left(\frac{1}{2}R_{min}+\rho +\sqrt{R_{min}\rho +\rho ^2}\right),$$
(57)
with $`R_{min}`$ the minimum coordinate location of the throat, given by Eq. (41). This coordinate representation of a black hole spacetime provides a foliation with $`K`$-constant, $`H`$-constant spacelike hypersurfaces. Each hypersurface is metrically equivalent to all others—the surfaces are independent of $`\overline{t}`$, and hence static. In addition, the hypersurfaces are asymptotically null ($`T^{}(r)1`$ as $`r\mathrm{}`$). Furthermore, the lapse, the shift, and all of the 3-metric and extrinsic curvature components are regular and well behaved, as illustrated in Fig. 5.
In addition to restricting ourselves to the singularity-avoiding family (HH) of $`K`$-surfaces, two additional conditions on the $`K`$-surfaces are demanded by the nature of our problem – the eventual simulation of the gravity-wave emission from two interacting black holes. First, the foliations must asymptote at large $`\overline{r}`$ to future null infinity. Therefore, we must restrict our attention to negative values for the trace of the extrinsic curvature tensor, $`K`$. Second, we require that such negative $`K`$ hypersurfaces enter the future singularity region. This will ensure proper coverage of the relevant region just above the future horizon, which is precisely where the gravity waves are produced. However, we expect the initial-data formulation for such surfaces to be involved, and this may guide our choices even more systemically. The two additional requirements limit us to the relatively narrow wedge of Fig. 3, formed by restricting to the “Future HH and SS ” shaded region with $`K<0`$.
A representative constant-crunch foliation generated by Eq. (56) is shown in Fig. 6. The avoidance of “grid stretching” is accomplished by a suitable choice of shift vector. To illustrate the non-zero shift we show the $`K=1`$, $`H3.11`$ foliation of Schwarzschild in Fig. 7, with the explicit misalignment of the $`r=\text{constant}`$ line segment and the normal vector.
Finally, we describe several preliminary numerical experiments using $`K`$-surfaces. A full treatment of a single black hole using $`K`$-constant foliations will be presented elsewhere Shoemaker2000 . Here we present several sample evolutions demonstrating the utility of constant mean curvature slicings. Figures 8 and 9 display results from the simplest possible test of $`K`$-constant foliation of the Schwarzschild geometry; the domain is taken to be a thin shell close to the horizon (in this case $`r[1,5]`$), analytic Dirichlet conditions are applied at both boundaries of the computational domain, and analytic lapse and shift conditions obtained from Eqs. (56) and (57) are used. The figures represent the singularity-avoiding foliation $`K=1`$ and $`H=3`$, for which the evolution was found to be stable and accurate over very long timescales. Using $`50`$ grid points, the code succesfully ran beyond $`t=50,000M`$ while maintaining high accuracy. The fractional error in the metric components was typically $`12\%`$.
Fig. 8 shows the convergence of the mean fractional error in the metric component $`\overline{a}=\sqrt{g_{\overline{r}\overline{r}}}`$ as a function of time. Each curve has been rescaled by a factor of $`4^p`$, where the number of grid points is given by $`400/2^p`$, with $`p=0,1,2,3`$. For $`t>2`$ the solution is approximately second order accurate. Initially, noise generated on the inner boundary causes fluctuations whose magnitude is largely independent of the number of grid points. Fig. 9 shows snapshots of the Hamiltonian constraint at various times in the evolution. Qualitatively, a wave which is triggered by truncation error is seen to propagate outwards from the inner boundary. The amplitude reduces rapidly, before growing once more as it is reflected off the outer boundary. The magnitude and speed of propagation of the wave quickly decay as the wave moves back into the domain, leaving a static solution which is stable beyond $`t=50,000M`$. The numerical error which initially propagates through the domain is caused entirely by the analytic Dirichlet boundary conditions, and can be largely eliminated by the use of more realistic conditions. Shoemaker2000
The numerical runs presented here evolve a $`K=1`$ singularity-avoiding hypersurface that asymptotes to null infinity, entering the past singularity region. This is not of the class of $`K`$-surfaces emphasized in this paper for use in numerical relativity. Ideally, we would have preferred presenting the evolution of a $`K=1`$ HH surface that enters the future singularity region. However, we were unable to find an HH surface in the future singularity region that evolved stably using the naive Dirichlet boundary conditions and analytic lapse and shift conditions. Nevertheless, we have evolved such surfaces stably by incorporating (1) area locking shift conditions, (2) isometery conditions at the throat, and (3) an outgoing boundary condition based on the difference between the numeric and analytic solutions. These numerical results, and the corresponding stability analysis, are not within the scope of this paper and will be presented elsewhere.Shoemaker2000
## V From One Black Hole to Two Spinning Black Holes
In this paper we considered a static spacetime metric for the Schwarzschild black hole, with spacelike hypersurfaces of constant (not necessarily zero) value of the trace of the extrinsic curvature tensor. This slicing provides a natural generalization of the maximal slicing scheme currently in use in some numerical approaches to the binary black hole problem.
An essential feature of our $`K`$-constant metric is a spatially-varying radial shift vector, which allows the surfaces to avoid the singularity while evading the grid stretching problems often encountered with other metrics. Work is in progress to develop a geometric handle on our shift condition, akin to that of the “minimal distortion” shift often discussed. The inner and outer boundary conditions are also particularly convenient. In the inner regions, the “horizon hugging” feature of the $`K`$-slices, together with their regularity, may remove the need to excise the grid within the apparent horizon, thus providing a natural “boundary without a boundary” avoidance of the singularity. In addition, at the outer boundary the surfaces are asymptotically null, which may aid in gravity wave extraction.
In examining the characteristics of the $`K`$-slices we reviewed the three families of $`K`$-surfaces, including the horizon-to-singularity (HS ) surfaces, as well as the more familiar horizon-to-horizon (HH ) and singularity-to-singularity (SS ) surfaces. All three families of surfaces either asymptote to the singularity, or to spatial infinity along a null surface. Thus the family of observers that make up such a foliation never “observe” the surfaces reach the singularity. This suggests a possible payoff in using the HS surfaces in the evolution of black hole spacetimes, as there is no danger of the surface reaching the singularity. Just as there are natural boundary conditions which “freeze out” gravitational radiation as it progresses to null infinity, so too can we expect boundary conditions where the dynamic space freezes out as it approaches the singularity. Are there such natural boundary conditions at the future singularity? It would be worthwhile to investigate such asymptotically-null boundary conditions at the singularity of an $`HS`$ surface.
The future utility of numerical simulations of black hole spacetimes hinges in large part on a suitable choice of coordinates for the initial data, and on the particular evolution of these coordinates through the four lapse and shift conditions. These conditions are the only handles by which to manage the growth of the metric and curvature components during evolution. It is through judicious choices of lapse and shift that one is able to effectively enable singularity avoidance, and allow for efficient extraction of gravitational radiation. The simplistic example presented here may provide some guidance as to how to proceed in the more general case. $`K`$-surfaces in the generic two black hole problem can be expected to preserve the singularity-avoidance “horizon hugging” behavior, as well as remain asymptotically null at the boundaries. It remains to be seen what shift vectors are required to manage the growth of the intrinsic and extrinsic properties of the metric in the more general case, though the constant-crunch shift presented here is a natural starting point.
We are currently pursuing three avenues towards furthering and generalizing the work presented here. First, we are analyzing the stability of this coordinatization to small perturbations. Second, we are using the (1+1)-dimensional code discussed in the last section to examine a greater portion of the parameter space of initial data, so as to fully explore the numerical stability of the slicing. From our preliminary numerical investigations, we expect to be able to do a full spacetime “evolution” of Schwarzschild, and have it run stably and accurately for extended periods of time over a wide range of $`H`$-$`K`$ parameter space. Finally, we are analyzing the Oppenheimer-Snyder collapse in such $`K`$-slicings, following the lead of Eardley and Smarr, Eardley1978 as this will further test our slicing in a non static setting. Thorne2000
###### Acknowledgements.
We wish to acknowledge the Los Alamos National Laboratory LDRD/ER program for financial support. WAM wishes to thank the Institute for Theoretical Physics at UCSB for providing a stimulating working environment in which to complete part of this research. DH was supported in part by the National Science Foundation under Grant No. PHY9907949 to the ITP. PL was supported in part by NSF grants PHY9800973 and PHY9800970. We wish to thank Richard Matzner for advice on handling the coordinate ambiguity at the throat. We are especially grateful to John A. Wheeler for encouraging us to examine these slices in the context of the numerical treatment of black hole spacetimes.
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# Massive skyrmions in quantum Hall ferromagnets
## I Introduction
In the last few years topological spin textures in the quantum Hall effect (QHE) have received considerable attention . The existence of skyrmions can be anticipated within the frame work of the Chern-Simon-Landau-Ginsburg mean field theory , i.e., integrating out the charge fluctuations of the composite bosons yield an effective model for the Chern-Simon’s gauge field. The transport properties of the skyrmions can be extracted through considering the fluctuating Chern-Simon’s gauge field, which can be derived by expanding the effective action about its minimum energy solution. One of the leading terms of this expansion is the Maxwell action of the Chern-Simon’s fluctuating field, e.g., $`(m^{}/\overline{\rho })_{k,\omega }|𝐉^s|^2`$. $`m^{}`$ is the electron effective mass in a host semiconductor, and $`\overline{\rho }=1/(2\pi \mathrm{}_0^2)`$ is the average electron density at $`\nu =1`$, and $`\mathrm{}_0`$ is the magnetic length.
Although a skyrmion mass is physically reasonable, in the usual minimal field theories skyrmions are considered as massless objects. ¿From the microscopic point of view, e.g., a microscopic Hartree-Fock approximation, it is not certain if the skyrmions are massive. The resolution of this ambiguity between the microscopic Hartree-Fock approximation and the Landau-Ginsburg-Chern-Simon’s theory (considered in this paper), is an open question.
In this paper we apply the theory of elasticity to investigate the effect of skyrmion mass on the thermodynamic properties of a skyrmion crystal. We derive the collective mode dispersion relations for the Skyrme lattices at long wave lengths. We then make a comparison with the massless microscopic Hartree-Fock calculations to reconcile the prediction of the Chern-Simon theory with the massless models . In addition, we suggest how the mass of skyrmions may be characterized within the microscopic picture. We also study the stability of the Skyrme lattices at low temperature and show that the low temperature phases of these Skyrme lattices are not affected by including a mass term in the effective action, unless if the mass is sufficiently large. We show the mass of skyrmions is suppressed by decreasing the Zeeman energy, and indicating that skyrmions are massless at zero Zeeman energy.
## II Skyrmion mass
If skyrmions have mass, this will affect a host of properties, from their tunneling through a constriction to the thermodynamics of a lattice of skyrmions. The rationale for their having mass is that in the starting point for many calculations, the Chern-Simons Lagrangian of Lee and Kane, the electrons have mass. After standard manipulations (introduction of a CP<sup>1</sup> field, changing variables to $`𝐦`$, the local spin field) one has a continuum theory in which gradients in the spin texture become associated with the charge distribution. If a skyrmion moves slowly across the system, this corresponds to the motion of one quasiparticle. It seems reasonable that the motion of such a texture will involve inertial terms. Questions similar in spirit arisen in determining the mass of vortices .
Recalling the duality relation between the topological 3-current of skyrmions and the Chern-Simon’s gauge field , $`J_\mu ^s=(\nu /2\pi )ϵ_{\mu \nu \lambda }_\nu A_\lambda `$, one might anticipate a second derivative term in time will appear in the effective action. After integrating out the statistical gauge field one obtains in the limit of low frequency and long wavelength:
$`S_E[𝐦]`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{𝐤,\omega }{}}\left(V(k)\right|J_0^s|^2+{\displaystyle \frac{m^{}}{\overline{\rho }}}|𝐉^s|^2+`$ (4)
$`i\alpha 𝐀^{(0)}(k)𝐉^s(k)`$
$`{\displaystyle \frac{2\pi \alpha }{k^2}}J_0^s(k)\widehat{𝐳}𝐤\times 𝐉^s(k))`$
$`+S_{NL\sigma M}+S_z+S_{\mathrm{fermi}}`$
where the last three terms are the nonlinear sigma action, the Zeeman contribution, and a term in the action that guarantees the skyrmions obey fermi statistics. The first term is electrostatic, and involves the fourier transform of an effective interaction potential $`V(k)`$ including screening by fluctuations in the texture, the second contains the kinetic energy and the third reflects that the skyrmion sees the original boson as a magnetic flux tube. The constant $`\alpha `$ is an odd integer and arises from a Chern-Simons term that makes the system fermionic.
The skyrmion current density is related to the spin texture via:
$$J_\lambda ^s=\frac{\nu }{8\pi }ϵ_{\lambda \mu \nu }\left(_\mu 𝐦\times _\nu 𝐦\right)𝐦$$
where the indices run over time and two spatial dimensions. The zeroth component is the topological charge density of the texture, which is proportional to the quasi-particle number density, $`\rho (r)J_0^s`$. If we consider a single skyrmion texture that moves uniformly, $`J_\lambda ^s=J_\lambda ^s(𝐱𝐯t)`$, then it is straightforward to show the relation between the skyrmionic current and their charge density satisfy the usual charge conservation law, i.e., $`𝐉^s=\rho 𝐯`$. Then the kinetic term simplifies to
$$\frac{1}{2}\underset{𝐤}{}\frac{m^{}}{\overline{\rho }}|𝐉^s(𝐤)|^2=\frac{1}{2}M_0v^2$$
which yields a transport mass for the skyrmions $`M_0=(m^{}/\overline{\rho })d^2𝐫\rho ^2(𝐫)/(2\pi )^2`$.
This mass is derived from kinetic considerations. It is possible that the correct “mass” to calculate will depend upon exactly what is being measured in a given experiment.
## III Collective modes
The long range order of the skyrme crystal is determined by the repulsive Coulomb interaction, and the topological, XY interaction of the hedge-hog fields. An antiferromagnetic ordering between the single skyrmions, within a square lattice minimizes the topological interaction. This can be realized after mapping of the topological hedge-hog fields of the charge one skyrmions onto a system of classical dipoles. However, a triangular lattice is favored by the Coulomb interaction, similar to the Wigner crystals. When the Zeeman energy is small enough, the skyrmions can pair up into charge two skyrmions and lower the total energy of the lattice. In contrast to single skyrmions where their topological hedge-hog fields are analog to a system of classical dipoles, the charge two skyrmions mimic a system of classical quadrapoles. This favors triangular lattice ordering of charge two skyrmions, i.e., a bi-skyrmion lattice.
The low-lying collective modes of a Skyrme lattice consist of phonons and spin waves. The dispersion relation of these collective modes can be obtained by adding a dynamical term to the effective Hamiltonian. This Hamiltonian has been derived in Ref. by a first principle calculation of a non-linear $`\sigma `$ model, assuming a specific skyrmion is localized and well separated from other skyrmions, i.e., they interact weakly. More precisely, we assume that we can divide the 2-dimensional configuration space into $`N`$ regions ($`N`$ is the number of skyrmions) such that a given skyrmion lives in the region that other skyrmions are close to their vacuum. This condition enables us to linearize the interaction potential energy among the isolated-skyrmions. This assumption may break down if the size of skyrmions, $`\lambda `$, becomes comparable with the distance between them. In this case the next to the linear terms in the potential energy may be significant, and one should take them into account. Roughly speaking, this happens if $`R2\lambda `$ ($`R`$ is the separation between two skyrmions). Here the skyrmion size $`\lambda `$ is defined as the radius at which the spin lies in the XY plane. The relevant potential energy functional of the excitations is obtained by introducing the lattice of equilibrium positions of the skyrmions $`R_{i\alpha }`$ is an initial lattice ansatz ($`\alpha `$ is the cartesian component of the position of the $`i`$th skyrmion), the displacement field $`u_\alpha (𝐑_i)`$, and the orientation field of the skyrmions $`\theta (𝐑_i)`$. The result is then
$`E[𝐮,\theta ]`$ $`=`$ $`{\displaystyle \underset{ij}{}}V_0(|𝐑_i+𝐮(𝐑_i)𝐑_j𝐮(𝐑_j)|)`$ (7)
$`+{\displaystyle \underset{<ij>}{}}J(|𝐑_i+𝐮(𝐑_i)𝐑_j𝐮(𝐑_j)|)`$
$`\times \mathrm{cos}(\chi (𝐑_i)+\theta (𝐑_i)\chi (𝐑_j)\theta (𝐑_j)).`$
Here the topological hedge-hog interaction for a single-skyrmion lattice is given by $`J(R)=c^2\stackrel{~}{g}/(4\pi ^2\mathrm{}_0^2)K_0(\kappa R)`$ where $`\stackrel{~}{g}=ge^2/(2ϵ\mathrm{}_0)`$ is the Zeeman energy, $`g`$ is the effective gyromagnetic ratio, $`\kappa ^2=\stackrel{~}{g}/(2\pi \mathrm{}_0^2\rho _s)`$, $`\rho _s`$ is the spin stiffness. $`K_0(x)`$ is the modified Bessel function. In addition, $`c`$ is a constant that can be obtained from the asymptotic form of an isolated skyrmions and equals $`c=30.4\mathrm{}_0`$. For a bi-skyrmion lattice $`J(R)=c^2\stackrel{~}{g}^2/(8\pi ^3\mathrm{}_0^4\rho _s)K_0(\kappa R)`$ and $`c=79\mathrm{}_0^2`$. One should note $`c`$, and therefore the topological interaction between skyrmions $`J(R)`$, depend on the local form of skyrmions.
### A Phonons
We find the spectrum of the phonons, using the standard technique of expanding the energy functional about its minima, i.e., $`𝐮=0`$, $`\theta =0`$, assuming the orientational field of skyrmions is frozen out. For the single-skyrmion case $`J(R)`$ is positive, hence $`\chi _i\chi _j=\pi `$ is the lowest energy state of the topological XY interaction. For the bi-skyrmion case $`J(R)`$ is negative and $`\chi _i\chi _j=0`$ is the lowest energy state. An estimate on the total energy of skyrmions show the bi-skyrmion configuration is likely if the Zeeman energy is very small ($`10^5`$). We therefore do not consider this configuration for the rest of this paper. Expanding the potential energy in terms of the displacement fields (up to quadratic order terms) gives
$$E[𝐮]=E_{\mathrm{classic}}+\frac{1}{2}\underset{𝐤BZ}{}\underset{\alpha \beta }{}𝐮_\alpha ^{}(𝐤)D_{\alpha \beta }(𝐤)𝐮_\beta (𝐤),$$
(8)
with $`E_{\mathrm{classic}}`$ the classical groundstate energy of the Skyrme lattice and $`D_{\alpha \beta }(𝐤)`$ is the dynamical matrix
$`D_{\alpha \beta }(𝐤)`$ $`=`$ $`{\displaystyle \frac{2\pi }{k}}{\displaystyle \frac{e^2}{ϵa_c}}k_\alpha k_\beta +a_c[(\mu +\lambda )k_\alpha k_\beta `$ (10)
$`+\mu k^2\delta _{\alpha \beta }+\gamma k_xk_y(1\delta _{\alpha \beta })+𝒪(k^4)].`$
The first term comes from the electrostatic interactions of a Wigner crystal, the others arise from standard 2D elasticity theory. The quantity $`a_c`$ is the area of a unit cell, while $`\mu ,\lambda `$, and $`\gamma `$ are the conventional Lamé coefficients for a square lattice.
Since the effective interaction consists of two terms, i.e., the direct and exchange interactions, the Lamé coefficients can be expressed by means of $`\mu =\mu _0+\mu _1`$, $`\lambda =\lambda _0+\lambda _1`$ and $`\gamma =\gamma _0+\gamma _1`$. Here $`0`$ and $`1`$ are labels associated with the direct and the exchange Coulomb energy. Expanding the exchange term (second term) in Eq.(7) up to the quadratic order in displacement field, and for the single-skyrmionic square lattices with only nearest neighbor exchange interactions we find
$$\mu _1=\frac{c^2\stackrel{~}{g}}{2\pi ^2a_c}\sqrt{\frac{\pi }{2}}\frac{0.5+\kappa R}{(\kappa R)^{1/2}}e^{\kappa R},$$
(12)
$$\lambda _1=\frac{c^2\stackrel{~}{g}}{2\pi ^2a_c}\sqrt{\frac{\pi }{2}}\frac{7/4+3\kappa R+(\kappa R)^2}{(\kappa R)^{1/2}}e^{\kappa R},$$
(13)
$$\gamma _1=\frac{c^2\stackrel{~}{g}}{2\pi ^2a_c}\sqrt{\frac{\pi }{2}}\frac{5/4+2\kappa R+(\kappa R)^2}{(\kappa R)^{1/2}}e^{\kappa R}.$$
(14)
We also find the contribution of the direct interaction to the elastic constants in units of $`(e^2/ϵ\mathrm{}_0a_c)\sqrt{|1\nu |/2\pi }`$ are $`\mu _0=0.146289`$, $`\lambda _0=0.536199`$ and $`\gamma _0=1.560224`$ (we take the nearest neighbor separation between the skyrmions $`R=\mathrm{}_0\sqrt{2\pi /|\nu 1|}`$ as the square lattice constant and therefore $`a_c=R^2`$). Note that $`\gamma =0`$ for a triangular lattices.
The time-derivative parts in the action of the displacement field $`𝐮`$, consists of the Wess-Zumino term and the kinetic energy of the skyrmions which are first and second derivatives in (imaginary) time, respectively. Combining these with our previous results gives the following effective Euclidean action for the low energy spectrum of the displacement field
$$S_{\mathrm{eff}}[𝐮]=\frac{1}{2}\underset{\alpha \beta }{}\underset{𝐤BZ}{}_0^\mathrm{}\beta 𝑑\tau \left[i\epsilon _{\alpha \beta }m\omega _cu_\alpha ^{}(𝐤,\tau )_\tau u_\beta (𝐤,\tau )M_0\delta _{\alpha \beta }u_\alpha ^{}(𝐤,\tau )_\tau ^2u_\beta (𝐤,\tau )+u_\alpha ^{}(𝐤,\tau )D_{\alpha \beta }(𝐤)u_\alpha (𝐤,\tau )\right],$$
(15)
where $`\omega _c`$ is the cyclotron frequency of the electron.The dependence of the skyrmion mass on the Zeeman energy can be obtained by solving the non-linear differential equation for the spin texture of a single skyrmion and is shown in Fig. 1. Note that unlike the total number of spins participating in the skyrmion, the skyrmion mass decreases with $`g`$. This arises because the mass is proportional to the square of the skyrmion density, and as $`g`$ decreases the skyrmion becomes more spread out. Expanding the time dependence of the complex field $`𝐮(𝐤,\tau )`$ in terms of the bosonic Matsubara frequencies, $`\omega _n=2\pi n/\beta `$, leads to the fluctuation matrix
$$𝒮_n(𝐤)=\left(\begin{array}{cc}M_0\omega _n^2+D_{xx}(𝐤)& m\omega _c\omega _n+D_{xy}(𝐤)\\ m\omega _c\omega _n+D_{yx}(𝐤)& M_0\omega _n^2+D_{yy}(𝐤)\end{array}\right).$$
(16)
The action in Eq. (15) is similar to the action of a Wigner crystals in the presence of a magnetic field. Although this is as expected, it should be noted that the magnetic field interaction is here exactly recovered by the topological Wess-Zumino term.
The contribution of the zero-point energy of the phonons to the total energy of the Skyrme lattices, may be obtained by integrating out the quadratic fluctuations in $`𝐮`$. It turns out that $`E=E_{\mathrm{classic}}+E_{\mathrm{flu}}`$ where
$`E_{\mathrm{flu}}=\underset{\beta \mathrm{}}{lim}`$ $`{\displaystyle \frac{1}{\beta }}{\displaystyle \underset{n}{}}{\displaystyle \underset{𝐤BZ}{}}\mathrm{ln}[\beta ^2\mathrm{}^2(\omega _{}(𝐤)+i\omega _n)`$ (18)
$`\times (\omega _+(𝐤)i\omega _n)],`$
and $`\omega _\pm (𝐤)`$ are the phonon frequencies
$`\omega _\pm ^2(𝐤)={\displaystyle \frac{1}{2M_0^2}}\left(m^2\omega _c^2+M_0(D_{xx}+D_{yy})\pm \sqrt{\left[m^2\omega _c^2+M_0(D_{xx}+D_{yy})\right]^24M_0^2(D_{xx}D_{yy}D_{xy}D_{yx})}\right),`$ (19)
which have been obtained by substituting the analytical continuation $`i\omega _n\omega (𝐤)+i\delta `$ into the fluctuation determinant. The dispersion relation of the phonons consists of a gapped mode, as well as a gapless modes. For the former, the gap starts at the cyclotron frequency $`\omega (𝐤0)=(m^{}/M_0)\omega _c`$ and at long wave lengths obeys
$$\omega _+^2(𝐤)=\left(\frac{m^{}}{M_0}\omega _c\right)^2+\frac{2\pi e^2}{ϵa_cM_0}k+a_c\frac{3\mu +\lambda }{M_0}k^2+𝒪(k^3),$$
(20)
where $`\phi `$ is the azimuthal angle of the wave vector $`𝐤`$ in the XY plane. The gapless mode on the other hand obeys
$`\omega _{}^2(𝐤)`$ $`=`$ $`{\displaystyle \frac{2\pi e^2}{ϵ(m^{}\omega _c)^2}}(\mu {\displaystyle \frac{\gamma }{2}}\mathrm{sin}^22\phi )k^3+𝒪(k^4).`$ (21)
Within a low Zeeman energy limit, where the mass of skyrmion is small, the gapped mode goes toward the higher energies and the effect of the mass of skyrmions becomes less significant. At this limit the prediction of the Chern-Simon theory, approaches to the current prediction of the microscopic Hartree-Fock approximation. However, for higher Zeeman energies, it is not clear if the prediction of the microscopic Hartree-Fock approximation leads to a single gapless mode or it can support the gapped mode and subsequently the possibility of the existence of a non-zero skyrmion mass too . The square lattice can be unstable against the lattice fluctuations. Moreover the contribution of the exchange interaction to the Lamé coefficients falls off exponentially (much faster than the contribution from the direct interaction) and when $`\nu 1`$ a triangular Wigner crystal configuration is more favorable than a square lattice. A similar situation arises when either $`\stackrel{~}{g}0`$ or $`\stackrel{~}{g}\mathrm{}`$. Our numerical studies show that this instability can be seen for values of the filling fraction and $`g`$-factors that are shown in Fig. 2. The curve of $`R=2\lambda `$ is also plotted. The overlap between skyrmions is significant below this curve and the topological interaction becomes strong. The assumption of weakly interacting skyrmions may fail within the region below this curve, and one should take the next to the linear topological XY terms into account. As it is shown in Fig. 2, for $`|1\nu |0.01`$, and for certain direction of $`𝐤`$, the gapless mode becomes imaginary, implying an instability of the square Skyrme lattice that is again consistent with the microscopic Hartree-Fock models . The stable region of the square lattices is characterized by dark area and the stable region of the triangular lattices is shown by the white area. It is seen the triangular lattice reappears when the Zeeman energy is small enough. Without Coulomb interactions a lattice configuration is always unstable against the attractive interaction between skyrmions, i.e., a single skyrmion with topological charge $`N`$ is the global minima of the skyrmionic energy functional.
As it is seen in Fig. 1, the quantum Hall skyrmions are massless at $`g=0`$. For the massless case the gapped mode goes to infinity and there is just only one gapless mode
$`(m\omega _c\omega `$ $`(𝐤))^2=\mu ({\displaystyle \frac{2\pi e^2}{ϵk}}+a_c^2[\lambda +2\mu ])k^4`$ (23)
$`2\gamma \left({\displaystyle \frac{2\pi e^2}{ϵk}}+a_c^2[\mu +\lambda +{\displaystyle \frac{1}{2}}\gamma ]\right)k_x^2k_y^2+𝒪(k^6),`$
which is identical to the gapless mode of the massive skyrmions. In fact, the long wave length power-law behavior of the gapless mode is totally unaffected by the mass of skyrmions. For both the massless and massive theories, we thus find that $`\omega k^{3/2}`$ at long wave lengths, which is consistent with the microscopic Hartree-Fock calculations . Experimentally, the mass of skyrmions may not be probed directly by phonon excitations unless we excite the gapped modes (the cyclotron modes), e.g., by optical measurements . Furthermore, the mass of skyrmions can not change the melting point of the 2D skyrmions, since $`T_m`$ is specified by the elastic constants of Skyrme lattices (the 2D melting point just depends on the interaction energy between skyrmions). Returning to the evaluation of the zero-point energy of the phonons, the dominant contribution at low temperatures is thus obtained as $`E_{\mathrm{flu}}=_k\mathrm{}\omega _{}(𝐤)`$. Similarly the average square displacement of the skyrmions is given by
$`u^2=\mathrm{}_0^4{\displaystyle \underset{𝐤BZ}{}}`$ $`{\displaystyle \frac{D_{xx}(𝐤)+D_{yy}(𝐤)}{\mathrm{}\omega _{}(𝐤)}}`$ (25)
$`\times \left[2n_B(\beta \mathrm{}\omega _{}(𝐤))+1\right],`$
where $`n_B(x)`$ is the Bose-Einstein distribution function.
### B Magnons
Next we consider the spin waves of the Skyrme lattices by expanding the energy functional in Eq. (7) also in the orientation field $`\theta `$. Up to the quadratic order, it gives the XY-energy contribution
$`E[𝐮,\theta ]=E[𝐮]+{\displaystyle \frac{1}{2}}K_{XY}{\displaystyle d^2𝐫|\theta (𝐫)|^2},`$ (26)
where $`K_{XY}=J(\kappa R)`$ is the effective stiffness associated with gradients in the skyrmion orientations. As is shown in Ref. , the effective action for the spin waves contain also a mass term ($`\mathrm{\Lambda }_0`$ is the moment of inertia) and we obtain finally
$`S_{\mathrm{eff}}[\theta ]`$ $`=`$ $`{\displaystyle _0^\mathrm{}\beta }d\tau {\displaystyle }d^2𝐫({\displaystyle \frac{\mathrm{\Delta }M}{a_c}}_\tau \theta `$ (28)
$`+{\displaystyle \frac{\mathrm{\Lambda }_0}{2a_c}}(_\tau \theta )^2+{\displaystyle \frac{K_{XY}}{2}}|\theta |^2).`$
The first term in Eq.(28) is the usual (dynamical) Berry’s phase of a quantum Hall ferromagnet, where $`\mathrm{\Delta }M`$ is the average change in the total magnetization induces by a single skyrmion texture. Note that in principle there is also a contribution from the Hopf term in the effective action for the quantum Hall ferromagnet, which at the quantum level ensures that the skyrmion obeys the correct spin-statistics relation . Since both these terms are a total derivative, however, the equation of motion is not affected by these terms and the long wave length dispersion relation that follows from action (28) turns out to be $`\omega (k)=c_sk`$, where $`c_s=\sqrt{a_cK_{XY}/\mathrm{\Lambda }_0}`$ is the velocity of the spin waves. The contribution from these fluctuations to the total energy of the crystal is again $`E_{fl\theta }=_k\mathrm{}\omega (𝐤)`$, and the mean square value of the associated fluctuations is
$$\theta ^2=\mathrm{}\underset{𝐤BZ}{}\frac{1}{\mathrm{\Lambda }_0\omega (𝐤)}\left[2n_B(\beta \mathrm{}\omega (𝐤))+1\right].$$
(29)
The coupling between the displacement and the orientational fields ($`𝐮`$ and $`\theta `$), turn out to be the next to leading order terms and are therefore negligible for our purposes. These terms lead to interactions between the phonons and spin waves, which we do not consider here. As a result we find for the zero-point energy of the phonons and spin waves simply $`E_{fl}=E_{\mathrm{flu}}+E_{fl\theta }`$. ¿From Eq. (25), we find $`u^2|1\nu |R^2/5`$ at zero temperature. We also find $`\theta ^2=\sqrt{2U/K_{XY}}`$ from Eq. (29) where $`U=1/(2\mathrm{\Lambda }_0)`$. For small values of $`K_{XY}`$ and/or large values of $`U`$ the quantum fluctuations are severe that the disordered phase can emerge. For very small Zeeman energy, where a phase transition from square single-skyrmion lattice into a triangular bi-skyrmion lattice can be observed , we obtain $`\theta ^2\stackrel{~}{g}^{0.46}`$. At the limit of small $`\stackrel{~}{g}`$, fluctuations are negligible and the long-range order of the bi-skyrmion lattices is not influenced by the quantum fluctuations.
## IV Conclusion
In this paper we have studied a system of 2-dimensional quantum Hall skyrmions, starting from a Chern-Simon-Landau-Ginsburg mean field theory. A Maxwell term can be obtained through the gradient expansion of the Chern-Simon action around its minimum energy solution. This term is responsible for generating of the skyrmionic inertial mass. Away from $`\nu =1`$ skyrmions stay in a crystal form. The long range order of crystal depends on the Landau level filling factor, and the Zeeman energy. Optical phonons which are out of phase fluctuations of the skyrmion lattices, are gapped since skyrmions carry an inertial mass. Therefore quantum Hall skyrmions behave like a Wigner crystal in the presence of external magnetic field. The inertial mass of skyrmions is vanished at zero Zeeman energy. In this situation the optical phonons are highly gapped, and they become inaccessible. We finish this paper with a final comment. At $`T=0`$, and far from $`\nu =1`$, the zero point quantum fluctuations of the phonons destroy the long-range order of skyrmion crystals. This is a high density limit of skyrmions where the Coulomb interaction screens from one $`1/r`$ to $`\mathrm{ln}(1/r)`$ by the Chern-Simon fluctuations. In this limit, the skyrmions behave like a gas of classical particles, i.e., they are crystalized under high pressures. This yields to the possibility of observing the re-entrance of the solid phase, followed by the disorder (liquid) phase at $`T=0`$ when $`\nu `$ is far enough from $`1`$.
## V acknowledgement
We acknowledge helpful conversations with Herbert Fertig, Steve Girvin, and Allan MacDonald. The work at the University of Oklahoma was supported by the NSF under grant No. EPS-9720651 and a grant from the Oklahoma State Regents for Higher Education.
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# 𝑇_c for dilute Bose gases: beyond leading order in 1/𝑁
## I Introduction
Second-order phase transitions have universal behavior, associated with long wavelength fluctuations, for which critical exponents and other universal quantities can often be successfully calculated using renormalization group techniques. For most such systems, the short distance physics is hopelessly complicated. In contrast, the phase transition of a dilute, interacting Bose gas provides a fascinating example where physics becomes simpler, and perturbative, at (relatively) small distance scales. For this system, it should be possible to marry techniques for treating long-distance critical fluctuations to a perturbative treatment of short distance physics, and so compute non-universal characteristics of the phase transition. A simple example of such a non-universal quantity is the phase transition temperature $`T_\mathrm{c}`$, and the effect of interactions on $`T_\mathrm{c}`$ has been explored by several authors , with a wide variety of theoretical results. In particular, the transition temperature has recently been calculated by Baym, Blaizot, and Zinn-Justin in the large $`N`$ approximation. For simplicity, they implicitly focus on the case of Bose gases with a single spin state, where the low-energy cross-section for atomic collisions can be parametrized by a single scattering length, $`a`$. As will be briefly reviewed below, the problem is first reduced to a calculation in a three dimensional O(2) scalar field theory at its critical point. Replacing that by an O($`N`$) theory with $`N=2`$, they find
$$T_\mathrm{c}=T_0\left[1+\frac{8\pi }{3\zeta (3/2)^{4/3}}an^{1/3}\left[1+O(N^1)\right]+O\left((an^{1/3})^2\right)\right]$$
(1)
in the dilute limit, where $`n`$ is the number density,<sup>*</sup><sup>*</sup>* For simplicity, we consider a uniform Bose gas, where $`n`$ is fixed. Alternatively, in an arbitrarily wide harmonic trap, $`n`$ should be interpreted as the actual density at the center of the trap at the transition temperature. and $`T_0`$ is the transition temperature of a non-interacting Bose gas,
$$T_0=\frac{2\pi \mathrm{}^2}{k_\mathrm{B}m}\left(\frac{n}{\zeta (3/2)}\right)^{2/3}.$$
(2)
Their result of $`\mathrm{\Delta }T_\mathrm{c}/T_0(T_\mathrm{c}T_0)/T_02.33an^{1/3}`$ is in good agreement with recent As alluded to earlier, there have been several different theoretical results and simulation results obtained by various methods (e.g. ), giving a large range of values for the coefficient of $`an^{1/3}`$ in $`\mathrm{\Delta }T_\mathrm{c}/T_0`$. There has also been some experimental data on the <sup>4</sup>He-Vycor system , which superficially seems to fit well an early theoretical estimate of Stoof , which is $`\mathrm{\Delta }T_\mathrm{c}/T_0(16\pi /3)\zeta (3/2)^{4/3}an^{1/3}4.66an^{1/3}`$. However, the detailed interpretation of this data is unclear. In that experiment, the Helium atoms are confined to an interconnected network of channels in the porous Vycor glass, and, for the low-density data that appears to fit Stoof, the interpaticle spacing is the same order of magnitude as the widths of the channels. Ref. simply assumes that the system can be modeled by a free Bose gas with (i) an effective mass for the atoms that is extracted experimentally, but (ii) the same scattering length as for bulk Helium, which is moreover taken from theoretical modeling. Because of these assumptions, the apparent agreement with Ref. should be treated with caution. numerical simulations that give $`\mathrm{\Delta }T_\mathrm{c}/T_0(2.2\pm 0.2)an^{1/3}`$. The result is surprising because it seems to work much better than the large $`N`$ expansion of critical exponents. For example, for $`O(N)`$ theory, the susceptibility critical exponent $`\gamma `$ is For Bose gases, a more physical example of a critical exponent is $`\nu =10.540(2/N)0.470(2/N)^2+O(N^3)`$, whose actual value is $`\nu 0.67`$ for $`N=2`$. The fact that O(2) critical exponents should be identified with Bose gas exponents is not completely trivial. A uniform, non-relativistic Bose gas is a constrained system: the particle density $`n`$ is fixed. This constraint causes the critical exponents $`\stackrel{~}{x}=(\stackrel{~}{\alpha },\stackrel{~}{\beta },\stackrel{~}{\gamma },\stackrel{~}{\nu })`$ of the actual system to be related to the standard exponents $`x=(\alpha ,\beta ,\gamma ,\nu )`$ of the field theory by (i) $`\stackrel{~}{\alpha }=\alpha /(1\alpha )`$, and $`\stackrel{~}{x}=x/(1\alpha )`$ for the others, if $`\alpha >0`$, or (ii) $`\stackrel{~}{x}=x`$ if $`\alpha <0`$. The actual value of $`\alpha `$ for the O(2) model is believed to be $`0.007\pm 0.006`$ . If negative, there is no difference between the exponents; if positive, there is in principle a very tiny difference. This relation explains, by the way, the difference between mean-field theory exponents for the O(2) model (e.g. $`\alpha =1/2`$) and the exponents of a non-interacting Bose gas (e.g. $`\stackrel{~}{\alpha }=1`$).
$`\gamma `$ $`=`$ $`2{\displaystyle \frac{24}{N\pi ^2}}+{\displaystyle \frac{64}{N^2\pi ^4}}\left({\displaystyle \frac{44}{9}}\pi ^2\right)+O(N^3)`$ (3)
$`=`$ $`21.216\left({\displaystyle \frac{2}{N}}\right)0.818\left({\displaystyle \frac{2}{N}}\right)^2+O(N^3).`$ (4)
This is not a marvelous expansion for $`N=2`$, for which the actual value is $`\gamma 1.32`$.
In this paper, we calculate the $`O(N^1)`$ relative correction to (1). We find
$$T_\mathrm{c}=T_0\left[1+\frac{8\pi }{3\zeta (3/2)^{4/3}}an^{1/3}\left[1\frac{0.527198}{N}+O(N^2)\right]+O\left((an^{1/3})^2\right)\right].$$
(5)
Setting $`N=2`$, this is only an 26% correction to the leading large $`N`$ result for $`\mathrm{\Delta }T_\mathrm{c}/T_0`$. We now have $`\mathrm{\Delta }T_\mathrm{c}/T_01.71an^{1/3}`$. Though this does not agree as well with the quoted simulation result, the moderatley small size of the correction supports the proposition that the large $`N`$ expansion works surprisingly well for $`T_\mathrm{c}`$.
In the remainder of this introduction, we review the long-distance O(2) effective theory for Bose condensation and then review the arguments of about how to calculate $`\mathrm{\Delta }T_\mathrm{c}/T_\mathrm{c}`$. In Sec. II, we review the leading-order calculation in large $`N`$ as done in . In Sec. III, we go on to calculate the next order in $`1/N`$. An appendix explains how to calculate some of the basic 3-dimensional integrals that appear at that order.
### A Review of effective theory
The basic assumption throughout will be that the average separation $`n^{1/3}`$ of atoms is large compared to the scattering length $`a`$. This can also be expressed as $`\lambda (T_0)a`$, where $`\lambda `$ is the thermal wavelength
$$\lambda (T)=\left(\frac{2\pi \mathrm{}^2}{mk_\mathrm{B}T}\right)^{1/2}.$$
(6)
It is well known that, at distance scales large compared to the scattering length $`a`$, an appropriate effective theory for a dilute Bose gas is the second-quantized Schrödinger equation, together with a chemical potential $`\mu `$ that couples to particle number density $`\psi ^{}\psi `$, and a $`|\psi |^4`$ contact interaction that reproduces low-energy scattering. The corresponding Lagrangian is
$$=\psi ^{}\left(i\mathrm{}_t\frac{\mathrm{}^2}{2m}^2\mu \right)\psi +\frac{2\pi \mathrm{}^2a}{m}(\psi ^{}\psi )^2.$$
(7)
In this context the corresponding mean-field equation of motion is called the Gross-Pitaevskii equation.<sup>§</sup><sup>§</sup>§ For a review, see . As with any effective theory, there are corrections represented by higher-dimensional, irrelevant interactions (in the sense of the renormalization group), At short distances, the $`_t`$ and $`^2`$ terms of the action $`𝑑td^3x`$ determine that times scales as (length)<sup>2</sup> and the scaling dimension of $`\psi `$ is (length)<sup>-3/2</sup>. such as $`(\psi ^{}\psi )^3`$ and $`\psi ^{}^4\psi `$. However, higher and higher dimension operators are parametrically less and less important if the distance scales of interest are large compared to the characteristic scales ($`a`$) of the atomic interactions. The $`(\psi ^{}\psi )^2`$ term in the Lagrangian (7) is in fact the lowest-dimension irrelevant interaction, and it is adequate for computing the leading-order effects of interactions in the diluteness expansion. For a discussion of analyzing corrections in this language, see ref. , which extended earlier work on corrections by refs. . A similar discussion for Fermi gases may be found in ref. .
Now treat the system at finite temperature using Euclidean time formalism. The field can then be decomposed into frequency modes with Matsubara frequencies $`\omega _n=2\pi nk_\mathrm{B}T/\mathrm{}`$. At sufficiently large distance scales ($`\lambda `$), and small chemical potential ($`|\mu |k_\mathrm{B}T`$), the $`(\mathrm{}^2/2m)^2\mu `$ terms in (7) become small compared to the $`O(\mathrm{}\omega _n)`$ time derivative term, provided $`n0`$. The non-zero Matsubara frequency modes then decouple from the dynamics, leaving behind an effective theory of only the zero-frequency modes $`\psi _0`$. Roughly,
$$\frac{1}{\mathrm{}}_0^\mathrm{}\beta 𝑑td^3x\beta d^3x\left[\psi _0^{}\left(\frac{\mathrm{}^2}{2m}^2\mu \right)\psi _0+\frac{2\pi \mathrm{}^2a}{m}(\psi _0^{}\psi _0)^2\right]$$
(8)
with $`\beta =1/k_BT`$. In detail, the parameters of (8) are renormalized by coupling to the non-zero modes, and there are again corrections in the form of irrelevant (and even marginal) interactions. However, these effects are all suppressed in the dilute limit<sup>\**</sup><sup>\**</sup>\** For the 3-dimensional effective theory (8), the short-distance scaling dimension of $`\psi _0`$ is (length)<sup>-1/2</sup>, the $`(\psi _0^{}\psi _0)^2`$ interaction is relevant, and a $`(\psi _0^{}\psi _0)^3`$ interaction would be marginal. Even though marginal, this last interaction can be ignored at the order of interest in the diluteness expansion because it has a small coefficient. For example, consider the term that would arise directly from the presence of a correction $`g_3(\psi ^{}\psi )^3`$ to the original Lagrangian (7). That would lead to a $`g_3(\psi _0^{}\psi _0)^3`$ term in (8) which, after rescaling, would become a term proportional to $`(mg_3/\mathrm{}^2\lambda ^4)(\varphi ^2)^3`$ in (9). Since $`\lambda \lambda (T_0)n^{1/3}`$ at the transition, the coefficient of this term is high order in the diluteness expansion in $`n^{1/3}`$. Similarly, an effective $`(\psi _0^{}\psi _0)^3`$ term arising from the 4-point interactions $`(\psi ^{}\psi )^2`$ and from integrating out physics at the scale $`\lambda `$ (due, for example, to non-zero Matsubara modes) would give rise to a $`(\varphi ^2)^3`$ term in (9) with coefficient proportional to $`u^3\lambda ^3\lambda ^3`$. and do not affect the computation of $`\mathrm{\Delta }T_\mathrm{c}/T_0`$ at leading order in $`an^{1/3}`$.
It is then convenient to write $`\psi _0=\mathrm{}^1(mk_\mathrm{B}T)^{1/2}(\varphi _1+i\varphi _2)`$ so that the effective action $`S=H/T`$ becomes a conventionally normalized O(2) field theory:
$$S=d^3x\left[\frac{1}{2}|\mathbf{}\varphi |^2+\frac{1}{2}r\varphi ^2+\frac{u}{4!}(\varphi ^2)^2\right],$$
(9)
where $`\varphi `$ is understood to be a 2-component real vector $`(\varphi _1,\varphi _2)`$ and
$$r=\frac{2m\mu }{\mathrm{}^2},u=\frac{96\pi ^2a}{\lambda ^2}.$$
(10)
### B Review of $`\mathrm{\Delta }T_\mathrm{c}/T_0`$
Our effective theory depends on two as yet undetermined parameters—$`r`$ and $`u`$, or equivalently $`\mu `$ and $`T`$. One constraint comes from fixing particle number density $`n`$:
$$n=\psi ^{}\psi =\frac{mk_\mathrm{B}T}{\mathrm{}^2}\varphi ^2.$$
(12)
At the critical temperature, a second requirement is that the system have infinite correlation length, which requires
$$\xi ^1=r+\mathrm{\Pi }(0)=0,$$
(13)
where $`\mathrm{\Pi }(p)`$ is the proper self-energy of the $`\varphi `$ field. The two equations (I B) determine the two unknowns $`r`$ and $`u`$, and hence $`T_\mathrm{c}`$. As noted by Baym et al., the density equation (12) can be rewritten as
$$\frac{\mathrm{}^2n}{mk_\mathrm{B}T}=_𝐩\frac{1}{p^2+r+\mathrm{\Pi }(p)}=_𝐩\frac{1}{p^2+[\mathrm{\Pi }(p)\mathrm{\Pi }(0)]},$$
(14)
where the last equality uses (13). Throughout this paper, we will use the notational shorthand
$$_𝐩\frac{d^3p}{(2\pi )^3}$$
(15)
for momentum integrals. (Technically, $`p`$ is a wave number rather than a momentum, but we will use conventional $`\mathrm{}=1`$ nomenclature, even though we have not set $`\mathrm{}`$ to 1.)
The expression (14) for the density is ultraviolet (UV) divergent and so receives contributions from short distance scales where the effective theory breaks down. This could be handled by appropriately regulating the effective theory and then perturbatively correcting the UV contribution. As pointed out by Baym et al., it is simpler to instead consider the difference $`nn_0(T)`$, where $`n_0(T)`$ is the same expression in the absence of interactions (i.e., with $`\mathrm{\Pi }`$ set to zero):
$$\frac{\mathrm{}^2[nn_0(T)]}{mk_\mathrm{B}T}=_𝐩\left[\frac{1}{p^2+r+\mathrm{\Pi }(p)}\frac{1}{p^2}\right]=_𝐩\left[\frac{1}{p^2+[\mathrm{\Pi }(p)\mathrm{\Pi }(0)]}\frac{1}{p^2}\right].$$
(16)
$`n_0(T)`$ represents the density a non-interacting Bose gas has if its transition temperature is $`T`$. It is given by inverting (2):
$$n_0(T)=\frac{\zeta (3/2)}{\lambda ^3(T)}.$$
(17)
This formula cannot be derived directly in the effective theory (9), but the difference $`nn_0`$ in (16) is insensitive to the UV and so can be.
The above constraints are entirely adequate to systematically determine $`\mathrm{\Delta }T_\mathrm{c}`$ in the large $`N`$ expansion, but there is a convenient way to simplify the bookkeeping a bit. Baym et al. give a simple argument that, to leading order in the density expansion,
$$\frac{\mathrm{\Delta }T_\mathrm{c}}{T_0}\frac{2}{3}\frac{[nn_0(T_\mathrm{c})]}{n},$$
(18)
where the factor of $`2/3`$ in (16) arises from the relation $`Tn_0^{2/3}`$. Combining (18) with (16), we can summarize as
$$\frac{\mathrm{\Delta }T_\mathrm{c}}{T_0}\frac{2mk_\mathrm{B}T_0}{3\mathrm{}^2n}_𝐩\left[\frac{1}{p^2+r+\mathrm{\Pi }(p)}\frac{1}{p^2}\right]=\frac{2mk_\mathrm{B}T_0}{3\mathrm{}^2n}_𝐩\left[\frac{1}{p^2+[\mathrm{\Pi }(p)\mathrm{\Pi }(0)]}\frac{1}{p^2}\right].$$
(19)
to leading order in $`an^{1/3}`$. It’s also useful to rephrase this, again in terms of the fields $`\varphi `$ of the effective theory, as
$$\frac{\mathrm{\Delta }T_\mathrm{c}}{T_0}\frac{2mk_\mathrm{B}T_0}{3\mathrm{}^2n}\mathrm{\Delta }\varphi ^2,$$
(20)
where
$$\mathrm{\Delta }\varphi ^2\varphi ^2\varphi ^2_{\mathrm{\Pi }0}.$$
(21)
Note that the problem of calculating $`\mathrm{\Delta }\varphi ^2`$ from the action (9), subject to the constraint (13), has only one dimensionful scale in it: $`u`$. The length scale of this problem, which will be the length scale of the physics that determines $`\mathrm{\Delta }T_\mathrm{c}/T_0`$, is therefore
$$u^1\frac{\lambda ^2}{a}$$
(22)
by dimensional analysis. In the dilute limit $`\lambda (T_\mathrm{c})a`$, this length scale is large compared to $`\lambda `$, which justifies use of the O(2) effective theory (9).
## II Review of leading order in $`1/N`$
We now review the leading large-$`N`$ calculation of $`\varphi ^2\varphi ^2_{\mathrm{\Pi }0}`$, and hence of $`\mathrm{\Delta }T_\mathrm{c}/T_0`$, by Baym et al. The details of our calculation are slightly different than theirs, and we will introduce techniques needed to proceed to higher order. We start with the standard large-$`N`$ generalization of the O(2) scalar field theory (9) to an $`O(N)`$ scalar theory: replace $`\varphi `$ by an $`N`$-component vector and treat $`Nu`$ as fixed in the $`N\mathrm{}`$ limit. The reader should keep in mind that $`u`$ is therefore order $`1/N`$. Standard $`N`$ power counting of Feynman diagrams consists of a power of $`u1/N`$ for each 4-point vertex and a power of $`N`$ for each flavor trace.
The set of diagrams that determine $`\mathrm{\Delta }\varphi ^2`$ at leading order in $`1/N`$ is depicted in Fig. 1, where the dashed line denotes bubble chains, as shown in Fig. 2. (For comparison, the diagram for $`\varphi ^2_{\mathrm{\Pi }0}`$ is shown in Fig. 3.) The cross denotes an insertion of the operator $`\varphi ^2`$, whose expectation we are computing. There is a simple way to summarize the effect on diagrammatic perturbation theory of the $`r\varphi ^2`$ term in the action (9) and the constraint (13) that $`r=\mathrm{\Pi }(0)`$.
> Rule 1: Use massless (gapless) scalar propagators $`1/p^2`$ when evaluating diagrams, ignoring the $`r\varphi ^2`$ term in the action. But whenever there is a one-particle irreducible sub-diagram $`X`$ that represents a contribution to the $`\varphi `$ proper self-energy $`\mathrm{\Pi }(p)`$, then<sup>††</sup><sup>††</sup>†† This rule is unambiguous for calculating expectations such as $`\varphi ^2`$. It is potentially ambiguous for calculating the free energy—for example, a diagram like Fig. 1 but without the cross on it. In that case, it is ambiguous which sub-diagrams would be considered self-energy insertions. A systematic way to treat the perturbation theory in all cases is to treat the $`r\varphi ^2`$ term in the action as a perturbation, include it in Feynman diagrams as a 2-point vertex, and then set $`r=\mathrm{\Pi }(0)`$ order by order in perturbation theory. replace $`X(p)`$ by $`X(p)X(0)`$.
We note for later reference that, for the purpose of this rule, a diagram that is cut in two pieces only by cutting a single internal dashed line is still one-particle irreducible, because cutting the bubble chain represented by a dashed line corresponds to cutting two $`\varphi `$ lines (Fig. 2).
The bubble chain sum shown in fig. 2 is given by
$$\text{}=N^1F_p,$$
(23)
where
$$F_p\frac{1}{\frac{3}{Nu}+\stackrel{~}{\mathrm{\Sigma }}_0(p)}$$
(25)
and $`\stackrel{~}{\mathrm{\Sigma }}_0(p)`$ represents the basic massless bubble integral (not summed over flavors)
$$\stackrel{~}{\mathrm{\Sigma }}_0(p)\frac{1}{2}_𝐥\frac{1}{l^2|𝐥+𝐩|^2}.$$
(26)
In $`d=3`$ dimensions,
$$\stackrel{~}{\mathrm{\Sigma }}_0(p)=\frac{1}{16p}.$$
(27)
Putting everything together, the diagram of fig. 1 gives
$$\mathrm{\Delta }\varphi ^2=_{\mathrm{𝐥𝐩}}\frac{F_p}{l^4}\left[\frac{1}{|𝐥+𝐩|^2}\frac{1}{p^2}\right]+O(N^1).$$
(28)
As pointed out by Baym et al. in , the above integral is not absolutely convergent in three dimensions, and one must be careful to consistently regulate the theory before proceeding. Integrals that are not absolutely convergent are at best ambiguous—they depend on the order one chooses to do the integrations. For example, if one evaluates (28) directly in three dimensions, doing the angular integrations first, then the $`l`$ integration, and then the $`p`$ integration, the result is zero. This is not in fact the correct answer. We will discuss this issue in some detail in order to justify the correctness of our procedure for later evaluating higher-order diagrams.
Baym et al.’s preferred method for the leading-order calculation is to use dimensional regularization and evaluate everything in $`d=3ϵ`$ dimensions. This is difficult at next order in $`1/N`$: the loop integrals we shall encounter are sufficiently complicated that evaluation in $`d=3ϵ`$ dimensions seems hard. Our strategy will be to instead always reduce diagrams to well-defined three-dimensional integrals, which are simpler to evaluate. We imagine starting with some consistent regularization scheme, like dimensional regularization, and will now discuss how to manipulate the integrals so that they will be absolutely convergent if we set $`d=3`$. We assume in what follows that the UV regulator respects parity and is invariant under shifts $`𝐩𝐩+𝐤`$ of loop momenta $`𝐩`$.
Let’s look at the divergences that cause absolute convergence of the integral (28) to fail in three dimensions. A simple one to correct is the behavior for $`𝐥`$ fixed and $`p\mathrm{}`$. The large $`p`$ piece of the $`𝐩`$ integration then behaves as $`_𝐩𝐩𝐥/p^4`$, which is logarithmically UV divergent (from the point of view of absolute convergence). This can be remedied by rewriting the regulated version of (28) by using the freedom to change the integration variable $`𝐩`$ to $`𝐩`$:
$$\mathrm{\Delta }\varphi ^2_{\mathrm{LO}}=_{\mathrm{𝐥𝐩}}\frac{F_p}{l^4}\left[\frac{1}{2|𝐥+𝐩|^2}+\frac{1}{2|𝐥𝐩|^2}\frac{1}{p^2}\right].$$
(29)
Now, if we throw away the UV regulator, the $`𝐥`$ fixed, large $`p`$ divergence is gone. This sort of divergence is trivial, easy to remedy, and won’t have much practical impact on our calculations (given the order in which we will eventually do integrations). We’ll simply acknowledge the issue in later calculations, without emphasizing it, by writing (29) as
$$\mathrm{\Delta }\varphi ^2_{\mathrm{LO}}=_{\mathrm{𝐥𝐩}}\frac{F_p}{l^4}\left[\frac{1}{|𝐥+𝐩|^2}\frac{1}{p^2}\right]_\pm ,$$
(30)
where the subscript $`\pm `$ means that one should average the expression with $`𝐩𝐩`$ (or equivalently with $`𝐥𝐥`$).
Unfortunately, even (29) is not absolutely convergent. There is still a logarithmic UV divergence associated with $`𝐥`$ and $`𝐩`$ simultaneously becoming large ($`lp\mathrm{}`$), as can be seen by simple power counting and the fact that $`F_p`$ approaches a non-zero constant for large $`p`$. Return to considering (29) with a UV regulator still in place. We can eliminate the UV divergence by rewriting
$$\mathrm{\Delta }\varphi ^2_{\mathrm{LO}}=_{\mathrm{𝐥𝐩}}\frac{(F_pF_{\mathrm{}})}{l^4}\left[\frac{1}{|𝐥+𝐩|^2}\frac{1}{p^2}\right]_\pm _{\mathrm{𝐥𝐩}}\frac{F_{\mathrm{}}}{l^4}\left[\frac{1}{|𝐥+𝐩|^2}\frac{1}{p^2}\right]_\pm .$$
(31)
The second integral vanishes, as can be seen by changing integration variable $`𝐩𝐩𝐥`$ in its first term. So
$$\mathrm{\Delta }\varphi ^2_{\mathrm{LO}}=_{\mathrm{𝐥𝐩}}\frac{(F_pF_{\mathrm{}})}{l^4}\left[\frac{1}{|𝐥+𝐩|^2}\frac{1}{p^2}\right]_\pm .$$
(32)
This is now UV convergent because $`F_pF_{\mathrm{}}0`$ as $`p\mathrm{}`$. However, we have traded the logarithmic UV divergence for a logarithmic infrared (IR) divergence, associated with $`pl0`$.
We now need some sort of infrared regulator. One physically motivated possibility for consistently regulating the infrared would be to consider the system infinitesimally above the critical temperature, so that all the massless scalar propagators $`1/p^2`$ should be replaced by massive ones $`1/(p^2+M^2)`$, where the mass $`M`$ represents the inverse correlation length $`\xi ^1`$. This defines an absolutely convergent integral in 3 dimensions, and the limit $`M0`$ would be taken only after the integrations.
Massless propagators $`1/p^2`$ will be much easier to deal with, however, in higher-order calculations. As a practical matter for computing diagrams, we prefer to introduce as few massive propagators as possible. It would be convenient, for example, to IR regulate (32) by introducing $`M`$ only in the $`1/l^2`$ propagators:
$$\mathrm{\Delta }\varphi ^2_{\mathrm{LO}}=\underset{M0}{lim}_{\mathrm{𝐥𝐩}}\frac{(F_pF_{\mathrm{}})}{(l^2+M^2)^2}\left[\frac{1}{|𝐥+𝐩|^2}\frac{1}{p^2}\right]_\pm ,$$
(33)
where $`F_p`$ is still defined in terms of the massless bubble integral, as in (23). One might worry that an ad hoc procedure of putting masses only on some propagators could be inconsistent, so let us argue more carefully. Return to the UV regulated version of (29) and note that the integral is not sensitive to the the region of integration where $`l`$ is infinitesimal, because this particular integral is IR convergent. There’s then no reason we can’t modify the infrared behavior of the integrand for infinitesimal $`l`$, without affecting the integral. So, for instance,
$$\mathrm{\Delta }\varphi ^2_{\mathrm{LO}}=\underset{M0}{lim}_{\mathrm{𝐥𝐩}}\frac{F_p}{(l^2+M^2)^2}\left[\frac{1}{|𝐥+𝐩|^2}\frac{1}{p^2}\right]_\pm .$$
(34)
But now, again rewriting $`F_p=(F_pF_{\mathrm{}})+F_{\mathrm{}}`$, the same steps as before reproduce (33).
Now that we have an absolutely convergent integral (33), we can do the integration in three dimensions and in any order we choose. It’s convenient to do the $`𝐥`$ integral first:
$$_𝐥\frac{1}{(l^2+M^2)^2}\left[\frac{1}{|𝐥+𝐩|^2}\frac{1}{p^2}\right]_\pm =\frac{1}{8\pi M(p^2+M^2)}\frac{1}{8\pi Mp^2}=\frac{M}{8\pi p^2(p^2+M^2)}.$$
(35)
The “$`\pm `$” prescription makes no difference to this particular integral, because the $`𝐥`$ integration by itself is completely convergent without it. Note that naively setting $`M`$ to zero at this stage would give the incorrect, zero result mentioned earlier. Instead, we have
$$\mathrm{\Delta }\varphi ^2_{\mathrm{LO}}=\underset{M0}{lim}_𝐩(F_pF_{\mathrm{}})\frac{M}{8\pi p^2(p^2+M^2)}.$$
(36)
The overall factor of $`M`$ in the numerator is canceled by a linear IR divergence in the $`𝐩`$ integration, which is cut off by $`M`$.
For small $`M`$, the integral (36) is dominated<sup>‡‡</sup><sup>‡‡</sup>‡‡ Some readers may worry that the integral (36) is dominated by arbitrarily small $`pM0`$. They may worry because at sufficiently small momentum our perturbative propagators are no longer good approximations to the full propagators. The full scalar propagators, for example, actually scale like $`1/l^{2+\eta }`$ rather than $`1/l^2`$ at small $`l`$ ($`Nu`$), where the critical exponent $`\eta `$ is $`O(N^1)`$. The difference becomes significant when $`lNu\mathrm{exp}(\eta )=Nu\mathrm{exp}[O(1/N)]`$. One might worry that the sensitivity of (36) to $`p0`$ is a sign that naive large $`N`$ perturbation theory must break down. It is important to realize, in the present case, that this infrared sensitivity is simply an artifact of our mathematical manipulations on the infrared-safe expression (33). Regardless of whether one used some sort of infrared-improved propagator in (33), that expression is not sensitive to far-infrared momenta. It is sensitive to momenta $`Nu`$, for which there is nothing wrong with a large-$`N`$ expansion based on perturbative propagators. by $`pM`$. So, in the limit of $`M0`$, we can simplify the calculation slightly by replacing $`F_pF_{\mathrm{}}`$ by $`F_0F_{\mathrm{}}`$. So
$$\mathrm{\Delta }\varphi ^2_{\mathrm{LO}}=(F_0F_{\mathrm{}})_𝐩\frac{M}{8\pi p^2(p^2+M^2)}=\frac{F_0F_{\mathrm{}}}{32\pi ^2}=\frac{Nu}{96\pi ^2}.$$
(37)
When combined with the formula (20) for $`\mathrm{\Delta }T_\mathrm{c}`$, this reproduce Baym et al.’s leading large $`N`$ result (1), in which $`N`$ has been set to 2.
## III Next order in $`1/N`$
The diagrams which contribute to $`\mathrm{\Delta }\varphi ^2`$ at next order in $`1/N`$ are shown in Figs. 4 and 5. The diagrammatic expansion comes from the standard introduction of an auxiliary field $`\sigma `$, represented by the dashed lines.<sup>\**</sup><sup>\**</sup>\** For a very quick review of standard large $`N`$, see, for example, section 2.1 of chapter 8 of . Some people might prefer to replace $`\sigma `$ by $`i\sigma `$ in the action (38), so that the Euclidean path integral for $`\sigma `$ is convergent, but it matters not at all for the purpose of large $`N`$ perturbation theory. The $`O(N)`$ action of (9) is rewritten as
$$S=d^3x\left[\frac{1}{2}|\mathrm{\Delta }\varphi |^2+\frac{1}{2}r\varphi ^2+\frac{1}{2}\varphi ^2\sigma \frac{1}{6u}\sigma ^2\right].$$
(38)
The $`\sigma `$ propagator is then turned into the bubble chain of fig. 2 by resumming the basic massless bubble of fig. 6 into the $`\sigma `$ propagator. Technically, this is accomplished by trivially rewriting (38) as
$`S`$ $`=`$ $`S_0+S_{\mathrm{subtractions}}+{\displaystyle d^3x\frac{1}{2}\varphi ^2\sigma },`$ (40)
$`S_0`$ $`=`$ $`{\displaystyle _𝐩}\left[\frac{1}{2}\varphi _𝐩p^2\varphi _𝐩+\frac{1}{2}\sigma _𝐩(NF_𝐩^1)\sigma _𝐩\right],`$ (41)
$`S_{\mathrm{subtractions}}`$ $`=`$ $`{\displaystyle _𝐩}\left[\frac{1}{2}r\varphi _𝐩\varphi _𝐩+\frac{1}{2}\sigma _𝐩N\stackrel{~}{\mathrm{\Sigma }}_0(p)\sigma _𝐩\right],`$ (42)
with $`F_p`$ and $`\stackrel{~}{\mathrm{\Sigma }}_0`$ given by (23). The terms designated $`S_{\mathrm{subtractions}}`$ may be ignored if one follows the previous Rule 1 as well as
> Rule 2: Do not include any diagrams that have the one-loop bubble, fig. 6, as a sub-diagram.
Note that Rule 1 eliminates any tadpole sub-diagrams, such as fig. 7. Formal large $`N`$ counting of diagrams is simply to count a factor of $`N^1`$ for each $`\sigma `$ propagator and a factor of $`N`$ for each $`\varphi `$ loop. The important momentum scale of the problem will be the scale $`pNu=O(N^0)`$, where the $`\sigma `$ propagator (23) makes the transition from its small $`p`$ behavior ($`F_pp`$) to its large $`p`$ behavior ($`F_p\text{constant}`$). Some authors like to completely integrate $`\varphi `$ out of (38), but we prefer to retain it, as there is then a more transparent relationship between Feynman diagrams and the corresponding Feynman integrals.
In evaluating the diagrams of fig. 4, we shall borrow techniques from ref. , where somewhat related diagrams were evaluated in gauge theories with large numbers of scalars. Our strategy will be to do the $`\varphi `$ loop integrals first, and then tackle the remaining integrals associated with $`\sigma `$ propagators.
### A Diagram a
Let’s start with Fig. 4a. The corresponding integral is
$$\mathrm{\Delta }\varphi ^2_\mathrm{a}=N^1_{\mathrm{𝐩𝐪}}F_pF_q_𝐥\frac{1}{l^6}\left[\frac{1}{|𝐥+𝐩|^2}\frac{1}{p^2}\right]_\pm \left[\frac{1}{|𝐥+𝐪|^2}\frac{1}{q^2}\right]_\pm .$$
(43)
As written, this integral is absolutely convergent and can be evaluated, without regularization, directly in three dimensions. To do the $`𝐥`$ integration, however, we find it convenient to temporarily introduce an IR regulator mass $`M`$. We may then separately integrate each of the terms of the integrand, which are not individually IR convergent. We can also use $`M`$ as a trick for reducing powers of $`l^2`$. Specifically, we rewrite the $`𝐥`$ integral as the $`M0`$ limit of
$`{\displaystyle _𝐥}{\displaystyle \frac{1}{(l^2+M^2)^3}}\left[{\displaystyle \frac{1}{|𝐥+𝐩|^2}}{\displaystyle \frac{1}{p^2}}\right]_\pm \left[{\displaystyle \frac{1}{|𝐥+𝐪|^2}}{\displaystyle \frac{1}{q^2}}\right]_\pm `$ (44)
$`={\displaystyle \frac{1}{2}}{\displaystyle \frac{d^2}{d(M^2)^2}}{\displaystyle _𝐥}{\displaystyle \frac{1}{(l^2+M^2)}}\left[{\displaystyle \frac{1}{|𝐥+𝐩|^2}}{\displaystyle \frac{1}{p^2}}\right]_\pm \left[{\displaystyle \frac{1}{|𝐥+𝐪|^2}}{\displaystyle \frac{1}{q^2}}\right]_\pm `$ (45)
$`={\displaystyle \frac{1}{2}}{\displaystyle \frac{d^2}{d(M^2)^2}}\left[I_1(𝐩,𝐪;M)p^2J_1(𝐪;M)q^2J_1(𝐩;M)+p^2q^2\left({\displaystyle \frac{M}{4\pi }}\right)\right]_\pm ,`$ (46)
(47)
where
$$I_1(𝐩,𝐪;M)_𝐥\frac{1}{(l^2+M^2)|𝐥+𝐩|^2|𝐥+𝐪|^2},$$
(48)
$$J_1(𝐩;M)_𝐥\frac{1}{(l^2+M^2)|𝐥+𝐩|^2},$$
(49)
and<sup>\*†</sup><sup>\*†</sup>\*† The integral $`_𝐥(l^2+M^2)^1`$ is $`M/4\pi `$ plus an $`M`$-independent UV divergence, and $`M^2`$ derivatives of the latter vanish. It’s of course not necessary to introduce $`_𝐥(l^2+M^2)^1`$ and this spurious UV divergence; one could simply evaluate $`_𝐥(l^2+M^2)^3`$ directly. But we find it convenient to consolidate the treatment of such integrals with that of the other terms in (47).
$$_𝐥\frac{1}{(l^2+M^2)^3}=\frac{1}{2}\frac{d^2}{d(M^2)^2}_𝐥\frac{1}{l^2+M^2}=\frac{1}{2}\frac{d^2}{d(M^2)^2}\left(\frac{M}{4\pi }\right).$$
(50)
The integral $`J_1`$ is straightforward to evaluate. A particularly simple way to evaluate $`I_1`$ is to make a conformal transformation which reduces it to the form of $`J_1`$. The results of both integrals, and the conformal transformation between them, are discussed in Appendix A. All we need here are the small $`M`$ expansions of those results, which turn out to be
$`I_1(𝐩,𝐪;M)`$ $`=`$ $`{\displaystyle \frac{1}{8pq|𝐩𝐪|}}{\displaystyle \frac{M}{4\pi p^2q^2}}{\displaystyle \frac{M^2𝐩𝐪}{8p^3q^3|𝐩𝐪|}}`$ (52)
$`+{\displaystyle \frac{M^3(p^2+4𝐩𝐪+q^2)}{12\pi p^4q^4}}+{\displaystyle \frac{M^4(3(𝐩𝐪)^2p^2q^2)}{16p^5q^5|𝐩𝐪|}}+O(M^5),`$
$$J_1(𝐩;M)=\frac{1}{8p}\frac{M}{4\pi p^2}+\frac{M^3}{12\pi p^4}+O(M^5).$$
(53)
Putting everything together,
$`{\displaystyle _𝐥}{\displaystyle \frac{1}{(l^2+M^2)^3}}\left[{\displaystyle \frac{1}{|𝐥+𝐩|^2}}{\displaystyle \frac{1}{p^2}}\right]_\pm \left[{\displaystyle \frac{1}{|𝐥+𝐪|^2}}{\displaystyle \frac{1}{q^2}}\right]_\pm `$ (54)
$`=\left[{\displaystyle \frac{\widehat{𝐩}\widehat{𝐪}}{8\pi Mp^3q^3}}+{\displaystyle \frac{3(\widehat{𝐩}\widehat{𝐪})^21}{16p^3q^3|𝐩𝐪|}}\right]_\pm +O(M)`$ (55)
$`=\left[{\displaystyle \frac{3(\widehat{𝐩}\widehat{𝐪})^21}{16p^3q^3|𝐩𝐪|}}\right]_\pm +O(M).`$ (56)
We can now set $`M=0`$. All that will matter in the integral (43) is the average $`\mathrm{}_\theta `$ over the angle between $`𝐩`$ and $`𝐪`$, which is
$$_𝐥\frac{1}{l^6}\left[\frac{1}{|𝐥+𝐩|^2}\frac{1}{p^2}\right]_\pm \left[\frac{1}{|𝐥+𝐪|^2}\frac{1}{q^2}\right]_\pm _\theta =\frac{1}{40p_>^6p_<},$$
(57)
where
$$p_>\mathrm{max}(p,q),p_<\mathrm{min}(p,q).$$
(58)
We are left with
$$\mathrm{\Delta }\varphi ^2_\mathrm{a}=\frac{1}{N}_{\mathrm{𝐩𝐪}}\frac{F_pF_q}{40p_>^6p_<}=\frac{1}{2\pi ^4N}_0^{\mathrm{}}p^2𝑑p_0^pq^2𝑑q\frac{F_pF_q}{40p^6q}.$$
(59)
The remaining integrals are easy to do, with the result
$$\mathrm{\Delta }\varphi ^2_\mathrm{a}=\frac{u}{15\pi ^4}\left(\frac{\pi ^2}{6}\frac{5}{4}\right).$$
(60)
### B Diagram b
Fig. 4b corresponds to
$$\mathrm{\Delta }\varphi ^2_\mathrm{b}=N^1_{\mathrm{𝐩𝐪}}F_pF_qB_{\mathrm{𝐩𝐪}},$$
(62)
$$B_{\mathrm{𝐩𝐪}}_𝐥\frac{1}{l^4}\left\{\frac{1}{|𝐥+𝐩|^4}\left[\frac{1}{|𝐥+𝐩+𝐪|^2}\frac{1}{q^2}\right]\frac{1}{p^4}\left[\frac{1}{|𝐩+𝐪|^2}\frac{1}{q^2}\right]\right\}_\pm .$$
(63)
This contribution to $`\mathrm{\Delta }\varphi ^2`$ is again absolutely convergent if the subscript $`\pm `$ is taken to mean averaging over $`𝐩𝐩`$ and also over $`𝐪𝐪`$. It is convenient to now rewrite the $`𝐥`$ integral as the $`M_1,M_20`$ limit of
$`{\displaystyle _𝐥}{\displaystyle \frac{1}{(l^2+M_1^2)^2}}\left\{{\displaystyle \frac{1}{(|𝐥+𝐩|^2+M_2^2)^2}}\left[{\displaystyle \frac{1}{|𝐥+𝐩+𝐪|^2}}{\displaystyle \frac{1}{q^2}}\right]{\displaystyle \frac{1}{(p^2+M_2^2)^2}}\left[{\displaystyle \frac{1}{|𝐩+𝐪|^2}}{\displaystyle \frac{1}{q^2}}\right]\right\}_\pm `$ (64)
$`={\displaystyle \frac{d}{d(M_1^2)}}{\displaystyle \frac{d}{d(M_2^2)}}\{I_2(𝐩+𝐪,𝐪,M_1,M_2){\displaystyle \frac{1}{q^2}}J_2(𝐩;M_1,M_2)`$ (65)
$`{\displaystyle \frac{1}{(p^2+M_2^2)}}[{\displaystyle \frac{1}{|𝐩+𝐪|^2}}{\displaystyle \frac{1}{q^2}}]({\displaystyle \frac{M_1}{4\pi }})\}_\pm ,`$ (66)
where
$$I_2(𝐩,𝐪;M_1,M_2)_𝐥\frac{1}{l^2(|𝐥+𝐩|^2+M_1^2)(|𝐥+𝐪|^2+M_2^2)},$$
(68)
$$J_2(𝐩;M_1,M_2)_𝐥\frac{1}{(l^2+M_1^2)(|𝐥+𝐩|^2+M_2^2)}.$$
(69)
The results for $`I_2`$ and $`J_2`$, and their small $`M_1,M_2`$ expansions, are given in Appendix A. The final result for the $`𝐥`$ integration, after taking the $`M_1,M_20`$ limit, is
$$B_{\mathrm{𝐩𝐪}}=\left[\frac{q2p(\widehat{𝐩}\widehat{𝐪})3q(\widehat{𝐩}\widehat{𝐪})^2}{8p^3q^2|𝐩+𝐪|^3}\right]_\pm ,$$
(70)
with angular average
$$B_{\mathrm{𝐩𝐪}}_\theta =\frac{\theta (pq)}{4p^6q},$$
(71)
where $`\theta (pq)`$ is the step function (1 for $`p>q`$; 0 for $`p<q`$). The remaining integrals are easy to do, giving
$$\mathrm{\Delta }\varphi ^2_\mathrm{b}=\frac{u}{3\pi ^4}\left(\frac{\pi ^2}{6}\frac{5}{4}\right).$$
(72)
### C Diagram c
Fig. 4c can be evaluated as the others, but the final integrals are a bit more complex. The diagram gives
$$\mathrm{\Delta }\varphi ^2_\mathrm{c}=N^1_{\mathrm{𝐩𝐪}}F_pF_qC_{\mathrm{𝐩𝐪}},$$
(74)
$$C_{\mathrm{𝐩𝐪}}_𝐥\frac{1}{l^4}\left[\frac{1}{|𝐥+𝐩|^2|𝐥+𝐪|^2|𝐥+𝐩+𝐪|^2}\frac{1}{p^2q^2|𝐩+𝐪|^2}\right]_\pm .$$
(75)
The $`𝐥`$ integration can be performed using methods similar to before:
$$C_{\mathrm{𝐩𝐪}}=\underset{M0}{lim}\frac{d}{d(M^2)}\left[H(𝐩,𝐪;M)\frac{1}{p^2q^2|𝐩+𝐪|^2}\left(\frac{M}{4\pi }\right)\right]_\pm ,$$
(76)
$$H(𝐩,𝐪;M)_𝐥\frac{1}{(l^2+M^2)|𝐥+𝐩|^2|𝐥+𝐪|^2|𝐥+𝐩+𝐪|^2}.$$
(77)
$`H`$ can be reduced to the basic integrals $`I_1`$ and $`J_1`$ encountered previously by rewriting the numerator 1 in (77) as
$$1=\frac{1}{2𝐩𝐪+M^2}\left[(l^2+M^2)+|𝐥+𝐩+𝐪|^2|𝐥+𝐩|^2|𝐥+𝐪|^2\right]$$
(78)
and then expanding the integrand into the corresponding four terms:
$$H(𝐩,𝐪;M)=\frac{1}{2𝐩𝐪+M^2}\left[I_1(𝐩,𝐪;0)+I_1(𝐩,𝐪;M)I_1(𝐩+𝐪,𝐪;M)I_1(𝐩+𝐪,𝐩;M)\right].$$
(79)
Using the expansion (52) of $`I_1`$, one obtains
$$C_{\mathrm{𝐩𝐪}}=\left[\frac{1}{16p^3q^3}\left(1+(\widehat{𝐩}\widehat{𝐪})^2\right)\left(\frac{1}{|𝐩𝐪|}\frac{1}{|𝐩+𝐪|}\right)+\frac{\widehat{𝐩}\widehat{𝐪}(\widehat{𝐩}\widehat{𝐪})^1}{8p^2q^2|𝐩+𝐪|^3}\right]_\pm .$$
(80)
This simplifies, after applying the $`[\mathrm{}]_\pm `$ prescription, to
$$C_{\mathrm{𝐩𝐪}}=\frac{\widehat{𝐩}\widehat{𝐪}(\widehat{𝐩}\widehat{𝐪})^1}{16p^2q^2}\left[\frac{1}{|𝐩+𝐪|^3}\frac{1}{|𝐩𝐪|^3}\right].$$
(81)
Angular averaging yields
$$C_{\mathrm{𝐩𝐪}}_\theta =\frac{1}{8p_>^7x^2(1+x^2)}\left[x+\frac{\mathrm{Sinh}^1x}{\sqrt{1+x^2}}\right],$$
(82)
where
$$xp_</p_>.$$
(83)
The remaining integrals over $`𝐩`$ and $`𝐪`$ are no longer so trivial. Notice first that for an arbitrary function $`f(x)`$ one can rewrite
$$\frac{1}{N}_{\mathrm{𝐩𝐪}}\frac{F_pF_q}{p_>^7}f(x)=\frac{1}{2\pi ^4N}_0^{\mathrm{}}𝑑p_0^1𝑑x\frac{F_pF_{xp}}{p^2}x^2f(x)=\frac{8u}{3\pi ^4}_0^1𝑑x\frac{x^3\mathrm{ln}x}{(1x)}f(x),$$
(84)
so that
$$\mathrm{\Delta }\varphi ^2_\mathrm{c}=\frac{u}{3\pi ^4}_0^1𝑑x\frac{x\mathrm{ln}x}{(1x)}\left[\frac{x}{1+x^2}+\frac{\mathrm{Sinh}^1x}{(1+x^2)^{3/2}}\right].$$
(85)
This is easiest to evaluate numerically, giving
$$\mathrm{\Delta }\varphi ^2_\mathrm{c}=\frac{cu}{3\pi ^4},$$
(86)
where $`c0.463715`$. We also have an analytic result:<sup>\*‡</sup><sup>\*‡</sup>\*‡ Our inelegant, brute force method for obtaining this result is borrowed from a footnote of ref. . The hard part is the $`\mathrm{Sinh}^1`$ term. We change variables from $`x`$ to $`y=x+\sqrt{1+x^2}`$. This turns $`\mathrm{Sinh}^1x`$ into $`\mathrm{ln}y`$. $`\mathrm{ln}x`$ can be written as a sum of terms of the form $`\mathrm{ln}(ya)`$, and the change of integration variable makes the rest of the integrand a rational function of $`y`$. We then split this rational function apart by partial fractions and do each integral, yielding di- and tri-logarithms of various arguments. Finally, we use a zoo of polylogarithm identities to simplify the answer.
$$c=\frac{\pi ^2}{48}\left[1+\frac{7}{\sqrt{2}}\mathrm{ln}(1+\sqrt{2})\right]\frac{2}{3}L(3,\chi _8),$$
(87)
where
$$L(s,\chi _8)=1\frac{1}{3^s}\frac{1}{5^s}+\frac{1}{7^s}+\frac{1}{9^s}\frac{1}{11^s}\frac{1}{13^s}+\frac{1}{15^s}+\frac{1}{17^s}\mathrm{}$$
(89)
is a particular case of Dirichlet’s $`L`$-function, and
$$L(3,\chi _8)=0.958380454563\mathrm{}.$$
(90)
### D Diagram d
Fig. 4d corresponds to
$$\mathrm{\Delta }\varphi ^2_\mathrm{d}=N^1_{\mathrm{𝐩𝐪}}F_pF_qF_{|𝐪𝐩|}D_{\mathrm{𝐩𝐪}},$$
(92)
$`D_{\mathrm{𝐩𝐪}}`$ $``$ $`{\displaystyle _𝐥^{}}{\displaystyle \frac{1}{l_{}^{}{}_{}{}^{2}|𝐥^{}+𝐩|^2|𝐥^{}+𝐪|^2}}{\displaystyle _𝐥}{\displaystyle \frac{1}{l^4}}\left[{\displaystyle \frac{1}{|𝐥+𝐩|^2|𝐥+𝐪|^2}}{\displaystyle \frac{1}{p^2q^2}}\right]_{\pm 𝐥}`$ (93)
$`=`$ $`I_1(𝐩,𝐪;0)\underset{M0}{lim}{\displaystyle \frac{d}{d(M^2)}}\left[I_1(𝐩,𝐪,M)p^2q^2\left({\displaystyle \frac{M}{4\pi }}\right)\right]`$ (94)
$`=`$ $`{\displaystyle \frac{\widehat{𝐩}\widehat{𝐪}}{64p^3q^3|𝐩𝐪|^2}}.`$ (95)
Here the subscript $`\pm 𝐥`$ means we implicitly average over $`𝐥𝐥`$ for absolute convergence. Doing the remaining integrals by brute force, we find
$$\mathrm{\Delta }\varphi ^2_\mathrm{d}=\frac{u}{\pi ^4}\left[\frac{7}{12}\left(\zeta (3)\frac{\pi ^2}{6}\right)+\frac{1}{6}\right],$$
(96)
where $`\zeta `$ is the Riemann zeta function. It’s interesting to note that $`\pi ^2/6`$ can also be written as $`\zeta (2)`$.
### E Diagram e
The final class of diagrams, fig. 4e, correspond to the leading order diagram, fig. 1, with the replacement
$$N^1F_𝐩(N^1F_𝐩)[\mathrm{\Sigma }(p)](N^1F_𝐩)N^1_𝐩,$$
(97)
where $`\mathrm{\Sigma }(p)`$ represents the contribution to the $`\sigma `$ self-energy at next-to-leading order, shown in the diagrams of fig. 5. Making this substitution in the leading-order calculation (37) gives
$$\mathrm{\Delta }\varphi ^2_\mathrm{e}=\frac{_0_{\mathrm{}}}{32\pi ^2}.$$
(98)
The nice feature of this relation is that we only need to calculate $`\mathrm{\Sigma }(p)`$ in the small and large $`p`$ limits.
The first self-energy diagram in fig. 5 contributes
$$\mathrm{\Sigma }_1(𝐩)=\frac{1}{2}_𝐪F_qF_{|𝐩+𝐪|}_{𝐥_1}\frac{1}{l_1^2|𝐥_1+𝐩|^2|𝐥_1+𝐩+𝐪|^2}_{𝐥_2}\frac{1}{l_2^2|𝐥_2+𝐩|^2|𝐥_2+𝐩+𝐪|^2}.$$
(99)
For small $`𝐩`$, the integration is dominated by $`l_1,l_2p`$ and $`qNu`$. So we can ignore $`𝐥_1`$, $`𝐥_2`$, and $`𝐩`$ compared to $`𝐪`$ and write For general $`p`$, the result is $`Nu\mathrm{\Sigma }_1(p)=24\pi ^2z^3\mathrm{Re}[2\mathrm{Li}_2(z)2\mathrm{Li}_2(2+z)+\frac{1}{2}\pi ^2]`$, where $`z48p/Nu`$, and $`\mathrm{Li}_2(z)_0^z𝑑x\mathrm{ln}(1x)/x`$ is the dilogarithm function. One may double check that the $`p0`$ limit agrees with (100).
$$\mathrm{\Sigma }_1(𝐩)\frac{1}{2}_𝐪q^4F_q^2\left[_𝐥\frac{1}{l^2|𝐥+𝐩|^2}\right]^2=\frac{Nu}{48\pi ^2p^2}.$$
(100)
This diagram has a quadratic IR divergence for $`𝐩=0`$. In contrast, the other two diagrams only diverge as linear$`\times `$log when $`𝐩=0`$, and so behave as $`p^1\mathrm{ln}p`$ for small $`p`$. In summary,
$$\mathrm{\Sigma }(𝐩)=\frac{Nu}{48\pi ^2p^2}+O(p^1\mathrm{ln}p)$$
(101)
for small $`𝐩`$. One may check by power counting diagrams and sub-diagrams that $`\mathrm{\Sigma }(\mathrm{})=0`$. Then
$$_0=\frac{16}{3\pi ^2}u,_{\mathrm{}}=0,$$
(102)
and One may also check this answer by direct, brute-force calculation of all the diagrams associated with fig. 4e.
$$\mathrm{\Delta }\varphi ^2_\mathrm{e}=\frac{u}{6\pi ^4}.$$
(103)
We conclude by mentioning one technical subtlety, glossed over above, concerning absolute convergence. The integration corresponding to the substitution (97) in the leading order analysis is
$$\mathrm{\Delta }\varphi ^2_\mathrm{e}=_{\mathrm{𝐥𝐩}}\frac{_p}{l^4}\left[\frac{1}{|𝐥+𝐩|^2}\frac{1}{p^2}\right]_\pm .$$
(104)
In contrast to the analogous leading-order expression (30), this integral is not absolutely convergent in the infrared ($`lp0`$), though it is convergent in the UV. One might therefore worry about the ad hoc introduction of an IR regulator $`M`$ in the calculation of this graph. However, this worry is easily bypassed by rewriting
$$\mathrm{\Delta }\varphi ^2_\mathrm{e}=_{\mathrm{𝐥𝐩}}\frac{(_p_0)}{l^4}\left[\frac{1}{|𝐥+𝐩|^2}\frac{1}{p^2}\right]_\pm ,$$
(105)
which should be understood as regulated in the UV. The UV-regulated integral of the $`_0`$ factor vanishes. Eq. (105) is now convergent in the IR, but logarithmically divergent in the UV, just as the original leading-order integral (30) was. One can now follow through the same argument as in the leading-order case to introduce an IR regulator and then remove the UV divergence, where $`F_p`$ in the leading-order analysis in now replaced by $`_p_0`$. The result is still (98).
### Summary
Summing all the diagrams then yields the total NLO contribution:
$$\mathrm{\Delta }\varphi ^2_{\mathrm{NLO}}=\frac{u}{3\pi ^4}\left[\frac{7}{4}\zeta (3)\frac{3}{2}\frac{17}{240}\pi ^2+\frac{7\pi ^2}{48\sqrt{2}}\mathrm{ln}(1+\sqrt{2})\frac{2}{3}L(3,\chi _8)\right],$$
(106)
with $`L(3,\chi _8)`$ given by (87). Combining with the leading-order result (37),
$$\frac{\mathrm{\Delta }\varphi ^2_{\mathrm{NLO}}}{\mathrm{\Delta }\varphi ^2_{\mathrm{LO}}}=\frac{0.527198}{N},$$
(107)
which is the relative NLO correction we presented for $`T_\mathrm{c}`$ in (5).
## ACKNOWLEDGMENTS
We are indebted to Eric Braaten for suggesting this project. We are also grateful to Tim Newman and Genya Kolemeisky for useful discussions. This work was supported, in part, by the U.S. Department of Energy under Grant No. DE-FG02-97ER41027.
## A Basic Integrals
Let’s begin with the integral $`J_2(𝐩;M_1,M_2)`$. This is quite easy to do by standard methods (for example, by introducing a Feynman parameter), and gives
$`J_2(𝐩;M_1,M_2)`$ $``$ $`{\displaystyle _𝐥}{\displaystyle \frac{1}{(l^2+M_1^2)(|𝐥+𝐩|^2+M_2^2)}}`$ (A1)
$`=`$ $`{\displaystyle \frac{1}{8\pi p}}\mathrm{Cos}^1\left({\displaystyle \frac{(M_1+M_2)^2p^2}{(M_1+M_2)^2+p^2}}\right)`$ (A2)
$`=`$ $`{\displaystyle \frac{1}{8p}}{\displaystyle \frac{(M_1+M_2)}{4\pi p^2}}+{\displaystyle \frac{(M_1+M_2)^3}{12\pi p^4}}+O(M^5).`$ (A3)
The integral $`J_1(𝐩;m)`$ of (49) is simply the special case $`J_1(𝐩;M)=J_2(𝐩;M,0)`$.
The integral $`I_2(𝐩,𝐪;M_1,M_2)`$ of (68) can be related to $`J_2(𝐩,𝐪;M_1,M_2)`$ by generalizing a trick presented in ref. . The idea is to change integration variables from $`𝐥`$ to its conformal inversion $`\stackrel{~}{𝐥}𝐥/l^2`$. The integration measure changes as $`(2\pi )^3d^3l=(2\pi )^3\stackrel{~}{l}^6d^3\stackrel{~}{l}`$. Propagators can be written in terms of the new variable $`\stackrel{~}{𝐥}`$ as
$`{\displaystyle \frac{1}{l^2}}`$ $`=`$ $`\stackrel{~}{l}^2,`$ (A4)
$`{\displaystyle \frac{1}{|𝐥+𝐩|^2+M^2}}`$ $`=`$ $`{\displaystyle \frac{\stackrel{~}{l}^2}{(p^2+m^2)[|\stackrel{~}{𝐥}+\stackrel{~}{𝐏}|^2+\stackrel{~}{M}_p^2]}},`$ (A5)
where
$$\stackrel{~}{𝐏}\frac{𝐩}{p^2+M^2},\stackrel{~}{M}_p\frac{M}{p^2+M^2}.$$
(A6)
Making this change of variables,
$`I_2(𝐩,𝐪;M_1,M_2)`$ $``$ $`{\displaystyle \frac{1}{(p^2+M_1^2)(q^2+M_2^2)}}{\displaystyle _{\stackrel{~}{l}}}\left[\left|\stackrel{~}{𝐥}+{\displaystyle \frac{𝐩}{p^2+M_1^2}}\right|^2+\left({\displaystyle \frac{M_1}{p^2+M_1^2}}\right)^2\right]^1`$ (A8)
$`\times \left[\left|\stackrel{~}{𝐥}+{\displaystyle \frac{𝐪}{q^2+M_2^2}}\right|^2+\left({\displaystyle \frac{M_2}{q^2+M_2^2}}\right)^2\right]^1`$
$`=`$ $`{\displaystyle \frac{1}{(p^2+M_1^2)(q^2+M_2^2)}}J_2[{\displaystyle \frac{𝐩}{p^2+M_1^2}}{\displaystyle \frac{𝐪}{q^2+M_2^2}};{\displaystyle \frac{M_1}{p^2+M_1^2}},{\displaystyle \frac{M_2}{q^2+M_2^2}}].`$ (A9)
For the application of this paper, the relevant terms in the small $`M_1,M_2`$ expansion of $`I_2`$ are
$`I_2(𝐩,𝐪;M_1,M_2)`$ $`=`$ $`{\displaystyle \frac{1}{8pq|𝐩𝐪|}}{\displaystyle \frac{1}{4\pi |𝐩𝐪|^2}}\left({\displaystyle \frac{M_1}{p^2}}+{\displaystyle \frac{M_2}{q^2}}\right)+O(M_1^2)+O(M_2^2)`$ (A12)
$`+{\displaystyle \frac{M_1M_2}{4\pi |𝐩𝐪|^4}}\left({\displaystyle \frac{M_1}{q^2}}+{\displaystyle \frac{M_2}{p^2}}\right)+O(M_1^4)+O(M_2^4)`$
$`+M_1^2M_2^2{\displaystyle \frac{[3pq2(p^2+q^2)(\widehat{𝐩}\widehat{𝐪})+pq(\widehat{𝐩}\widehat{𝐪})^2]}{8p^2q^2|𝐩𝐪|^5}}+O(M^5).`$
The integral $`I_1`$ of (48) is related by $`I_1(𝐩,𝐪;M)=I_2(𝐩,𝐩𝐪;M,0)`$ and gives
$$I_1(𝐩,𝐪;M)=\frac{\mathrm{Cos}^1(2\omega _{\mathrm{𝐩𝐪}}^21)}{8\pi M|𝐩𝐪|^2\sqrt{\omega _{\mathrm{𝐩𝐪}}^21}},$$
(A13)
where
$$\omega _{\mathrm{𝐩𝐪}}\frac{M|𝐩𝐪|}{\sqrt{(p^2+M^2)(q^2+M^2)}}.$$
(A14)
The small $`M`$ expansion is given in (52).
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# 1 Introduction
## 1 Introduction
In 1978, Lepowsky and Wilson introduced the idea of free field representation in the theory of affine Lie algebras. In the first paper , a representation of the affine Lie algebra $`\widehat{𝔰𝔩}_2`$ was found for the special case of level one. Soon later this construction was extended to more general cases by introducing an associative algebra called the $`𝒵`$-algebra , and a connection between the Rogers-Ramanujan identities and the affine Lie algebras was uncovered. This study, however, was restricted to the case of the level $`k`$ being a non-negative integer, because their idea of representing the $`𝒵`$-algebra was to use $`k`$-copies of Heisenberg algebras.
We now know that the Wakimoto construction affords a way of realizing representations of an arbitrary level $`k`$. The Wakimoto representation plays an important role, both as a powerful computational tool in mathematical physics, as well as for theoretical studies in pure representation theory. There is, however, one basic difference between the Wakimoto and Lepowsky-Wilson’s original constructions. While the former is based on the homogeneous Heisenberg subalgebra, the latter is based on the principal Heisenberg subalgebra. Though strange as it may seem, we have not seen in the literature a construction which respects the principal gradation and works for an arbitrary level.
In this note we revisit this problem. We observe here a simple but peculiar fact that the $`𝒵`$-algebra arises as a certain specialization of the deformed Virasoro algebra . The representation of $`\widehat{𝔰𝔩}_2`$ mentioned above is obtained as an immediate consequence of this observation. We also touch upon the screening currents, vertex operators and a connection to the elliptic Knizhnik-Zamolodchikov (KZ) equation studied by Etingof .
## 2 Free field realization of $`\widehat{𝔰𝔩}_2`$
### 2.1 $`\widehat{𝔰𝔩}_2`$ in the principal picture
In order to fix the notation, let us recall the realization of $`\widehat{𝔰𝔩}_2`$ in the principal picture. Let $`e,f,h`$ be the standard generators of $`𝔤=𝔰𝔩_2`$ with $`[h,e]=2e`$, $`[h,f]=2f`$, $`[e,f]=h`$. The invariant bilinear form on $`𝔰𝔩_2`$ is chosen as $`(e,f)=1,(h,h)=2`$. Let $`𝔤_0=h`$, $`𝔤_1=ef`$. The affine Lie algebra $`\widehat{𝔰𝔩}_2`$ is realized as a vector space
$`\widehat{𝔰𝔩}_2=𝔤_0[t^2,t^2]𝔤_1t[t^2,t^2]c\rho ,`$
endowed with the Lie bracket
$`[at^m,bt^n]=[a,b]t^{m+n}+{\displaystyle \frac{1}{2}}(a,b)mc\delta _{m+n,0},`$
$`c:\text{ central},[\rho ,at^m]=mat^m,`$
where $`a,b𝔤`$. We set
$`\beta _n=(e+f)t^n(n\text{ odd}),x_n=\{\begin{array}{cc}ht^n\hfill & \text{(}n\text{ even}\text{),}\hfill \\ (e+f)t^n\hfill & \text{(}n\text{ odd}\text{).}\hfill \end{array}`$
In terms of the generating series (currents)
$`\beta (\zeta )={\displaystyle \underset{n:\mathrm{odd}}{}}\beta _n\zeta ^n,x(\zeta )={\displaystyle \underset{n}{}}x_n\zeta ^n,`$
the commutation relations read as
$`[\beta (\zeta ),\beta (\xi )]={\displaystyle \frac{c}{2}}\left((D\delta )\left({\displaystyle \frac{\xi }{\zeta }}\right)(D\delta )\left({\displaystyle \frac{\xi }{\zeta }}\right)\right),`$ (1)
$`[\beta (\zeta ),x(\xi )]=\left(\delta \left({\displaystyle \frac{\xi }{\zeta }}\right)\delta \left({\displaystyle \frac{\xi }{\zeta }}\right)\right)x(\xi ),`$ (2)
$`[x(\zeta ),x(\xi )]=2\delta \left({\displaystyle \frac{\xi }{\zeta }}\right)\beta (\xi )+c(D\delta )\left({\displaystyle \frac{\xi }{\zeta }}\right).`$ (3)
Here $`\delta (\zeta )=_m\zeta ^m`$ and $`D=D_\zeta `$ stands for $`\zeta \frac{d}{d\zeta }`$.
In the sequel we fix a complex number $`k2`$ (the level), and focus attention to representations of $`\widehat{𝔰𝔩}_2`$ on which the central element $`c`$ acts as $`k`$ times the identity. Our aim is to find a free field realization of the relations (1)-(3). For that purpose let us introduce three kinds of bosonic fields
$`\varphi _1(\zeta )={\displaystyle \underset{\genfrac{}{}{0pt}{}{n:\mathrm{even}}{n0}}{}}{\displaystyle \frac{\varphi _{1,n}}{n}}\zeta ^n+\varphi _{1,0}\mathrm{log}\zeta +Q,`$
$`\varphi _i(\zeta )={\displaystyle \underset{n:\mathrm{odd}}{}}{\displaystyle \frac{\varphi _{i,n}}{n}}\zeta ^n(i=0,2).`$
The fields $`\varphi _0(\zeta )`$, $`\varphi _2(\zeta )`$ are odd, $`\varphi _i(\zeta )=\varphi _i(\zeta )`$, while the derivative of $`\varphi _1(\zeta )`$ is even, $`(D\varphi _1)(\zeta )=(D\varphi _1)(\zeta )`$. We set the commutation relations for their Fourier modes as
$`[\varphi _{1,m},\varphi _{1,n}]=4(k+2)m\delta _{m+n,0},[\varphi _{1,0},Q]=4(k+2),`$
$`[\varphi _{0,m},\varphi _{0,n}]=4km\delta _{m+n,0},`$
$`[\varphi _{2,m},\varphi _{2,n}]=4km\delta _{m+n,0},`$
all other commutators being $`0`$. For $`j`$, we denote by $`_{j,k}`$ the Fock space for three bosons
$`_{j,k}=[\varphi _{0,n},\varphi _{2,n}(n=1,3,5,\mathrm{}),\varphi _{1,n}(n=2,4,6,\mathrm{})]|j,k`$
generated on the Fock vacuum $`|j,k`$,
$`\varphi _{i,n}|j,k=0(n>0,i=0,1,2),\varphi _{1,0}|j,k=2j|j,k,e^{\frac{Q}{2(k+2)}}|j,k=|j+1,k.`$
We denote by $`d:_{j,k}_{j,k}`$ the grading operator
$`[d,\varphi _{i,n}]=n\varphi _{i,n},[d,Q]=\varphi _{1,0},d|j,k={\displaystyle \frac{2j^2+k}{4(k+2)}}|j,k.`$
We adopt the conventional ‘normal ordering rule’ and the ‘normal ordering symbol’ $`:\mathrm{}:`$ for our bosonic fields. For example,
$`:e^{\varphi _1(\zeta )}:=e^{_{n1}\frac{\varphi _{1,2n}}{2n}\zeta ^{2n}}e^{_{n1}\frac{\varphi _{1,2n}}{2n}\zeta ^{2n}}e^Q\zeta ^{\varphi _{1,0}}.`$
We can now state the main result of this note.
###### Proposition 2.1
Let $`j,k`$ be complex numbers with $`k0,2`$. Then the following gives a level $`k`$ representation of $`\widehat{𝔰𝔩}_2`$ on the Fock space $`_{j,k}`$:
$`\beta (\zeta )={\displaystyle \frac{1}{2}}D\varphi _0(\zeta ),`$ (4)
$`x(\zeta )={\displaystyle \frac{1}{2}}:\left(D\varphi _1(\zeta )+D\varphi _2(\zeta )\right)e^{\frac{\varphi _2(\zeta )}{k}+\frac{\varphi _0(\zeta )}{k}}:,`$ (5)
$`c=k,\rho =d.`$ (6)
This representation is highest weight in the sense that
$`\beta _n|j,k=0,x_n|j,k=0\mathrm{for}n>0,`$
$`x_0|j,k=j|j,k.`$
The highest weight is $`\frac{k}{2}(\mathrm{\Lambda }_1+\mathrm{\Lambda }_0)+j(\mathrm{\Lambda }_1\mathrm{\Lambda }_0)`$ where $`\mathrm{\Lambda }_0,\mathrm{\Lambda }_1`$ are the fundamental weights of $`\widehat{𝔰𝔩}_2`$. The character of this representation (counted according to the principal gradation)
$`\text{tr}_{_{j,k}}(q^\rho )=q^{\frac{2j^2+k}{4(k+2)}}{\displaystyle \frac{1}{(q;q)_{\mathrm{}}(q;q^2)_{\mathrm{}}}},`$
is the same as that of the Verma module of $`\widehat{𝔰𝔩}_2`$. Here $`(z;p)_{\mathrm{}}=_{n=0}^{\mathrm{}}(1p^nz)`$. In particular, the above representation is irreducible for generic values of $`j,k`$.
### 2.2 Connection with the deformed Virasoro algebra
Because of the commutation relation (1), the current $`\beta (\zeta )`$ can tautologically be identified with the bosonic field $`(1/2)D\varphi _0(\zeta )`$. As Lepowsky and Wilson have shown, the other current $`x(\zeta )`$ can be realized as
$`x(\zeta )=z(\zeta ):e^{\frac{\varphi _0(\zeta )}{k}}:,`$ (7)
provided $`z(\zeta )=_nz_n\zeta ^n`$ commutes with $`\varphi _0(\zeta )`$ and satisfies the relation of the $`𝒵`$-algebra
$`\left({\displaystyle \frac{\zeta _1\zeta _2}{\zeta _1+\zeta _2}}\right)^{2/k}z(\zeta _1)z(\zeta _2)=\left({\displaystyle \frac{\zeta _2\zeta _1}{\zeta _2+\zeta _1}}\right)^{2/k}z(\zeta _2)z(\zeta _1)+k(D\delta )\left({\displaystyle \frac{\zeta _2}{\zeta _1}}\right).`$ (8)
As we explain below, this algebra is related with the deformed Virasoro algebra (DVA).
The DVA is an associative algebra generated by $`T_n`$ ($`n`$) (see ). In terms of $`T(\zeta )=_nT_n\zeta ^n`$ the defining relations read
$`f(\zeta _2/\zeta _1)T(\zeta _1)T(\zeta _2)T(\zeta _2)T(\zeta _1)f(\zeta _1/\zeta _2)={\displaystyle \frac{(1q)(1t^1)}{1p}}\left[\delta \left({\displaystyle \frac{p\zeta _2}{\zeta _1}}\right)\delta \left({\displaystyle \frac{p^1\zeta _2}{\zeta _1}}\right)\right],`$ (9)
where $`q`$ and $`t`$ are parameters, $`p=q/t`$, and
$`f(\zeta )=\mathrm{exp}\left\{{\displaystyle \underset{n1}{}}{\displaystyle \frac{1}{n}}{\displaystyle \frac{(1q^n)(1t^n)}{1+p^n}}\zeta ^n\right\}.`$
Now set
$`q=e^h,t=q^{\frac{k+2}{2}},`$
and consider the limit $`h0`$. Suppose that the expansion
$`T(\zeta )=0+hT^{(1)}(\zeta )+h^2T^{(2)}(\zeta )+\mathrm{}`$ (10)
takes place. Under this assumption we find that, at the second order in $`h`$, the relation (9) reduces to the $`𝒵`$-algebra relation (8) with the identification $`\sqrt{1}T^{(1)}(\zeta )=z(\zeta )`$. We have verified the expansion (10) using the known bosonization for DVA . This leads to the formula
$`z(\zeta )={\displaystyle \frac{1}{2}}:(D\varphi _1(\zeta )+D\varphi _2(\zeta ))e^{\frac{\varphi _2(\zeta )}{k}}:.`$
The formula (5) for $`x(\zeta )`$ follows from this and (7). In fact this bosonization formula for the $`𝒵`$-algebra was first obtained by a guesswork based on the one for the ordinary Wakimoto realization which respects the homogeneous gradation . The relationship between the $`𝒵`$-algebra and the deformed Virasoro algebra was noticed only afterwards.
Thus, somewhat unexpectedly, one can view the DVA with the parameters $`q`$ and $`t=q^{\frac{k+2}{2}}`$ as a quantum deformation of Lepowsky-Wilson’s $`𝒵`$-algebra at level $`k`$.
The DVA admits two kinds of “screening operators” $`S_\pm (\xi )`$ commuting with $`T(\zeta )`$ up to a total difference. In the next section we shall show that one of them in the limit, $`S(\xi )`$, becomes the screening operator for the present bosonization of $`\widehat{𝔰𝔩}_2`$. The limit of the other one $`\eta (\xi )`$, after a modification by zero-mode, plays a role of the operator “$`B`$” which appears in the construction of the elliptic KZ equation .
Remark. It seems likely that a similar construction persists in the case of the deformed $`W_N`$ algebra . Let $`\omega `$ be a primitive $`N`$-th root of unity, and set $`t=\omega q^{\frac{k+N}{N}}`$. Suppose that in the limit $`q1`$ we have the expansion of the $`W`$-currents
$`W_i(\zeta )=0+hW_i^{(1)}(\zeta )+h^2W_i^{(2)}(\zeta )+\mathrm{}.`$ (11)
We have checked for $`N=3`$ (and partially for all $`N`$) that the $`W_i^{(1)}`$ ($`i=1,\mathrm{},N1`$) then satisfy the relations of the $`𝒵`$-algebra for $`\widehat{𝔰𝔩}_N`$. However, for $`N3`$ we have not been able to show (11), which seems to hold not in the free field realization but only at the level of correlation functions.
## 3 Vertex operators and KZ equations
### 3.1 Screening current $`S(\xi )`$
Our representation of $`\widehat{𝔰𝔩}_2`$ on the Fock space $`_{j,k}`$ may become reducible for some specific values of the parameters $`j,k`$. In the usual Wakimoto construction the object which controls this phenomenom is the screening current. Let us consider its analog. Define $`S(\zeta ):_{j,k}_{j2,k}`$ by
$`S(\zeta )={\displaystyle \frac{1}{2}}\zeta ^{\frac{2}{k+2}}:D\varphi _2(\zeta )e^{\frac{\varphi _1(\zeta )}{k+2}}:.`$
It enjoys the expected properties
$`[\beta (\zeta ),S(\xi )]=0,`$
$`[x(\zeta ),S(\xi )]`$
$`={\displaystyle \frac{k+2}{2}}D_\xi (\delta \left({\displaystyle \frac{\xi }{\zeta }}\right)\xi ^{\frac{2}{k+2}}:e^{\frac{\varphi _1(\xi )}{k+2}+\frac{\varphi _2(\xi )}{k}+\frac{\varphi _0(\xi )}{k}}:+\delta ({\displaystyle \frac{\xi }{\zeta }})\xi ^{\frac{2}{k+2}}:e^{\frac{\varphi _1(\xi )}{k+2}\frac{\varphi _2(\xi )}{k}\frac{\varphi _0(\xi )}{k}}:).`$
These equations imply that the screening charge
$`Q^m={\displaystyle \mathrm{}\frac{d\xi _1}{\xi _1}\mathrm{}\frac{d\xi _m}{\xi _m}S(\xi _1)\mathrm{}S(\xi _m)}:_{j,k}_{j2m,k},`$
commutes with the action of $`\beta _n,x_n`$, if closed contours for all the $`\xi _i`$’s exist.
### 3.2 Vertex Operators
Let us consider the vertex operator (VO) associated with the two-dimensional representation. Set $`V=u_+u_{}`$, and let $`V^{}(\zeta )=V[\zeta ,\zeta ^1]`$ be the $`\widehat{𝔰𝔩}_2`$-module given by
$`\beta _n(u_\pm \zeta ^m)=u_\pm \zeta ^{m+n},x_n(u_\pm \zeta ^m)=\{\begin{array}{cc}u_{}\zeta ^{m+n}\hfill & \text{(}n\text{ even),}\hfill \\ u_{}\zeta ^{m+n}\hfill & \text{(}n\text{ odd).}\hfill \end{array}`$
We define $`\mathrm{\Phi }(\zeta ):_{j,k}_{j+1,k}V^{}(\zeta )`$ by
$`\mathrm{\Phi }(\zeta )v=\mathrm{\Phi }_+(\zeta )vu_++\mathrm{\Phi }_{}(\zeta )vu_{},`$
$`\mathrm{\Phi }_\pm (\zeta )=\zeta ^{\frac{1}{2(k+2)}}:e^{\frac{\varphi _1(\zeta )}{2(k+2)}\pm \frac{\varphi _2(\zeta )}{2k}\pm \frac{\varphi _0(\zeta )}{2k}}:.`$
Then we have the intertwining property
$`(x\mathrm{id}+\mathrm{id}x)\mathrm{\Phi }(\zeta )=\mathrm{\Phi }(\zeta )x({}_{}{}^{}x\widehat{𝔰𝔩}_2).`$
Remark.$`V^{}(\zeta )`$ contains a proper submodule $`W=\text{span}\{u_+\zeta ^mu_{}(\zeta )^mm\}`$. Its quotient $`V(\zeta )=V^{}(\zeta )/W`$ is isomorphic to the irreducible evaluation module of $`\widehat{𝔰𝔩}_2`$ associated with $`V`$. The above VO naturally gives rise to the intertwiner $`_{j,k}_{j+1,k}V(\zeta )`$.
### 3.3 Elliptic Knizhnik-Zamolodchikov equation
As usual, the highest-to-highest matrix elements of VO’s satisfy the KZ equation. Let us consider the elliptic KZ equation studied by Etingof .
Let $`M_{j,k}`$ denote the Verma module over $`\widehat{𝔰𝔩}_2`$ with highest weight $`\frac{k}{2}(\mathrm{\Lambda }_1+\mathrm{\Lambda }_0)+j(\mathrm{\Lambda }_1\mathrm{\Lambda }_0)`$ and highest weight vector $`v_{j,k}`$. Denote by $`\mathrm{\Psi }(\zeta ):M_{j,k}M_{j^{},k}V(\zeta )`$ the intertwining operator. Let $`B:M_{j,k}M_{j,k}`$ denote the linear isomorphism characterized by
$`Bv_{j,k}=v_{j,k},`$
$`B\beta (\zeta )=\beta (\zeta )B,Bx(\zeta )=x(\zeta )B.`$
It was shown in that the trace function
$`F(\zeta _1,\mathrm{},\zeta _n)=\mathrm{tr}_{M_{j,k}}\left(\mathrm{\Psi }(\zeta _1)\mathrm{}\mathrm{\Psi }(\zeta _n)Bq^\rho \right)`$
satisfies the following elliptic KZ equation
$`(k+2)D_{\zeta _i}F(\zeta _1,\mathrm{},\zeta _n)={\displaystyle \underset{j(i)}{}}r^{ij}(\zeta _i/\zeta _j)F(\zeta _1,\mathrm{},\zeta _n),`$
where $`r(z)`$ denotes the elliptic classical $`r`$ matrix
$`r(z)=\left({\displaystyle \underset{l}{}}{\displaystyle \frac{q^lz}{(q^lz)^21}}\right)(e+f)(e+f)\left({\displaystyle \underset{l}{}}{\displaystyle \frac{(q)^lz}{(q^lz)^21}}\right)(ef)(ef)`$
$`+\left({\displaystyle \underset{l}{}}()^l{\displaystyle \frac{(q^lz)^2+1}{(q^lz)^21}}\right){\displaystyle \frac{1}{2}}hh,`$
and $`r^{ij}(z)`$ signifies $`r(z)`$ acting nontrivially on the $`(i,j)`$-th components.
Let us look for a bosonic realization of the operator $`B`$. In the usual Wakimoto construction, we have the ‘$`\xi `$-$`\eta `$ system’ after bosonizing the ‘$`\beta `$-$`\gamma `$ ghost system’. An analog of $`\eta `$ in the present case is
$`\eta (\zeta )={\displaystyle \underset{n}{}}\eta _n\zeta ^n=\zeta ^{\frac{k+2}{2}}:e^{\frac{\varphi _1(\zeta )}{2}+\frac{\varphi _2(\zeta )}{2}}:,`$
which satisfies
$`[\beta (\zeta ),\eta (\xi )]=0,`$
$`[x(\zeta ),\eta (\xi )]_+=2\xi _\xi (\xi ^{\frac{k+2}{2}}\delta \left({\displaystyle \frac{\zeta }{\xi }}\right):e^{\frac{\varphi _1(\xi )}{2}+\frac{(k+2)\varphi _2(\xi )}{2k}+\frac{\varphi _0(\xi )}{k}}:),`$
where $`[A,B]_+=AB+BA`$. From the above, we find that the zero-th Fourier mode $`\eta _0:_{j,k}_{j+k+2,k}`$ satisfies
$`\eta _0\beta (\zeta )=\beta (\zeta )\eta _0,\eta _0x(\zeta )=x(\zeta )\eta _0.`$
Therefore, if $`j+k+2=j+2m`$ for some $`m_0`$, the combination of $`\eta _0`$ and the screening operator $`Q^m`$ implements $`B`$.
In particular, when $`m=0`$ and $`2j_{>0}`$, the screening operators do not appear and the traces can be expressed without using integrals. For example, when $`(j,k)=(1/2,3)`$ and $`(j,k)=(1,4)`$ we find
$`\mathrm{tr}_{_{1/2,3}}\left(\mathrm{\Phi }_\pm (\zeta _1)\eta _0q^\rho \right)=q^{\frac{5}{8}}(q^2;q^2)_{\mathrm{}}^{\frac{3}{2}},`$
$`\mathrm{tr}_{_{1,4}}\left(\mathrm{\Phi }_{\epsilon _1}(\zeta _1)\mathrm{\Phi }_{\epsilon _2}(\zeta _2)\eta _0q^\rho \right)`$
$`=q^{\frac{1}{4}}(q^2;q^2)_{\mathrm{}}^{\frac{9}{4}}\zeta ^{1/4}\mathrm{\Theta }_{q^2}(\zeta ^2)^{\frac{1}{4}}\mathrm{\Theta }_{q^2}(\epsilon _1\epsilon _2q\zeta )`$
$`=q^{\frac{1}{2}}\left\{{\displaystyle \frac{\sqrt{1}\epsilon _1\epsilon _2}{8\pi ^3}}\mathrm{}^{}\left({\displaystyle \frac{\mathrm{ln}(\epsilon _1\epsilon _2q\zeta )}{2\pi \sqrt{1}}}\right|1,2\tau )\right\}^{\frac{1}{4}},`$
where $`\zeta =\zeta _2/\zeta _1`$, $`q=\mathrm{exp}(2\pi \sqrt{1}\tau )`$, and $`\mathrm{}(z|\omega ,\omega ^{})`$ denotes the Weierstrass elliptic function with fundamental periods $`\omega ,\omega ^{}`$. These formulas have been discussed in (there are minor errors in Section 5, equation (5.2) of ).
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# Acknowledgments
## Acknowledgments
SD and AG thank Abhay Ashtekar for discussions and various suggestions. The works of SD and AG were supported by the National Science Foundation grant PHY95-14240 and the Eberly Research Funds of Penn State. PM thanks the Theory Division of CERN, where the collaboration started.
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