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warning/0007/cond-mat0007103.html	ar5iv	text	# Fermi and non-Fermi Liquid Behavior in Quantum Impurity Systems: Conserving Slave Boson Theory ## 0.1 Introduction Highly correlated electron systems are characterized by a strong repulsion between electrons on the same lattice site, effectively restricting the dynamics to the Fock subspace of states without double occupancy of sites. The prototype model for such systems is the Anderson impurity model, which consists of an electron in a localized level $`\epsilon _d<0`$ (called d-level in the following) with on-site repulsion $`U`$, hybridizing via a transition matrix element $`V`$ with one or several degenerate conduction electron bands or channels . Depending on the number of channels $`M`$, the model exhibits the single- or the multi-channel Kondo effect, where at temperatures $`T`$ below the Kondo temperature $`T_K`$ the local electron spin is screened ($`M=1`$) or overscreened ($`M2`$) by the conduction electrons, leading to Fermi liquid (FL) or to non-Fermi liquid (NFL) behavior with characteristic low-temperature singularities, respectively. As perhaps the simplest model to investigate the salient features of correlations induced by short-range repulsion, the Anderson model plays a central role for the description of strongly correlated electron systems: In the limit of large spatial dimensions strongly correlated lattice systems reduce in general to a single Anderson impurity hybridizing with a continuum of conduction electron states whose properties are determined from a self-consistency condition imposed by the translational invariance of the system . Quantum impurity models have received further interest due to their relevance for mesoscopic systems like single electron transistors or quantum point contacts . Nonlinear conductance anomalies observed in the latter systems have provided one of the strongest cases for the physical realization of the two-channel Kondo effect generated by two-level systems with electron assisted tunneling. The above-mentioned systems call for the development of accurate and flexible theoretical methods, applicable to situations where exact solution methods are not available. We here present a general, well-controlled auxiliary boson technique which correctly describes the FL as well as the NFL case of the generalized SU(N)$`\times `$SU(M) Anderson impurity model. As a standard diagram technique it has the potential to be generalized for correlated lattice problems as well as for non-equilibrium situations in mesoscopic systems. In section 2 we describe several exact properties of the auxiliary particle representation, while the conserving slave boson theory is developed and evaluated in section 3. ## 0.2 Exact Auxiliary Particle Representation ### 0.2.1 The SU(N)$`\times `$SU(M) Anderson Impurity Model The auxiliary or slave boson method is a powerful tool to implement the effective restriction to the sector of Fock space with no double occupancy imposed by a large on-site repulsion $`U`$. The creation operator for an electron with spin $`\sigma `$ in the $`d`$-level is written in terms of fermionic operators $`f_\sigma `$ and bosonic operators $`b`$ as $`d_\sigma ^{}=f_\sigma ^{}b`$. This representation is exact, if the constraint that the total number operator of auxiliary fermions $`f_\sigma `$ and bosons $`b`$ is equal to unity is obeyed. $`f_\sigma ^{}`$ and $`b^{}`$ may be envisaged as creating the three allowed states of the impurity: singly occupied with spin $`\sigma `$ or empty. In view of the possibility of both FL and NFL behavior in quantum impurity systems mentioned in the introduction it is useful to introduce $`M`$ degenerate channels for the conduction electron operators $`c_{\sigma \mu }^{}`$, labeled $`\mu =1,2,\mathrm{},M`$, in such a way that in the limit of impurity occupation number $`n_d1`$ (Kondo limit) the $`M`$-channel Kondo model is recovered, i.e. the model obeys an SU(M) channel symmetry. The slave bosons then form an SU(M) multiplet $`b_{\overline{\mu }}`$ which transforms according to the conjugate representation of SU(M), so that $`\mu `$ is a conserved quantum number. Generalizing, in addition, to arbitrary spin degeneracy $`N`$, $`\sigma =1,2,\mathrm{},N`$, one obtains the SU(N)$`\times `$SU(M) Anderson impurity model in slave boson representation $$H=\underset{\stackrel{}{k},\sigma ,\mu }{}\epsilon _\stackrel{}{k}c_{\stackrel{}{k}\mu \sigma }^{}c_{\stackrel{}{k}\mu \sigma }+E_d\underset{\sigma }{}f_\sigma ^{}f_\sigma +V\underset{\stackrel{}{k},\sigma ,\mu }{}(c_{\stackrel{}{k}\mu \sigma }^{}b_{\overline{\mu }}^{}f_\sigma +h.c.),$$ (0.1) where the local operator constraint $`\widehat{Q}_\sigma f_\sigma ^{}f_\sigma +_\mu b_{\overline{\mu }}^{}b_{\overline{\mu }}=1`$ must be fulfilled at all times. ### 0.2.2 Gauge Symmetry and Exact Projection onto the Physical Fock Space The system described by the auxiliary particle Hamiltonian (0.1) is invariant under simultaneous, local $`U(1)`$ gauge transformations, $`f_\sigma f_\sigma \mathrm{e}^{i\varphi (\tau )}`$, $`b_{\overline{\mu }}b_{\overline{\mu }}\mathrm{e}^{i\varphi (\tau )}`$, with $`\varphi (\tau )`$ an arbitrary, time dependent phase. While the gauge symmetry guarantees the conservation of the local, integer charge $`Q`$, it does not single out any particular $`Q`$, like $`Q=1`$. In order to effect the projection onto the $`Q=1`$ sector of Fock space, one may use the following procedure : Consider first the grand-canonical ensemble with respect to $`Q`$ and the associated chemical potential $`\lambda `$. The expectation value in the $`Q=1`$ subspace of any physical operator $`\widehat{A}`$ acting on the impurity states is then obtained as $$\widehat{A}=\underset{\lambda \mathrm{}}{lim}\frac{\frac{}{\zeta }\text{tr}\left[\widehat{A}e^{\beta (H+\lambda Q)}\right]_G}{\frac{}{\zeta }\text{tr}\left[e^{\beta (H+\lambda Q)}\right]_G}=\underset{\lambda \mathrm{}}{lim}\frac{\widehat{A}_G}{Q_G},$$ (0.2) where the index $`G`$ denotes the grand canonical ensemble and $`\zeta `$ is the fugacity $`\zeta =\mathrm{e}^{\beta \lambda }`$. In the second equality of Eq. (0.2) we have used the fact that any physical operator $`\widehat{A}`$ acting on the impurity is composed of the impurity electron operators $`d_\sigma `$, $`d_\sigma ^{}`$, and thus annihilates the states in the $`Q=0`$ sector, $`\widehat{A}\|Q=0=0`$. It is obvious that the grand-canonical expectation value involved in Eq. (0.2) may be factorized into auxiliary particle propagators using Wick’s theorem, thus allowing for the application of standard diagrammatic techniques. It is important to note that, in general, $`\lambda `$ plays the role of a time dependent gauge field. In Eq. (0.2) a time independent gauge for $`\lambda `$ has been chosen. In this way, the projection is only performed at one instant of time, explicitly exploiting the conservation of the local charge $`Q`$. Thus, choosing the time independent gauge means that in the subsequent development of the theory, the $`Q`$ conservation must be implemented exactly. This is achieved in a systematic way by means of conserving approximations , i.e. by deriving all self-energies and vertices by functional derivation from one common Luttinger-Ward functional $`\mathrm{\Phi }`$ of the fully renormalized Green’s functions, $$\mathrm{\Sigma }_{b,f,c}=\delta \mathrm{\Phi }\{G_b,G_f,G_c\}/\delta G_{b,f,c}.$$ (0.3) This amounts to calculating all quantities of the theory in a self-consistent way, but has the great advantage that gauge field fluctuations need not be considered. ### 0.2.3 Infrared Threshold Behavior of Auxilary Propagators The projection onto the physical subspace, Eq. (0.2), implies that the pseudofermion and slave boson Green’s functions $`G_f`$, $`G_b`$ are definied as the usual time-ordered, grand canonical expectation values of a pair of creation and annihilation operators, however evaluated in the limit $`\lambda \mathrm{}`$. It follows that the traces involved in $`G_f`$, $`G_b`$ are taken purely over the the $`Q=0`$ sector of Fock space, and thus the backward-in-time contribution to the auxiliary particle propagators vanishes. Consequently, the auxiliary particle propagators are formally identical to the core hole propagators appearing in the well-known X-ray problem , and the long-time behavior of $`G_f`$ ($`G_b`$) is determined by the orthogonality catastrophe of the overlap of the Fermi sea without impurity ($`Q=0`$) and the fully interacting conduction electron sea in the presence of a pseudofermion (slave boson) ($`Q=1`$). It may be shown that the auxiliary particle spectral functions have threshold behavior with vanishing spectral weight at $`T=0`$ for energies $`\omega `$ below a threshold $`E_o`$, and power law behavior above $`E_o`$, $`A_{f,b}(\omega )\mathrm{\Theta }(\omega E_o)\omega ^{\alpha _{f,b}}`$. For the single-channel Anderson model, which is known to have a FL ground state, the threshold exponents may be deduced from an analysis in terms of scattering phase shifts, using the Friedel sum rule, since in the spin screened FL state the impurity acts as a pure potential scatterer , $`\alpha _f={\displaystyle \frac{2n_dn_d^2}{N}},\alpha _b=1{\displaystyle \frac{n_d^2}{N}}(N1,M=1)`$ (0.4) These results have been confirmed by numerical renormalization group (NRG) calculations and by use of the Bethe ansatz solution in connection with boundary conformal field theory (CFT) . On the contrary, in the NFL case of the multi-channel Kondo model the threshold exponents have been deduced by a CFT solution as $`\alpha _f={\displaystyle \frac{M}{M+N}},\alpha _b={\displaystyle \frac{N}{M+N}}(N2,MN)`$ (0.5) Since the dependence of $`\alpha _f`$, $`\alpha _b`$ on the impurity occupation number $`n_d`$ shown above originates from pure potential scattering, it is characteristic for the FL case. The auxiliary particle threshold exponents are, therefore, indicators for FL or NFL behavior in quantum impurity models of the Anderson type. ## 0.3 Conserving Slave Particle T-Matrix Approximation ### 0.3.1 Non-Crossing Approximation (NCA) The conserving formulation discussed in section 2.2 precludes mean field approximations which break the $`U(1)`$ gauge symmetry, like slave boson mean field theory. Although the latter can in some cases successfully describe the low $`T`$ behavior of models with a FL ground state, it leads to a spurious phase transition at finite $`T`$ and, in particular, fails to describe NFL systems. Rather, the approximation should be generated from a Luttinger-Ward functional $`\mathrm{\Phi }`$. Using the hybridization $`V`$ as a small parameter, one may generate successively more complex approximations. The lowest order conserving approximation generated in this way is the Non-crossing Approximation (NCA) , defined by the first diagram in Fig. 0.2, labeled “NCA”. The NCA is successful in describing Anderson type models at temperatures above and around the Kondo temperature $`T_K`$, and even reproduces the threshold exponents Eq. (0.5) for the NFL case of the Anderson impurity model. However, it fails to describe the FL regime at low temperatures. This may be traced back to the failure to capture the spin-screened Kondo singlet ground state of the model, since coherent spin flip scattering is not included in NCA, as seen below. ### 0.3.2 Dominant Contributions at Low Energy In order to eliminate the shortcomings of the NCA mentioned above, we may use as a guiding principle to look for contributions to the vertex functions which renormalize the auxiliary particle threshold exponents to their correct values, since this is a necessary condition for the description of FL and NFL behavior, as discussed in section 2.3. As shown by power counting arguments , there are no corrections to the NCA exponents in any finite order of perturbation theory. Thus, any renormalization of the NCA exponents must be due to singularities arising from an infinite resummation of terms. In general, the existence of collective excitations leads to a singular behavior of the corresponding two–particle vertex function. In view of the tendency of Kondo systems to form a collective spin singlet state, we expect a singularity in the spin singlet channel of the pseudofermion–conduction electron vertex function. It is then natural to perform a partial resummation of those contributions which, at each order in the hybridization $`V`$, contain the maximum number of spin flip processes. This amounts to calculating the conduction electron–pseudofermion vertex function in the “ladder” or T-matrix approximation, $`T^{(cf)}`$, where the irreducible vertex is given by $`V^2G_b`$. The Bethe–Salpeter equation for $`T^{(cf)}`$ reads (Fig. 0.1), $`T_{\sigma \tau ,\sigma ^{}\tau ^{}}^{(cf)\mu }(i\omega _n,i\omega _n^{},i\mathrm{\Omega }_n)=`$ $`+`$ $`V^2G_{b\overline{\mu }}(i\omega _n+i\omega _n^{}i\mathrm{\Omega }_n)\delta _{\sigma \tau ^{}}\delta _{\tau \sigma ^{}}`$ $``$ $`V^2T{\displaystyle \underset{\omega _n^{\prime \prime }}{}}G_{b\overline{\mu }}(i\omega _n+i\omega _n^{\prime \prime }i\mathrm{\Omega }_n)\times `$ $`G_{f\sigma }(i\omega _n^{\prime \prime })G_{c\mu \tau }^0(i\mathrm{\Omega }_ni\omega _n^{\prime \prime })T_{\tau \sigma ,\sigma ^{}\tau ^{}}^{(cf)\mu }(i\omega _n^{\prime \prime },i\omega _n^{},i\mathrm{\Omega }_n),`$ where $`\sigma `$, $`\tau `$, $`\sigma ^{}`$, $`\tau ^{}`$ represent spin indices and $`\mu `$ a channel index. A similar integral equation holds for the charge fluctuation T-matrix $`T^{(cb)}`$; it is obtained from $`T^{(cf)}`$ by interchanging $`f_\sigma b_\mu `$ and $`c_{\sigma \mu }c_{\sigma \mu }^{}`$. Inserting NCA Green’s functions for the intermediate state propagators of Eq. (0.3.2), we find at low temperatures and in the Kondo regime $`(n_d\stackrel{>}{}0.7)`$ a pole of $`T^{(cf)}`$ in the singlet channel as a function of the center–of–mass (COM) frequency $`\mathrm{\Omega }`$, at a frequency which scales with the Kondo temperature, $`\mathrm{\Omega }=\mathrm{\Omega }_{cf}T_K`$. Similarly, the corresponding $`T`$-matrix $`T^{(cb)}`$ in the conduction electron–slave boson channel, evaluated within the analogous approximation, develops a pole at negative values of $`\mathrm{\Omega }`$ in the empty orbital regime $`(n_d\stackrel{<}{}0.3)`$. In the mixed valence regime ($`n_d0.5)`$ the poles in both $`T^{(cf)}`$ and $`T^{(cb)}`$ coexist. The appearance of poles in the two–particle vertex functions $`T^{(cf)}`$ and $`T^{(cb)}`$, which signals the formation of collective states, may be expected to influence the behavior of the system in a major way. ### 0.3.3 Self-consistent Formulation The approximation considered so far is not yet consistent: The spectral weight produced by the poles of $`T^{(cf)}`$ and $`T^{(cb)}`$ at negative frequencies $`\mathrm{\Omega }`$ is strictly prohibited in the limit $`T0`$ by the threshold property of auxiliary particle vertex functions (compare section 2.3). However, recall that a minimum requirement on the approximation used is the conservation of gauge symmetry. This requirement is not met when the integral kernel of the $`T`$-matrix equation is approximated by the NCA result. Rather, the approximation should be generated from the Luttinger-Ward functional shown in Fig. 0.2. It is defined as the infinite series of all vacuum skeleton diagrams which consist of a single ring of auxiliary particle propagators, where each conduction electron line spans at most two hybridization vertices. The first diagram of the infinite series of CTMA terms corresponds to NCA. By functional differentiation with respect to the conduction electron Green’s function and the pseudofermion or the slave boson propagator, respectively, the shown $`\mathrm{\Phi }`$ functional generates the ladder approximations $`T^{(cf)}`$, $`T^{(cb)}`$ for the total conduction electron-pseudofermion vertex function and for the total conduction electron-slave boson vertex function (Fig. 0.1). The auxiliary particle self-energies are obtained in the conserving scheme as the functional derivatives of $`\mathrm{\Phi }`$ with respect to $`G_f`$ or $`G_b`$, respectively (Eq. (0.3)), and are, in turn, nonlinear functionals of the full, renormalized auxiliary particle propagators. This defines a set of self-consistency equations, which we term conserving T-matrix approximation (CTMA). The CTMA is justified on formal grounds by a cancellation theorem for all diagrams not included in the generating functional $`\mathrm{\Phi }`$ : As seen in Fig. 0.2, $`\mathrm{\Phi }`$ includes all contributions where a conduction electron line spans up to two hybridization vertices; contributions not included contain at least one conduction electron arch spanning four vertices. These contributions cancel pairwise at arbitrary loop order to leading and subleading order in the external frequency, as illustrated in Fig. 0.3. ### 0.3.4 Results We have solved the CTMA equations numerically for a wide range of impurity occupation numbers $`n_d`$ from the Kondo to the empty impurity regime both for the single-channel and for the two-channel Anderson model down to temperatures of the order of at least $`10^2T_K`$. The solution of the CTMA equations forces the T-matrices to have vanishing spectral weight at negative COM frequencies $`\mathrm{\Omega }`$. Indeed, the numerical evaluation shows that the poles of $`T^{(cf)}`$ and $`T^{(cb)}`$ are shifted to $`\mathrm{\Omega }=0`$ by self-consistency, where they merge with the continuous spectral weight present for $`\mathrm{\Omega }>0`$, and thus renormalize the threshold exponents of the auxiliary spectral functions. For $`N=2`$, $`M=1`$ the threshold exponents $`\alpha _f`$, $`\alpha _b`$ extracted from the numerical solutions are shown in Fig. 0.4. In the Kondo limit of the multi–channel case ($`N2`$, $`M=2,4`$) the CTMA solutions are found not to alter the NCA values and reproduce the the correct threshold exponents, $`\alpha _f=M/(M+N)`$, $`\alpha _b=N/(M+N)`$. The good agreement of the CTMA exponents with their exact values over the complete range of $`n_d`$ for the single–channel model and in the Kondo regime of the multi–channel model may be taken as evidence that the T-matrix approximation correctly describes both the FL and the non–FL regimes of the SU(N)$`\times `$SU(M) Anderson model (N=2, M=1,2,4). The static spin susceptibility $`\chi `$ of the single- and of the two-channel Anderson model in the Kondo regime calculated within CTMA as the derivative of the magnetization $`M=\frac{1}{2}g\mu _Bn_fn_f`$ with respect to a magnetic field $`H`$ is shown in Fig. 0.5. Good quantitative agreement with exact solutions is found for $`N=2`$, $`M=1`$ (FL). For $`N=2`$, $`M=2`$ (NFL) CTMA correctly reproduces the exact logarithmic temperature dependence below the Kondo scale $`T_K`$. In contrast, the NCA solution recovers the logarithmic behavior only far below $`T_K`$. ## 0.4 Conclusion We have reviewed a novel technique to describe correlated quantum impurity systems with strong onsite repulsion, which is based on a conserving formulation of the auxiliary boson method. The conserving scheme allows to implement the conservation of the local charge $`Q`$ without taking into account time dependent fluctuations of the gauge field $`\lambda `$. By including the leading infrared singular contributions (spin flip and charge fluctuation processes), physical quantities, like the magnetic susceptibility, are correctly described both in the Fermi and in the non–Fermi liquid regime, over the complete temperature range, including the crossover to the correlated many–body state at the lowest temperatures. As a standard diagram technique this method has the potential to be applicable to problems of correlated systems on a lattice as well as to mesoscopic systems out of equilibrium via the Keldysh technique. We wish to thank S. Böcker, T.A. Costi, S. Kirchner, A. Rosch, A. Ruckenstein and Th. Schauerte for stimulating discussions. This work is supported by DFG through SFB 195 and by the Hochleistungsrechenzentrum Stuttgart.
warning/0007/cond-mat0007173.html	ar5iv	text	# Phonon anomalies due to strong electronic correlations in layered organic metals ## Abstract We show how the coupling between the phonons and electrons in a strongly correlated metal can result in phonon frequencies which have a non-monotonic temperature dependence. Dynamical mean-field theory is used to study the Hubbard-Holstein model that describes the $`\kappa `$-(BEDT-TTF)<sub>2</sub>X family of superconducting molecular crystals. The crossover with increasing temperature from a Fermi liquid to a bad metal produces phonon anomalies that are consistent with recent Raman scattering and acoustic experiments. An important problem concerning strongly correlated metals such as heavy fermions, cuprates and organic superconductors is understanding the interplay of the strong interactions between electrons and the interactions between the electrons and phonons. This is particularly relevant to understanding the question of whether the phonons play any role in superconductivity . In this Letter we show how strong electronic correlations can lead to phonon frequencies varying with temperature. Although, our calculations focus on explaining recent experiments on a particular family of organic superconductors, the physics involved is relevant to other strongly correlated metals in which there is a significant redistribution of the electronic spectral weight as the temperature is varied. The quasi-two-dimensional organic superconductors, $`\kappa `$-(BEDT-TTF)<sub>2</sub>X, are strongly correlated electron systems. Recent Raman scattering experiments find that in the metallic state the frequency of certain phonons associated with the BEDT-TTF molecules have a non-monotonic temperature dependence below 200 K. Acoustic experiments find that around 40 K there is a significant softening of the speed of longitudinal sound propogating perpendicular to the layers. This softening, of the order of a few per cent, is more than an order of magnitude larger than the softening associated with the superconducting transition. We will show that such a temperature dependence can arise due to the destruction of Fermi liquid quasiparticles that occurs above the coherence temperature, $`T^{}`$ . Anomalies in acoustic phonons in the heavy fermion UPt<sub>3</sub> have also been seen at temperatures of the order of $`T^{}`$. The simplest possible strongly correlated electron model for the $`\kappa `$-(BEDT-TTF)<sub>2</sub>X family is a single band Hubbard model on an anisotropic triangular lattice at half-filling. We have found that many of the transport properties of the metallic phase are consistent with the predictions of dynamical mean-field theory (DMFT) which captures exactly the dynamical fluctuations at each lattice site, but neglects all non-local spatial correlations. In order to understand the effect of electronic correlations on the molecular phonon modes we have considered the relevant Holstein-Hubbard model where the phonon amplitude couples linearly to the local charge density. The phonon self-energy is proportional to the electronic density-density correlation function which we calculate using DMFT. When the electron-electron interactions are sufficiently large we find that phonon frequencies can have a non-monotonic temperature dependence. This is related to the crossover from a Fermi liquid to a bad metal that occurs with Increasing temperature and has a significant effect on the electronic transport properties. The actual form of the temperature dependence of the shift in phonon frequency varies significantly with the phonon frequency. We consider the Hubbard-Holstein Hamiltonian $`H`$ $`=`$ $`t_1{\displaystyle \underset{ij,\sigma }{}}(c_{i\sigma }^{}c_{j\sigma }+h.c.)+t_2{\displaystyle \underset{ik,\sigma }{}}(c_{i\sigma }^{}c_{k\sigma }+h.c.)`$ (1) $`+`$ $`U{\displaystyle \underset{i}{}}n_in_i\mu {\displaystyle \underset{i\sigma }{}}n_{i\sigma }`$ (2) $`+`$ $`{\displaystyle \frac{g}{\sqrt{2}}}{\displaystyle \underset{i\sigma }{}}(a_i^{}+a_i)n_{i\sigma }+\omega _0{\displaystyle \underset{i}{}}a_i^{}a_i`$ (3) where the electronic part describes electrons on the antibonding orbitals of each dimer of BEDT-TTF molecules which are located on an anisotropic triangular lattice. $`t_1`$ and $`t_2`$ are the nearest and next-nearest-neighbour hoppings. $`U`$ is the Coulomb repulsion between two electrons on the same site and $`\mu `$ is the chemical potential. The operator, $`c_i^{}`$, creates an electron on the anti-bonding orbital of the dimer. The operator, $`a_i^{}`$, creates a phonon at site $`i`$, which describes a molecular vibration of frequency $`\omega _0`$. We consider only the in-phase vibrations of the two phonon modes associated with the dimer; these can be activated by Raman scattering (see below). $`g`$ is the coupling between the Raman active phonon and the electron density on a dimer. We focus on the parameter regime where the electron- electron interaction is dominant and we are well away from any instability (superconducting or charge-density-wave) due to the electron-phonon coupling. We can then decouple the set of Dyson’s equations in the electron-phonon problem, so that the electron self-energy contains the electron-electron scattering mechanism only; any effects coming from the interaction of the electrons with phonons on the electron propagator are neglected. The electron Green’s function is given by $$G_\sigma (𝐤,i\omega _n)=\frac{1}{i\omega _nϵ_𝐤\mathrm{\Sigma }(i\omega _n)}$$ (4) where $`\omega _n=(2n+1)\pi T`$ is a Matsubara fermion frequency for temperature $`T`$. $`ϵ_𝐤=t_1\mathrm{cos}(k_x)+t_2\mathrm{cos}(k_x+k_y)`$ is the dispersion relation for the anisotropic triangular lattice. $`\mathrm{\Sigma }(i\omega _n)`$ is the momentum independent self-energy computed within DMFT from the associated Anderson impurity problem, using the iterative perturbation theory. Details can be found elsewhere . The phonon problem is solved separately through the associated Dyson equation $$D(𝐪,\omega )=\frac{\omega _0^2}{\omega ^2\omega _0^2g^2\omega _0^2\mathrm{\Pi }(𝐪,\omega )/2}$$ (5) where $`\mathrm{\Pi }(𝐪,\omega )`$ is the electronic density-density correlation function (or polarization) which includes the full effect of the electron-electron interactions at the level of DMFT. For $`\mathrm{\Pi }(𝐪,\omega )`$ we take the particle-hole bubble, which, in terms of Matsubara frequencies is given by $$\mathrm{\Pi }(𝐪,i\nu _n)=T\underset{𝐤,\sigma ,\omega _n}{}G_\sigma (𝐤,i\omega _n)G_\sigma (𝐤+𝐪,i\omega _n+i\nu _n)$$ (6) and $`G_\sigma (𝐤,i\omega _n)`$ is the Green’s function of the electrons obtained from the solution of (4). The temperature dependence of the polarization predominantly comes from the temperature dependence of the individual one-electron Greens functions. Fig. 1 shows how the Fermi liquid quasiparticle peak in the spectral density of states is strongly temperature dependent. We show results for the case, $`t_1=t_2=t`$ , which corresponds to a triangular lattice. Due to the electron-phonon interactions, the phonon frequency is shifted from its bare value $`\omega _0`$ and has a finite lifetime. This shift can be obtained from the poles of (5), which we denote by $`\stackrel{~}{\omega }_𝐪=\omega _𝐪+i\mathrm{\Gamma }_𝐪`$. For weak electron-phonon coupling the frequency shift is $$\frac{\mathrm{\Delta }\omega }{\omega _0}\frac{\omega _𝐪\omega _0}{\omega _0}=\frac{g^2}{4}\mathrm{\Pi }^R(𝐪,\omega _0)$$ (7) where $`\mathrm{\Pi }^R(𝐪,\omega _0)`$ denotes the real part of the polarization. The phonon damping is proportional to the imaginary part of the polarization. Fig. 2 shows the temperature dependence of the real part of the polarization, $`\mathrm{\Pi }^R(𝐪=0,\omega _0)`$, for different bare phonon frequencies, $`\omega _0`$. Our DMFT calculations show that for frequencies close to $`U/2`$ and sufficiently strong interactions the phonon frequency can have a non-monotonic temperature dependence near the coherence temperature, $`T^{}`$. From Fig. 2, we find that $`T^{}0.1t100K`$, for $`t0.1`$ eV. However, this behaviour disappears gradually as $`U`$ decreases and becomes smaller than the bandwidth, $`W=4.5t`$. This is clearly seen for $`U=2t`$. The effective mass enhancement for $`U=5.5t`$ is $`m^{}/m3.8`$; this is in the range of the effective masses found in $`\kappa `$-(BEDT-TTF)<sub>2</sub>X . Lin et al. have measured an anomalous softening below about 100 K of the Raman frequency shifts in the phonons, $`\nu _9=`$ 505 cm<sup>-1</sup> and $`\nu _{60}(B_{3g})=`$ 890 cm<sup>-1</sup> for $`\kappa `$-(BEDT-TTF)<sub>2</sub> Cu(SCN)<sub>2</sub> (see Fig. 5 in ref.). The higher frequency mode, $`\nu _3=`$ 1478 cm<sup>-1</sup>, also exhibits a strong (but monotonic) temperature dependence. Table I gives the magnitude of the total temperature dependence between about 10 and 300 K. Other phonons do not exhibit such a strong temperature dependence. Fig. 2 and equation (7) implies that for a coupling $`gt`$ the temperature dependence can be as large as 5 % for $`\omega _0=U/2`$ and of the order of 1 % for larger or smaller frequencies. The values of $`g`$ deduced below (see also Table 1) are consistent with this if $`t0.1`$ eV which is also reasonable . However, if $`gt\omega _0`$ this raises questions about vertex corrections to (6) and the contribution of electron phonon coupling to the electronic self energy. Yet the effects shown in Fig. 2 are at most a few per cent and so we consider that effects that are higher order in $`g`$ will be very small. We find even stronger effects for low-frequency phonons. In Fig. 3 the real part of the phonon self-energy is plotted for $`𝐪=0`$ and $`\omega _0=0.05t`$. For strong correlations ($`U>W`$), a dip in the real part of the self-energy appears at the coherence temperature, $`T^{}`$. For decreasing values of $`U`$, the position of the dip moves to higher temperatures (because $`T^{}`$ increases) and the dip becomes smaller. These results could be relevant to understanding recent acoustic experiments . The velocity of ultrasonic waves which are propagating perpendicular to the layers was found to have a non-monotonic temperature dependence. These waves have frequencies of 100 MHz and velocities of about 2000 cm/s. The velocity versus temperature shows a broad dip of a few per cent around 40 K. This softening becomes larger as the pressure is decreased and is about three times larger for $`\kappa `$-(BEDT-TTF)<sub>2</sub>Cu\[N(CN)<sub>2</sub>\]Br than $`\kappa `$-(BEDT-TTF)<sub>2</sub> Cu(SCN)<sub>2</sub>. Decreasing the pressure or changing the anion from Cu(SCN)<sub>2</sub> to Cu\[N(CN)<sub>2</sub>\]Br corresponds to increasing the electronic correlations or increasing of $`U/t`$. For example, it has been observed that as the pressure decreases $`m^{}/m_e`$ increases and the metal-insulator transition is approached . Our calculated variation of the position of the dip with $`U`$ is in qualitative agreement with the observed variation of the position of the dip with pressure (compare Fig. 4 in Ref. ). However, caution is in order because further experiments find that the softening only occurs for waves propogating parallel to the layers when their polarisation is perpendicular to the layers. Our model can only explain this if modulation of the interlayer spacing has a much stronger coupling to the electronic charge density within the dimers than modulation of the interdimer spacing. The family $`\kappa `$-(BEDT-TTF)<sub>2</sub>X are particulary amenable to study the effect of electronic correlations on the molecular phonons because the dimer structure of the molecular crystal allows us to extract the electron-phonon coupling strength $`g`$ from experimental data. The crystal structure is such that the BEDT-TTF molecules are arranged in pairs that are reasonably well separated from one another. There are three electrons per dimer. For each phonon mode on a molecule there is a symmetric and an anti-symmetric combination on the dimer. By parity conservation, the anti-symmetric modes are infra-red active and the symmetric modes are Raman active. If there is a Holstein model coupling for a single molecule the Hamiltonian for a dimer is $`H_{ep}=t_0\left(c_1^{}c_2+c_2^{}c_1\right)+{\displaystyle \underset{i=1}{\overset{2}{}}}gc_i^{}c_i\left(a_i+a_i^{}\right)+\omega _0a_i^{}a_i`$ where $`t_0`$ is the transfer integral for hopping between molecules within the dimer. This is much larger than the hoppings $`t_1`$ and $`t_2`$ between dimers that appears in (3) . The bonding and anti-bonding orbitals associated with the dimer are separated by $`2t_0`$, even in the presence of a Hubbard term. The Hamiltonian can be re-written as $`H_{ep}=t_0\left(c_1^{}c_2+c_2^{}c_1\right)+{\displaystyle \underset{\alpha =\pm }{}}{\displaystyle \frac{g}{\sqrt{2}}}n_\alpha \left(a_\alpha +a_\alpha ^{}\right)+\omega _0a_\alpha ^{}a_\alpha `$ where $`a_\pm =\frac{1}{\sqrt{2}}(a_1\pm a_2)`$ and $`n_\pm =c_1^{}c_1\pm c_2^{}c_2`$. The $`n_+`$ term is the same as the number operator, $`n_i`$, appearing in (3). It is the fluctuations in this term that produce the temperature dependence shown in Figures 2 and 3. In contrast, the infrared active mode will not couple to these fluctuations and so should have no such temperature dependence. If the number of electrons on the dimer is fixed then the symmetric mode has frequency $`\omega _0`$. A second-order perturbative calculation predicts that for $`\omega _0t_0`$, the infra-red frequency $`\omega _{IR}`$ will be smaller than the Raman frequency ($`\omega _R=\omega _0`$), $$\frac{\omega _{IR}\omega _R}{\omega _R}=\frac{g^2}{2\omega _0t_0}.$$ (8) Table I lists the details of the measured frequency shifts for three different modes for $`\kappa `$-(BEDT-TTF)<sub>2</sub>Cu(SCN)<sub>2</sub>. For $`t_0=0.3`$ eV these values are used together with (8) to evaluate the electron-phonon coupling $`g`$. The values obtained are in reasonable agreement with those found in a MNDO frozen phonon quantum chemistry calculation on a single BEDT-TTF molecule. Zeyher and Zwicknagl considered the frequency shift of phonons when they enter the superconducting phase. Within the framework of BCS theory, phonons with frequency $`\omega _0`$ much smaller than the Fermi energy and much larger than the superconducting gap $`\mathrm{\Delta }`$ will harden by an amount $$\frac{\omega _s\omega _R}{\omega _0}=\frac{8g^2N(0)}{\omega _0}\left(\frac{\mathrm{\Delta }}{\omega _0}\right)^2\mathrm{ln}\left(\frac{\omega _0}{\mathrm{\Delta }}\right)$$ (9) where $`N(0)`$ is the density of states per spin at the Fermi energy. There should be no shift in the frequency of the infra-red active mode on entering the superconducting phase. Using $`\mathrm{\Delta }2k_BT_c0.002`$ eV and the values for $`g`$ in Table I we have evaluated the estimated shift in the phonon frequency for three different modes in Table I. Eldridge et al. found that in $`\kappa `$-(BEDT-TTF)<sub>2</sub>Cu\[N(CN)<sub>2</sub>\]Br the $`\nu _{60}(B_{3g})`$ mode at 890 cm<sup>-1</sup> hardened by about 0.2 % on entering the superconducting phase. A shift in this mode was not detected in $`\kappa `$-(BEDT-TTF)<sub>2</sub>Cu(SCN)<sub>2</sub>. No frequency shift in the other high-frequency modes was detected to within about 0.1 %. However, one should be cautious about making quantitative comparisons between experiment and Table I since (9) is only valid for $`\omega _0t`$ and we have $`\omega _0t`$. Pedron et al. did observe the hardening of modes with frequencies ranging from 27 to 134 cm<sup>-1</sup>. For acoustic phonons the softening will scale like $`(\omega _0/\mathrm{\Delta })^2`$ and so be much smaller than the effects shown in Fig. 3. In conclusion, we have shown how in a strongly correlated metal the redistribution of spectral weight over the scale of the band width with varying temperature can result in phonon frequencies with anomalous temperature dependence. The effects involved are larger than those associated with the superconducting transition because the latter only involves a redistribution of spectral weight over energies of the order of the superconducting gap which is much less than the band width. This work was supported by the Australian Research Council. We thank J. E. Eldridge, J. B. Marston and M. Poirier for helpful discussions.
warning/0007/quant-ph0007084.html	ar5iv	text	# Noise-free scattering of the quantized electromagnetic field from a dispersive linear dielectric ## I Introduction A fundamental problem in quantum optics is how the properties of light change as it propagates through a medium. If the medium is nonlinear, new frequencies can be produced and the quantum noise properties of the field can be altered. This leads to such interesting phenomena as solitons, squeezing, and quantum phase diffusion, all of which have been observed . If the medium is linear, the situation is not as dramatic, but linear media serve as a first step in the description of their nonlinear brethren, and they present problems in their own right, such as the inclusion of dispersion. It is generally thought that an accurate, first-principles treatment of dispersion necessitates the inclusion of absorption, and consequently additional noise or reservoir operators. Here we analyse a quantum field theoretic model that demonstrates dispersion-induced absorption, but without any additional noise operators appearing in the scattering relations. Dielectric media can be described in a number of different ways. They can be characterized by their susceptibilities, a procedure which, when fields are included leads to the macroscopic Maxwell equations. Consequently, we shall call this the macroscopic approach. It leads to difficulties when one wants to include dispersion in a Lagrangian or Hamiltonian formulation, because dispersion is a consequence of the fact that the response of the medium is not instantaneous, but depends on the values of the field over a range of times . For certain kinds of fields, in particular narrow-band ones, these problems can be overcome by using an approximate Lagrangian which is local in time . Another approach is to construct a microscopic model for the medium and to include the degrees of freedom of the medium in the theory -. This, for obvious reasons, we shall call the microscopic approach. It has the advantage that the inclusion of dispersion is not a problem, but the disadvantage is that for each new medium, a new model must be constructed. There are also intermediate approaches which use frequency-dependent susceptibilities, but also add quantum noise operators to the equations of motion for the fields , . Most of the work on quantized fields in media has concentrated on what happens inside the medium. However, there has been a steady stream of research which has considered fields entering and leaving a medium as well. This is essential if one wants to describe real experiments, in which the fields are generated outside the medium, then pass through it, and are finally measured in free space. Perhaps the first to investigate this question were Lang, et. al., who examined the connection between the field inside and outside a laser cavity while studying why the laser linewidth is so narrow , . Their model consisted of a cavity bounded on one end by a perfectly reflecting wall and on the other by a thin dielectric slab, and the cavity itself is filled with an active medium. This cavity is embedded in a larger cavity which represents the universe. They showed how the modes of the universe are related to the cavity quasi-modes, which makes it possible to find the output field in terms of the cavity field. In their analysis the field was classical, but soon thereafter the quantum version of their model was constructed and used to investigate the relation between the field inside and outside the cavity for a laser in the linear regime by Ujihara . This approach was later used by Gea-Banacloche, et. al. to study the relationship between the squeezing generated inside a cavity to that outside the cavity . This last problem had first been considered by Yurke who based his approach on an earlier paper by Denker and himself . They developed a quantum theory of electronic networks in which the network itself is located at $`x=0`$ and transmission lines extending from there to $`x=+\mathrm{}`$ bring input signals to the network and carry output signals away. Fields propagating toward $`x=0`$ are input fields and those propagating away are output fields, and the object is to find the output fields in terms of the input ones for a given network. A related input-output theory was developed by Collett and Gardiner and was put on a firmer footing by Carmichael . This theory considers a cavity containing a medium, active or passive, linear or nonlinear, which is coupled to a reservoir. The reservoir operators serve as the input and output fields. The dynamics of the system inside the cavity is described by a master equation, and its solution is used to find the time-dependent reservoir operators and, thereby, the output field. More recently a number of groups have examined the scattering of the quantized electromagnetic field from inhomogeneous linear dielectrics. Glauber and Lewenstein considered the case of a non-dispersive, lossless dielectric which is described by a real position-dependent susceptibility . The scattering from dispersive media was studied by Knöll and Leonhardt who considered a medium consisting of damped harmonic oscillators . Rather than use a formal scattering approach they used time-dependent Green’s functions to solve the field equations. A different treatment of dispersive media was given by Matloob, et. al. . They considered an arbitrary complex, frequency-dependent dielectric function and quantized the theory at the level of the equations of motion rather than starting with a Lagrangian. Working if frequency space they found the fields emerging from a dielectric slab in terms of those entering it. A similar analysis was carried out by Gruner and Welsch . A final approach is based on polaritons in finite media \- . For an infinite dielectric interacting with the electromagnetic field the eigenstates of the Hamiltonian are mixed matter-field modes known as polaritons. If the medium is finite, the polaritons acquire a finite lifetime. By looking at the the electromagnetic parts of the polariton modes, scattering of the field from the medium can be described. In this paper we shall apply the quantum scattering formalism for fields to describe an electromagnetic wave scattering from a finite medium, which behaves as a mirror or beam-splitter. The medium is treated microscopically, and it is dispersive. The final result is an explicit expression for the out operators in terms of those of the input field. The calculation starts from first principles, and, consequently, shows how some of the relations between in and out operators, which are often used in quantum optics, follow from an underlying scattering theory. An important feature of the model used here is that it includes a dielectric medium with multiple bare resonances, leading to a number of discrete absorption bands. This is typical of real dielectric materials, and leads to a dielectric constant that can be modeled with the widely used Sellmeir expansion in frequency. As we noted in our discussion of previous results, there are three treatments of scattering from a dispersive, linear dielectric. Two are to some extent phenomenological in that they do not start from a Hamiltonian describing the field-medium system , . The third approach treats time-dependent fields rather than finding the asymptotic in and out fields which are the basic objects in a scattering treatment . We believe that this leaves room for a more fundamental approach which can place the theory on a firmer foundation. The results we find from field theoretic scattering theory are similar to those found in the approach pioneered by Yurke, and have the virtue that they are simple and intuitively clear. By employing a fundamental approach, we have the advantage that the meaning of all of the operators which we employ is well-defined, which is not always the case in the more phenomenological treatments. The results presented here can be viewed as a justification of earlier phenomenological theories. ## II Model We shall consider a one-dimensional model of the electromagnetic field and the medium which was developed in reference . This model can be used to describe the normal incidence of an electromagnetic wave on a medium, where the wave travels in the $`x`$ direction and is polarized in the $`z`$ direction. The field can be represented by means of the dual potential, $`\mathrm{\Lambda }(x,t)`$, which is appropriate if there are no free charges. In the case of a $`z`$-polarized normally incident plane wave, $`\mathrm{\Lambda }(x,t)`$ is the $`y`$ component of the dual potential. The fields are given by $$D=\frac{\mathrm{\Lambda }}{x}B=\mu _0\frac{\mathrm{\Lambda }}{t}.$$ (1) The medium consists of dipoles which are harmonic oscillators with masses $`m_\nu `$ and bare frequencies $`\mathrm{\Omega }_\nu `$, where $`\nu =1,\mathrm{},N`$. The un-renormalized oscillator frequencies can be chosen to correspond to transition frequencies of atoms or molecules making up an actual material. Each oscillator is described by a field, $`r_\nu (x)`$, which gives the displacement of the oscillator at position $`x`$ and with frequency $`\mathrm{\Omega }_\nu `$. It is convenient to represent the oscillators in terms of the polarization fields $$p_\nu (x)=q_\nu \rho _\nu (x)r_\nu (x),$$ (2) where $`\rho _\nu (x)`$ is the density of oscillators with frequency $`\mathrm{\Omega }_\nu `$, and the dipole corresponding to oscillators of type $`\nu `$ consists of charges $`q_\nu `$. We shall work in the multi-polar gauge so that the coupling between the electromagnetic field and the medium is proportional to $`_\nu p_\nu (x)D(x)`$. The medium self-interaction terms proportional to the square of the total polarization are incorporated into the frequencies, $`\mathrm{\Omega }_\nu `$. For a volume such as a waveguide of cross-sectional area $`A`$, the Lagrangian density for the medium-field system is given by $``$ $`=`$ $`{\displaystyle \frac{A}{2ϵ_0}}\{{\displaystyle \frac{1}{c^2}}\dot{\mathrm{\Lambda }}^2(x)(_x\mathrm{\Lambda }(x))^2+{\displaystyle \underset{\nu }{}}[{\displaystyle \frac{1}{g_\nu (x)}}(\dot{p}_\nu ^2(x)\mathrm{\Omega }_\nu ^2p_\nu ^2(x))`$ (4) $`+2p_\nu (x)_x\mathrm{\Lambda }(x)]\},`$ where $$g_\nu (x)=\frac{q_\nu ^2\rho _\nu (x)}{m_\nu ϵ_0}.$$ (5) ### A Refractive index ¿From the Lagrangian density we find the equations of motion for the fields $`_t^2\mathrm{\Lambda }c^2_x^2\mathrm{\Lambda }`$ $`=`$ $`c^2_x{\displaystyle \underset{\nu }{}}p_\nu `$ (6) $`_t^2p_\nu +\mathrm{\Omega }_\nu ^2p_\nu `$ $`=`$ $`g_\nu _x\mathrm{\Lambda }.`$ (7) For a medium of constant density, (i. e. $`g_\nu (x)`$ is independent of $`x`$), we can solve the above equations by assuming that both $`p_\nu `$ and $`\mathrm{\Lambda }`$ are proportional to $`e^{i(kx\omega t)}`$. The values of $`\omega `$ are the frequencies of the modes of the system and are given by the solutions of equation : $$\omega ^2=(kc)^2\left[1\underset{\nu }{}\frac{g_\nu }{\mathrm{\Omega }_\nu ^2\omega ^2}\right].$$ (8) Defining the index of refraction, $`n(\omega )`$, to be $`kc/\omega `$, we find $$n(\omega )=\left[1\underset{\nu }{}\frac{g_\nu }{\mathrm{\Omega }_\nu ^2\omega ^2}\right]^{1/2}.$$ (9) This is very similar to the classical Sellmeir expansion for the refractive index. Note that this expansion is not identical to the Sellmeir expansion, but can be converted into the commonly used Sellmeir form through a renormalization of the bare resonant frequencies of the oscillators. The characteristic property of this type of equation is that it possesses solutions for the refractive index that are either purely real (transmission bands) or purely imaginary (absorption bands). At the bare resonance frequency, the refractive index is zero. Near a resonance, where $`\omega \mathrm{\Omega }_\nu `$, the refractive index is real for $`\omega >\mathrm{\Omega }_\nu `$, and imaginary for $`\omega <\mathrm{\Omega }_\nu `$. At a finite detuning below a resonance, the refractive index goes to infinity just below the start of the corresponding absorption band. ### B Lagrangian quantization ¿From the Lagrangian density we can find the canonical momenta corresponding to $`\mathrm{\Lambda }`$ and $`p_\nu `$, which we shall denote by $`\mathrm{\Pi }`$ and $`\pi _\nu `$, respectively. These are given by $$\mathrm{\Pi }(x)=\mu _0\dot{\mathrm{\Lambda }}(x)\pi _\nu (x)=\frac{\dot{p}_\nu (x)}{ϵ_0g_\nu (x)}.$$ (10) The theory is quantized by imposing the commutation relations $$[\widehat{\mathrm{\Lambda }}(x,t),\widehat{\mathrm{\Pi }}(x^{},t)]=i\mathrm{}\delta (xx^{})/A$$ (11) and $$[\widehat{p}_\nu (x,t),\widehat{\pi }_\nu ^{}(x^{},t)]=i\mathrm{}\delta _{\nu ,\nu ^{}}\delta (xx^{})/A.$$ (12) The canonical momenta and the Lagrangian density can now be used to find the Hamiltonian density for the quantized theory $`(x)`$ $`=`$ $`{\displaystyle \frac{A}{2ϵ_0}}:\{{\displaystyle \frac{ϵ_0}{\mu _0}}\widehat{\mathrm{\Pi }}^2(x)+(_x\widehat{\mathrm{\Lambda }}(x))^2+{\displaystyle \underset{\nu }{}}[ϵ_0^2g_\nu (x)\widehat{\pi }_\nu ^2(x)`$ (14) $`+{\displaystyle \frac{\mathrm{\Omega }_\nu ^2}{g_\nu (x)}}\widehat{p}_\nu ^2(x)2\widehat{p}_\nu (x)_x\widehat{\mathrm{\Lambda }}(x)]\}:.`$ This can be put in a different form if we define annihilation and creation operators, $`\widehat{\xi }_\nu (x)`$ and $`\widehat{\xi }_\nu ^{}(x)`$, for the oscillators, where $$\widehat{\xi }_\nu (x)=\frac{1}{\sqrt{2\mathrm{}}}\left(\sqrt{\frac{\mathrm{\Omega }_\nu }{ϵ_0g_\nu (x)}}\widehat{p}_\nu +i\sqrt{\frac{ϵ_0g_\nu (x)}{\mathrm{\Omega }_\nu }}\widehat{\pi }_\nu \right),$$ (15) so that $$[\widehat{\xi }_\nu (x),\widehat{\xi }_\nu ^{}^{}(x^{})]=\delta _{\nu ,\nu ^{}}\delta (xx^{})/A.$$ (16) We finally have for the Hamiltonian density $`(x)`$ $`=`$ $`{\displaystyle \frac{A}{2ϵ_0}}:\{{\displaystyle \frac{ϵ_0}{\mu _0}}\widehat{\mathrm{\Pi }}^2(x)+(_x\widehat{\mathrm{\Lambda }}(x))^2+{\displaystyle \underset{\nu }{}}[2ϵ_0\mathrm{}\mathrm{\Omega }_\nu \widehat{\xi }_\nu ^{}(x)\widehat{\xi }_\nu (x)`$ (18) $`2\sqrt{{\displaystyle \frac{\mathrm{}ϵ_0g_\nu (x)}{2\mathrm{\Omega }_\nu }}}(\widehat{\xi }_\nu (x)+\widehat{\xi }_\nu ^{}(x))_x\widehat{\mathrm{\Lambda }}(x)]\}:`$ ## III Scattering Theory In order to determine what happens when an electromagnetic wave scatters off of the medium we shall apply the standard formulation of scattering for quantum fields . This is done in the Heisenberg picture so that it is the field operators which are time dependent. Because we shall consider a medium which is bounded in the $`x`$ direction, the interaction is bounded in time. This can be seen either by considering the incoming waves to be wave packets, so that the interaction takes place only while the packet is inside the medium, or by using plane waves and turning the interaction on and off adiabatically. In either approach, the fields will go to free fields both as $`t\mathrm{}`$ and as $`t\mathrm{}`$. The free fields as $`t\mathrm{}`$ are the in fields, and those as $`t\mathrm{}`$ are the out fields. The time dependent field operators which carry the full time dependence of the Hamiltonian, including the interaction, are also known as the interpolating fields, because they interpolate between the in and the out fields. Our goal is to use the interpolating fields to find an expression for the out fields in terms of the in fields. This will give us a complete description of the scattering process.We note here that a related Heisenberg-picture approach to quantum scattering theory relevant to quantum optics measurements was recently developed by Dalton et al.. They developed a similar basic formalism, but did not consider specific examples. ### A In and out fields To find the relationship between the in and out-fields, we need to express the equations of motion of the interpolating fields as integral equations. From the Hamiltonian for our model we find $`(_t^2c^2_x^2)\widehat{\mathrm{\Lambda }}`$ $`=`$ $`c^2_x{\displaystyle \underset{\nu }{}}\sqrt{{\displaystyle \frac{\mathrm{}ϵ_0g_\nu }{2\mathrm{\Omega }_\nu }}}\left[\widehat{\xi }_\nu +\widehat{\xi }_\nu ^{}\right]`$ (19) $`(_t+i\mathrm{\Omega }_\nu )\widehat{\xi }_\nu `$ $`=`$ $`{\displaystyle \frac{i}{ϵ_0\mathrm{}}}\sqrt{{\displaystyle \frac{\mathrm{}ϵ_0g_\nu }{2\mathrm{\Omega }_\nu }}}_x\widehat{\mathrm{\Lambda }}.`$ (20) In order to express these as integral equations we define the Green’s functions $`\mathrm{\Delta }^{(ret)}(x,t)`$, $`\mathrm{\Delta }^{(adv)}(x,t)`$, $`\mathrm{\Gamma }_\nu ^{(ret)}(x,t)`$, and $`\mathrm{\Gamma }_\nu ^{(adv)}(x,t)`$. They satisfy the equations $`(_t^2c^2_x^2)\mathrm{\Delta }^{(ret)}(x,t)`$ $`=`$ $`\delta (x)\delta (t)`$ (21) $`(_t+i\mathrm{\Omega }_\nu )\mathrm{\Gamma }_\nu ^{(ret)}(x,t)`$ $`=`$ $`\delta (x)\delta (t)`$ (22) $`(_t^2c^2_x^2)\mathrm{\Delta }^{(adv)}(x,t)`$ $`=`$ $`\delta (x)\delta (t)`$ (23) $`(_t+i\mathrm{\Omega }_\nu )\mathrm{\Gamma }_\nu ^{(ret)}(x,t)`$ $`=`$ $`\delta (x)\delta (t),`$ (24) and the boundary conditions $`\mathrm{\Delta }^{(ret)}(x,t)=\mathrm{\Gamma }_\nu ^{(ret)}(x,t)`$ $`=`$ $`0\text{for}t<0`$ (25) $`\mathrm{\Delta }^{(adv)}(x,t)=\mathrm{\Gamma }_\nu ^{(adv)}(x,t)`$ $`=`$ $`0\text{for}t>0.`$ (26) The retarded Green’s functions can be expressed as $`\mathrm{\Delta }^{(ret)}(x,t)`$ $`=`$ $`{\displaystyle \frac{1}{(2\pi )^2}}{\displaystyle 𝑑k𝑑\omega \frac{e^{i(kx\omega t)}}{(ck)^2(\omega +iϵ)^2}}`$ (27) $`\mathrm{\Gamma }_\nu ^{(ret)}(x,t)`$ $`=`$ $`{\displaystyle \frac{\delta (x)}{2\pi }}{\displaystyle 𝑑\omega \frac{e^{i\omega t}}{i(\mathrm{\Omega }iϵ\omega )}},`$ (28) where $`ϵ0^+`$, and the advanced Green’s functions are given by almost identical expressions, the only difference being that $`ϵ`$ is replaced by $`ϵ`$. The integral equations corresponding to the differential equations, Eqs. (20), are $`\widehat{\mathrm{\Lambda }}(x,t)`$ $`=`$ $`\widehat{\mathrm{\Lambda }}_{in}(x,t)c^2{\displaystyle 𝑑x^{}𝑑t^{}\mathrm{\Delta }^{(ret)}(xx^{},tt^{})}`$ (30) $`_x^{}{\displaystyle \underset{\nu }{}}\sqrt{{\displaystyle \frac{\mathrm{}ϵ_0g_\nu (x^{})}{2\mathrm{\Omega }_\nu }}}\left[\widehat{\xi }_\nu (x^{},t)+\widehat{\xi }_\nu ^{}(x^{},t)\right]`$ $`\widehat{\xi }_\nu (x,t)`$ $`=`$ $`\widehat{\xi }_\nu ^{(in)}(x,t)+{\displaystyle \frac{i}{ϵ_0\mathrm{}}}{\displaystyle 𝑑x^{}𝑑t^{}\mathrm{\Gamma }^{(ret)}(xx^{},tt^{})}`$ (32) $`\sqrt{{\displaystyle \frac{\mathrm{}ϵ_0g_\nu (x^{})}{2\mathrm{\Omega }_\nu }}}_x^{}\widehat{\mathrm{\Lambda }}(x^{},t^{}).`$ The corresponding expression involving the out-fields is: $`\widehat{\mathrm{\Lambda }}(x,t)`$ $`=`$ $`\widehat{\mathrm{\Lambda }}_{out}(x,t)c^2{\displaystyle 𝑑x^{}𝑑t^{}\mathrm{\Delta }^{(adv)}(xx^{},tt^{})}`$ (34) $`_x^{}{\displaystyle \underset{\nu }{}}\sqrt{{\displaystyle \frac{\mathrm{}ϵ_0g_\nu (x^{})}{2\mathrm{\Omega }_\nu }}}\left[\widehat{\xi }_\nu (x^{},t)+\widehat{\xi }_\nu ^{}(x^{},t)\right]`$ $`\widehat{\xi }_\nu (x,t)`$ $`=`$ $`\widehat{\xi }_\nu ^{(out)}(x,t)+{\displaystyle \frac{i}{ϵ_0\mathrm{}}}{\displaystyle 𝑑x^{}𝑑t^{}\mathrm{\Gamma }^{(adv)}(xx^{},tt^{})}`$ (36) $`\sqrt{{\displaystyle \frac{\mathrm{}ϵ_0g_\nu (x^{})}{2\mathrm{\Omega }_\nu }}}_x^{}\widehat{\mathrm{\Lambda }}(x^{},t^{}).`$ Note that the integral equations incorporate the boundary conditions for the fields. The first set of equations implies that $`\widehat{\mathrm{\Lambda }}(x,t)`$ and $`\widehat{\xi }_\nu (x,t)`$ will go to $`\widehat{\mathrm{\Lambda }}_{in}(x,t)`$ and $`\widehat{\xi }_\nu ^{(in)}(x,t)`$, respectively, as $`t\mathrm{}`$, and the second set implies that they will go to $`\widehat{\mathrm{\Lambda }}_{out}(x,t)`$ and $`\widehat{\xi }_\nu ^{(out)}(x,t)`$, respectively, as $`t\mathrm{}`$. What we shall do is solve the first set of equations for $`\widehat{\mathrm{\Lambda }}(x,t)`$ and $`\widehat{\xi }_\nu (x,t)`$ in terms of the in fields, and then insert this solution into the second set to find the out fields in terms of the in fields. ### B Fourier decomposition We begin solving Eqs. (32) by taking the time Fourier transform of both sides. Defining $`\widehat{\mathrm{\Lambda }}(x,\omega )`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2\pi }}}{\displaystyle 𝑑te^{i\omega t}\widehat{\mathrm{\Lambda }}(x,t)}`$ (37) $`\widehat{\xi }_\nu (x,\omega )`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2\pi }}}{\displaystyle 𝑑te^{i\omega t}\widehat{\xi }_\nu (x,t)},`$ (38) and similarly for the in and out-fields, we find that $`\widehat{\mathrm{\Lambda }}(x,\omega )`$ $`=`$ $`\widehat{\mathrm{\Lambda }}_{in}(x,\omega ){\displaystyle \frac{ic}{2\omega }}{\displaystyle 𝑑x^{}e^{i\omega \|xx^{}\|/c}_x^{}\underset{\nu }{}\sqrt{\frac{\mathrm{}ϵ_0g_\nu (x^{})}{2\mathrm{\Omega }_\nu }}}`$ (40) $`\left[\widehat{\xi }_\nu (x^{},\omega )+\widehat{\xi }_\nu ^{}(x^{},\omega )\right],`$ and $$\widehat{\xi }_\nu (x,\omega )=\widehat{\xi }_\nu ^{(in)}(x,\omega )+\frac{1}{ϵ_0\mathrm{}}\sqrt{\frac{\mathrm{}ϵ_0g_\nu (x)}{2\mathrm{\Omega }_\nu }}\frac{1}{\mathrm{\Omega }_\nu iϵ\omega }_x\widehat{\mathrm{\Lambda }}(x,\omega ).$$ (41) In deriving these equations we made use of the fact that $$𝑑te^{i\omega t}\mathrm{\Delta }^{(ret)}(x,t)=\frac{i}{2\omega c}e^{i\omega \|x\|/c}.$$ (42) We can derive an equation for only the field $`\widehat{\mathrm{\Lambda }}(x,k_0)`$ by substituting from Eq. (41) into Eq. (40). We find that $`\widehat{\mathrm{\Lambda }}(x,\omega )`$ $`=`$ $`\widehat{\mathrm{\Lambda }}_{in}(x,\omega ){\displaystyle \frac{ic}{2\omega }}{\displaystyle 𝑑x^{}e^{i\omega \|xx^{}\|/c}_x^{}\underset{\nu }{}\sqrt{\frac{\mathrm{}ϵ_0g_\nu (x^{})}{2\mathrm{\Omega }_\nu }}}`$ (45) $`[\widehat{\xi }_\nu ^{(in)}(x^{},\omega )+\widehat{\xi }_\nu ^{(in)}(x^{},\omega )`$ $`+{\displaystyle \frac{1}{ϵ_0\mathrm{}}}\sqrt{{\displaystyle \frac{\mathrm{}ϵ_0g_\nu (x^{})}{2\mathrm{\Omega }_\nu }}}{\displaystyle \frac{2\mathrm{\Omega }_\nu }{\mathrm{\Omega }_\nu ^2(\omega +iϵ)^2}}_x^{}\widehat{\mathrm{\Lambda }}(x^{},\omega )].`$ Our next step is to turn this into a differential equation, but before doing so we shall make a simplifying assumption. The field $`\widehat{\xi }_\nu ^{(in)}(x,t)`$ is a free field which oscillates at the frequency $`\mathrm{\Omega }_\nu `$, and this implies that $`\widehat{\xi }_\nu ^{(in)}(x,\omega )`$ is nonzero only when $`\omega =\mathrm{\Omega }_\nu `$ . We are mainly interested in cases where the incoming light is not resonant with the medium, so we shall initially assume that $`\omega \mathrm{\Omega }_\nu `$ for $`\nu =1,\mathrm{}N`$. This implies that we can drop $`\widehat{\xi }_\nu ^{(in)}(x,\omega )`$ and $`\widehat{\xi }_\nu ^{(in)}(x,\omega )`$ from the above equation and set $`ϵ=0`$. We return to the resonant case later. Next, we then apply the differential operator $`c^2_x^2+\omega ^2`$ to both sides. This annihilates the in-field term and converts the integral equation into a homogeneous differential equation. The result is $$_x\left(\frac{c^2}{n^2(x,\omega )}_x\widehat{\mathrm{\Lambda }}(x,\omega )\right)+\omega ^2\widehat{\mathrm{\Lambda }}(x,\omega )=0.$$ (46) Here, $`n(x,\omega )`$, the space and frequency dependent index of refraction of the medium, is given by $$n(x,\omega )=\left(1\underset{\nu }{}\frac{g_\nu (x)}{\mathrm{\Omega }_\nu ^2\omega ^2}\right)^{1/2}.$$ (47) In this form, the equations have a rather classical appearance, and the matter operators no longer appear in the formulation, which gives rise to a substantial simplification. ## IV Dielectric layer We now want to specialize our equations to the case of a dielectric layer with a uniform density of oscillators. This corresponds to the important case of a beam-splitter or mirror, although we make no restrictions as to the size of the layer. The medium extends from $`x=L`$ to $`x=L`$. Inside the medium, $`n(x,\omega )`$ has a value of $`n_0(\omega )`$, and outside the medium it has a value of $`1`$. The solutions to Eq. (46) should be continuous and $`n^2(x,\omega )`$ times their derivative should be continuous. These correspond to a continuous magnetic and electric field, respectively. ### A Classical Case In order to find the solution of the operator equation, Eq. (46), we first find solutions of the corresponding c-number equation, which we shall denote as $`u(x,\omega )`$. We begin by dividing the line into three regions, region I for $`x<L`$, region II for $`LxL`$, and region III for $`x>L`$. In regions I and III, $`u(x,\omega )`$ satisfies $$\left[c^2_x^2+\omega ^2\right]u(x,\omega )=0,$$ (48) and in region II $$\left[c^2_x^2+\omega ^2n_0^2(\omega )\right]u(x,\omega )=0,$$ (49) where $`n_0`$ is the value of $`n(x,\omega )`$ in region II. A solution incident from the left, $`u_l(x,\omega )`$, which satisfies the equation and has the proper continuity properties, is given by $$u_l(x,\omega )=\{\begin{array}{cc}e^{ik(\omega )x}+R(\omega )e^{ik(\omega )x}& \text{in\hspace{0.17em} \hspace{0.17em} region\hspace{0.17em} \hspace{0.17em} I}\hfill \\ B_r^{(l)}(\omega )e^{i\kappa (\omega )x}+B_l^{(l)}(\omega )e^{i\kappa (\omega )x}& \text{in\hspace{0.17em} \hspace{0.17em} region\hspace{0.17em} \hspace{0.17em} II}\hfill \\ T(\omega )e^{ik(\omega )x}& \text{in\hspace{0.17em} \hspace{0.17em} region\hspace{0.17em} \hspace{0.17em} III\hspace{0.17em} \hspace{0.17em} ,}\hfill \end{array}$$ (50) where $`k(\omega )=\omega /c`$ , and $`\kappa (\omega )=n_0k(\omega )`$ . It is to be remembered that waves proportional to $`e^{ikx}`$ or $`e^{i\kappa x}`$ are propagating to the right, and those proportional to $`e^{ikx}`$ or $`e^{i\kappa x}`$ are propagating to the left. Suppressing the frequency arguments for clarity, the coefficients in the above equation are given by $`R={\displaystyle \frac{i(n_0^21)\mathrm{sin}(2\kappa L)}{D}}e^{2ikL}`$ $`T={\displaystyle \frac{2n_0}{D}}e^{2ikL}`$ (51) $`B_r^{(l)}={\displaystyle \frac{n_0(n_0+1)}{D}}e^{i(\kappa +k)L}`$ $`B_l^{(l)}={\displaystyle \frac{n_0(n_01)}{D}}e^{i(\kappa k)L},`$ (52) where $`D=2n_0\mathrm{cos}(2\kappa L)i(n_0^2+1)\mathrm{sin}(2\kappa L)`$. Note that $`R`$ and $`T`$ are, respectively, the reflection and transmission coefficients for the medium, and that $`\|R\|^2+\|T\|^2=1`$. For a solution incident from the right we have $$u_r(x,\omega )=\{\begin{array}{cc}T(\omega )e^{ik(\omega )x}& \text{in\hspace{0.17em} \hspace{0.17em} region\hspace{0.17em} \hspace{0.17em} I}\hfill \\ B_r^{(r)}(\omega )e^{i\kappa (\omega )x}+B_l^{(r)}(\omega )e^{i\kappa (\omega )x}& \text{in\hspace{0.17em} \hspace{0.17em} region\hspace{0.17em} \hspace{0.17em} II}\hfill \\ e^{ik(\omega )x}+R(\omega )e^{ik(\omega )x}& \text{in\hspace{0.17em} \hspace{0.17em} region\hspace{0.17em} \hspace{0.17em} III\hspace{0.17em} \hspace{0.17em} ,}\hfill \end{array}$$ (53) where $`k`$ , $`\kappa `$, $`R`$ and $`T`$ are as before, while $`B_r^{(r)}=B_l^{(l)}`$, and $`B_l^{(r)}=B_r^{(l)}`$. ### B Asymptotic fields Both $`u_r`$ and $`u_l`$ are solutions of the differential equation, Eqs. (48) and (49), and this implies that they are also solutions to the corresponding integral equation $`\mathrm{\Lambda }(x,\omega )=`$ $`\mathrm{\Lambda }_{in}(x,\omega )`$ $`{\displaystyle \frac{ic}{2\omega }}{\displaystyle _L^L}𝑑x^{}e^{i\omega \|xx^{}\|/c}_x^{}`$ (55) $`\left[\left(1{\displaystyle \frac{1}{n^2(x,\omega )}}\right)_x^{}\mathrm{\Lambda }(x^{},\omega )\right],`$ for particular choices of the field $`\mathrm{\Lambda }_{in}(x,\omega )`$. We can find $`\mathrm{\Lambda }_{in}(x,\omega )`$ for both solutions simply by substituting them into Eq. (55). We must be careful, however, because the expression inside the square brackets is not continuous at $`x=\pm L`$, and it is being differentiated, so that the discontinuities will lead to finite contributions after being integrated. One way to find these contributions is to consider a refractive index which is continuous, but which goes to the desired one as a limit. For example, let us suppose that $`n(x,\omega )`$ is $`1`$ for $`x<L\delta `$ and $`x>L+\delta `$, is equal to $`n_0`$ for $`LxL`$, goes continuously from $`1`$ to $`n_0`$ as $`x`$ goes from $`L\delta `$ to $`L`$, and goes continuously from $`n_0`$ to $`1`$ as $`x`$ goes from $`L`$ to $`L+\delta `$. We can then take the limit $`\delta 0`$. Let us examine what happens in the interval between $`L\delta `$ and $`L`$; the interval between $`L`$ and $`L+\delta `$ is similar. As $`\delta 0`$ we have that $`{\displaystyle _{L\delta }^L}𝑑x^{}e^{i\omega \|xx^{}\|/c}_x^{}\left[\left(1{\displaystyle \frac{1}{n^2(x,\omega )}}\right)_x^{}\mathrm{\Lambda }(x^{},\omega )\right]`$ (56) $`e^{i\omega \|x+L\|/c}{\displaystyle _{L\delta }^L}𝑑x^{}_x^{}\left[\left(1{\displaystyle \frac{1}{n^2(x,\omega )}}\right)_x^{}\mathrm{\Lambda }(x^{},\omega )\right]`$ (57) $`=e^{i\omega \|x+L\|/c}\left(1{\displaystyle \frac{1}{n_0^2(\omega )}}\right)_x\mathrm{\Lambda }(x,\omega )\|_{x=L^+},`$ (58) where $`x=L^+`$ denotes the limit as $`xL`$ from the positive direction ($`x=L^{}`$ is defined in an analogous fashion). The other limit of the integral contributes zero, due to the refractive index term approaching unity. Explicitly putting in the terms resulting from the boundaries of the medium gives $`\mathrm{\Lambda }(x,\omega )=`$ $`\mathrm{\Lambda }_{in}(x,\omega )`$ $`{\displaystyle \frac{ic}{2\omega }}{\displaystyle _{L^+}^L^{}}𝑑x^{}e^{i\omega \|xx^{}\|/c}\left(1{\displaystyle \frac{1}{n_0^2(\omega )}}\right)_x^{}^2\mathrm{\Lambda }(x^{},\omega )`$ (61) $`{\displaystyle \frac{ic}{2\omega }}[e^{i\omega \|x+L\|/c}(1{\displaystyle \frac{1}{n_0^2(\omega )}})_x\mathrm{\Lambda }(x,\omega )\|_{x=L^+}`$ $`e^{i\omega \|xL\|/c}(1{\displaystyle \frac{1}{n_0^2(\omega )}})_x\mathrm{\Lambda }(x,\omega )\|_{x=L^{}}].`$ If we now substitute $`u_l(x,\omega )`$ into this equation instead of $`\mathrm{\Lambda }(x,\omega )`$, we find that $$\mathrm{\Lambda }_{in}(x,\omega )=e^{i\omega x/c},$$ (62) and if we substitute $`u_r(x,\omega )`$, we find $$\mathrm{\Lambda }_{in}(x,\omega )=e^{i\omega x/c}.$$ (63) ### C Quantum Case In the quantum case we have the usual expansion of a free field in terms of annihilation and creation operators. This leads, in the present case, to: $$\widehat{\mathrm{\Lambda }}_{in}(x,t)=𝑑k\sqrt{\frac{\mathrm{}cϵ_0}{4\pi A\|k\|}}\left[\widehat{a}_k^{(in)}e^{i(kx\|k\|ct)}+(\widehat{a}_k^{(in)})^{}e^{i(kx\|k\|ct)}\right],$$ (64) which implies that for $`\omega >0`$ $$\widehat{\mathrm{\Lambda }}_{in}(x,\omega )=\sqrt{\frac{\mathrm{}ϵ_0}{2cAk(\omega )}}\left[\widehat{a}_{k(\omega )}^{(in)}e^{i\omega x/c}+\widehat{a}_{k(\omega )}^{(in)}e^{i\omega x/c}\right].$$ (65) The results of the previous paragraph allow us to see that if $`\widehat{\mathrm{\Lambda }}(x,\omega )`$ is given by $$\widehat{\mathrm{\Lambda }}(x,\omega )=\sqrt{\frac{\mathrm{}ϵ_0}{2cAk(\omega )}}[\widehat{a}_{k(\omega )}^{(in)}u_l(x,\omega )+\widehat{a}_{k(\omega )}^{(in)})u_r(x,\omega )],$$ (66) then it is a solution of Eq. (61) with $`\widehat{\mathrm{\Lambda }}_{in}`$ given by Eq. (65). This gives us the interpolating field in terms of the in field. Our remaining task is to use the expression for the interpolating field to find the out-field in terms of the in-field. This can be done by substituting the expression for $`\widehat{\mathrm{\Lambda }}(x,\omega )`$ given in the previous paragraph into the equation which relates the interpolating field to the out-field $`\widehat{\mathrm{\Lambda }}(x,\omega )=`$ $`\widehat{\mathrm{\Lambda }}_{out}(x,\omega )`$ $`+{\displaystyle \frac{ic}{2\omega }}{\displaystyle _{L^+}^L^{}}𝑑x^{}e^{i\omega \|xx^{}\|/c}\left(1{\displaystyle \frac{1}{n_0^2(\omega )}}\right)_x^{}^2\widehat{\mathrm{\Lambda }}(x^{},\omega )`$ (69) $`+{\displaystyle \frac{ic}{2\omega }}[e^{i\omega \|x+L\|/c}(1{\displaystyle \frac{1}{n_0^2(\omega )}})_x\widehat{\mathrm{\Lambda }}(x,\omega )\|_{x=L^+}`$ $`e^{i\omega \|xL\|/c}(1{\displaystyle \frac{1}{n_0^2(\omega )}})_x\widehat{\mathrm{\Lambda }}(x,\omega )\|_{x=L^{}}],`$ which follows from Eqs. (36). The derivation is almost identical to that of Eq. (61), so we do not give it explicitly. Making this substitution we find that, for $`\omega >0`$ $`\widehat{\mathrm{\Lambda }}_{out}(x,\omega )`$ $`=`$ $`\sqrt{{\displaystyle \frac{\mathrm{}ϵ_0}{2cAk(\omega )}}}[(T(\omega )e^{ik(\omega )x}+R(\omega )e^{ik(\omega )x})\widehat{a}_{k(\omega )}^{(out)}`$ (71) $`+(R(\omega )e^{ik(\omega )x}+T(\omega )e^{ik(\omega )x})\widehat{a}_{k(\omega )}^{(out)}].`$ The out-field can also be expressed in terms of out creation and annihilation operators, $$\widehat{\mathrm{\Lambda }}_{out}(x,t)=𝑑k\sqrt{\frac{\mathrm{}cϵ_0}{4\pi A\|k\|}}\left[\widehat{a}_k^{(out)}e^{i(kx\|k\|ct)}+(\widehat{a}_k^{(out)})^{}e^{i(kx\|k\|ct)}\right].$$ (72) Taking the Fourier transform of this equation with respect to time, for $`\omega >0`$, gives $$\widehat{\mathrm{\Lambda }}_{out}(x,\omega )=\sqrt{\frac{\mathrm{}ϵ_0}{2cAk(\omega )}}[\widehat{a}_{k(\omega )}^{(out)}e^{ik(\omega )x}+\widehat{a}_{k(\omega )}^{(out)}e^{ik(\omega )x})].$$ (73) Comparing Eqs. (71) and (73) we see that for $`\omega >0`$ $`\widehat{a}_{k(\omega )}^{(out)}`$ $`=`$ $`T(\omega )\widehat{a}_{k(\omega )}^{(in)}+R(\omega )\widehat{a}_{k(\omega )}^{(in)}`$ (74) $`\widehat{a}_{k(\omega )}^{(out)}`$ $`=`$ $`R(\omega )\widehat{a}_{k(\omega )}^{(in)}+T(\omega )\widehat{a}_{k(\omega )}^{(in)}.`$ (75) These equations are the solution to the scattering problem. We note that the transmission and reflection coefficients satisfy the usual relation of $`\|T(\omega )\|^2+\|R(\omega )\|^2=1`$. This holds even in the bandgap regions, where transmission occurs via an evanescent field. Thus, the scattering problem is explicitly unitary, as we would expect. It is important to notice that unitarity holds even inside the band-gaps of the problem, indicating that the absorption bands simply modify the reflection and transmission coefficients, without removing photons. Another way to think of this, is that even when a photon is removed through virtual excitation of an atomic resonance, the photon will eventually be re-radiated in either the forward of backward directions. ### D Resonances Let us now consider what happens when $`\omega =\pm \mathrm{\Omega }_\nu .`$ This case was excluded from our earlier treatment, since it must involve some matter-operator contribution, which we have neglected so far. Of course, as the resonances are discrete, the frequencies involved are essentially a set of measure zero, lying on the upper edge of each band-gap. In this case, terms proportional to either $`\widehat{\xi }_\nu ^{(in)}(x,\omega )`$ or $`\widehat{\xi }_\nu ^{(in)}(x,\omega )`$ will be present in Eq. (45). This, in turn, means that the solution for the interpolating field given in Eq.(66) must be modified. In particular, a term proportional to the matter fields must be added. This is not surprising; it is usually accepted that in a dispersive medium there must be absorption, and this in turn generally requires coupling to a reservoir field. However, it would be rather surprising to find quantum noise only occurring at discrete frequencies corresponding to the band edge. We have already established, in particular, that absorption in the band at frequencies different from the resonances simply changes the transmission and reflection coefficients, without adding any noise source. We will now show that even with the matter terms included, our earlier conclusions still hold; there are no extra noise sources for the out-fields. In order to see this, we first note that $$\widehat{\xi }_\nu ^{(in)}(x,t)=e^{i\mathrm{\Omega }_\nu t}\widehat{\xi }_\nu ^{(in)}(x),$$ (76) where: $`\widehat{\xi }_\nu ^{(in)}(x)=e^{i\mathrm{\Omega }_\nu t}\widehat{\xi }_\nu ^{(in)}(x,0)`$. This implies that: $`\widehat{\xi }_\nu ^{(in)}(x,\omega )`$ $`=`$ $`\sqrt{2\pi }\delta (\omega \mathrm{\Omega }_\nu )\widehat{\xi }_\nu ^{(in)}(x)`$ (77) $`\widehat{\xi }_\nu ^{(in)}(x,\omega )`$ $`=`$ $`\sqrt{2\pi }\delta (\omega +\mathrm{\Omega }_\nu )\widehat{\xi }_\nu ^{(in)}(x).`$ (78) Application of the differential operator $`c^2_x^2+\omega ^2`$ to Eq. (45), this time keeping the matter terms, gives: $$_x\left(\frac{c^2}{n^2(x,\omega )}_x\widehat{\mathrm{\Lambda }}(x,\omega )\right)+\omega ^2\widehat{\mathrm{\Lambda }}(x,\omega )=_x\widehat{F}_{in}(x,\omega ),$$ (79) where we now include an inhomogeneous term defined as: $$\widehat{F}_{in}(x,\omega )=c^2\underset{\nu }{}\sqrt{\frac{\mathrm{}\pi ϵ_0g_\nu (x)}{\mathrm{\Omega }_\nu }}\left(\delta (\omega \mathrm{\Omega }_\nu )\widehat{\xi }_\nu ^{(in)}(x)+\delta (\omega +\mathrm{\Omega }_\nu )\widehat{\xi }_\nu ^{(in)}(x)\right).$$ (80) This equation can be solved by means of a Green’s function which satisfies: $$_x\left(\frac{c^2}{n^2(x,\omega )}_xG(x,x^{})\right)+\omega ^2G(x,x^{})=\delta (xx^{}),$$ (81) together with the proper boundary conditions. For the boundary conditions, we shall choose $`G(x,x^{})`$ to have only outgoing waves at $`x=\pm \mathrm{}`$. Any fields produced by these matter terms are generated in a finite region and propagate outward, so that these boundary conditions are the appropriate ones. The Green’s function can be found by standard techniques, and is given by: $`G(x,x^{})`$ $`=`$ $`{\displaystyle \frac{1}{2i\omega cT(\omega )}}u_{\mathrm{}}(x,\omega )u_r(x^{},\omega )[x>x^{}]`$ (82) $`=`$ $`{\displaystyle \frac{1}{2i\omega cT(\omega )}}u_r(x,\omega )u_{\mathrm{}}(x^{},\omega )[x<x^{}].`$ (83) The solution to Eq. (79) with only outgoing waves at $`x=\pm \mathrm{}`$ , which we shall call $`\widehat{\mathrm{\Lambda }}_s(x,\omega )`$ is then: $$\widehat{\mathrm{\Lambda }}_s(x,\omega )=_L^L𝑑x^{}G(x,x^{})_x^{}\widehat{F}_{in}(x^{},\omega ).$$ (84) Substitution of $`\widehat{\mathrm{\Lambda }}_s(x,\omega )`$ into Eq.(45) \- the integral equation for $`\widehat{\mathrm{\Lambda }}(x,\omega )`$ \- shows that it is, as expected, a solution with $`\widehat{\mathrm{\Lambda }}_{in}(x,\omega )=0`$. A complete solution for general $`\widehat{\mathrm{\Lambda }}_{in}(x,\omega )`$ can then be obtained by adding to $`\widehat{\mathrm{\Lambda }}_s`$ a solution of the homogeneous equation with the proper $`\widehat{\mathrm{\Lambda }}_{in}`$, as discussed in the previous section. As has already been noted, $`\widehat{F}_{in}(x,\omega )`$ is only nonzero if $`\omega =\pm \mathrm{\Omega }_\nu `$ , for some index $`\nu `$. At these values, which correspond to the upper boundary of each band-edge, the index of refraction vanishes inside the medium, and: $`T(\omega )`$ $`=`$ $`{\displaystyle \frac{c}{ci\omega L}}e^{2i\omega L/c}`$ (85) $`R(\omega )`$ $`=`$ $`{\displaystyle \frac{i\omega L}{ci\omega L}}e^{2i\omega L/c}.`$ (86) For $`L<x<L`$ we have that: $$u_{\mathrm{}}(x,\omega )=u_r(x,\omega )=\frac{c}{ci\omega L}e^{i\omega L/c}.$$ (87) Therefore, for $`\omega =\pm \mathrm{\Omega }_\nu `$ , it follows that the Green’s function is constant, and $`\widehat{\mathrm{\Lambda }}_s(x,\omega )`$ is proportional to: $$_L^L𝑑x^{}_x^{}\widehat{F}_{in}(x^{},\omega )=\widehat{F}_{in}(L,\omega )\widehat{F}_{in}(L,\omega ).$$ (88) Here, however, we must be careful since the source function $`\widehat{F}_{in}(x,\omega )`$ is proportional to $`\sqrt{g_\nu (x)}`$, which is discontinuous at $`x=\pm L`$ , so that the right-hand side of the above equation is not well-defined. In order to resolve this difficulty, we must use the same technique as before, and consider a continuously changing refractive index over a small boundary region. That is, we suppose that $`n(x,\omega )`$ is $`1`$ for $`x<L\delta `$ and $`x>L+\delta `$, is equal to $`0`$ for $`LxL`$, goes continuously from $`1`$ to $`0`$ as $`x`$ goes from $`L\delta `$ to $`L`$, and goes continuously from $`0`$ to $`1`$ as $`x`$ goes from $`L`$ to $`L+\delta `$. A similar behaviour is assumed for$`\sqrt{g_\nu (x)}`$ . We can then take the limit $`\delta 0`$. Let us examine what happens in the interval between $`L\delta `$ and $`L`$; the interval between $`L`$ and $`L+\delta `$ is similar. For $`x>L+\delta `$ , the integral we wish to consider is then: $`{\displaystyle _{L\delta }^{L+\delta }}𝑑x^{}u_r(x^{},\omega )_x^{}\widehat{F}_{in}(x^{},\omega )`$ $`=`$ $`{\displaystyle _{L\delta }^L}𝑑x^{}u_r(x^{},\omega )_x^{}\widehat{F}_{in}(x^{},\omega )`$ (89) $`+`$ $`u_r(L,\omega )\left(\widehat{F}_{in}(L,\omega )\widehat{F}_{in}(L,\omega )\right)`$ (90) $`+`$ $`{\displaystyle _L^{L+\delta }}𝑑x^{}u_r(x^{},\omega )_x^{}\widehat{F}_{in}(x^{},\omega ).`$ (91) For $`x`$ in other intervals, the situation is similar. Note that $`u_r(x^{},\omega )`$ is now a solution of the homogeneous version of Eq. (81), including the modified index of refraction. From this equation it is relatively straightforward to show that, for $`\delta `$ small, $$_L^{L+\delta }𝑑x^{}u_r(x^{},\omega )_x^{}\widehat{F}_{in}(x^{},\omega )u_r(L,\omega )\left(\widehat{F}_{in}(L+\delta ,\omega )\widehat{F}_{in}(L,\omega )\right),$$ (92) and this becomes an equality as $`\delta 0`$ . A similar relationship holds for the integral from $`L\delta `$ to $`L`$ : $$_{L\delta }^L𝑑x^{}u_r(x^{},\omega )_x^{}\widehat{F}_{in}(x^{},\omega )u_r(L,\omega )\left(\widehat{F}_{in}(L\delta ,\omega )\widehat{F}_{in}(L,\omega )\right).$$ (93) Adding these contributions up, and noting that if $`n(x,\omega )=0`$ for $`LxL`$, then $`u_r(L,\omega )=u_r(L,\omega )`$, we find that: $$\underset{\delta 0}{lim}_{L\delta }^{L+\delta }𝑑x^{}u_r(x^{},\omega )_x^{}\widehat{F}_{in}(x^{},\omega )=0.$$ (94) This implies that the matter operators do not contribute to the interpolating field solution, even at the resonance frequencies which we did not consider in detail previously. Thus, the previous relation between the in and out operators still holds at the resonances where $`\omega =\pm \mathrm{\Omega }_\nu `$ . At first sight, this seems difficult to understand, since in general one would need to include noise operators to conserve commutation relations, and hence unitarity. However, this is consistent because, as can be seen from Eq.(86), we have $`\|T(\omega )\|^2+\|R(\omega )\|^2=1`$, even when $`\omega =\pm \mathrm{\Omega }_\nu `$ for $`\nu =1,..N`$ . In summary, we reach the somewhat surprising conclusion that no additional noise operators are needed in the asymptotic properties of the present model - even at resonance. ## V Photodetection example Given the state of the in-field, these equations allow us to calculate the properties of the out-field, and hence calculate observable scattering properties. In order to see how this works let us consider an example. We shall find the probability that a photo-detector located at $`x`$, where $`x>0`$ and is far from the medium, will fire at time $`t`$. At the long times required for propagation to this location, the fields will asymptotically become out-fields. The photo-detection probability is therefore proportional to $$in\|\widehat{D}_{out}^{()}(x,t)\widehat{D}_{out}^{(+)}(x,t)\|in=in\|(_x\widehat{\mathrm{\Lambda }}_{out}^{()}(x,t))(_x\widehat{\mathrm{\Lambda }}_{out}^{(+)}(x,t))\|in,$$ (95) where $`\|in`$ is the in state, $$\widehat{\mathrm{\Lambda }}_{out}^{(+)}(x,t)=𝑑k\sqrt{\frac{\mathrm{}cϵ_0}{4\pi A\|k\|}}\widehat{a}_k^{(out)}e^{i(kx\|k\|ct)},$$ (96) and $`\widehat{\mathrm{\Lambda }}_{out}^{()}(x,t)=(\widehat{\mathrm{\Lambda }}_{out}^{(+)}(x,t))^{}`$. Now let $`f(k)`$ be a function which is zero if $`k<0`$ . The Fourier transform of $`f(k)`$ is closely related to the shape of the pulse which is being sent into the medium. Define $$\widehat{a}_{in}^{}[f]=𝑑kf(k)(\widehat{a}_k^{(in)})^{},$$ (97) and let $$\|in=\mathrm{exp}(\widehat{a}_{in}[f]\widehat{a}_{in}^{}[f])\|0_{in}.$$ (98) This is a coherent state composed of wave packets with the intensity of the field and the shape of the wave packet determined by $`f(k)`$. For this state the correlation function in Eq. (95) is given by $$in\|\widehat{D}_{out}^{()}(x,t)\widehat{D}_{out}^{(+)}(x,t)\|in=\frac{\mathrm{}cϵ_0}{4\pi A}\left\|𝑑kf(k)T(ck)e^{i(kx\|k\|ct)}\right\|^2,$$ (99) where we have used Eq. (75) to relate the in and out operators, and we have explicitly indicated the $`k`$ dependence of the transmission coefficient. As expected, this equation demonstrates explicitly that photodetection rates are suppressed for frequency components that correspond to the dielectric absorption bands, where $`T(\omega )0`$. At these frequencies, the predominant effect will be a strong reflection, with no photodetection occurring at the detector location on the other side of the mirror. ## VI Conclusion We have presented an analysis of a quantized electromagnetic wave scattering off a linear, dispersive medium of finite extent. The medium consists of harmonic oscillators whose energy-level spacings can be chosen to match those of an actual medium in the spirit of the classical Sellmeir expansion. What emerges is a relation between the in and out fields which is most simply stated in terms of their annihilation operators. The overall results are exact, and simply expressed in terms of linear transmission and reflection coefficients. The medium has both transmission and absorption bands. Since the model is a full quantum-field version of the widely used Lorenz model that leads to the Sellmeir expansion, it has a wide area of applicability to realistic dielectric media with a variety of dispersion relations. The final results are very similar to the classical expressions which relate the amplitudes of the incoming and outgoing waves. The transmission and reflection coefficients which one finds from the quantum and classical analyses are identical, as they should be, given that the model is a linear one, and must reduce to the classical theory in the correspondence limit. Expressions such as those appearing in Eq. (75) are often used in quantum optics, and we believe it is useful to see how they emerge from the underlying scattering theory. An important feature of the theory treated here is that it has absorption bands, without any corresponding noise terms in the field equation. This is due to the dielectric model used here, in which the dielectric constant is always either purely real or purely imaginary. In this type of model, all photons absorbed are re-emitted. Thus, absorption simply results in a strong reflection, with purely evanescent fields inside the dielectric. In a related phenomenological treatment , one finds similar behaviour: if the dielectric constant is always either purely real or purely imaginary, it is possible to have dispersion without any additional noise terms. The reason for this is that the scattering terms alone are sufficient to ensure that the input-output relations remain unitary, with no change in the commutators. Thus, in the present model, there is no need for any additional source terms. The theory therefore provides a justification for the use of simple input-output relations to describe idealized dielectric or metallic mirrors, even when the dielectric response is dispersive. However, for realistic media it is generally the case that absorption can also occur even in the transmission bands. Treating this would require the use of more sophisticated models, including a complex refractive index. ###### Acknowledgements. This research was supported by the USA National Science Foundation under grant INT-9602515, and by the Australian Research Council.
warning/0007/cond-mat0007085.html	ar5iv	text	# 1/f Noise in a Coulomb Glass ## Abstract At low temperatures electron hopping in a three dimensional Coulomb glass produces fluctuations in the single particle density of states and hence in the resistivity. This results in a low frequency resisitivity noise spectrum which goes as $`f^\alpha `$ where $`\alpha `$ is very close to 1. This holds down to extremely low frequencies. Low frequency $`1/f`$ noise is ubiquitous; it is found in a wide variety of conducting systems such as metals, semiconductors, tunnel junctions , and even superconducting SQUIDs . Yet the microscopic mechanisms are still not well understood. In some cases the electrons are in a Coulomb glass which is an insulator with randomly placed electrons that have Coulomb interactions . Lightly doped semiconductors and disordered metals are examples of such systems. This paper focuses on $`1/f`$ noise in Coulomb glasses. Experimental studies on doped silicon inversion layers have shown that low frequency $`1/f`$ noise is produced by hopping conduction . More recent experiments have observed 1/f noise down to 0.1 Hz in boron–doped silicon . Because the systems are glassy, electron hopping can occur on very long time scales and this produces low frequency noise. In this paper we show that the resulting noise spectrum goes as $`f^\alpha `$ where $`f`$ is frequency and the exponent $`\alpha 1`$. Shklovskiĭ developed the first theory of 1/f noise in Coulomb glasses. He suggested that it is produced by fluctuations in the number of electrons in an infinite percolating cluster . These fluctuations are caused by the slow exchange of electrons between the infinite conducting cluster and small isolated donor clusters. A more rigorous calculation combined with numerical simulations of Shklovskii’s model found a noise spectrum that went as $`f^\alpha `$ where $`\alpha `$ was considerably lower than 1. Furthermore, below a minimum frequency of order 1–100 Hz, the noise spectral density saturated and became a constant independent of frequency. A similar conclusion holds for a model suggested by Kozub in which electron hops within finite clusters produce fluctuations in the potential seen by hopping conduction electrons that contribute to the current. Kogan has argued that transitions between valleys in the energy landscape produces 1/f noise because high barriers result in slow fluctuations in hopping conduction . In our approach electron hopping shifts the single particle energies $`\epsilon `$ because they depend on Coulomb interactions with other sites. This leads to fluctuations in the single particle density of states $`g(\epsilon )`$ which, in turn, produces fluctuations in the conductivity. The conductivity depends on the density of states $`g(\epsilon \mu )`$ in the vicinity of the Fermi energy $`\mu `$. Note that $`g(\epsilon \mu )`$ can be affected by hops between sites $`i`$ and $`j`$ even if the energies on these sites are not near the Fermi energy because an electron or hole on site $`i`$ or $`j`$ can interact with other sites whose energy is (or was) near the Fermi energy. We start with a model of the Coulomb glass that follows that of Baranovskiĭ, Shklovskiĭ, and Éfros (BSE) . In this model, the electrons occupy the sites of a periodic lattice, and the number of electrons is half the number of sites. Each site has a random onsite energy $`\varphi _i`$ chosen from a uniform distribution extending from $`A`$ to $`A`$. Thus, $`g_o`$, the density of states without interactions, is flat. A site can contain 0 or 1 electron. The Hamiltonian can be written as $$H=\underset{i}{}\varphi _in_i+\underset{i>j}{}\frac{e^2}{\kappa r_{ij}}n_in_j$$ (1) where the occupation number $`n_i`$ equals $`\frac{1}{2}`$ if site $`i`$ is occupied and $`\frac{1}{2}`$ if site $`i`$ is unoccupied, $`e`$ is the electron charge, $`\kappa `$ is the dielectric constant and $`r_{ij}`$ is the distance between sites $`i`$ and $`j`$. The single site energy is $`\epsilon _i=\varphi _i+_j\frac{e^2}{\kappa r_{ij}}n_j`$. At zero temperature Coulomb interactions between localized electrons result in a so–called Coulomb gap in the single particle density of states that is centered at the Fermi energy . We will use Mott’s argument for variable range hopping to relate fluctuations in the density of states to fluctuations in the resistivity. One can regard a Coulomb glass as a random resistor network with a transition between sites $`i`$ and $`j`$ associated with a resistance $`R_{ij}`$ given by $$R_{ij}=R_{ij}^o\mathrm{exp}(\xi _{ij})$$ (2) where the prefactor $`R_{ij}^o=kT/(e^2\gamma _{ij}^o)`$ with $`\gamma _{ij}^o`$ given by $$\gamma _{ij}^o=\frac{D^2\|\mathrm{\Delta }_i^j\|}{\pi ds^5\mathrm{}^4}\left[\frac{2e^2}{3\kappa a}\right]^2\frac{r_{ij}^2}{a^2}\left[1+\left(\frac{\mathrm{\Delta }_i^ja}{2\mathrm{}s}\right)^2\right]^4$$ (3) where $`D`$ is the deformation potential, $`s`$ is the speed of sound, $`d`$ is the mass density, and $`\mathrm{\Delta }_i^j=\epsilon _j\epsilon _ie^2/\kappa r_{ij}`$. $`\mathrm{\Delta }_i^j`$ is the change in energy that results from hopping from $`i`$ to $`j`$. $`a=\kappa a_B`$ is the effective Bohr radius of a donor, and $`a_B`$ is the usual Bohr radius ($`a_B=\mathrm{}^2/me^2`$). We will set the mass $`m`$ equal to the electron mass so that $`a_B=0.529\AA `$. In eq. (2), the exponent is given by $$\xi _{ij}=\frac{2r_{ij}}{a}+\frac{\epsilon _{ij}}{kT}$$ (4) The exponent reflects the thermally activated hopping rate between $`i`$ and $`j`$ as well as the wavefunction overlap between the sites. $`\epsilon _{ij}`$ is given by : $$\epsilon _{ij}=\{\begin{array}{cc}\|\epsilon _j\epsilon _i\|\frac{e^2}{\kappa r_{ij}},\hfill & (\epsilon _i\mu )(\epsilon _j\mu )<0\hfill \\ \mathrm{max}[\|\epsilon _i\mu \|,\|\epsilon _j\mu \|],\hfill & (\epsilon _i\mu )(\epsilon _j\mu )>0\hfill \end{array}$$ (5) At both high and low compensations, electron hopping usually occurs on one side of the Fermi level $`\mu `$ and the lower expression applies. At intermediate compensations and in the regime of variable range hopping, hopping electrons often cross the Fermi level and the upper expression applies. In the regime of variable range hopping Mott pointed out that hopping conduction at low temperatures comes from states near the Fermi energy. Let $`\stackrel{~}{\epsilon }=\epsilon \mu `$. If we consider states within $`\epsilon _o`$ of the Fermi energy, then the concentration of states in this band is $`N(\epsilon _o)=_{\epsilon _o}^{\epsilon _o}g(\stackrel{~}{\epsilon })𝑑\stackrel{~}{\epsilon }`$ where $`g(\stackrel{~}{\epsilon })`$ is the density of states measured from the Fermi energy. So the typical separation between sites is $`R=[N(\epsilon _o)]^{1/3}`$. To estimate the resistance corresponding to hopping between two typical states of the band, we replace $`r_{ij}`$ with $`R`$ and $`\|\epsilon _j\epsilon _i\|`$ with $`\epsilon _o`$ in eqs. (4) and (5) to obtain $`\xi (\epsilon _o)`$. Minimizing $`\xi (\epsilon _o)`$ yields $`\overline{\epsilon }_o`$. Plugging this into eqs. (4) and (2) yields the variable range hopping formula for the resistivity $`\overline{\rho }(T)=\rho _o(T)\mathrm{exp}(\xi (\overline{\epsilon }_o))`$. In our model the noise results from electron hopping which produces fluctuations in the density of states $`g(\epsilon )=\overline{g}(\epsilon )+\delta g(\epsilon )`$, where $`\overline{g}(\epsilon )`$ is the average density of states. This in turn creates fluctuations in $`N(\epsilon _o)`$, $`\xi (\epsilon _o)`$, $`\overline{\epsilon }_o`$, and $`\rho (T)`$. We can calculate these fluctuations by applying classical perturbation theory to the derivation of the variable range formula. To first order, $`\delta \xi (\epsilon _o)=\delta \rho (T)/\overline{\rho }(T)=(2kTg(T,\overline{\epsilon }_o))^1_{\overline{\epsilon }_o}^{\overline{\epsilon }_o}\delta g(T,\stackrel{~}{\epsilon })𝑑\stackrel{~}{\epsilon }`$. We have included the temperature dependence of the density of states because at finite temperatures the Coulomb gap fills in and the density of states no longer vanishes at the Fermi energy . The autocorrelation function for the fluctuations in the resistivity is $`{\displaystyle \frac{<\delta \rho (T,t_2)\delta \rho (T,t_1)>}{\overline{\rho }^2(T)}}={\displaystyle \frac{1}{4k^2T^2g^2(T,\overline{\epsilon }_o)}}`$ (6) $`{\displaystyle _{\overline{\epsilon }_o}^{\overline{\epsilon }_o}}𝑑\stackrel{~}{\epsilon }{\displaystyle _{\overline{\epsilon }_o}^{\overline{\epsilon }_o}}𝑑\stackrel{~}{\epsilon }^{}<\delta g(T,\stackrel{~}{\epsilon },t_2)\delta g(T,\stackrel{~}{\epsilon }^{},t_1)>`$ (7) We assume that there is no correlation between the fluctuations in the density of states at different energies, so $`<\delta g(T,\stackrel{~}{\epsilon },t_2)\delta g(T,\stackrel{~}{\epsilon }^{},t_1)>`$ $`=`$ $`E<\delta g(T,\stackrel{~}{\epsilon },t_2)\delta g(T,\stackrel{~}{\epsilon },t_1)>`$ (9) $`\delta (\stackrel{~}{\epsilon }\stackrel{~}{\epsilon }^{})`$ where $`E`$ is an energy of order $`2\overline{\epsilon }_o`$. Furthermore we assume that the time and energy dependence of the density of states autocorrelation function are separable, allowing us to write $$_{\overline{\epsilon }_o}^{\overline{\epsilon }_o}𝑑\stackrel{~}{\epsilon }<\delta g(T,\stackrel{~}{\epsilon },t_2)\delta g(T,\stackrel{~}{\epsilon },t_1)>=C(\overline{\epsilon }_o,T)f(T,t_2t_1)$$ (10) where we are assuming translational invariance in time (stationary processes). $`C(\overline{\epsilon }_o,T)`$ is a function of $`\overline{\epsilon }_o`$ and temperature. The function $`f(T,t)`$ characterizes the time dependence of the return to equilibrium by the system after it is perturbed by a fluctuation in the density of states. Inserting eqns. (9) and (10) in (7) yields $$\frac{<\delta \rho (T,t_2)\delta \rho (T,t_1)>}{\overline{\rho }^2(T)}=\frac{EC(\overline{\epsilon }_o,T)}{4k^2T^2g^2(T,\overline{\epsilon }_o)}f(T,t_2t_1)$$ (11) To relate this to the spectral density of the noise $`S(\omega )`$, let $`\psi _\rho (t_2t_1)=<\delta \rho (T,t_2)\delta \rho (T,t_1)>`$ and let $`\psi _\rho (\omega )`$ be the Fourier transform of $`\psi _\rho (t_2t_1)`$. According to the Wiener–Khintchine theorem , for a stationary process the spectral density of fluctuations is given by $`S_\rho (\omega )`$ $`=`$ $`2\psi _\rho (\omega )`$ (12) $`=`$ $`{\displaystyle \frac{E\overline{\rho }^2(T)C(\overline{\epsilon }_o,T)}{2k^2T^2g^2(T,\overline{\epsilon }_o)}}f(T,\omega )`$ (13) We do not know the temperature dependence of $`f(T,t)`$, so for the moment we will suppress this and just refer to $`f(t)`$. Theoretical calculations find that after large deviations from equilibrium, the density of states returns to equilibrium with a time dependence given by $`g(\mu ,t)\mathrm{ln}t`$ or $`g(\mu ,t)t^\theta `$ where $`\theta 1`$ . This agrees with experiments done at low temperatures . If we assume that these functional forms are also valid for $`f(t)`$ which applies to small perturbations, then we obtain $`1/f`$ noise. We now describe the calculation leading to this conclusion . One starts with the Hamiltonian (1) but assumes that the Coulomb interactions are turned on at time $`t=0`$: $$H=\underset{i}{}\varphi _in_i+\underset{i>j}{}\frac{e^2}{\kappa r_{ij}}n_in_j\theta (t)$$ (14) where the step function $`\theta (t)`$ is 0 for $`t<0`$ and 1 for $`t0`$. So for $`t<0`$ the noninteracting density of states is a constant $`g_o`$. Once the interactions are turned on, one follows the subsequent time development of the Coulomb gap. The Coulomb gap arises because the stability of the ground state with respect to single electron hopping from an occupied site $`i`$ to an unoccupied site $`j`$ requires $`\mathrm{\Delta }_i^j>0`$. So we need to subtract from the density of states those states which violate this stability condition. This leads to a self–consistent equation for the density of states : $`g(\stackrel{~}{\epsilon },`$ $`t)`$ $`=g_o{\displaystyle \underset{j>i}{}}(1a_o^3{\displaystyle _A^A}d\stackrel{~}{\epsilon }^{}g(\stackrel{~}{\epsilon }^{},t)\theta ({\displaystyle \frac{e^2}{\kappa r_{ij}}}+\stackrel{~}{\epsilon }\stackrel{~}{\epsilon }^{})`$ (16) $`F(n_i^{}=1,n_j^{}=0)\theta (t\tau _{ij}(\stackrel{~}{\epsilon }^{},\stackrel{~}{\epsilon },r_{ij})))`$ where the single–site energy $`\stackrel{~}{\epsilon }_i=\stackrel{~}{\epsilon }`$, $`\stackrel{~}{\epsilon }_j=\stackrel{~}{\epsilon }^{}`$, and $`a_o`$ is the lattice constant. $`n_i^{}=n_i+1/2`$; so $`n_i^{}=1`$ if site i is occupied and $`0`$ if site $`i`$ is unoccupied. $`F(n_i^{},n_j^{})`$ is the probability that donors $`i`$ and $`j`$ have occupation numbers $`n_i^{}`$ and $`n_j^{}`$, respectively, while all other sites have their ground state occupation numbers $`\stackrel{~}{n}_k^{}`$. $`\tau _{ij}^1`$ is the number of electrons which jump from site $`i`$ to site $`j`$ per unit time. $`\theta (t\tau _{ij})`$ represents the fact that at time $`t`$, the primary contributions to the change in the density of states will be from those hops for which $`\tau _{ij}<t`$ . In writing eq. (16), we assume that these hops together with phonons have equilibrated the system as much as is possible at time $`t`$. The hopping rate $`\tau _{ij}^1`$ is given by $$\tau _{ij}^1=\gamma _{ij}^o\mathrm{exp}(\frac{2r_{ij}}{a})[1+N(\mathrm{\Delta }_i^j)]F(n_i^{}=1,n_j^{}=0)$$ (17) where $`N(\mathrm{\Delta }_i^j)`$ is the phonon occupation factor and reflects the contribution of phonon assisted hopping. We are also allowing for spontaneous emission of phonons since we are considering a nonequilibrium situation in which electrons hop in order to lower their energy. Following we can rewrite the self–consistent equation $`g(\stackrel{~}{\epsilon },t)`$: $`g(\stackrel{~}{\epsilon },t)=g_o\mathrm{exp}\{{\displaystyle \frac{1}{2}}{\displaystyle _A^A}d\stackrel{~}{\epsilon }^{}g(\stackrel{~}{\epsilon }^{},t){\displaystyle _{a_o}^{\mathrm{}}}dr4\pi r^2`$ (18) $`F(n(\stackrel{~}{\epsilon })=1,n(\stackrel{~}{\epsilon }^{})=0)\theta ({\displaystyle \frac{e^2}{\kappa r}}+\stackrel{~}{\epsilon }\stackrel{~}{\epsilon }^{})\theta (t\tau (\stackrel{~}{\epsilon }^{},\stackrel{~}{\epsilon },r))\}`$ (19) At low energies large distances play an important role and so we have replaced the sum by an integral over $`r`$ in the exponent. The origin is at site $`i`$. $`n(\stackrel{~}{\epsilon })`$ is the occupation probability of a site with energy $`\stackrel{~}{\epsilon }`$. $`\tau (\stackrel{~}{\epsilon }^{},\stackrel{~}{\epsilon },r)`$ is given by (17) with $`r_{ij}`$ replaced by $`r`$, $`\stackrel{~}{\epsilon }_i`$ replaced by $`\stackrel{~}{\epsilon }`$, and $`\stackrel{~}{\epsilon }_j`$ replaced by $`\stackrel{~}{\epsilon }^{}`$. Since it is not clear how the stability condition $`\mathrm{\Delta }_i^j>0`$ can be applied to finite temperatures, we confine our calculations to the case of $`T=0`$. In this case the phonon occupation factor $`N(\mathrm{\Delta }_i^j)=0`$ and the electron occupation factor $`F(n_i=1,n_j=0)=1`$ if $`\stackrel{~}{\epsilon }_i<0`$ and $`\stackrel{~}{\epsilon }_j>0`$. Otherwise $`F(n_i=1,n_j=0)=0`$. We can solve eq. (19) iteratively on the computer. After a few iterations the typical difference between successive iterations is typically less than 1 part in $`10^5`$. We find that the Coulomb gap develops slowly over many decades in time . After an infinite amount of time, the density of states at the Fermi energy $`\mu `$ goes to zero and $`g(\stackrel{~}{\epsilon })\stackrel{~}{\epsilon }^2`$. The functional form of the time dependence of $`g(\stackrel{~}{\epsilon },t)`$ varies with the energy $`\stackrel{~}{\epsilon }`$ and with $`g_o`$. For conduction noise we are interested in the time dependence of the density of states at the Fermi energy which is shown in Figure 1 for $`10^8`$ s $`<t<10^8`$ s. For $`g_o=2\times 10^5`$ states/K–Å<sup>3</sup>, we can fit our results to the form $`g(\mu ,t)=B_1\mathrm{ln}(t_o/t)`$ where $`t<t_o`$, $`t_o3\times 10^{43}`$ sec, and for $`g_o=6.25\times 10^5`$ states/K–Å<sup>3</sup>, $`g(\mu ,t)B_2t^\theta `$ where $`\theta 0.05`$. The values of $`B_1`$, $`B_2`$, and the other parameters used to obtain these results are given in the caption of Figure 1. These fits change slightly for longer times. For example, for $`10^8`$ s $`<t<10^{100}`$ s and $`g_o=2\times 10^5`$ states/K–Å<sup>3</sup>, the fit to our results has the form $`g(\mu ,t)t^\theta `$ where $`\theta 0.01`$. Still we see that the density of states at the Fermi energy approaches its equilibrium value roughly logarithmically in time. This is consistent with recent experiments on thin semiconducting and metallic films which have shown that the system adjusted to changes in the Fermi energy approximately logarithmically in time. These films were grown on insulating substrates which separated them from a gate electrode that regulated the electron density, and hence the chemical potential, of the film. The conductance was measured as a function of the gate voltage. If the gate voltage was changed suddenly from, say, $`V_o`$ to $`V_1`$, the conductance had a very fast initial rise, followed by a period of rapid relaxation, which in turn was followed by a long period of very slow relaxation. The relaxation could be described by $`\mathrm{ln}t`$ or $`t^\theta `$ with $`\theta `$ being small and varying slowly with time. This is consistent with our view that when the gate voltage is changed, the Fermi energy changes, and time dependent relaxations arise because the system must dig a new Coulomb gap in the density of states at the new Fermi energy. So both theory and experiment indicate that the nonequilibrium density of states approaches its equilibrium value roughly logarithmically in time. Returning to the original model described by (1), we assume that this time dependence holds true in the linear response regime at low temperatures. If a fluctuation $`\delta g(\mu ,t=0)`$ at $`t=0`$ pushes the density of states away from its mean equilibrium value at the Fermi energy, then this perturbation will decay according to $`f(t)`$ which enters into eqs. (10) and (11). Our nonequilibrium calculation indicates that $`f(t)`$ can have the form: $$f_1(t)=B_1\mathrm{ln}(\frac{t_o}{t})$$ (20) where $`t<t_o`$, and $`t_o`$ is on the order of the age of the universe or longer, or $$f_2(t)=B_2t^\theta $$ (21) where $`\theta 1`$, and $`B_1`$ and $`B_2`$ are positive constants of order $`g_o`$. In both cases $`t`$ is greater than some $`t_{\mathrm{min}}`$ of order 10<sup>-8</sup> s, say. The time dependence is a function of the energy, so here we set $`\epsilon =\mu `$. Fourier transforming $`f_1(t)`$ and keeping the real part, we find that $$f_1(\omega )\frac{\pi B_1}{2}\frac{1}{\omega }$$ (22) This implies that the noise spectral density $`S(\omega )1/\omega `$. Fourier transforming $`f_2(t)`$ and keeping the real part yields $$f_2(\omega )\frac{\pi B_2\theta }{2}\frac{1}{\omega ^{1\theta }}$$ (23) for $`\theta 1`$. This implies $`S(\omega )1/\omega ^{1\theta }`$. To summarize, electron hopping leads to fluctuations in the density of states that relax back to equilibrium roughly logarithmically in time. This leads to 1/f noise in the spectral density $`S(\omega )`$ of the noise in the resistivity. In particular we find that $`S(\omega )1/\omega ^\alpha `$ where $`\alpha =1`$ if the relaxation is logarithmic in time, and $`\alpha =1\theta `$ if the relaxation is a power law that goes as $`t^\theta `$ where $`\theta 1`$. In general $`\alpha `$ depends on temperature and is weakly dependent on the noninteracting density of states $`g_o`$ and on the times scales. As eq. (13) indicates, the noise amplitude also depends on the temperature. Unfortunately we cannot ascertain these temperature dependences because we do not know the temperature dependence of the fluctuations $`\delta g(T,\stackrel{~}{\epsilon },t)`$ in the density of states. However we believe that our mechanism for 1/f noise should be valid at low temperatures ($`T\stackrel{<}{}20`$ K) where the logarithmic time dependence of the conductance is observed after the Coulomb glass has been pushed out of equilibrium by the sudden application of a gate voltage . I thank Hervé Carruzzo for helpful discussions. This work was supported in part by ONR grant N00014-00-1-0005 and CULAR funds provided by the University of California for the conduct of discretionary research by Los Alamos National Laboratory.
warning/0007/astro-ph0007110.html	ar5iv	text	# Bizarre Hard X-ray Outbursts of Cygnus X-1 ## 1 Introduction Cyg X-1 is a persistent X-ray source believed to be powered by accretion onto a black hole from a massive companion. Most of the time Cyg X-1 spends in the hard state and sometimes it switches to the soft state (see, e.g., Liang & Nolan 1984; Zhang et al. 1997; Poutanen 1998; Gierliński et al. 1999). The X-ray luminosity above $`2\mathrm{keV}`$ is estimated to be about $`3\times 10^{37}\mathrm{erg}\mathrm{s}^1`$ (Gierliński et al. 1997), assuming the distance to the source, $`D=2\mathrm{kpc}`$ (Massey, Johnson, & DeGioia-Eastwood 1995; Malysheva 1997). Variability of Cyg X-1 was extensively studied with various instruments (see, e.g., Ling et al. 1987; Gilfanov et al. 1995; Phlips et al. 1996; Kuznetsov et al. 1997; Paciesas et al. 1997; Wen et al. 1999; Brocksopp et al. 1999; Bałucińska-Church et al. 2000). These studies show the stability of the hard state of Cyg X-1. The photon flux above 50 keV is normally $`0.1\mathrm{ph}\mathrm{cm}^2\mathrm{s}^1`$ with variations by a factor 2. In this Letter, we analyze the period of the unusually strong hard X-ray activity of Cyg X-1 on 1999 April 21. Two strong outbursts activated the BATSE/CGRO (Fishman et al. 1989) onboard trigger. We analyze the temporal structure of the outbursts in the four LAD energy channels 1–4, estimate the peak luminosity of Cyg X-1, and discuss possible implications for theoretical models of accretion in this object. ## 2 The Outbursts The unusual activity of Cyg X-1 can be traced back to 1999 April 19 (TJD 11287). Two events recorded on April 19 were found in the BATSE data during search for non-triggered gamma-ray bursts (GRBs) (Stern et al. 1999, 2000a). The best fit locations were within $`4^\mathrm{o}5^\mathrm{o}`$ from Cyg X-1 ($`1\sigma `$ errors exceeded $`10^\mathrm{o}`$ for those events). We use the location procedure described in Stern et al. (2000b) which is similar to that used for GRB location by Pendleton et al. (1999). The estimated peak fluxes above 50 keV were about 0.3 and $`0.5\mathrm{ph}\mathrm{cm}^2\mathrm{s}^1`$ which is high as compared with the normal flux from Cyg X-1. On April 20, BATSE detectors were triggered by another event with right ascension, $`\alpha =305\stackrel{}{\mathrm{.}}1`$, and declination, $`\delta =26\stackrel{}{\mathrm{.}}2`$ (BATSE estimate, note also that this event is identified as a GRB in the BATSE data base). Again, the location was close to Cyg X-1 (slightly beyond 1$`\sigma `$ error circle).<sup>1</sup><sup>1</sup>1The coordinates of Cyg X-1 are $`\alpha =299\stackrel{}{\mathrm{.}}5`$ and $`\delta =35\stackrel{}{\mathrm{.}}2`$. The peak flux of this event was about $`0.5\mathrm{ph}\mathrm{cm}^2\mathrm{s}^1`$. On April 21, the two brightest outbursts occurred, with the interval $`2.5`$ hours. Between and after the two outbursts, Cyg X-1 demonstrated nothing unusual, being in the hard state with normal luminosity. The summary of all the five events is given in Table 1. We will concentrate on the two brightest events on April 21 (two last lines in Table 1). The locations of both events coincide with the location of Cyg X-1 with 2<sup>o</sup> accuracy. In principle, it might be accidental projections of GRBs on Cyg X-1. The probability that one of $`10^3`$ detected GRBs with comparable brightness will appear within 2<sup>o</sup> from Cyg X-1 is about 0.3. However, the probability of appearance of two such GRBs from the same location within 2.5 hours is low, $`10^5`$. Besides, the long duration of the events, $`10^3`$ s, is very unusual for GRBs. One concludes that the events are outbursts of Cyg X-1. TABLE 1 Summary of Outbursts on 1999 April 19 - 21 ST2000<sup>a</sup> #<sup>b</sup> time<sup>c</sup>, s $`\alpha ^\mathrm{d}`$ $`\delta ^\mathrm{d}`$ $`1\sigma ^\mathrm{d}`$ 11287a 16816 297.8 38.2 18.4 11287b 18071 296.6 38.0 12.9 11288e 7523 83437 305.6 27.6 6.0 11289c 7524 54530 299.9 34.4 2.4 11289e 7525 63100 298.4 35.2 1.2 <sup>a</sup> Name in the catalog of Stern et al. (1999), consisting of TJD and a letter. <sup>b</sup> BATSE trigger number. <sup>c</sup> Time of the day in seconds. <sup>d</sup> The best fit location and its $`1\sigma `$ error (deg). ### 2.1 Outburst 11289c The first outburst on April 21 has similar light curves in LAD energy channels 1–3 (see top panel in Fig. 1). The signal in channel 4 (above 300 keV) was low, practically undetectable. The signal was cut off by Earth occultation at $`55.23`$ ks. We model the background during the outburst using a linear fit extrapolating the background after the occultation. In order to get the photon and energy fluxes, we fit the signal count rate in the three energy channels using the detector response matrix computed with the code of Pendleton (see Pendleton et al. 1999 and references therein). An exponentially cutoff power-law is assumed as a spectral hypothesis. The photon and energy fluxes are then obtained by integrating the best fitting spectrum in the proper energy band. The error is estimated to be about 15%. (Note that the overall systematic error in the BATSE flux normalization may be about 20%, see, e.g., Much et al. 1996.) TABLE 2 Peak Fluxes and Luminosities Outburst $`F_{>50}`$<sup>a</sup> $`F_{>30}`$<sup>b</sup> $`L_{>50}`$<sup>c</sup> $`L_{>30}`$<sup>d</sup> 11289c 1.1 2.4 $`0.75\times 10^{38}`$ $`1.2\times 10^{38}`$ 11289e 1.3 3.9 $`0.8\times 10^{38}`$ $`1.6\times 10^{38}`$ <sup>a</sup> Peak flux ($`\mathrm{ph}\mathrm{cm}^2\mathrm{s}^1`$) above 50 keV. <sup>b</sup> Peak flux above 30 keV. <sup>c</sup> Peak luminosity ($`\mathrm{erg}\mathrm{s}^1`$) above 50 keV (assuming distance $`D=2`$ kpc). <sup>d</sup> Peak luminosity above 30 keV. The peak flux in 50–300 keV band (channels 2–3) is about $`1.1\mathrm{ph}\mathrm{cm}^2\mathrm{s}^1`$ which is $`10`$ times higher than the normal flux from Cyg X-1. The flux averaged over hundred seconds (54.55-54.65 ks interval) is $`0.75\mathrm{ph}\mathrm{cm}^2\mathrm{s}^1`$. The peak luminosity above 30 keV exceeds $`10^{38}\mathrm{erg}\mathrm{s}^1`$ (see Table 2 for details). Note that the strong increase in the luminosity was not accompanied by crucial changes in the spectrum. At least, the hardness ratios during the outburst stayed similar to the normal hard state of Cyg X-1 (see § 3). ### 2.2 Outburst 11289e The second outburst on April 21 was detected $``$ 8 ks after the first outburst (see bottom panel in Fig. 2.1). Unfortunately, there is a gap in the data at the time when Cyg X-1 rose above the horizon ($``$62.8 ks). The data records start at $`62.9`$ ks. Note that there was no occultation of Cyg X-1 in the end of the outburst. The sharp fall off in the light curves in each channel is the intrinsic behavior of the source. It allows one to estimate the peak flux and luminosity of the outburst which we give in Table 2. The corresponding luminosities averaged over interval E (see Fig. 1) are 20% smaller than the peak values. The ionospheric background showed smooth latitude variations during the outburst (see bottom panel of Fig. 1) on the time-scale of order 1 ks. The background level is different in different detectors, therefore, it is difficult to extract exactly the signal seen in the detectors looking at Cyg X-1. The signal uncertainty caused by the background variations is of order 20%. One can get a good location fit for the signal by analyzing the variable part of the light curve. At a given time interval, we fit the light curve by a linear function and take it as a “reference” level. (It thus includes both the background and the linear part of the true signal.) We then take the difference between the count rate (at this or another interval) and the reference level as a “signal”. The location fits demonstrate that almost all variable part of the signal comes from the direction of Cyg X-1 with $`5^\mathrm{o}`$ accuracy at all time intervals A-F marked in Figure 1 (see Table 3). The residual signal can be fitted by 200 s – 400 s long linear fragments with a good $`\chi ^2`$ (lines 1-3 in Table 3). This means that no fraction of the variable signal can be mimicked by magnetospheric phenomena like particle precipitation which either produces a diffuse flux of photons or can be a localized source at some distance from the spacecraft. In the first case the location fit is so bad that the residual $`\chi ^2`$ is comparable to the initial one, and in the second case one observes a fast change in the direction of the source due to the satellite motion. In our case, the satellite moved $`6000`$ km between intervals B and E, while the source’s position on the sky did not change. TABLE 3 Location Fits for Different Time Intervals Signal<sup>a</sup> Reference<sup>b</sup> Channel $`\alpha ^\mathrm{c}`$ $`\delta ^\mathrm{c}`$ $`1\sigma ^\mathrm{c}`$ $`\chi _0^2/\chi _r^2`$ <sup>d</sup> B B 2+3 302.1 37.1 6.6 11.4 / 0.93 D+E D+E 1 302.3 38.4 2.0 34.0 / 1.13 D+E D+E 2+3 301.3 40.2 6.0 16.9 / 1.05 D+E G 2+3 298.4 35.2 1.2 447 / 10.8 E E 1 303.9 40.1 4.9 15.3 / 0.90 E+F G 3 298.8 33.5 2.1 64.4 / 2.35 F G 2+3 302.4 40.3 4.1 51.0 / 1.66 <sup>a</sup> The time interval in which the location fit is done for the residuals after subtracting a linear (“reference”) function from the count rate. <sup>b</sup> The interval at which the linear reference function was fitted. <sup>c</sup> The best fit location and its $`1\sigma `$ error (deg). <sup>d</sup> $`\chi _0^2`$ and $`\chi _r^2`$ are $`\chi ^2`$ per dof for the count rate relative to the reference level and for the residual signal after the location fit, respectively. While the first April 21 outburst is unusual only for its strength, the second one is even more surprising. The behavior of the light curve in different energy channels suggests the presence of two independent emission components. The first (highly variable) component dominates the signal in channels 1 and 2, and the second component (with low amplitude of variability) dominates the signal in channel 3. The soft component terminated at 64.05 ks and then one observes only the hard component in the three channels. Unfortunately, the spectral data for this outburst were lost, which did not allow us to confirm the two-components by spectral analysis. Since we cannot exactly separate the signal from the background, we study the outburst using the strong non-Poisson variability produced by the signal in the count rate. In each time interval, A-F, we compute the root-mean-square (rms) of the signal variations, $`S_i`$, around the linear fits ($`i=1,2,3`$ is the channel number). We then compute the “fluctuation hardness ratio”, fhr, that is the ratio of the rms in different channels. The rms (with subtracted Poisson component) and $`\mathrm{fhr}_{21}\mathrm{rms}_2/\mathrm{rms}_1`$ are given in Table 4. Assuming the rms to be roughly proportional to the average level of the signal, the fhr may give an estimate for the true hardness ratio. The decrease of the $`\mathrm{fhr}_{21}`$ as the outburst progresses indicates that the variable component gets softer. TABLE 4 Characteristics of the Count Rate Variability Interval $`\mathrm{rms}_1`$ $`\mathrm{rms}_2`$ $`\mathrm{rms}_3`$ $`c_{23}`$ $`\mathrm{fhr}_{21}`$ A 97 109 70 0.63 $`\pm `$ 0.09 1.12 $`\pm `$ 0.06 B 300 316 169 0.92 $`\pm `$ 0.02 1.05 $`\pm `$ 0.02 C 211 213 96 0.77 $`\pm `$ 0.05 1.01 $`\pm `$ 0.03 D 499 370 100 0.71 $`\pm `$ 0.04 0.74 $`\pm `$ 0.01 E 568 286 90 0.15 $`\pm `$ 0.07 0.50 $`\pm `$ 0.02 F 109 196 106 0.92 $`\pm `$ 0.07 1.32 $`\pm `$ 0.10 The $`1\sigma `$ errors of $`c_{23}`$ and $`\mathrm{fhr}_{21}`$ are of Poisson nature; they are estimated with the bootstrap method. The measured hardness ratios are not affected by the background uncertainty because they are based on the signal dispersion in a narrow time interval where the background is almost linear. We then compute the cross-correlation coefficient between channels 2 and 3, $`c_{23}(S_2S_3S_2S_3)/(\mathrm{rms}_2\mathrm{rms}_3`$) (see Table 4). The cross-correlation decreases during the outbursts and becomes very low in interval E. It confirms that the soft component gets so soft that it practically does not contribute to the signal in channel 3. In interval F, the cross-correlation is high, confirming that the soft component has disappeared and we see only the second (hard) component in all the three channels. The luminosities of the soft and hard components can be roughly estimated from the step-like cut offs at 64.05 ks and 64.1 ks, respectively. The average luminosity of Cyg X-1 in time interval E above 50 keV (30 keV) is $`L_{>50}2.5\times 10^{37}\mathrm{erg}\mathrm{s}^1`$ ($`L_{>30}7.0\times 10^{37}\mathrm{erg}\mathrm{s}^1`$) in the soft component and $`L_{>50}4.0\times 10^{37}\mathrm{erg}\mathrm{s}^1`$ ($`L_{>30}5.5\times 10^{37}\mathrm{erg}\mathrm{s}^1`$) in the hard component. ## 3 Comparison with the Normal Hard State of Cyg X-1 The normal hard state of Cyg X-1 was studied using the Earth occultations which occur on each orbit of the CGRO. The step in the light curve at the moment of occultation shows the amplitude of the signal from Cyg X-1 in each energy channel. The study of Cyg X-1 with this method was done by Ling et al. (1997) and Paciesas et al. (1997). We performed a similar analysis for $`60`$ occultations and evaluated the flux from Cyg X-1 and the hardness ratio $`\mathrm{hr}_{32}`$ (ratio of the signal count rates in channels 3 and 2) for each occultation. The results are shown in Figure 2. The observed large dispersion in the hardness ratio may be caused by the measurement errors rather than intrinsic variability of Cyg X-1. To estimate the errors we did a similar analysis of 80 occultations of the Crab nebula which is known to be a steady source. We found $`\mathrm{hr}_{32}=0.66\pm 0.13`$ and the count rate $`316\pm 110`$, where the errors represent the standard deviations. The major cause of the errors is the uncertainty in the changing background (to have good statistics one has to fit the signal on a relatively long period $`100`$ s before and after the occultation, and the background curvature plays a role). One should note that the error in the hardness ratio of Cyg X-1 is likely to be smaller than that for the Crab since its flux is larger. In the outbursts, these errors are even smaller. For comparison we also show the estimates of the flux and hardness ratio for the two April 21 outbursts. Note that outburst 11289c and interval F of outburst 11289e have hardness ratios similar to the normal hard state. Interval E of outburst 11289e is significantly softer. The presence of two emission components (soft, highly variable, and hard, with lower variability) in outburst 11289e is intriguing. Are they present in the normal state of Cyg X-1? In the sample of 63 occultations we found that the fractional rms of the count rate (i.e. the ratio of the rms to the mean count rate) in channel 2 was higher than that in channel 3 in 47 cases. This indicates higher variability in channel 2 (the probability of such an accidental excess is $`0.510^4`$). The decreasing of the variability with energy was also observed in the 2–40 keV band (Nowak et al. 1999; Revnivtsev et al. 2000). These facts are consistent with the presence of a hard component with a low variability though this interpretation is not unique. Note also that Gierliński et al. (1997) got the best fit to the broad-band X-ray spectrum of Cyg X-1 with two thermal Comptonization components of different temperatures. ## 4 Discussion The hard state of Cyg X-1 has been a puzzle since its discovery. The standard accretion disk model (Shakura & Sunyaev 1973) was not able to explain the X-ray spectrum, and two modifications of the model were suggested: a two-temperature hot disk and an active corona atop the standard disk (see Beloborodov 1999a for a recent review). The advective hot-disk models (e.g. Esin et al. 1998) are consistent with the observed spectrum and luminosity of Cyg X-1 if the accretion rate has a specific value $`\dot{M}\dot{M}_{\mathrm{max}}`$ $`10\alpha ^2L_E/c^2`$ and the viscosity parameter $`\alpha 0.20.3`$. Here $`L_\mathrm{E}=`$ $`4\pi cGMm_p/\sigma _T=`$ $`1.3\times 10^{39}`$ $`(M/10M_{})\mathrm{erg}\mathrm{s}^1`$ is the Eddington luminosity. Small variations in the accretion rate, $`\mathrm{\Delta }\dot{M}/\dot{M}10\%`$, were predicted to destroy the hard state (Esin et al. 1998; but see also Zdziarski 1998). By contrast, the luminosity of Cyg X-1 is known to vary by a factor of two without substantial changes in the spectral shape (e.g., Paciesas et al. 1997; Gierliński et al. 1997). Such fluctuations already challenged the model, and the hard outbursts analyzed in this Letter are even more difficult to explain. The model would work only assuming a specific dependence $`\alpha \dot{M}^{1/2}`$ that keeps $`\dot{M}\dot{M}_{\mathrm{max}}`$; $`\alpha `$ then should be about unity at the peak of the outburst. In the context of the disk-corona model, the outbursts can be interperted as an enhanced coronal activity of the accretion disk. In this model, the X-ray spectral slope is controlled by one parameter, the feedback factor due to X-ray reprocessing by the disk. The hard-state spectrum of Cyg X-1 is well explained if the coronal plasma is ejected away from the disk with a mildly relativistic velocity $`\beta =v/c0.3`$ (Beloborodov 1999b; Malzac, Beloborodov, & Poutanen 2001). Alternatively, the observed emission may be produced by a static corona atop a strongly ionized disk (e.g., Ross, Fabian, & Young 1999; Nayakshin 1999). The corona becomes $`e^\pm `$–dominated at high luminosities and its temperature decreases (e.g., Svensson 1984; Stern et al. 1995; Poutanen & Svensson 1996). Pair creation may cause the shift of the spectral break in outburst 11289e to smaller energies. Cyg X-1 is a massive X-ray binary and it may be fed mainly by the donor wind. Then the pattern of accretion can change completely compared to the standard viscous $`\alpha `$disk or its modifications. The captured wind matter has a low angular momentum, just about critical for disk formation (Illarionov & Sunyaev 1975). Under such conditions, a small-scale inviscid disk forms, which accretes super-sonically (Beloborodov & Illarionov 2001). The disk forms in the ring-like caustic of the accretion flow where energy is liberated in inelastic collision of gas streams. A Comptonized power-law spectrum is then emitted with a standard break at $`100`$ keV, and it has appearance of a normal hard state of Cyg X-1. The two-component emission in outburst 11289e probably requires a two-zone emission model. For instance, the soft component may be associated with a variable coronal emission atop the disk, and the hard component with an inner relativistic jet. A corona-jet model was recently proposed by Brocksopp et al. (1999) based on correlations between the X-rays and the radio emission observed in Cyg X-1. The April 1999 outbursts were, however, too short to have a significant impact on the flux in radio and no substantial changes were detected (R. Fender and G. Pooley, private communication). In the context of the wind-fed accretion model, the two-component emission can naturally appear if viscous accretion disk and wind-disk collision are both operating in the source. Concluding, the 1999 April 21 outbursts are the brightest events detected from Cyg X-1 by BATSE during 9 years of its operation. The hard X-ray luminosity above 30 keV was in excess of $`10^{38}\mathrm{erg}\mathrm{s}^1`$. Unfortunately, there are no simultaneous observations in the soft X-rays, so that one can only guess what was the total luminosity. The large luminosity and unusual spectral behavior challenge the theoretical models of accretion in Cyg X-1. We thank Rob Fender, Guy Pooley, and Jerry Bonnell for useful discussions. We are grateful to Rob Preece and Geoff Pendleton for providing the code for computing the detector response matrix and an anonymous referee for useful comments that significantly improved the paper. This research has made use of data obtained through the High Energy Astrophysics Science Archive Research Center Online Service, provided by the NASA/Goddard Space Flight Center. This work was supported by the Swedish Natural Science Research Council, the Swedish Royal Academy of Sciences, the Wenner-Gren Foundation for Scientific Research, the Anna-Greta and Holger Crafoord Fund, and RFBR grant 00-02-16135.
warning/0007/gr-qc0007034.html	ar5iv	text	# Hyperbolic formulations and numerical relativity II: Asymptotically constrained systems of the Einstein equations ## I Introduction Numerical relativity, an approach to solve the Einstein equations numerically, is supposed to be the only way to study highly non-linear gravitational phenomena. Although the attempt has already decades of history, we still do not have a definite recipe for integrating the Einstein equations that will give us us accurate and long-term stable time evolutions. Here and hereafter, we mean “stable evolution” that the system keeps the violation of the constraints within a suitable small value in its free numerical evolution. As the authors discussed in our preceding paper (Paper I) , one direction for obtaining a more stable system is to apply a set of dynamical equations which have manifest hyperbolic form (or first-order form). The standard Arnowitt-Deser-Misner (ADM) formulation does not have this feature, but there are many alternative proposals for constructing a hyperbolic set of equations (, see also references in ). However, we showed in Paper I that a symmetric hyperbolic form (mathematically, the “ultimate” level of hyperbolicity) does not necessary give the best performance for stable numerical evolution compared with weakly and strongly hyperbolic systems. This experiment was performed using Ashtekar’s connection variables, since this formulation enables us to compare three levels of hyperbolic formulations keeping the same fundamental dynamical variables. In this article, we discuss different (but somewhat related) approaches to obtaining stable evolution of the Einstein equations. The idea is to construct a system robust against the perturbative error produced during numerical time integration. We discuss the following two systems. The first one is the so-called “$`\lambda `$-system”, which was proposed originally by Brodbeck, Frittelli, Hübner and Reula (BFHR) . The idea of this approach is to introduce additional variables, $`\lambda `$, which indicates the violation of the constraints, and to construct a symmetric hyperbolic system for both the original variables and $`\lambda `$ together with imposing dissipative dynamical equations for $`\lambda `$s. BFHR constructed their $`\lambda `$-system based on Frittelli-Reula’s symmetric hyperbolic formulation of the Einstein equations , and we have also presented a similar system for Ashtekar’s connection formulation based on its symmetric hyperbolic expression . In §II, we review this system and present numerical examples which show this system behaves as expected. The second one has the same motivation but turns to be more practical, which we call “adjusted-system”. The essential procedure is to to add constraint terms to the right-hand-side of the dynamical equations with multipliers, and to choose the multipliers so as to decrease the violation of the constraint equations. This second step will be explained by obtaining non-positive or pure-imaginary eigenvalues of the adjusted constraint propagation equations. We remark that adjusting the dynamical equation using the constraints is not a new idea. This can be seen for example in a remedial ADM system by Detweiler , in a conformally decoupled trace-free re-formulation of ADM by Nakamura et al , and also in constructing hyperbolic formulations . We also remark that this eigenvalue criterion is also the core part of the theoretical support of the above $`\lambda `$-system. In §III, we describe this approach and present numerical examples again in the Maxwell system and in the Ashtekar’s system. This “adjusted-system” does not change the number of dynamical variables, and does not require hyperbolicity in the original set of equations. Therefore we think our results promote further applications in numerical relativity. We will not repeat our explanation of Ashtekar’s connection formulation in our notation, nor of our detailed numerical procedures, since they are described in our Paper I . ## II Asymptotically constrained system 1: $`\lambda `$-system We begin by reviewing the fundamental procedures of the “$`\lambda `$-system” proposed by Brodbeck, Frittelli, Hübner and Reula (BFHR) . We, then, demonstrate how this system works in Maxwell’s equations, and Ashtekar’s connection formulation of the Einstein equations in the following subsections. ### A The “$`\lambda `$ system” The actual procedures for constructing a $`\lambda `$ system are followings. 1. Prepare a symmetric hyperbolic evolution system which describe the problem; say $$_tu^\gamma =A^{i\gamma }{}_{\delta }{}^{}_{i}^{}u^\delta +B^\gamma ,$$ (1) where $`u^\gamma `$ ($`\gamma =1,\mathrm{},N`$) is a set of dynamical variables, $`A(u(x^i))`$ forms a symmetric matrix (Hermitian matrix when $`u`$ is complex variables) and $`B(u(x^i))`$ is a vector, where $`A`$ and $`B`$ do not include any further spatial derivatives in these components. The system may have constraint equations, which should be the first class. Ideally, we expect that the evolution equation of the set of constraints $`C^\rho `$ $`(\rho =1,\mathrm{},M)`$, which hereafter we denote constraint propagation equation, forms a first order hyperbolic system (cf. ), say $$_tC^\rho =D^{i\rho }{}_{\sigma }{}^{}_{i}^{}C^\sigma +E^\rho {}_{\sigma }{}^{}C_{}^{\sigma },$$ (2) (where $`D,E`$ are the same with $`A,B`$ above) but this hyperbolicity may not be necessary. 2. Introduce $`\lambda ^\rho `$ as a measure of violation of the constraint equation, $`C^\rho 0`$. ($``$ denotes “weakly equal”.) Here $`C^\rho `$ is a given function of $`u`$ and is assumed to be linear in its first-order space derivatives. We impose that $`\lambda `$ obeys a dissipative equation of motion $$_t\lambda ^\rho =\alpha _{(\rho )}C^\rho \beta _{(\rho )}\lambda ^\rho \text{ (we do not sum over }\rho \text{ and }(\rho )\text{ on right hand side)}$$ (3) with the initial data $`\lambda ^\rho =0`$, and by setting $`\alpha 0,\beta >0`$. We remark that $`\lambda ^\rho `$ remains zero during the time evolution if there is no violation of the constraints. 3. Take a set $`(u,\lambda )`$ of dynamical variables, and modify the evolution equations so as to form a symmetric hyperbolic system. That is, the set of equations $$_t\left(\begin{array}{c}u^\gamma \\ \lambda ^\rho \end{array}\right)\left(\begin{array}{cc}A^{i\gamma }_\delta & 0\\ F^{i\rho }_\delta & 0\end{array}\right)_i\left(\begin{array}{c}u^\delta \\ \lambda ^\sigma \end{array}\right),$$ (4) ($``$ means that we have extracted only the term which appears in the principal part of the system) can be modified as $$_t\left(\begin{array}{c}u^\gamma \\ \lambda ^\rho \end{array}\right)\left(\begin{array}{cc}A^{i\gamma }_\delta & \overline{F}^i{}_{\sigma }{}^{}^\gamma \\ F^{i\rho }_\delta & 0\end{array}\right)_i\left(\begin{array}{c}u^\delta \\ \lambda ^\sigma \end{array}\right),$$ (5) where the additional terms will not disturb the hyperbolicity of equations of $`u^\gamma `$, rather they make the whole system symmetric hyperbolic, which guarantees the well-posedness of the system. Therefore the derived system, (5), should have unique solution. If a perturbative violation of constraints, $`\lambda ^\rho 0`$, occurs during the evolution, by choosing appropriate $`\alpha `$s and $`\beta `$s in (3), $`\lambda `$s can be made decaying to zero, which means the total system evolves into the constraint surface asymptotically. We note that this procedure requires that the original system $`u`$ forms a symmetric hyperbolic system, so that applications to the Einstein equations are somewhat restricted. BFHR constructed this $`\lambda `$-system using a Frittelli-Reula’s formulation . We also applied this system to the symmetric hyperbolic version of Ashtekar’s formulation . We next review a brief proof why the system (5) ensures that the evolution is constrained asymptotically. We first remark again that we only consider perturbative violations of constraints in our evolving system. The steps are following. 1. Since we modify the equations for $`u^\gamma `$, the propagation equation of the constraints are also modified; write them schematically as $$_tC^\rho =D^{i\rho }{}_{\sigma }{}^{}_{i}^{}C^\sigma +E^\rho {}_{\sigma }{}^{}C_{}^{\sigma }+G^{ij\rho }{}_{\sigma }{}^{}_{i}^{}_j\lambda ^\sigma +H^{i\rho }{}_{\sigma }{}^{}_{i}^{}\lambda ^\sigma +I^\rho {}_{\sigma }{}^{}\lambda _{}^{\sigma }.$$ (6) 2. In order to see the asymptotic behaviors of $`(\lambda ^\rho ,C^\rho )`$, we write them using their Fourier components so that their evolution equations take an homogenous form. That is, we transform $`(\lambda ^\rho ,C^\rho )`$ to $`(\widehat{\lambda }^\rho ,\widehat{C}^\rho )`$ as $$\lambda (x,t)^\rho =\widehat{\lambda }(k,t)^\rho \mathrm{exp}(ikx)d^3k,C(x,t)^\rho =\widehat{C}(k,t)^\rho \mathrm{exp}(ikx)d^3k.$$ (7) Then we see the evolution equations (3) and (6) become $$_t\left(\begin{array}{c}\widehat{\lambda }^\rho \\ \widehat{C}^\rho \end{array}\right)=\left(\begin{array}{cc}\beta _{(\rho )}\delta _\sigma ^\rho & \alpha _{(\rho )}\delta _\sigma ^\rho \\ G^{ij\rho }{}_{\sigma }{}^{}k_{i}^{}k_j+iH^{i\rho }{}_{\sigma }{}^{}k_{i}^{}+I^\rho _\sigma & iD^{i\rho }{}_{\sigma }{}^{}k_{i}^{}+E^\rho _\sigma \end{array}\right)\left(\begin{array}{c}\widehat{\lambda }^\sigma \\ \widehat{C}^\sigma \end{array}\right)=:P\left(\begin{array}{c}\widehat{\lambda }^\sigma \\ \widehat{C}^\sigma \end{array}\right).$$ (8) 3. If all eigenvalues of this coefficient matrix $`P`$ have negative real part, a pair $`(\widehat{\lambda },\widehat{C})`$ evolves as $`\mathrm{exp}(\mathrm{\Lambda }t)`$ asymptotically where $`\mathrm{\Lambda }`$ is the diagonalized matrix of $`P`$, which indicates that the original variables $`(\lambda ,C)`$ evolves similarly. It would be best if we could determine the $`\alpha `$ and $`\beta `$ in such a way in general, but it is not possible. Therefore we extract the principal order of $`P`$ and examine the condition for $`\alpha `$ and $`\beta `$ so that $`P`$ only has negative (real) eigenvalues. We remark again that this procedure is justified when we only consider a perturbative error from the constraint surface. ### B Example 1: Maxwell equations As a first example, we present the Maxwell equations in a form of $`\lambda `$-system. The Maxwell equations form linear and symmetric hyperbolic dynamical equations, together with two constraint equations, which might be the best system to start with. #### 1 $`\lambda `$-system The Maxwell equations for an electric field $`E^i`$ and a magnetic field $`B^i`$ in the vacuum consist of two constraint equations, $`C_E`$ $`:=`$ $`_iE^i0,`$ (9) $`C_B`$ $`:=`$ $`_iB^i0,`$ (10) and a set of dynamical equations, $$_t\left(\begin{array}{c}E^i\\ B^i\end{array}\right)=\left(\begin{array}{cc}0& cϵ^i{}_{j}{}^{}^l\\ cϵ^i{}_{j}{}^{}^l& 0\end{array}\right)_l\left(\begin{array}{c}E^j\\ B^j\end{array}\right),$$ (11) which satisfies symmetric hyperbolicity. The constraint evolutions become $`_tC_E=0`$ and $`_tC_B=0`$, which indicate (trivial) symmetric hyperbolicity. According to the above procedure, we introduce $`\lambda `$s which obey $`_t\lambda _E`$ $`=`$ $`\alpha _1C_E\beta _1\lambda _E,`$ (12) $`_t\lambda _B`$ $`=`$ $`\alpha _2C_B\beta _2\lambda _B,`$ (13) with the initial data $`\lambda _E=\lambda _B=0`$ and take $`(E,B,\lambda _E,\lambda _B)`$ as a set of variables to evolve: $$_t\left(\begin{array}{c}E^i\\ B^i\\ \lambda _E\\ \lambda _B\end{array}\right)=\left(\begin{array}{cccc}0& cϵ^i{}_{j}{}^{}^l& 0& 0\\ cϵ^i{}_{j}{}^{}^l& 0& 0& 0\\ \alpha _1\delta ^l_j& 0& 0& 0\\ 0& \alpha _2\delta ^l_j& 0& 0\end{array}\right)_l\left(\begin{array}{c}E^j\\ B^j\\ \lambda _E\\ \lambda _B\end{array}\right)+\left(\begin{array}{c}0\\ 0\\ \beta _1\lambda _E\\ \beta _2\lambda _B\end{array}\right).$$ (14) We obtain immediately an expected symmetric form as $$_t\left(\begin{array}{c}E^i\\ B^i\\ \lambda _E\\ \lambda _B\end{array}\right)=\left(\begin{array}{cccc}0& cϵ^i{}_{j}{}^{}^l& \alpha _1\delta ^{li}& 0\\ cϵ^i{}_{j}{}^{}^l& 0& 0& \alpha _2\delta ^{li}\\ \alpha _1\delta _j^l& 0& 0& 0\\ 0& \alpha _2\delta _j^l& 0& 0\end{array}\right)_l\left(\begin{array}{c}E^j\\ B^j\\ \lambda _E\\ \lambda _B\end{array}\right)+\left(\begin{array}{c}0\\ 0\\ \beta _1\lambda _E\\ \beta _2\lambda _B\end{array}\right).$$ (15) #### 2 Analysis of eigenvalues Now the evolution equations for the constraints $`C_E`$ and $`C_B`$ become $`_tC_E=\alpha _1(\mathrm{\Delta }\lambda _E),_tC_B=\alpha _2(\mathrm{\Delta }\lambda _B)`$ (16) where $`\mathrm{\Delta }=_i^i`$. We take the Fourier integrals for constraints $`C`$s \[(16)\] and $`\lambda `$s \[(12), (13)\], in the form of (7), to obtain $`_t\left(\begin{array}{c}\widehat{C}_E\\ \widehat{C}_B\\ \widehat{\lambda }_E\\ \widehat{\lambda }_B\end{array}\right)`$ $`=`$ $`\left(\begin{array}{cccc}0& 0& \alpha _1k^2& 0\\ 0& 0& 0& \alpha _2k^2\\ \alpha _1& 0& \beta _1& 0\\ 0& \alpha _2& 0& \beta _2\end{array}\right)\left(\begin{array}{c}\widehat{C}_E\\ \widehat{C}_B\\ \widehat{\lambda }_E\\ \widehat{\lambda }_B\end{array}\right),`$ (17) where $`k^2=k_ik^i`$. We find the matrix is constant. Note that this is exact expression. Since the eigenvalues are $`(\beta _1\pm \sqrt{\beta _1^24\alpha _1^2k^2})/2`$ and $`(\beta _2\pm \sqrt{\beta _2^24\alpha _2^2k^2})/2`$, the negative eigenvalue requirement becomes $`\alpha _1,\alpha _20`$ and $`\beta _1,\beta _2>0`$. #### 3 Numerical demonstration We present a numerical demonstration of the above Maxwell “$`\lambda `$-system”. We prepare a code which produces electromagnetic propagation in $`xy`$-plane, and monitor the violation of the constraint during time integration. Specifically we prepare the initial data with a Gaussian packet at the origin, $`E^i(x,y,z)`$ $`=`$ $`(Aye^{B(x^2+y^2)},Axe^{B(x^2+y^2)},0),`$ (18) $`B^i(x,y,z)`$ $`=`$ $`(0,0,0),`$ (19) where $`A`$ and $`B`$ are constants, and let it propagate freely, under the periodic boundary condition. The code itself is quite stable for this problem. In Fig.1, we plot L2 norm of the error ($`C_E`$ over the whole grid) as a function of time. The solid line (constant) in Fig.1 (a) is of the original Maxwell equation. If we introduce $`\lambda `$s, then we see the error will be reduced by a particular choice of $`\alpha `$ and $`\beta `$. Fig.1 (a) is for changing $`\alpha `$ with $`\beta =2.0`$, while Fig.1 (b) is for changing $`\beta `$ with $`\alpha =0.5`$. Here, we simply use $`\alpha :=\alpha _1=\alpha _2`$ and $`\beta :=\beta _1=\beta _2`$. We see better performance for $`\beta >0`$ \[Fig.1 (b)\], which is the case of negative eigenvalues of the constraint propagation equation. We also see the system will diverge for large $`\alpha `$ \[Fig.1 (b)\]. The upper bound of $`\alpha `$ can be explained by the violation of the Courant-Friedrich-Lewy (CFL) condition, where the characteristic speed comes from the flux term of the dynamical equations (15). ### C Example 2: Einstein equations (Ashtekar equations) The second demonstration is of the vacuum Einstein equations in Ashtekar’s connection formalism . Before going through the $`\lambda `$-system, we will briefly outline the equations. The fundamental Ashtekar variables are the densitized inverse triad, $`\stackrel{~}{E}_a^i`$, and the SO(3,C) self-dual connection, $`𝒜_i^a`$, where the indices $`i,j,\mathrm{}`$ indicate the 3-spacetime, and $`a,b,\mathrm{}`$ is for SO(3) space. The total four-dimensional spacetime is described together with the gauge variables $`\underset{^{}}{𝑁},N^i,𝒜_0^a`$, which we call the densitized lapse function, shift vector and the triad lapse function. Since the Hilbert action takes the form $`S`$ $`=`$ $`{\displaystyle \mathrm{d}^4x[(_t𝒜_i^a)\stackrel{~}{E}_a^i+\underset{^{}}{𝑁}𝒞_H+N^i𝒞_{Mi}+𝒜_0^a𝒞_{Ga}]},`$ (20) the system has three constraint equations, $`𝒞_H𝒞_{Mi}𝒞_{Ga}0`$, which are called the Hamiltonian, momentum, and Gauss constraint equation, respectively. They are written as $`𝒞_H`$ $`:=`$ $`(i/2)ϵ^{ab}{}_{c}{}^{}\stackrel{~}{E}_{a}^{i}\stackrel{~}{E}_b^jF_{ij}^c,`$ (21) $`𝒞_{Mi}`$ $`:=`$ $`F_{ij}^a\stackrel{~}{E}_a^j,`$ (22) $`𝒞_{Ga}`$ $`:=`$ $`𝒟_i\stackrel{~}{E}_a^i,`$ (23) where $`F_{\mu \nu }^a:=2_{[\mu }𝒜_{\nu ]}^aiϵ^a{}_{bc}{}^{}𝒜_{\mu }^{b}𝒜_\nu ^c`$ is the curvature 2-form and $`𝒟_i\stackrel{~}{E}_a^j:=_i\stackrel{~}{E}_a^jiϵ_{ab}{}_{}{}^{c}𝒜_{i}^{b}\stackrel{~}{E}_c^j`$. The original dynamical equation for $`(\stackrel{~}{E}_a^i,𝒜_i^a)`$ constitutes a weakly hyperbolic form, $`_t\stackrel{~}{E}_a^i`$ $`=`$ $`i𝒟_j(ϵ^{cb}{}_{a}{}^{}\underset{^{}}{𝑁}\stackrel{~}{E}_c^j\stackrel{~}{E}_b^i)+2𝒟_j(N^{[j}\stackrel{~}{E}_a^{i]})+i𝒜_0^bϵ_{ab}{}_{}{}^{c}\stackrel{~}{E}_{c}^{i},`$ (24) $`_t𝒜_i^a`$ $`=`$ $`iϵ^{ab}{}_{c}{}^{}\underset{^{}}{𝑁}\stackrel{~}{E}_b^jF_{ij}^c+N^jF_{ji}^a+𝒟_i𝒜_0^a`$ (25) where $`𝒟_jX_a^{ji}:=_jX_a^{ji}iϵ_{ab}{}_{}{}^{c}𝒜_{j}^{b}X_c^{ji},`$ for $`X_a^{ij}+X_a^{ji}=0`$. It is also possible to express (24) and (25) to reveal symmetric hyperbolicity . For more detailed definitions and our notation, please see Appendix A of our Paper I . #### 1 $`\lambda `$-system for controlling constraint violations Here, we only consider the $`\lambda `$-system which controls the violation of the constraint equations. In , we have also discussed an advanced version of the $`\lambda `$-system which controls the violations of the reality condition. We introduce new variables ($`\lambda ,\lambda _i,\lambda _a`$), obeying the dissipative evolution equations, $`_t\lambda `$ $`=`$ $`\alpha _1𝒞_H\beta _1\lambda ,`$ (26) $`_t\lambda _i`$ $`=`$ $`\alpha _2\stackrel{~}{𝒞}_{Mi}\beta _2\lambda _i,`$ (27) $`_t\lambda _a`$ $`=`$ $`\alpha _3𝒞_{Ga}\beta _3\lambda _a,`$ (28) where $`\alpha _i0`$ (possibly complex) and $`\beta _i>0`$ (real) are constants. If we take $`y_\alpha :=(\stackrel{~}{E}_a^i,𝒜_i^a,\lambda ,\lambda _i,\lambda _a)`$ as a set of dynamical variables, then the principal part of (26)-(28) can be written as $`_t\lambda `$ $``$ $`i\alpha _1ϵ^{bcd}\stackrel{~}{E}_c^j\stackrel{~}{E}_d^l(_l𝒜_j^b),`$ (29) $`_t\lambda _i`$ $``$ $`\alpha _2[e\delta _i^l\stackrel{~}{E}_b^j+e\delta _i^j\stackrel{~}{E}_b^l](_l𝒜_j^b),`$ (30) $`_t\lambda _a`$ $``$ $`\alpha _3_l\stackrel{~}{E}_a^l.`$ (31) The characteristic matrix of the system $`u_\alpha `$ is not Hermitian. However, if we modify the right-hand-side of the evolution equation of ($`\stackrel{~}{E}_a^i,𝒜_i^a`$), then the set becomes a symmetric hyperbolic system. This is done by adding $`\overline{\alpha }{}_{3}{}^{}\gamma _{}^{il}(_l\lambda _a)`$ to the equation of $`_t\stackrel{~}{E}_a^i`$, and by adding $`i\overline{\alpha }{}_{1}{}^{}ϵ_{}^{a}{}_{c}{}^{}{}_{}{}^{d}\stackrel{~}{E}_{i}^{c}\stackrel{~}{E}_d^l(_l\lambda )+\overline{\alpha }{}_{2}{}^{}(e\gamma ^{lm}\stackrel{~}{E}_i^a+e\delta _i^m\stackrel{~}{E}^{la})(_l\lambda _m)`$ to the equation of $`_t𝒜_i^a`$. The final principal part, then, is written as $$_t\left(\begin{array}{c}\stackrel{~}{E}_a^i\\ 𝒜_i^a\\ \lambda \\ \lambda _i\\ \lambda _a\end{array}\right)\left(\begin{array}{ccccc}^l{}_{a}{}^{}{}_{}{}^{bi}_j& 0& 0& 0& \overline{\alpha }{}_{3}{}^{}\gamma _{}^{il}\delta _a^b\\ 0& 𝒩^l{}_{i}{}^{a}{}_{b}{}^{}^j& i\overline{\alpha }{}_{1}{}^{}ϵ_{}^{a}{}_{c}{}^{}{}_{}{}^{d}\stackrel{~}{E}_{i}^{c}\stackrel{~}{E}_d^l& \overline{\alpha }{}_{2}{}^{}e(\delta _i^j\stackrel{~}{E}^{la}\gamma ^{lj}\stackrel{~}{E}_i^a)& 0\\ 0& i\alpha _1ϵ_b{}_{}{}^{cd}\stackrel{~}{E}_{c}^{j}\stackrel{~}{E}_d^l& 0& 0& 0\\ 0& \alpha _2e(\delta _i^j\stackrel{~}{E}_b^l\delta _i^l\stackrel{~}{E}_b^j)& 0& 0& 0\\ \alpha _3\delta _a^b\delta _j^l& 0& 0& 0& 0\end{array}\right)_l\left(\begin{array}{c}\stackrel{~}{E}_b^j\\ 𝒜_j^b\\ \lambda \\ \lambda _j\\ \lambda _b\end{array}\right).$$ (32) where $`^{labij}`$ $`=`$ $`iϵ^{abc}\underset{^{}}{𝑁}\stackrel{~}{E}_c^l\gamma ^{ij}+N^l\gamma ^{ij}\delta ^{ab},`$ (33) $`𝒩^{labij}`$ $`=`$ $`i\underset{^{}}{𝑁}(ϵ^{abc}\stackrel{~}{E}_c^j\gamma ^{li}ϵ^{abc}\stackrel{~}{E}_c^l\gamma ^{ji}e^2\stackrel{~}{E}^{ia}ϵ^{bcd}\stackrel{~}{E}_c^j\stackrel{~}{E}_d^le^2ϵ^{acd}\stackrel{~}{E}_d^i\stackrel{~}{E}_c^l\stackrel{~}{E}^{jb}`$ (34) $`+e^2ϵ^{acd}\stackrel{~}{E}_d^i\stackrel{~}{E}_c^j\stackrel{~}{E}^{lb})+N^l\delta ^{ab}\gamma ^{ij},`$ (35) Clearly, the solution $`(\stackrel{~}{E}_a^i,𝒜_i^a,\lambda ,\lambda _i,\lambda _a)=(\stackrel{~}{E}_a^i,𝒜_i^a,0,0,0)`$ represents the original solution of the Ashtekar system. #### 2 Analysis of eigenvalues After linearizing and taking the Fourier transformation (7), the propagation equation of the constraints $`(𝒞_H,\stackrel{~}{𝒞}_{Mi},𝒞_{Ga})`$ and $`(\lambda ,\lambda _i,\lambda _a)`$ are written as, $`_t\left(\begin{array}{c}\widehat{𝒞}_H\\ \widehat{\stackrel{~}{𝒞}}_{Mi}\\ \widehat{𝒞}_{Ga}\\ \widehat{\lambda }\\ \widehat{\lambda }_i\\ \widehat{\lambda }_a\end{array}\right)`$ $`=`$ $`\left(\begin{array}{cccccc}0& ik_j& 0& 2\overline{\alpha }{}_{1}{}^{}k_{m}^{}k^m& 0& 0\\ ik_i& k_mϵ^m{}_{i}{}^{}^j& 0& 0& \overline{\alpha }{}_{2}{}^{}(k_ik^j+k_mk^m\delta _i^j)& 0\\ 0& 2\delta _a^b& ϵ^{mb}{}_{a}{}^{}k_{m}^{}& 2i\overline{\alpha }{}_{1}{}^{}k_{a}^{}& \overline{\alpha }{}_{2}{}^{}ϵ_{}^{amj}k_m& \overline{\alpha }{}_{3}{}^{}k_{m}^{}k^m\delta _a^b\\ \alpha _1& 0& 0& \beta _1& 0& 0\\ 0& \alpha _2\delta _i^j& 0& 0& \beta _2\delta _i^j& 0\\ 0& 0& \alpha _3\delta _a^b& 0& 0& \beta _3\delta _a^b\end{array}\right)\left(\begin{array}{c}\widehat{𝒞}_H\\ \widehat{\stackrel{~}{𝒞}}_{Mj}\\ \widehat{𝒞}_{Gb}\\ \widehat{\lambda }\\ \widehat{\lambda }_j\\ \widehat{\lambda }_b\end{array}\right),`$ (36) In order to link the discussion with our later demonstration in the plane symmetric spacetime, we here consider only the Fourier component of $`k_i=(1,0,0)`$ for simplicity. The eigenvalues, $`E_i(i=1,\mathrm{},14)`$, of the characteristic matrix of (36) can be written explicitly as $`(E_1,\mathrm{},E_{10})`$ $`=`$ $`(1/2)\beta _3\pm (1/2)\sqrt{\beta _3^24\|\alpha _3\|^2},`$ $`(1/2)(i+\beta _3)\pm (1/2)\sqrt{14\|\alpha _3\|^22i\beta _3+\beta _3^2},`$ $`(1/2)(i+\beta _3)\pm (1/2)\sqrt{14\|\alpha _3\|^22i\beta _3+\beta _3^2},`$ $`(1/2)(i+\beta _2)\pm (1/2)\sqrt{14\|\alpha _2\|^22i\beta _2+\beta _2^2},`$ $`(1/2)(i+\beta _2)\pm (1/2)\sqrt{14\|\alpha _2\|^22i\beta _2+\beta _2^2}`$ and as solutions ($`E_{11},\mathrm{},E_{14}`$) of the quartic equation $$x^4+(\beta _2+\beta _1)x^3+(2\|\alpha _1\|^2+2\|\alpha _2\|^2+1+\beta _1\beta _2)x^2+(2\|\alpha _2\|^2\beta _1+\beta _2+\beta _1+2\|\alpha _1\|^2\beta _2)x+(\beta _1\beta _2+4\|\alpha _1\|^2\|\alpha _2\|^2)=0,$$ (37) where $`\|\alpha _i\|^2=\alpha _i\overline{\alpha }_i`$. We omit the explicit expressions of $`E_{11},\mathrm{},E_{14}`$ in order to save space. A possible set of conditions on $`\alpha _\rho ,\beta _\rho ,(\rho =1,2,3)`$ for $`\mathrm{}e(E_i)<0`$ are $$\alpha _\rho 0\text{and}\beta _\rho >0.$$ (38) This is true (necessary and sufficient) for $`E_1,\mathrm{},E_{10}`$, and also plausible for $`E_{11},\mathrm{},E_{14}`$ as far as our numerical evaluation tells (see Fig.2). #### 3 Numerical demonstration In this subsection, we demonstrate that the $`\lambda `$-system for the Ashtekar equation actually works as expected. The model we present here is gravitational wave propagation in a planar spacetime under periodic boundary condition. We perform a full numerical simulation using Ashtekar’s variables. We prepare two $`+`$-mode strong pulse waves initially by solving the ADM Hamiltonian constraint equation, using York-O’Murchadha’s conformal approach. Then we transform the initial Cauchy data (3-metric and extrinsic curvature) into the connection variables, $`(\stackrel{~}{E}_a^i,𝒜_i^a)`$, and evolve them using the dynamical equations. For the presentation in this article, we apply the geodesic slicing condition (ADM lapse $`N=1`$, with zero shift and zero triad lapse). We have used both the Brailovskaya integration scheme, which is a second order predictor-corrector method, and the so-called iterative Crank-Nicholson integration scheme for numerical time evolutions. The details of the numerical method are described in the Paper I , where we also described how our code shows second order convergence behaviour. In order to show the expected “stabilization behaviour” clearly, we artificially add an error in the middle of the time evolution. More specifically, we set our initial guess 3-metric as $$\widehat{\gamma }_{ij}=\left(\begin{array}{ccc}1& 0& 0\\ sym.& 1+K(e^{(xL)^2}+e^{(x+L)^2})& 0\\ sym.& sym.& 1K(e^{(xL)^2}+e^{(x+L)^2})\end{array}\right),$$ (39) in the periodically bounded region $`x=[5,+5]`$, and added an artificial inconsistent rescaling once at time $`t=6`$ for the $`𝒜_y^2`$ component as $`𝒜_y^2𝒜_y^2(1+\mathrm{error})`$. Here $`K`$ and $`L`$ are constants and we set $`K=0.3`$ and $`L=2.5`$ for the plots. Fig.3 (a) shows how the violation of the Hamiltonian constraint equation, $`𝒞_H`$, become worse depending on the term $`\mathrm{error}`$. The oscillation of the L2 norm $`𝒞_H`$ in the figure due to the pulse waves collide periodically in the numerical region. We, then, fix the error term as a 20% spike, and try to evolve the same data in different equations of motion, i.e., the original Ashtekar’s equation \[solid line in Fig.3 (b)\], strongly hyperbolic version of Ashtekar’s equation (dotted line) and the above $`\lambda `$-system equation (other lines) with different $`\beta `$s but the same $`\alpha `$. As we expected, all the $`\lambda `$-system cases result in reducing the Hamiltonian constraint errors. ### D Remarks for the $`\lambda `$-system In the previous subsections, we showed that $`\lambda `$-system works as we expected. The system evolves into a constraint surface asymptotically even if we added an error artificially. However, the $`\lambda `$-system can not be introduced generally, because (i) the construction of the $`\lambda `$-system requires that the original dynamical equation be in symmetric hyperbolic form, which is quite restrictive for the Einstein equations, (ii) the system requires many additional variables and we also need to evaluate all the constraint equations at every time steps, which is hard task in computation. Moreover, it is not clear that the $`\lambda `$-system can control constraint equations which do not have any spatial differential terms. (e.g., the primary metric reality condition in the Ashtekar formulation.) <sup></sup><sup></sup>* This statement is not inconsistent with our previous work, in which we also proposed a $`\lambda `$-system that can control the secondary triad reality condition. We, next, propose an alternative system which also enable us to control the violation of constraint equations, but is robust for the above points. ## III Asymptotically constrained system 2: Adjusted system We here propose another approach for obtaining stable evolutions, which we name the “adjusted-system”. The essential procedure is to add constraint terms to the right-hand-side of the dynamical equations with multipliers, and to choose the multipliers so as the adjusted equations decrease the violation of constraints during time evolution. This system has several advantages than the previous $`\lambda `$-system. ### A “Adjusted system” The actual procedure for constructing an adjusted system is as follows. 1. Prepare a set of evolution equations for dynamical variables and the first class constraints which describe the problem. It is not required that the system is in the first order form nor hyperbolic form. However here we start from the same form with (1) and (2). We repeat them as $`_tu^\gamma `$ $`=`$ $`A^{i\gamma }{}_{\delta }{}^{}_{i}^{}u^\delta +B^\gamma ,`$ (40) $`_tC^\rho `$ $`=`$ $`D^{i\rho }{}_{\sigma }{}^{}_{i}^{}C^\sigma +E^\rho {}_{\sigma }{}^{}C_{}^{\sigma },`$ (41) where $`A(u(x^i))`$ is not required to form a symmetric or Hermitian matrix. 2. Add the constraint terms, $`C^\rho `$, (and/or their derivatives) to the dynamical equation (40) with multipliers $`\kappa `$, $$_tu^\gamma =A^{i\gamma }{}_{\delta }{}^{}_{i}^{}u^\delta +B^\gamma +\kappa _\rho ^\gamma C^\rho +\kappa _\rho ^{\gamma i}_iC^\rho .$$ (42) We call the added terms, $`\kappa _\rho ^\gamma C^\rho `$ and/or $`\kappa _\rho ^{\gamma i}_iC^\rho `$, “adjusted terms”, and leave $`\kappa _\rho ^\gamma `$ and $`\kappa _\rho ^{\gamma i}`$ unspecified for the moment. Because of these adjusted terms, the original constraint propagation equations, (41), must also be adjusted: $$_tC^\rho =D^{i\rho }{}_{\sigma }{}^{}_{i}^{}C^\sigma +E^\rho {}_{\sigma }{}^{}C_{}^{\sigma }+F^{ij\rho }{}_{\sigma }{}^{}_{i}^{}_jC^\sigma +G^{i\rho }{}_{\sigma }{}^{}_{i}^{}C^\sigma +H^\rho {}_{\sigma }{}^{}C_{}^{\sigma }.$$ (43) The last three terms are due to the adjusted terms. 3. Specify the multipliers $`\kappa `$, by evaluating the eigenvalues that appear in the RHS of (43). Practically, by taking the Fourier transformation (7), we can reduce (43) to homogeneous form, $$_t\widehat{C}^\rho =(ik_iD^{i\rho }{}_{\sigma }{}^{}+E^\rho {}_{\sigma }{}^{}k_ik_jF^{ij\rho }{}_{\sigma }{}^{}+ik_iG^{i\rho }{}_{\sigma }{}^{}+H^\rho {}_{\sigma }{}^{})\widehat{C}^\sigma .$$ (44) We, then, take a linearization against a certain background spacetime, $$_t{}_{}{}^{(1)}\widehat{C}_{}^{\rho }=(ik_i{}_{}{}^{(0)}D_{}^{i\rho }{}_{\sigma }{}^{}+{}_{}{}^{(0)}E_{}^{\rho }{}_{\sigma }{}^{}k_ik_j{}_{}{}^{(0)}F_{}^{ij\rho }{}_{\sigma }{}^{}+ik_i{}_{}{}^{(0)}G_{}^{i\rho }{}_{\sigma }{}^{}+{}_{}{}^{(0)}H_{}^{\rho }{}_{\sigma }{}^{}){}_{}{}^{(1)}\widehat{C}_{}^{\sigma }.$$ (45) (here ($`n`$) indicates the order in linearization) and evaluate the eigenvalues of the coefficient matrix in (45). For this process, we propose two guidelines. 1. The first one is to obtain negative real-part of the eigenvalues. This is from the same principle in $`\lambda `$-system when we specified $`\alpha `$ and $`\beta `$, in order to force the system approach the constraint surface asymptotically. Provided that we obtain $`\kappa `$ which produce all the negative-real-part eigenvalues, the Fourier component $`\widehat{C}`$ decays to zero in time evolution, and the original constraint term $`C`$ also. 2. An alternative guideline is to obtain as many non-zero eigenvalues as one can. More precisely, this case is supposed to have pure imaginary eigenvalues. In such a case, the constraint propagation equations (e.g. $`_t\widehat{C}=\pm ik\widehat{C}`$) behave like the normal wave equations in its original component (e.g. $`_tC=\pm _xC`$), and its stability can be discussed using von Neumann stability analysis. As is well known, stability depends on the choice of numerical integration scheme, but it is also certain that we can control (or decrease) the amplitude of the constraint terms. The advantage of this adjusted system is that we do not need to add variables to the fundamental set, while the above first guideline (3a) is the same mechanism which is applied for the $`\lambda `$-system. We note that the non-zero eigenvalue feature was conjectured in Alcubierre et. al. in order to show the advantage of the conformally-scaled ADM system, but the discussion there is of dynamical equations and not of constraint propagation equations. The guideline (3b) is obtained heuristically as we will show in Fig.5 that a system with three zero eigenvalues is more stable than one with five. We, however, conjecture that systems with non-zero (or pure-imaginary) eigenvalues in their constraint propagation equations have more dissipative features than that of zero-eigenvalue system. This is from the von Neumann’s stability analysis, evaluating dynamical variables with the finite-differenced quantities. See Appendix B for more details. We remark that adding constraint terms to the dynamical equations is not a new idea. For example, Detweiler applied this procedure to the ADM equations and used the finiteness of the norm to obtain a new system. This is also one of the standard procedures for constructing a symmetric hyperbolic system (e.g. ). We believe, however, that the above guidelines yield the essential mechanism for our purpose, to constructing a stable dynamical system. In the following subsections and the Appendix A, we demonstrate that this adjusted system actually works as desired in the Maxwell system and in the Ashtekar system of the Einstein equations, in which above two guidelines are applied respectively. ### B Example 1: Maxwell equations #### 1 adjusted system We here again consider the Maxwell equations (9)-(11). We start from the adjusted dynamical equations $`_tE_i`$ $`=`$ $`cϵ_i{}_{}{}^{jk}_{j}^{}B_k+P_iC_E+p^j{}_{i}{}^{}(_jC_E)+Q_iC_B+q^j{}_{i}{}^{}(_jC_B),`$ (46) $`_tB_i`$ $`=`$ $`cϵ_i{}_{}{}^{jk}_{j}^{}E_k+R_iC_E+r^j{}_{i}{}^{}(_jC_E)+S_iC_B+s^j{}_{i}{}^{}(_jC_B),`$ (47) where $`P,Q,R,S,p,q,r`$ and $`s`$ are multipliers. These dynamical equations adjust the constraint propagation equations as $`_tC_E`$ $`=`$ $`(_iP^i)C_E+P^i(_iC_E)+(_iQ^i)C_B+Q^i(_iC_B)`$ (49) $`+(_ip^{ji})(_jC_E)+p^{ji}(_i_jC_E)+(_iq^{ji})(_jC_B)+q^{ji}(_i_jC_B),`$ $`_tC_B`$ $`=`$ $`(_iR^i)C_E+R^i(_iC_E)+(_iS^i)C_B+S^i(_iC_B)`$ (51) $`+(_ir^{ji})(_jC_E)+r^{ji}(_i_jC_E)+(_is^{ji})(_jC_B)+s^{ji}(_i_jC_B).`$ This will be expressed using Fourier components by $`_t\left(\begin{array}{c}\widehat{C}_E\\ \widehat{C}_B\end{array}\right)`$ $`=`$ $`\left(\begin{array}{cc}_iP^i+iP^ik_i+ik_j(_ip^{ji})k_ik_jp^{ji}& _iQ^i+iQ^ik_i+ik_j(_iq^{ji})k_ik_jq^{ji}\\ _iR^i+iR^ik_i+ik_j(_ir^{ji})k_ik_jr^{ji}& _iS^i+iS^ik_i+ik_j(_is^{ji})k_ik_js^{ji}\end{array}\right)\left(\begin{array}{c}\widehat{C}_E\\ \widehat{C}_B\end{array}\right)=:T\left(\begin{array}{c}\widehat{C}_E\\ \widehat{C}_B\end{array}\right).`$ (52) Assuming the multipliers are constants or functions of $`E`$ and $`B`$, we can truncate the principal matrix as $${}_{}{}^{(0)}T=\left(\begin{array}{cc}iP^ik_ik_ik_jp^{ji}& iQ^ik_ik_ik_jq^{ji}\\ iR^ik_ik_ik_jr^{ji}& iS^ik_ik_ik_js^{ji}\end{array}\right),$$ (53) with eigenvalues $$\mathrm{\Lambda }^\pm =\frac{p+s\pm \sqrt{p^2+4qr2ps+s^2}}{2},$$ (54) where $`p:=iP^ik_ik_ik_jp^{ji},q:=iQ^ik_ik_ik_jq^{ji},r:=iR^ik_ik_ik_jr^{ji},s:=iS^ik_ik_ik_js^{ji}.`$ If we fix $`q=r=0`$, then $`\mathrm{\Lambda }^\pm =p,s`$. Further if we assume $`p^{ji},s^{ji}>0`$, and set everything else to zero, then $`\mathrm{\Lambda }^\pm <0`$, that is we can get the all eigenvalues which have negative real part. That is, our guideline (a) is satisfied. (Conversely, if we choose $`q=r=0`$ and $`p^{ji},s^{ji}<0`$, then $`\mathrm{\Lambda }^\pm >0`$.) #### 2 Numerical Demonstration We applied the above adjusted system to the same wave propagation problem as in §II B 3. For simplicity, we fix $`\kappa =p^{ij}=s^{ij}`$ and set other multipliers equal to zero. In Fig.4, we show the L2 norm of constraint violation as a function of time, with various $`\kappa `$. As was expected, we see better performance for $`\kappa >0`$ (of the system with negative real part of constraint propagation equation), while diverging behavior for $`\kappa <0`$ (of the system with positive real part of constraint propagation equation). ### C Example 2: Einstein equations (Ashtekar equations) #### 1 Adjusted system for controlling constraint violations We here only consider the adjusted system which controls the departures from the constraint surface. In the Appendix, we present an advanced system which controls the violation of the reality condition together with numerical demonstration. Even if we restrict ourselves to adjusted equations of motion for $`(\stackrel{~}{E}_a^i,𝒜_i^a)`$ with constraint terms (no their derivatives), generally, we could adjust them as $`_t\stackrel{~}{E}_a^i`$ $`=`$ $`i𝒟_j(ϵ_a^{cb}\underset{^{}}{𝑁}\stackrel{~}{E}_c^j\stackrel{~}{E}_b^i)+2𝒟_j(N^{[j}\stackrel{~}{E}_a^{i]})+i𝒜_0^bϵ_{ab}^c\stackrel{~}{E}_c^i+X_a^i𝒞_H+Y_a^{ij}𝒞_{Mj}+P_a^{ib}𝒞_{Gb},`$ (55) $`_t𝒜_i^a`$ $`=`$ $`iϵ_c^{ab}\underset{^{}}{𝑁}\stackrel{~}{E}_b^jF_{ij}^c+N^jF_{ji}^a+𝒟_i𝒜_0^a+\mathrm{\Lambda }\underset{^{}}{𝑁}\stackrel{~}{E}_i^a+Q_i^a𝒞_H+R_i{}_{}{}^{ja}𝒞_{Mj}^{}+Z_i^{ab}𝒞_{Gb},`$ (56) where $`X_a^i,Y_a^{ij},Z_i^{ab},P_a^{ib},Q_i^a`$ and $`R_i^{aj}`$ are multipliers. However, in order to simplify the discussion, we restrict multipliers so as to reproduce the symmetric hyperbolic equations of motion , i.e., $$X=Y=Z=0,P_a^{ib}=\kappa _1(N^i\delta _a^b+i\underset{^{}}{𝑁}ϵ_a{}_{}{}^{bc}\stackrel{~}{E}_{c}^{i}),Q_i^a=\kappa _2(e^2\underset{^{}}{𝑁}\stackrel{~}{E}_i^a),R_i{}_{}{}^{ja}=\kappa _3(ie^2\underset{^{}}{𝑁}ϵ^{ac}{}_{b}{}^{}\stackrel{~}{E}_{i}^{b}\stackrel{~}{E}_c^j).$$ (57) Here $`\kappa _1=\kappa _2=\kappa _3=1`$ is the case of symmetric hyperbolic equation for $`(\stackrel{~}{E}_a^i,𝒜_i^a)`$, while $`\kappa _1=\kappa _2=\kappa _3=0`$ is the (Ashtekar’s original) weakly hyperbolic equation, and other choices of $`\kappa `$s let the equation satisfy the level of strongly hyperbolic form. With these adjusted terms, the constraint propagation equations become $`_t{}_{}{}^{(1)}𝒞_{H}^{}`$ $`=`$ $`(12\kappa _3)(_j{}_{}{}^{(1)}𝒞_{Mj}^{}),`$ (58) $`_t{}_{}{}^{(1)}𝒞_{Mi}^{}`$ $`=`$ $`(12\kappa _2)(_i{}_{}{}^{(1)}𝒞_{H}^{})+i\kappa _3ϵ^{mj}{}_{i}{}^{}(_m{}_{}{}^{(1)}𝒞_{Mj}^{}),`$ (59) $`_t{}_{}{}^{(1)}𝒞_{Ga}^{}`$ $`=`$ $`2\kappa _3{}_{}{}^{(1)}𝒞_{Ma}^{}+i\kappa _1ϵ_a{}_{}{}^{bm}(_m{}_{}{}^{(1)}𝒞_{Gb}^{}).`$ (60) against the Minkowskii background. The eigenvalues of the coefficient matrix after the Fourier-transformation are $`(0,\pm i\kappa _1\sqrt{k^2},\pm i\kappa _3\sqrt{k^2},\pm i(2\kappa _21)(2\kappa _31)\sqrt{k^2})`$ (61) where $`k^2:=k_ik^i`$. For example, $`(0(\mathrm{multiplicity}5),\pm i\sqrt{k^2})`$ $`\mathrm{for}\kappa _1=\kappa _2=\kappa _3=0:\mathrm{original}\mathrm{system}`$ (62) $`(0(\mathrm{multiplicity}3),\pm i\sqrt{k^2}(\mathrm{multiplicity}3))`$ $`\mathrm{for}\kappa _1=\kappa _2=\kappa _3=1:\mathrm{symmetric}\mathrm{hyperbolic}\mathrm{system}.`$ (63) That is, our guideline (b) is obtained. The above adjustment, (55)-(57), will not produce negative-real-part eigenvalues, so our guideline (a) cannot be applied here. If we adjust the dynamical equation using the spatial derivatives of constraint terms, then it is possible to get all negative eigenvalues like in the Maxwell system (though this is complicated). However, since we found that this adjustment, (55)-(57), gives us an example of controlling the violation of constraint equations for our purpose, we only show this simpler version here. #### 2 Numerical Demonstration As a demonstration, we use here the same model as in §II C 3, that is, gravitational wave propagation in the plane symmetric spacetime, with an artificial error in the middle of time evolution. We examine how the adjusted multipliers contribute to the system’s stability. In Fig.5, we show the results of this experiment. We plot the violation of the constraint equations both $`𝒞_H`$ and $`𝒞_{Mx}`$. An artificial error term was added at $`t=6`$, as a kick of $`𝒜_y^2𝒜_y^2(1+\text{ error})`$, where the error amplitude is +20% as before. We set $`\kappa \kappa _1=\kappa _2=\kappa _3`$ for simplicity. The solid line is the case of $`\kappa =0`$, that is the case of “no adjusted” original Ashtekar equation (weakly hyperbolic system). The dotted line is for $`\kappa =1`$, equivalent to the symmetric hyperbolic system. We see other line ($`\kappa =2.0`$) shows better performance than the symmetric hyperbolic case. ## IV Discussion With the purpose of searching for an evolution system of the Einstein equations which is robust against perturbative errors for the free evolution of the initial data, we studied two “asymptotically constrained” systems. First, we examined the previously proposed “$`\lambda `$-system”, which introduces artificial flows to constraint surfaces based on the symmetric hyperbolic formulation. We showed that this system works as expected for the wave propagation problem in the Maxwell system and in Ashtekar’s system of general relativity. However, the $`\lambda `$-system cannot be applied to general dynamical systems in general relativity, since the system requires the base system to be symmetric hyperbolic form. Alternatively, we proposed a new mechanism to control the stability, which we named the “adjusted system”. This is simply obtained by adding constraint terms in the dynamical equations and adjusting the multipliers. We proposed two guidelines for specifying multipliers which reduce the numerical errors; that is, non-positive-real-part or pure-imaginary eigenvalues of the adjusted constraint propagation equations. This adjusted system was also tested in the Maxwell system and in Ashtekar’s system. As we denoted earlier, the idea of adding constraint terms is not new. However, we think that our guidelines for controlling the decay of constraint equations are appropriate for our purposes, and were not suggested before. Up to our numerical experiments, our guidelines give us clear indications whether the constraints decay (i.e. stable system) or not for perturbative errors, though we also think that this is not a complete explanation for all cases. This feature may be explained or proven in different ways, such as finiteness of the norm (of evolution equations or of constraint propagation equations), or by another mechanism in future. Secondary conclusion is that the symmetric hyperbolic equation is not always the best one for controlling stable evolution. As we show in the wave propagation model in the adjusted Ashtekar’s equation, our eigenvalue guidelines affect more than the system’s hyperbolicity. (We found a similar conclusion in .) We think this result opens a new direction to numerical relativists for future treatment of the Einstein equations. We are now applying our idea to the standard ADM and conformally scaled ADM system to explain these differences. Results will be reported elsewhere . ## Acknowledgements HS appreciates helpful comments by Abhay Ashtekar, Jorge Pullin, Douglas Arnold and L. Samuel Finn, and the hospitality of the CGPG group. We thank Bernard Kelly for careful reading the manuscript. Numerical computations were performed using machines at CGPG. This work was supported in part by the NSF grants PHYS95-14240, and the Everly research funds of Penn State. HS was supported by the Japan Society for the Promotion of Science as a research fellow abroad. ## A Controlling Reality condition by adjusted system We demonstrate here that our adjusted system in the Ashtekar formulation also works for controlling reality conditions. As a model problem, we concern the degenerate point passing problem which we considered previously in . In §A 1, we review this background briefly, and in §A 2 we show our numerical demonstrations. ### 1 Degenerate point passing problem In , the authors had examined the possibility of dynamical passing of the degenerate point in the spacetime. There the authors found that we are able to pass (i.e. continue time evolutions) if we could foliate the time-constant hypersurface into complex plane assuming that such a degenerate point exists on the real plane. Such foliations are available within Ashtekar’s original formulation, since the fundamental variables are complex quantities. The trick is to violate the reality condition locally, only in the vicinity of a degenerate point. As a model, we construct a metric, $`{}_{}{}^{(4)}g`$, which possesses a degenerate point ($`det{}_{}{}^{(3)}g=0`$) at the origin $`t=x=0`$ in Minkowskii background metric: $`ds^2`$ $`=`$ $`[1(2tx\mathrm{exp}(t^2x^2))^2]dt^2+4tx\mathrm{exp}(t^2x^2)[1(12x^2)\mathrm{exp}(t^2x^2)]dtdx`$ (A2) $`+[1(12x^2)\mathrm{exp}(t^2x^2)]^2dx^2+dy^2+dz^2.`$ We consider the time evolution, which initial data is described by a particular time slice $`t<0`$ of (A2), and whose time-constant hypersurfaces are foliated by the gauge condition, $`N`$ $`=`$ $`1,(\underset{^{}}{𝑁}=e^1),`$ (A3) $`N_x`$ $`=`$ $`2tx\mathrm{exp}(t^2x^2)[1(12x^2)\mathrm{exp}(t^2x^2)]+iat\mathrm{exp}(b(t^2+x^2)),`$ (A4) $`𝒜_0^a`$ $`=`$ $`0,`$ (A5) which enables to detour into the complex plane. Our goal is to demonstrate that the time evolution comes back to the real plane without any divergence in variables and curvatures. Such a “recovering condition” can be described by $`{\displaystyle _t_{}^{t_+}}\mathrm{}N(t,𝒙)𝑑t=0,{\displaystyle _t_{}^{t_+}}\mathrm{}N^i(t,𝒙)𝑑t=0,`$ Foliation recovering condition (A6) $`\mathrm{}N(t,𝒙)0,\mathrm{}N^i(t,𝒙)0,\mathrm{}\left[\stackrel{~}{E}_a^i\stackrel{~}{E}^{ja}/\mathrm{det}\stackrel{~}{E}(t,𝒙)\right]0,`$ Asymptotic reality condition (A7) for all four limits $`𝒙𝒙_{}\pm \mathrm{\Delta }𝒙`$, $`tt_{}\pm \mathrm{\Delta }t`$. Numerically, this problem becomes an eigenvalue problem, since our boundary conditions, (A6) and (A7), specify much freedom. To see if the evolution satisfies the criteria or not, we introduced two measures $`F(t_{final})`$ $`:=`$ $`\underset{x}{\mathrm{max}}\left\|\mathrm{}\left(e(t=t_{final},x)1\right)\right\|\text{(asymptotically flat)}`$ (A8) $`R(t_{final})`$ $`:=`$ $`\underset{x}{\mathrm{max}}\left\|\mathrm{}\left(e(t=t_{final},x)\right)\right\|\text{(asymptotically real)}`$ (A9) and searched the parameters $`a`$ and $`b`$ in (A4). If we apply our adjusted system to this model, then we expect that the allowed range for the parameters $`a`$ and $`b`$ becomes more general, since the real-surface-recovering feature is in the flow of the adjusted system’s foliation. ### 2 Application of the adjusted system As was shown in the previous section, for this purpose, we have to foliate our hypersurface in the complex-valued region and foliate back to the real-valued surface. That is, we can treat the reality condition, both primary and secondary, as a part of the constraint equations. For the above degenerate point-passing problem, we need to control only the violation of $`\mathrm{}m(\stackrel{~}{E}_a^i\stackrel{~}{E}_a^j)`$. Therefore, similar to the proposal of the adjusted system discussed in §II C, our adjusted dynamical equations can be written as $`_t\stackrel{~}{E}_a^i`$ $`=`$ $`i𝒟_j(ϵ_a^{cb}\underset{^{}}{𝑁}\stackrel{~}{E}_c^j\stackrel{~}{E}_b^i)+2𝒟_j(N^{[j}\stackrel{~}{E}_a^{i]})+i𝒜_0^bϵ_{ab}^c\stackrel{~}{E}_c^i+X_a^i𝒞_H+Y_a^{ij}𝒞_{Mj}+P_a^{ib}𝒞_{Gb}+T^i{}_{ajk}{}^{}\mathrm{}m(\stackrel{~}{E}_b^j\stackrel{~}{E}_b^k),`$ (A10) $`_t𝒜_i^a`$ $`=`$ $`iϵ_c^{ab}\underset{^{}}{𝑁}\stackrel{~}{E}_b^jF_{ij}^c+N^jF_{ji}^a+𝒟_i𝒜_0^a+\mathrm{\Lambda }\underset{^{}}{𝑁}\stackrel{~}{E}_i^a+Q_i^a𝒞_H+R_i^{aj}𝒞_{Mj}+Z_i^{ab}𝒞_{Gb}+V^a{}_{ijk}{}^{}\mathrm{}m(\stackrel{~}{E}_b^j\stackrel{~}{E}_b^k),`$ (A11) where $`X_a^i,Y_a^{ij},Z_i^{ab},P_a^{ib},Q_i^a,R_i^{aj},T^i_{ajk}`$ and $`V^a_{ijk}`$ are adjusted multipliers. If we simply set $`X_a^i=Y_a^{ij}=Z_i^{ab}=P_a^{ib}=Q_i^a=R_i^{aj}=V^a{}_{ijk}{}^{}=0`$ and $`T^i{}_{ajk}{}^{}=i\kappa \delta _j^i\delta _{ak}`$, (where $`\kappa `$ is real constant), then we obtain the constraint propagation equation $`_t({}_{}{}^{(0)}\mathrm{}(\stackrel{~}{E}_a^i\stackrel{~}{E}_a^j))`$ $`=`$ $`2\kappa ({}_{}{}^{(0)}\mathrm{}(\stackrel{~}{E}_a^i\stackrel{~}{E}_a^j))+\text{other constraint terms}.`$ (A12) The eigenvalue of the Fourier-transformed RHS is $`2\kappa `$. That is, if we set $`\kappa >0(<0)`$ then the eigenvalue is negative (positive), while $`\kappa =0`$ recovers the original non-adjusted system. The results of numerical demonstration are shown in Fig.6. We plot the L2 norm of the violation of the reality condition as a function of time, $`t`$ (this evolution is from $`t=5`$ to 5 ). Around the time $`t=0`$ the error appears due to our “detour” slicing condition, and the original system ($`\kappa =0`$) will not recover the reality surface with the choice of $`a`$ and $`b`$ in (A4) for this plot. However, for the positive $`\kappa `$ case, the foliation will be forced to recover the reality surface, while for negative $`\kappa `$ case it will not. Therefore this example again supports our guidelines, i.e. negative eigenvalue of constraint propagation equation will guarantee the evolution to the constraint surface. ## B von Neumann analysis of constraint propagation equations Here we show von Neumann’s stability analysis for the constraint propagation equations, in order to support our guideline (3b) for the adjusted system (§III A). The von Neumann analysis (see e.g. ) gives us powerful predictions for the stability of a finite difference approximation. Briefly, the analysis consists from the Fourier decomposition in the spatial directions of the dynamical variables and its one-step time evolution with a particular time integration scheme. If we wrote the fundamental variable $`\varphi (x,t)`$, then the criteria for the stability is $`\|\lambda _i\|1`$ where $`\lambda _i`$ are the eigenvalues of the amplification matrix $`G`$, which is in the expression of the evolution equations in the form of $`\varphi (x,t+\mathrm{\Delta }t)=G\varphi (x,t)`$. In our discussion, the constraint propagation equations are not directly used for numerical integrations, but are used as a guideline for the stability. The application of von Neumann analysis, however, is also allowed for the constraint propagation equations, as far as substituting the finite derivatives in the analysis using those of the fundamental dynamical variables. Here we show the most simplest cases for the adjusted Maxwell system and the adjusted Ashtekar system. ##### a Adjusted Maxwell system We start from choosing $`\kappa :=P_1=P_2=P_3`$ and other multipliers zero in the system (46) and (47). The Fourier component of the propagation equation for $`C_E`$ (52) becomes $`_t\widehat{C}_E=i\kappa (k_x+k_y+k_z)\widehat{C}_E,`$ which eigenvalue (54) is $`i\kappa (k_x+k_y+k_z)`$. That is, non-zero $`\kappa `$ gives us a pure-imaginary eigenvalue. By applying von Neumann analysis, we obtain the amplitude $`G`$s for FTCS (forward time and center space difference), Brailovskaya and (2-iteration) iterative Crank Nicholson schemes as $`\|G_{FTCS}\|^2`$ $`=`$ $`1+(\kappa \sigma )^2,`$ (B1) $`\|G_{Br}\|^2`$ $`=`$ $`1(\kappa \sigma )^2+(\kappa \sigma )^4,`$ (B2) $`\|G_{CN2}\|^2`$ $`=`$ $`1(\kappa \sigma )^4/4+(\kappa \sigma )^6/16,`$ (B3) respectively, where $`\sigma =(\mathrm{\Delta }t/\mathrm{\Delta }x)(\mathrm{sin}(k_x\mathrm{\Delta }x)+\mathrm{sin}(k_y\mathrm{\Delta }x)+\mathrm{sin}(k_z\mathrm{\Delta }x))`$ and we assume 3-dimensional finite grid of equal space $`\mathrm{\Delta }x`$ in all directions. Except for the FTCS scheme, we see that non-zero $`\|\kappa \|`$ (near $`\kappa =0`$) yields $`\|G\|<1`$. The bigger $`\|\kappa \|`$ (near $`\kappa =0`$) gives less $`\|G\|<1`$. The simulation we showed in §III B is not this case (since we tried to show the one which satisfy the guideline (3a)), but we also obtained the numerical results which confirm our conjecture here. ##### b Adjusted Ashtekar system Similarly, for the constraint propagation equations (58)- (60), with $`\kappa :=\kappa _1=\kappa _2=\kappa _3`$, we obtain the eigenvalues $`\lambda _i`$ of the amplification matrix $`G`$ for the above three schemes, $`\|\lambda \|_{FTCS}^2`$ $`=`$ $`1,1+(\kappa \sigma )^2,1+\{(12\kappa )(\kappa \sigma )\}1^2,`$ (B4) $`\|\lambda \|_{Br}^2`$ $`=`$ $`1,1(\kappa \sigma )^2+(\kappa \sigma )^4,1\{(12\kappa )\sigma \}^2+\{(12\kappa )\sigma \}^4,`$ (B5) $`\|\lambda \|_{CN2}^2`$ $`=`$ $`1,1(\kappa \sigma )^4/4+(\kappa \sigma )^6/16,1\{(12\kappa )\sigma \}^4/4+\{(12\kappa )\sigma \}^6/16,`$ (B6) with multiplicity 1, 4 and 2, respectively. Here again we see that non-zero $`\|\kappa \|`$ makes the system $`\|G\|<1`$ for Brailovskaya and 2-iteration Crank-Nicholson schemes. This analysis supports why the guideline (3b) works for our results shown in Fig. 5.
warning/0007/nucl-th0007017.html	ar5iv	text	# SIMULTANEOUS HEAVY ION DISSOCIATION AT ULTRARELATIVISTIC ENERGIES ## 1 Single and mutual electromagnetic dissociation Let us consider a collision of heavy ultrarelativistic nuclei with the masses and charges $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{Z}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}`$ and $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{Z}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}`$ at the impact parameter $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{b}$}`$ exceeding the sum of nuclear radii, $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{b}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{>}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{R}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{R}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}`$. According to the Weizsäcker-Williams (WW) method, the impact of the Lorentz-boosted Coulomb field of the nucleus $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}`$ on the collision partner $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}`$ is treated as the absorption of equivalent photons $`^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{?}}$}`$. The mean number of photons absorbed by the nucleus $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}`$ in such collision is: $$\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{m}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{b}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{N}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{Z}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{b}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\sigma }$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{𝑑}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}$$ (1) where the spectrum of virtual photons, $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{N}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{Z}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{b}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}`$ $`^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{?}}$}`$, and the total photoabsorption cross section, $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\sigma }$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}`$, are used. Assuming that the probability of multiphoton absorption is given by the Poisson distribution with the mean multiplicity $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{m}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{b}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}`$, one has the cross section for the single electromagnetic dissociation to a given channel $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{i}$}`$ $`^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{?}}$}`$: $$\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\sigma }$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{i}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\pi }$}\underset{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{b}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{c}$}}{\overset{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{}}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{b}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{𝑑}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{b}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{P}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{b}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{b}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{c}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1.34}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\left(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{/}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{3}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{/}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{3}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{0.75}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{/}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{3}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{/}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{3}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\right)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}$$ (2) where $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{b}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{c}$}`$ is in fm and the probability of dissociation at impact parameter $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{b}$}`$ is given by: $$\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{P}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{b}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{e}$}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{m}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{b}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{𝑑}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{N}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{Z}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{b}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\sigma }$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{f}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{i}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{.}$}$$ (3) Here $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{f}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{i}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}`$ is the branching ratio for the considered channel $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{i}$}`$ in the absorption of a photon with the energy $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}`$ on the nucleus $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}`$. These values are calculated by photonuclear reaction models $`^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{?}}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{?}}$}}`$ or taken from experiments. In the WW method, the graph for the mutual electromagnetic dissociation, Fig. 1, may be constructed from two graphs of the single dissociation by interchanging the roles of “emitter” and “absorber” at the secondary photon exchange. Such procedure is possible since the first emitted photon with $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{m}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{a}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{x}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\gamma }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{/}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{R}$}`$ does not change essentially the total energy, $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\gamma }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{M}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}`$, of the emitting nucleus, $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{m}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{a}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{x}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{/}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{/}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{R}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{M}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{10}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{4}$}`$, and there are no correlations between the energies $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}`$ and $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}`$. In other words, both photon exchanges may be considered as independent processes. Moreover, at ultrarelativistic energies the collision time is much shorter then a typical deexcitation period when a nucleus loses its charge via the proton emission or fission. It means that the equivalent photon spectrum from the excited nucleus, $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}}`$, is equal to the spectrum from the nucleus in its ground state, $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}`$, see Fig. 1. Therefore the cross section for the mutual electromagnetic dissociation of the nuclei $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}`$ and $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}`$ to given channels $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{i}$}`$ and $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{j}$}`$ is given by: $$\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\sigma }$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{m}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{i}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{j}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\pi }$}\underset{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{b}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{c}$}}{\overset{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{}}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{b}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{𝑑}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{b}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{P}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{b}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{P}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{b}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{.}$}$$ (4) Substituting $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{P}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{b}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}`$ for each of the nuclei one has: $$\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\sigma }$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{m}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{i}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{j}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{𝑑}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{𝑑}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{𝒩}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{m}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\sigma }$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\sigma }$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{f}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{i}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{f}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{j}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}$$ (5) with the spectral function $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{𝒩}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{m}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}`$ defined for the mutual dissociation: $$\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{𝒩}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{m}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\pi }$}\underset{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{b}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{c}$}}{\overset{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{}}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{b}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{𝑑}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{b}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{e}$}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{m}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{b}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{N}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{Z}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{b}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{N}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{Z}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{b}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{.}$}$$ (6) ## 2 Nucleon removal in grazing nuclear collisions The cross section for the abrasion of $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{n}$}`$ neutrons and $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{z}$}`$ protons from the projectile $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{Z}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}`$ in a peripheral collision with the target $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{Z}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}`$ may be derived from the Glauber multiple scattering theory $`^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{?}}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{?}}$}}`$. A simple parameterization exists for the single neutron removal cross section $`^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{?}}$}`$: $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\sigma }$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{n}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{u}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{c}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{n}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}{\displaystyle \frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{Z}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\sigma }$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{G}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{P}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{e}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{s}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{c}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\sigma }$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{G}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\pi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\left(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{b}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{c}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}{\displaystyle \frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{\Delta }}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{b}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\right)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{\Delta }}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{b}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}`$ where $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{\Delta }}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{b}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{0.5}$}`$ fm and $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{P}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{e}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{s}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{c}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{0.75}$}`$ is the probability for a neutron to escape without suffering the interaction with a spectator fragment. This can be extended to the case of mutual (1n,1n) emission: $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\sigma }$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{n}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{u}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{c}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{n}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{n}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\left(}$}{\displaystyle \frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{Z}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\right)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\left(}$}{\displaystyle \frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{Z}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\right)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\sigma }$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{G}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{P}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{e}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{s}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{c}$}}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{.}$}`$ (7) ## 3 Results and discussion The results of the abrasion model for the charge changing cross sections of the single dissociation of 158A GeV $`{}_{}{}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{208}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{Pb}}$}`$ ions are shown in Fig. 2 in addition to the electromagnetic contribution calculated by the RELDIS code. As one can see, the electromagnetic contribution dominates for few nucleon removal process. The interaction of knocked out nucleons with spectators and spectator de-excitation process itself were neglected in this version of the abrasion model. However, good agreement with the experimental data $`^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{?}}$}`$ is found. After this verification the model can be extrapolated to the energies of RHIC and LHC heavy-ion colliders. There is a proposal to use the simultaneous neutron emission for beam luminosity monitoring via the correlated registration of forward neutrons in zero degree calorimeters at RHIC $`^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{?}}$}`$. The model predicts $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\sigma }$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{m}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{n}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{n}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{750}$}`$ and 970 mb for AuAu and PbPb collisions at RHIC and LHC, respectively, for the correlated single neutron emission. We found that an important part of the dissociation events leading to the single neutron emission is accompanied by the emission of charged particles and nuclear fragments. Such events are due to the equivalent photon absorption above the Giant Resonance region. For grazing nuclear collisions Eq. (7) gives $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\sigma }$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{n}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{u}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{c}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{n}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{n}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{100}$}`$ mb. Our detailed consideration of the mutual nuclear dissociation is in progress. We are grateful to A.S. Botvina, G. Dellacasa, J.J. Gaardhøje, G. Giacomelli, A.B. Kurepin and S. White for useful discussions. I.A.P. and I.N.M. are indebted to the Organizing Committee of Bologna 2000 Conference for the kind hospitality and financial support. ## References
warning/0007/physics0007037.html	ar5iv	text	# Captions ## Captions Fig. 1 Phase portrait of the system for $`n_o=1.6`$ and homogeneous ion density, following from Eqs. (4) and (8), see also Fig. 2. The dashed lines stand for vacuum regions, the continuous lines stand for plasma regions, the actual trajectory is given by the thick line and it runs clockwise. The large dashed line denotes the regions where the electron density is negative. Fig. 2 The continuous line represents the plasma-field structures in a semi-infinite plasma initially occupying the region $`x0`$, for $`n_o=1.6`$, the unperturbed electron density being represented by the dashed line. The dotted line represents the resulting electron density distribution. All quantities are dimensionless.
warning/0007/math0007112.html	ar5iv	text	# Introduction. ## Introduction. In quantum physics one is frequently given a unital C\–algebra $``$ for the observables and a distinguished proper subset of states $`S_0S()`$ of its state space together with the constraint that the physical system can only realise these states and no others. Examples. (1) Constrained systems;– here one has a distinguished set of unitaries $`𝒰_u`$ and for the set of physically realisable states $`S_0`$ we have the Dirac states $`\left\{\omega S()\right\|\omega (𝒰)=1\},`$ cf. Grundling and Hurst \[GH\]. (2) The algebra of the canonical commutation relations $`=\overline{\mathrm{\Delta }(S,B)}`$ over a symplectic space $`(S,B)`$ (cf. Manuceau \[Ma\]), where $`S_0`$ is required to be the set of regular states. (3) In finite dimensional quantum mechanics, $``$ is taken as a factor of Type I and the physically relevant states, $`S_0,`$ as its set of normal states. (4) In algebraic quantum field theory, let $``$ be the inductive limit of a net of local algebras, and let $`S_0`$ be the set of locally normal states with respect to some distinguished state (cf. Haag \[Ha\]). In such a situation one can object that the given system $`(,S_0)`$ is not satisfactory because it leads naturally to nonphysical objects, for instance the weak\–closure of $`S_0`$ can be the full state space $`S().`$ Indeed, from the point of view that the states and observables should be in some kind of duality (a Heisenberg–Schrödinger picture equivalence), one can argue that if $`S_0`$ is the physical state space, then $``$ is not the correct algebra of observables. From the mathematical point of view there are also problems, e.g. we may not have a decomposition theory of states in $`S_0`$ in terms of other states in $`S_0.`$ Another problem, is that if we have two pairs $`(^{(i)},S_0^{(i)}),`$ $`i=1,\mathrm{\hspace{0.33em}2}`$ and an imprimitivity bimodule for the algebras, then when we induce a representation from one algebra to the other, we can easily move out of the class of allowed states. Our idea here is to replace $``$ by an algebra $``$ which has precisely $`S_0`$ as its state space, in a sense to be made precise below. For some examples of $`(,S_0),`$ we do in fact have a more convenient algebra which in some sense has precisely $`S_0`$ as its state space, e.g. if we take $`=(),`$ and let $`S_0`$ be its set of normal states, then the algebra of compact operators $`=𝒦()`$ is an algebra with state space $`S_0`$ in the sense that $`S_0𝒦()=S(𝒦())`$ and the states on $`𝒦()`$ extend uniquely. (We use here the notation $`S()`$ for the state space of its argument). Group algebras provide another set of examples – see the next section. So, inspired by these examples, we will study here the following situation. Given a pair $`(,S_0),`$ find a pair of C\–algebras $``$ and an embedding $`\xi :`$ such that the states on $``$ extend uniquely to $`\xi (),`$ and $`\theta (S())=S_0`$ where $`\theta `$ denotes the extension map $`\theta :S()S()`$. Naturally, there are existence questions to be answered, and we will address these first. In the next section we develop our basic theory, in Sect. 2 we construct a host up to a central algebra, and in Sect. 3 we do a few applications. ## 1. Basic Concepts. To prepare the ground, we first recall some background material. Recall that a hereditary subalgebra $``$ of a C\–algebra $`𝒜`$ is a C\–subalgebra such that $`0<A<B`$ for $`A𝒜,`$ $`B`$ implies that $`A`$. Equivalently (cf. Murphy \[Mu\]) $``$ is hereditary if $`𝒜.`$ Such algebras are plentiful, and are in bijection with the set of closed left ideals of $`𝒜.`$ All closed two–sided ideals are hereditary. For us, the most important property is: a C\–subalgebra $``$ is hereditary iff each state $`\omega S()`$ has a unique extension to a state $`\theta (\omega )S(𝒜)`$ (cf. Kusuda \[Ku\]). If we define a projection $`P𝒜^{\prime \prime }`$ as the unit of $`^{\prime \prime }𝒜^{\prime \prime },`$ then the map $`\theta :^{}𝒜^{}`$ defined by $$\theta (\phi )(A):=\phi (PAP)=\underset{\alpha }{lim}\phi (E_\alpha AE_\alpha )$$ $`(1)`$ for all $`A𝒜`$ and some approximate identity $`\{E_\alpha \}`$ of $`,`$ is precisely the unique extension map on the states $`S()^{}.`$ Moreover $`\theta `$ is an isometry, hence its range is norm closed. Given a representation of $`,`$ we can always induce a representation on $`𝒜`$ from it (cf. Fell and Doran, XI.7.6 \[FD\]), but usually this will be on a different space than the original representation. In the case that $``$ is a two–sided ideal of $`𝒜`$ we have that $`P𝒜^{}𝒜^{\prime \prime }`$, and so $$\theta (\phi )(A):=\phi (PA)=\underset{\alpha }{lim}\phi (E_\alpha A)$$ $`(2)`$ and now, in addition, representations of $``$ also extend uniquely on the same space to $`𝒜.`$ For a representation $`\pi :(),`$ the unique extension is $$\stackrel{~}{\pi }(A):=\pi (PA)=\underset{\alpha }{\text{s-lim}}\pi (E_\alpha A)$$ $`(3)`$ In this paper we will always use the notation $`(,S_0),`$ to denote a unital C\–algebra $``$ and a distinguished subset of its state space. We choose $``$ to be unital, since this ensures that its state space is w\–closed, hence norm closed (cf. Pedersen \[Pe\] 3.2.1). We are now ready for our basic definitions: Def. Given a pair $`(,S_0)`$ consisting of a unital C\–algebra and a proper subset of its states, we say a C\–algebra $``$ is a host for the pair if there is a unital C\–algebra $``$ (faithful embedding) and a unital \-homomorphism $`\xi :`$ such that: (i) $``$ is hereditary in $`,`$ (ii) $``$ is generated by $``$ and $`\xi (),`$ (iii) the map $`\theta :S()S()`$ defined by unique extension $$\theta (\omega )(F)=\underset{\alpha }{lim}\omega (E_\alpha \xi (F)E_\alpha ),F$$ $`(4)`$ for any approximate identity $`\{E_\alpha \}`$ of $`,`$ is injective and has range $`\theta (S())=S_0.`$ $``$In the case that $``$ is in addition a two sided ideal of $`,`$ we call it an ideal host. Remarks (1) The requirement (ii) that $``$ be generated by $``$ and $`\xi (),`$ is not essential;– if we start from some larger algebra $``$ satisfying the other requirements, we can always replace $``$ by $`C^{}(\xi ())`$ inside $``$ because $``$ is hereditary for any subalgebra of $``$ containing it. An important point in the definition, is that we allow $``$ to be outside $`.`$ (2) If $`(,S_0)`$ has an ideal host $`,`$ then there is a natural homomorphism of $``$ into the multiplier algebra $`M()^{\prime \prime },`$ and so we can equivalently define an ideal host for $`(,S_0)`$ as a C\–algebra $``$ together with a homomorphism $`\xi :M()`$ such that $`\theta :S()S()`$ is injective, with image $`S_0.`$ (3) If $`(,S_0)`$ has only a host $`,`$ then we may want to try a similar embedding than in the last remark. The point is that $``$ acts as a set of quasi–multipliers of $``$ (i.e. $`E`$ for each $`E`$) and we know by C. Akemann and G. Pedersen, Prop. 4.2 \[AP\], that there is a linear bijection between the quasi–multipliers and elements $`A^{\prime \prime }`$ such that $`A.`$ However, in general this bijection is not a homomorphism for C\–algebras of quasi–multipliers, so we can not exploit this bijection as in the previous remark. (4) Ideal hosts are of course more useful than hosts because then one has also unique extensions of representations on the same space. As we shall see however, their existence place strong structural restrictions on $`S_0.`$ Most of our analysis here will concern ideal hosts. First we do a list of examples, and then some structural analysis. Examples (1) For the pair $`(,S_0)=((),S_N)`$ where $`S_N`$ denotes the set of normal states, an ideal host algebra is $`=𝒦()`$ with the identity map $`\xi :()()=M(𝒦())=`$. (2) An important example for physics, is the following. Let the pair $`(,S_0)`$ consists of the CCR–algebra over a finite dimensional symplectic space, and its set of regular states. To be more concrete, consider the CCR–algebra on $`R^2,`$ i.e. the unique simple C\–algebra $``$ generated by unitaries $`\left\{\delta _𝐱\right\|𝐱R^2\}`$ satisfying the Weyl relations: $$\delta _𝐱\delta _𝐲=\rho (𝐱,𝐲)\delta _{𝐱+𝐲},\text{where}\rho (𝐱,𝐲):=\mathrm{exp}[i(x_1y_2x_2y_1)].$$ Define another C\–algebra $``$ as the C\–envelope of the twisted convolution algebra, where the latter consists of $`L^1(R^2)`$ equipped with the multiplication and involution: $$fg(𝐱)=_{R^2}f(𝐲)g(𝐱𝐲)\rho (𝐲,𝐱)𝑑𝐲,f^{}(𝐱)=\overline{f(𝐱)}.$$ This algebra $``$ is known to be isomorphic to $`𝒦(L^2(R))`$ (cf. I.E. Segal \[Se\]). Then $`M()`$ by the action $`\delta _𝐱f(𝐲)=\rho (𝐱,𝐲)f(𝐲𝐱).`$ The unique extensions of states on $``$ to $``$ produce precisely the set of regular states $$S_0:=\left\{\omega S()\right\|𝐱\omega (\delta _𝐱)\text{is continuous}\}.$$ The extension map $`\theta :S()S_0`$ is injective, since for each $`\varphi S_0`$ we can reconstruct the $`\omega S()`$ such that $`\theta (\omega )=\varphi `$ via the formula $$\omega (f):=_{R^2}f(𝐱)\phi (\delta _𝐱)𝑑𝐱$$ for all $`fL^1(R^2).`$ (Put in other words, here $``$ is the twisted discrete group algebra for $`R^2,`$ and $``$ is the usual twisted group algebra for $`R^2.)`$ Since $``$ is simple, and $`S_0S(),`$ we see that $`=\{0\}.`$ Of course $``$ is a far better behaved C\–algebra than $`,`$ it is even separable. See \[Gr\] for generalisations of this example. (3) Let $`G`$ be a nondiscrete topological group, and denote $`G`$ equipped with the discrete topology by $`G_d.`$ Let the pair $`(,S_0)`$ be the discrete group algebra $`=C^{}(G_d)`$ and its set of states continuous with respect to the topology of $`G`$ (i.e. $`g\omega (\delta _g)`$ is continuous, where $`\delta _g`$ denotes the Dirac point measure at $`g`$). When $`G`$ is locally compact, the usual group algebra $`C^{}(G)`$ is an ideal host algebra if we use the imbedding $`C^{}(G_d)M\left(C^{}(G)\right)`$ obtained by convolution of measures. Inspired by this example, we define Def. Let $`G`$ be a topological group, not necessarily locally compact. Then a group algebra for it, is any ideal host for the pair $`(C^{}(G_d),S_0)`$ where $`S_0`$ denotes the states $`\omega `$ on $`C^{}(G_d)`$ such that the map $`g\omega (\delta _g)`$ is continuous. In \[Gr\] we took this definition for a group algebra in order to construct a group algebra for inductive limit groups. We can also adapt it for covariance algebras. In this paper, however, we will not construct any such group algebras for more general groups. Next we analyze some structural consequences for the existence of a host for the pair $`(,S_0).`$ We will usually consider the homomorphism $`\xi :`$ as an embedding, to save on notation. Theorem 1.1. If a pair $`(,S_0)`$ has a host $`,`$ then (i) $`S_0`$ is a norm–closed face in $`S().`$ (ii) If $``$ is an ideal host, then the norm–closed face $`S_0`$ is also invariant, i.e. if $`\omega S_0`$ then $`\omega _BS_0`$ for all $`B`$ with $`\omega (B^{}B)=1,`$ and where $`\omega _B(F):=\omega (B^{}FB)`$ for $`F.`$ (iii) $`\theta :S()S_0`$ is an isomorphism, i.e. it is affine and a homeomorphism w.r.t. the norm topology. (iv) In the case that $`S_0`$ is the face obtained by extending the states from a hereditary subalgebra $`𝒜`$ to $`,`$ then $`=𝒜.`$ Proof: (i) Norm closure: We first show that $`\theta :^{}^{}`$ is norm continuous. Recall that since $``$ is hereditary in $``$ we have $`\theta (\phi )(F):=\phi (PFP)`$ with $`P^{\prime \prime }^{\prime \prime }^{\prime \prime }.`$ Now $$\begin{array}{ccc}\hfill \theta (\phi )& =sup\left\{\right\|\theta (\phi )(F)\left\|\right\|F,F1\}\hfill & \\ & =sup\left\{\right\|\phi (PFP)\left\|\right\|F,F1\}\hfill & \end{array}$$ Now we know from Kusuda \[Ku\] Theorem 2.2 that $`P^{\prime \prime }P=^{\prime \prime }`$ since $``$ is hereditary, and so $$\left\{\right\|\phi (PFP)\left\|\right\|F,F1\}\left\{\right\|\phi (A)\left\|\right\|A^{\prime \prime },A1\}$$ and hence, since the supremum of the last set is just $`\phi ,`$ we find that $`\theta (\phi )\phi .`$ Thus $`\theta :^{}^{}`$ is norm continuous. By assumption $`\theta `$ is injective on $`S();`$ we prove that it is also injective on $`^{}.`$ If it were not, there would be $`\phi _i^{}`$ such that $`\theta (\phi _1\phi _2)=0.`$ Then for $`\psi :=\phi _1\phi _2`$ do a Jordan decomposition, $`\psi =\rho _+\rho _{}+i(\mu _+\mu _{})`$ and then $`\theta (\psi )=0`$ implies $`\theta (\rho _+\rho _{})=0=\theta (\mu _+\mu _{}),`$ i.e. $`\theta (\rho _+)=\theta (\rho _{})`$ and $`\theta (\mu _+)=\theta (\mu _{}).`$ But $`\theta `$ is injective on states, so $`\rho _+=\rho _{}`$ and $`\mu _+=\mu _{},`$ i.e. $`\psi =0.`$ Thus we know that $`\theta `$ is both norm continuous and invertible on $`^{}`$ so by a corollary to the open mapping theorem, its inverse must also be norm continuous, and hence by Theorem 5.8, p216 of A. Taylor \[Ta\] we conclude that $`\theta (^{})`$ is norm closed. Since $`S()^{}`$ is also norm closed (recall that $``$ is unital), so is $`\theta (^{})S()`$ and as this is just the image under $`\theta `$ of $`S(),`$ we conclude that $`S_0`$ is norm closed. Now (iii) also follows from the preceding. That $`S_0`$ is a convex set, follows from the fact that $`\theta `$ is linear and $`S()`$ is convex, so we just need to prove the facial property. Let $`\omega =\lambda \phi +(1\lambda )\psi S_0`$ where $`\lambda [0,1].`$ We need to show that $`\phi ,\psi S_0.`$ Since $`\omega S_0`$ there is a unique $`\omega ^{}S()`$ such that $`\omega =\theta (\omega ^{}).`$ By the hereditary property, $`\omega ^{}`$ extends uniquely to a state $`\widehat{\omega ^{}}`$ on $`^{\prime \prime }`$ (and this extension of course restricts to $`\omega `$ on $``$). By definition of $`P`$ we have that $`\widehat{\omega ^{}}(P)=1,`$ and conversely, given any state $`\gamma `$ on $`^{\prime \prime }`$ with $`\gamma (P)=1`$ we have that $`\gamma S(),`$ hence $`\theta (\gamma )=\gamma S_0.`$ Now $`0=\widehat{\omega ^{}}(IP)=\lambda \widehat{\phi }(IP)+(1\lambda )\widehat{\psi }(IP)`$ where $`\widehat{\phi },\widehat{\psi }`$ are extensions of $`\phi ,\psi `$ to $`^{\prime \prime }.`$ Thus by positivity of all terms in the sum, we conclude that $`\widehat{\phi },\widehat{\psi }`$ vanish on $`IP`$ and hence $`\phi =\theta (\widehat{\phi }),`$ $`\psi =\theta (\widehat{\psi }).`$ Thus $`\phi ,\psi S_0,`$ i.e. $`S_0`$ is a face. For (ii), assume that $``$ is an ideal of $`,`$ then we want to prove invariance of the face. Let $`\omega =\theta (\omega ^{})S_0,`$ then clearly $`\omega _B^{}():=\omega ^{}(B^{}B)`$ defines a state on $``$ using the fact that $``$ is an ideal (here we took $`B`$ with $`\omega ^{}(B^{}B)=1`$). Thus $`S()`$ is invariant, and by the definitions $`\omega _B=\theta (\omega _B^{})S_0`$ because extension and conjugation of a state commutes when the projection $`P`$ in Equation (1) commutes with $`,`$ and this it does since $``$ is an ideal. (iv) Since $``$ is hereditary in $`,`$ we see that $``$ is hereditary in $`.`$ Moreover, the face $`S_0S()`$ is precisely the states which extend from the states on $`,`$ which are uniquely determined by their values on $``$ by the hereditary property. Since $`S_0`$ are also the states which extend from $`𝒜,`$ by the bijection between such faces and hereditary subalgebras (cf. Pedersen 3.10.7 \[Pe\] and Murphy 3.2.1 \[Mu\]) we conclude that $`=𝒜.`$ . Remarks. (1) It is known that an invariant convex norm–closed set of states $`S_0`$ is also a face (cf. footnote in \[HKK\]), so for an ideal host it suffices to say that $`S_0`$ is an invariant convex norm–closed set. In other words, this says that the cone which $`S_0`$ generates, $`R_+S_0,`$ is a folium. We will extend the term “folium” to also mean invariant norm–closed convex sets of state spaces. (2) Since $`\theta :S()S_0`$ is an affine bijection, it restricts to a bijection between the pure states on $``$ and the extreme points of $`S_0.`$ This immediately limits the class of faces and folia for which hosts exist, because there are many folia without extreme points, e.g. the folium of normal states of $`L^{\mathrm{}}(R)`$ (measures absolutely continuous w.r.t. the Lebesgue measure). We will sharpen this observation below. Note also that as $`\theta `$ involves both an extension and a restriction, it need not a priori take pure states to pure states. Since $`S()`$ is generated as the w\–closed convex hull of its pure states, we can use $`\theta `$ to transfer this weak\–topology to $`S_0`$ (but this is different from the weak\–topology of $`),`$ to conclude that w.r.t. this topology $`S_0`$ is a compact convex set, hence via Choquet theory, there are integral decompositions of states w.r.t. measures on $`S_0.`$ In the next section we will exploit these decompositions. It would be nice to have some intrinsic definition of the w\–topology induced by $`\theta `$ on $`S_0`$ but we do not have this yet. Theorem 1.2. Let $`f(^{})_+`$ be a folium. Then it is the set of normal positive forms of the von Neumann algebra $`\pi ()^{\prime \prime }`$ where $`\pi =\underset{\phi f}{}\pi _\phi .`$ Conversely, the set of positive normal forms of any von Neumann algebra is a folium. Proof: See Haag, Kadison, Kastler in \[HKK\]. . This is quite useful, in that any folium can now be analyzed as the normal state space of some concrete C\–algebra. Remark 1.3. For later use, we need to know about projections associated with faces and folia. Start with a pair $`(,S_0),`$ where by Theorem 1.1 we now assume that $`S_0`$ is a norm–closed face. Corresponding to this, we know from Pedersen 3.6.11 \[Pe\] that there is a projection $`P^{\prime \prime }`$ which we now show how to construct. (To use Pedersen 3.6.11, we need to know that a norm closed face $`S_0`$ generates a cone $`R_+S_0`$ which is norm–closed and hereditary and in $`(^{\prime \prime })_{},`$ but it is quite straightforward to verify this). First define $$\mathrm{\Gamma }:=\left\{\phi (^{\prime \prime })_{}\right\|\|\phi \|R_+S_0\},$$ and this is in fact a left invariant vector space, by the proof in Pedersen 3.6.11. Then its annihilator $`\mathrm{\Gamma }^{}^{\prime \prime }`$ is a $`\sigma \text{–weakly}`$ closed left ideal, hence $`\mathrm{\Gamma }^{}(\mathrm{\Gamma }^{})^{}`$ is a weak–operator closed hereditary subalgebra of $`^{\prime \prime }.`$ If we denote its unit (which is a projection in $`^{\prime \prime })`$ by $`Q,`$ then the desired projection we want is $`P=IQ.`$ To recover $`R_+S_0`$ from $`P,`$ we just take the set of $`\phi (_{}^{\prime \prime })_+=(^{})_+`$ such that $`\phi (P)=1.`$ In the case that $`S_0`$ is a folium (i.e. also invariant), we find that $`\mathrm{\Gamma }`$ is a two–sided invariant space, hence $`\mathrm{\Gamma }^{}`$ is a two–sided $`\sigma \text{–weakly}`$ closed ideal. It then follows from Pedersen 2.5.4 \[Pe\] that its unit $`Q^{}^{\prime \prime },`$ and hence $`P^{}^{\prime \prime }.`$ An obvious method by which one may think one can construct a host algebra, is to take the algebra $`\stackrel{~}{}:=C^{}(PP)^{\prime \prime }.`$ Whilst this is certainly hereditary in $`^{\prime \prime },`$ and the states which uniquely extend from $`\stackrel{~}{}`$ to $`:=C^{}(\stackrel{~}{})^{\prime \prime }`$ will satisfy $`\omega (P)=1,`$ this is not enough to guarantee that their restrictions to $``$ will be in $`S_0.`$ This is because given a $`\phi S(),`$ one can only conclude that $`\phi S_0`$ if its normal extension to $`^{\prime \prime }`$ satisfies $`\stackrel{~}{\phi }(P)=1,`$ and whilst for an $`\omega S()`$ which extended from one on $`\stackrel{~}{}`$ we have $`\omega (P)=1,`$ we do not know that $`\omega `$ is the normal extension of its restriction $`\omega .`$ Thus $`\theta `$ may not map onto $`S_0`$ for this choice $`\stackrel{~}{}.`$ By the previous remark we know there are folia without hosts, so that we know the above procedure must sometimes fail. From Takesaki Prop. 2.17 (p129) \[Tak\] we know that if $``$ is a closed two–sided ideal of $`𝒜,`$ then $`\pi ()^{\prime \prime }=\pi (𝒜)^{\prime \prime }`$ for any representation $`\pi `$ which is nondegenerate on $`.`$ This implies that if we have an ideal host $``$ for $`(,S_0),`$ then $`\pi ()^{\prime \prime }=\pi ()^{\prime \prime }`$ for the representation $`\pi =\underset{\omega S()}{}\pi _\omega .`$ This fact leads us to suspect that $`\pi ()^{\prime \prime }=\pi ()^{\prime \prime }`$ and this is what we now want to prove, but we need a lemma first. We use Pedersen’s notation $`[]`$ for “closed linear span.” Lemma 1.4. Let $`(,S_0)`$ have an ideal host $`,`$ and let $`\pi :()`$ be a representation. Let $`_e`$ be the essential subspace of $`\pi ().`$ Then for any vector $`\mathrm{\Omega }_e`$ we have $$\left[\pi ()\mathrm{\Omega }\right]=\left[\pi ()\mathrm{\Omega }\right].$$ Proof: Denote $`_\mathrm{\Omega }:=\left[\pi ()\mathrm{\Omega }\right]`$ and $`\pi _\mathrm{\Omega }:=\pi _\mathrm{\Omega },`$ then due to the fact that $``$ is an ideal host, $`\pi _\mathrm{\Omega }`$ extends uniquely on the same space $`_\mathrm{\Omega }`$ to $`.`$ Hence $`\left[\pi ()\mathrm{\Omega }\right]=\left[\pi _\mathrm{\Omega }()\mathrm{\Omega }\right]_\mathrm{\Omega }=\left[\pi ()\mathrm{\Omega }\right].`$ We prove the reverse inclusion by contradiction. Assume it is not true, then there exists some nonzero $`\psi [\pi ()\mathrm{\Omega }]_e`$ such that $`\psi [\pi ()\mathrm{\Omega }].`$ We have $$\pi ()\psi \pi ()\mathrm{\Omega }$$ $`(5)`$ because $`(\pi (A)\psi ,\pi (B)\mathrm{\Omega })=(\psi ,\pi (A^{}B)\mathrm{\Omega })=0`$ for all $`A,B`$. Normalise: $`\psi =1=\mathrm{\Omega }`$ and choose $`\alpha ,\beta C`$ such that $`\|\alpha \|^2+\|\beta \|^2=1`$ and $`\alpha 0\beta `$ and define: $$\begin{array}{ccc}\hfill \phi :=& \alpha \mathrm{\Omega }+\beta \psi ,\omega _\phi (A):=(\phi ,\pi (A)\phi )\hfill & \\ \hfill \omega _\psi (A):=& (\psi ,\pi (A)\psi ),\omega (A):=(\mathrm{\Omega },\pi (A)\mathrm{\Omega }).\hfill & \end{array}$$ Now since $`\phi _e,`$ $`\omega _\phi `$ is nondegenerate on $`,`$ hence $`\theta (\omega _\phi )=\omega _\phi `$ and so for all $`A`$ we have: $$\begin{array}{ccc}\hfill \theta (\omega _\phi )(A)& =\omega _\phi (A)=(\alpha \mathrm{\Omega }+\beta \psi ,\pi (A)\left(\alpha \mathrm{\Omega }+\beta \psi \right))\hfill & \\ & =\|\alpha \|^2\omega (A)+\|\beta \|^2\omega _\psi (A)=\lambda \omega (A)+(1\lambda )\omega _\psi (A)\hfill & \end{array}$$ where we made use of the orthogonality (5) and we have set $`\lambda :=\|\alpha \|^2(0,\mathrm{\hspace{0.17em}1}).`$ Thus $$\theta (\omega _\phi )=\theta \left(\lambda \omega +(1\lambda )\omega _\psi \right).$$ $`(6)`$ However, for an element $`L`$ we cannot use the orthogonality (5), and so we get $$\omega _\phi (L)=\lambda \omega (L)+(1\lambda )\omega _\psi (L)+\overline{\alpha }\beta (\mathrm{\Omega },\pi (L)\psi )+\alpha \overline{\beta }(\psi ,\pi (L)\mathrm{\Omega })$$ $`(7)`$ We show that the last two terms can always be made nonzero by some choice of $`L.`$ If this were not the case, they must be zero for all $`L=\gamma R`$ where $`\gamma C,`$ $`R=R^{}.`$ Then we have for the last two terms of Equation (7) that $$2\mathrm{R}\mathrm{e}\left[\alpha \overline{\beta }\gamma (\psi ,\pi (R)\mathrm{\Omega })\right]=0$$ for all $`\gamma `$ and $`R,`$ i.e. $`(\psi ,\pi (R)\mathrm{\Omega })=0`$ for all $`R=R^{}.`$ However, $``$ is spanned by its selfadjoint elements, thus $`\psi \pi ()\mathrm{\Omega },`$ and so $$\psi \left[\pi ()\mathrm{\Omega }\right]\psi $$ and thus $`\psi =0`$ which is a contradiction with our initial assumption. Thus the last two terms of Equation (7) are nonzero for some $`L,`$ i.e. on $``$ we have $$\omega _\phi \lambda \omega +(1\lambda )\omega _\psi .$$ This, together with Equation (6) contradicts the assumption that $`\theta `$ is injective on $`S().`$ Thus our initial assumption is wrong, so $`\left[\pi ()\mathrm{\Omega }\right]\left[\pi ()\mathrm{\Omega }\right],`$ and in fact we have equality $`\left[\pi ()\mathrm{\Omega }\right]=\left[\pi ()\mathrm{\Omega }\right].`$ . Remarks. (1) Until now, we have used the standard notation $`^{\prime \prime }`$ for the universal von Neumann algebra of $``$ (not unique for concrete C\–algebras). To avoid confusion in subsequent arguments, we will sometimes explicitly indicate the universal representations, and our notation is that $`\pi _{}:(_{})`$ is the universal representation of $`,`$ i.e. $`\pi _{}=\underset{\omega S()}{}\pi _\omega `$ and $`^{\prime \prime }\pi _{}()^{\prime \prime }.`$ Note that when $``$ is an ideal host for $`(,S_0),`$ then $`\pi _{}`$ extends uniquely on the same space to a representation of $`,`$ and this implies that $`\pi _{}`$ of $``$ is a subrepresentation of $`\pi _{}.`$ (2) One may try to generalise this lemma away from ideal hosts to hosts, in which case we suspect that for any vector $`\mathrm{\Omega }_e`$ we have $$\left[\pi ()\mathrm{\Omega }\right]\left[\pi ()\mathrm{\Omega }\right].$$ However, the proof so far eludes us. If one starts as in the proof by assuming some nonzero $`\psi [\pi ()\mathrm{\Omega }]\backslash [\pi ()\mathrm{\Omega }],`$ then $`\psi =\psi _0+\psi _1`$ where $`\psi _0[\pi ()\mathrm{\Omega }]\psi _1`$ but we may have that $`\psi _i_e,`$ even though $`\psi _e,`$ and this causes problems. Now we are ready to prove: Theorem 1.5. Let $`(,S_0)`$ have an ideal host $`,`$ then $`\pi _{}()^{\prime \prime }=\pi _{}()^{\prime \prime }(^{\prime \prime })`$ where we use the same symbol $`\pi _{}`$ for the unique extension of it from $``$ to $`^{\prime \prime }M().`$ Proof: In $`(_{})`$ let $`A\pi _{}()^{}`$ and let $`B,`$ and recall that $`\pi _{}(B)\psi =\underset{\alpha }{lim}\pi _{}(BE_\alpha )\psi `$ for all $`\psi _{}`$ and any approximate identity $`\{E_\alpha \}`$ of $`.`$ Thus for all $`\psi _{}`$ we have $$[A,\pi _{}(B)]\psi =\underset{\alpha }{lim}[A,\pi _{}(BE_\alpha )]\psi =0$$ because $`A\pi _{}()^{}`$ and $``$ is an ideal. This is true for all $`B,`$ and so $`\pi _{}()^{}\pi _{}()^{}.`$ We now prove the reverse inclusion, and we do it by contradiction. Assume that $`\pi _{}()^{}\pi _{}()^{}.`$ Then since von Neumann algebras are spanned by their projections, we can find a nontrivial projection $`P\pi _{}()^{}\backslash \pi _{}()^{}`$ (otherwise, if all the projections of $`\pi _{}()^{}`$ were in $`\pi _{}()^{}`$ the algebras would be equal). Recall that $`\pi _{}=\underset{\omega S()}{}\pi _\omega `$, so there must be some state in $`S()`$, say $`\omega _0`$, such that $$[\pi _{}(),P]P_{\omega _0}0=[\pi _{}(),P]$$ $`(8)`$ where $`P_{\omega _0}`$ denotes the projection onto the subspace $`_{\omega _0}_{}`$ of the subrepresentation $`\pi _{\omega _0}:(_{\omega _0}).`$ Let $`\mathrm{\Omega }_{\omega _0}`$ be the normalised cyclic vector for this representation. We claim that $$P\mathrm{\Omega }_{\omega _0}0(IP)\mathrm{\Omega }_{\omega _0}.$$ If $`P\mathrm{\Omega }_{\omega _0}=0,`$ then $$0=\left[\pi _{}()P\mathrm{\Omega }_{\omega _0}\right]=\left[P\pi _{}()\mathrm{\Omega }_{\omega _0}\right]=\left[P\pi _{\omega _0}()\mathrm{\Omega }_{\omega _0}\right]=P_{\omega _0}$$ where we made use of Lemma 1.4, that $`\mathrm{\Omega }_{\omega _0}`$ is cyclic for $`\pi _{\omega _0}().`$ But if $`P`$ annihilates $`_{\omega _0},`$ it must commute with $`\pi _{}()`$ on $`_{\omega _0},`$ and this contradicts Equation (8), hence $`P\mathrm{\Omega }_{\omega _0}0.`$ Similarly, if $`(IP)\mathrm{\Omega }_{\omega _0}=0,`$ then by the same argument $`(IP)`$ commutes with $`\pi _{}()`$ on $`_{\omega _0},`$ hence so does $`P.`$ So, also $`(IP)\mathrm{\Omega }_{\omega _0}0.`$ Thus we can write: $$\begin{array}{ccc}\hfill \mathrm{\Omega }_{\omega _0}& =P\mathrm{\Omega }_{\omega _0}+(IP)\mathrm{\Omega }_{\omega _0}\hfill & \\ & =\alpha \mathrm{\Omega }_P+\beta \mathrm{\Omega }_P^{}\text{where:}\hfill & \\ \hfill \mathrm{\Omega }_P& :=\frac{P\mathrm{\Omega }_{\omega _0}}{P\mathrm{\Omega }_{\omega _0}},\mathrm{\Omega }_P^{}:=\frac{(IP)\mathrm{\Omega }_{\omega _0}}{(IP)\mathrm{\Omega }_{\omega _0}}\hfill & \end{array}$$ and where $`\|\alpha \|^2+\|\beta \|^2=1`$ and $`\alpha 0\beta .`$ Adapting now the proof in Lemma 1.4, since we have that $`[\pi _{}(),P]=0`$ we find for all $`A`$: $$\begin{array}{ccc}\hfill \omega _0(A)& =\lambda \omega _P(A)+(1\lambda )\omega _P^{}(A)\text{where:}\hfill & \\ \hfill \omega _P(N)& =(\mathrm{\Omega }_P,\pi _{}(N)\mathrm{\Omega }_P),\text{and}\hfill & \\ \hfill \omega _P^{}(N)& =(\mathrm{\Omega }_P^{},\pi _{}(N)\mathrm{\Omega }_P^{}),\hfill & \end{array}$$ and we set $`\lambda :=\|\alpha \|^2.`$ That is, we have $$\theta (\omega _0)=\theta \left(\lambda \omega _P+(1\lambda )\omega _P^{}\right).$$ $`(9)`$ Now for an $`L`$ we have similar to before: $$\begin{array}{ccc}\hfill \omega _0(L)=& (\alpha \mathrm{\Omega }_P+\beta \mathrm{\Omega }_P^{},\pi _{}(L)\left(\alpha \mathrm{\Omega }_P+\beta \mathrm{\Omega }_P^{}\right))\hfill & \\ \hfill =& \lambda \omega _P(L)+(1\lambda )\omega _P^{}(L)\hfill & \\ & +\overline{\alpha }\beta (\mathrm{\Omega }_P,\pi _{}(L)\mathrm{\Omega }_P^{})+\alpha \overline{\beta }(\mathrm{\Omega }_P^{},\pi _{}(L)\mathrm{\Omega }_P).\hfill & (10)\hfill \end{array}$$ We prove that we can always find an $`L`$ to make the last line nonzero. If it is always zero, it is zero for all $`L=\gamma R`$ where $`\gamma C`$ and $`R=R^{}.`$ Thus $$2\mathrm{R}\mathrm{e}\left[\alpha \overline{\beta }\gamma (\mathrm{\Omega }_P^{},\pi _{}(R)\mathrm{\Omega }_P)\right]=0\gamma C,R=R^{}$$ and thus $`(\mathrm{\Omega }_P^{},\pi _{}(R)\mathrm{\Omega }_P)=0`$ for all $`R=R^{}.`$ But $``$ is spanned by its selfadjoint elements, and so $$(\mathrm{\Omega }_P^{},\pi _{}()\mathrm{\Omega }_P)=0\text{i.e.}\pi _{}()\mathrm{\Omega }_P\pi _{}()\mathrm{\Omega }_P^{}.$$ Thus $`\pi _{}()`$ restricted to $`[\pi _{}()\mathrm{\Omega }_P][\pi _{}()\mathrm{\Omega }_P^{}]`$ decomposes into two cyclic representations $`\pi _P\pi _P^{}.`$ Since $`\mathrm{\Omega }_P`$ is cyclic for $`\pi _P(),`$ by Lemma 1.4 it is also cyclic for $`\pi _P().`$ Thus, $$[\pi _{}()\mathrm{\Omega }_P]=[\pi _{}()\mathrm{\Omega }_P]=[\pi _{}()P\mathrm{\Omega }]=[P\pi _{}()\mathrm{\Omega }]=P_{\omega _0}$$ and likewise we get $`[\pi _{}()\mathrm{\Omega }_P^{}]=(IP)_{\omega _0}.`$ So for all $`L`$ and $`A`$ we have $$\begin{array}{ccc}\hfill P\pi _{}(L)\pi _{}(A)\mathrm{\Omega }& =P\pi _{}(LA)(\alpha \mathrm{\Omega }_P+\beta \mathrm{\Omega }_P^{})=\alpha \pi _{}(LA)\mathrm{\Omega }_P\hfill & \\ & =\pi _{}(LA)P\mathrm{\Omega }=\pi _{}(L)P\pi _{}(A)\mathrm{\Omega }.\hfill & \end{array}$$ Thus $`[P,\pi _{}(L)]P_{\omega _0}=0`$ for all $`L.`$ But this contradicts Equation (8), hence the last line in Equation (10) is nonzero for some $`L.`$ Thus $$\omega _0\lambda \omega _P+(1\lambda )\omega _P^{}$$ and this, together with Equation (9) now contradicts the assumption that $`\theta `$ is injective on $`S().`$ Thus, the initial assumption was wrong, and we conclude $`\pi _{}()^{}=\pi _{}()^{},`$ and hence $`\pi _{}()^{\prime \prime }=\pi _{}()^{\prime \prime }.`$ . The unique extension of $`\pi _{}`$ from $``$ to $``$ is of course just the representation $`\pi _{S_0}:=\underset{\omega S_0}{}\pi _\omega \mathrm{Rep}(),`$ using the definition of an ideal host. Corollary 1.6. Let $`𝒩`$ be a von Neumann algebra, and let $`S_0`$ be its set of normal states. If $``$ is an ideal host for the pair $`(𝒩,S_0),`$ then $`\pi _{S_0}(𝒩)=^{\prime \prime }`$ and hence $`\pi _{S_0}(𝒩)`$ contains $``$ as an ideal. Proof: Recall the embedding $`𝒩M()^{\prime \prime }.`$ Note that if $`\omega S(),`$ then the unique extension of $`\pi _\omega `$ to $`𝒩`$ is (unitarily equivalent to) $`\pi _{\theta (\omega )},`$ and this we see from $$\theta (\omega )(N)=\underset{\alpha }{lim}\omega (NE_\alpha )=\underset{\alpha }{lim}(\mathrm{\Omega }_\omega ,\pi _\omega (NE_\alpha )\mathrm{\Omega }_\omega )=(\mathrm{\Omega }_\omega ,\stackrel{~}{\pi }_\omega (N)\mathrm{\Omega }_\omega ).$$ Thus, recalling that the universal representation of $``$ is $`\pi _{}=\underset{\omega S()}{}\pi _\omega `$ and that $`\theta (S())=S_0,`$ we conclude that the unique extension of $`\pi _{}`$ to $`𝒩`$ is $`\pi _{}𝒩=\underset{\omega S_0}{}\pi _\omega =\pi _{S_0}.`$ Thus $`\pi _{S_0}`$ on $`𝒩`$ is a normal representation, and so $$\pi _{S_0}(𝒩)=\pi _{S_0}(𝒩^{\prime \prime })=\pi _{S_0}(𝒩)^{\prime \prime }=^{\prime \prime }$$ where we used Theorem 1.5 for the last equality. . Corollary 1.7. If $``$ is an ideal host for a pair $`(,S_0),`$ then $`\theta :S()S_0`$ maps the pure states of $``$ to pure states on $`.`$ Hence all extreme points of $`S_0`$ are pure. Proof: By Theorem 1.5 we have $`\pi _{S_0}()^{\prime \prime }=\pi _{S_0}()^{\prime \prime }.`$ Let $`\omega _0S_0,`$ then since $`\pi _{S_0}=\underset{\omega S_0}{}\pi _\omega ,`$ we have that the restriction map $`R:\pi _{S_0}()\pi _{\omega _0}()`$ by $`R(\pi _{S_0}(F)):=\pi _{\omega _0}(F)`$ for all $`F`$ is a normal \–homomorphism, hence it extends to $`\pi _{S_0}()^{\prime \prime }=\pi _{}()^{\prime \prime }`$ and so $$\begin{array}{ccc}\hfill R\left(\pi _{S_0}()^{\prime \prime }\right)& =R\left(\pi _{S_0}()\right)^{\prime \prime }=\pi _{\omega _0}()^{\prime \prime }\hfill & \\ & =R\left(\pi _{}()\right)^{\prime \prime }=\pi _{\omega _0}()^{\prime \prime }.\hfill & \end{array}$$ Or to be more notationally precise, $`\pi _{\omega _0}()^{\prime \prime }=\pi _{\theta ^1(\omega _0)}()^{\prime \prime }.`$ Now $`\theta `$ is an affine bijection, so it restricts to a bijection between the pure states of $`S()`$ and the extreme points of $`S_0.`$ If $`\phi `$ is a pure state on $`,`$ then $`(_\phi )=\pi _\phi ()^{\prime \prime }=\pi _{\theta (\phi )}()^{\prime \prime },`$ i.e. $`\theta (\phi )`$ is also pure on $`.`$ . This last corollary now severely limits the class of folia for which ideal hosts exist, indeed, we can quickly prove many obstruction theorems, e.g. the next one. Corollary 1.8. If $`𝒩`$ is a simple factor and $`S_0`$ is its folium of normal states, then there is no ideal host for the pair $`(𝒩,S_0).`$ Proof: it suffices by Corollary 1.7 to observe that $`S_0`$ contains no pure states. For if an $`\omega S_0`$ were pure, then using the fact that $`\pi _\omega `$ is normal, we see that $`\pi _\omega (𝒩)=\pi _\omega (𝒩)^{\prime \prime }=(_\omega ).`$ Since $`𝒩`$ is simple, $`\pi _\omega `$ is an isomorphism hence $`𝒩(_\omega )`$ and the latter is not simple. This is a contradiction, so $`S_0`$ has no pure states. . Remarks. (1) Thus, since we know from Kadison and Ringrose 6.8.4 and 6.6.5 \[KR\], that every finite factor and each countably decomposable type III factor is simple, Corollary 1.8 shows that there are many von Neumann algebras for which there is no ideal host for its normal states. From Corollary 1.6 we see that we can only expect ideal hosts for a von Neumann algebra which has a norm closed proper ideal which is weak operator dense. If $``$ is a C\–algebra, but not a von Neumann algebra, Theorem 1.5 just tells us that we should be looking for an ideal host $``$ in the von Neumann algebra $`\pi _{S_0}()^{\prime \prime },`$ but by Corollary 1.7 we also know that this may fail, unless all the extreme points of $`S_0`$ are pure. (2) Theorem 1.5 states a “weak” uniqueness, in that it claims that all ideal hosts for the same pair must have the same universal algebra $$^{\prime \prime }=\pi _{S_0}()^{\prime \prime }.$$ Theorem 1.5 also tells us where all ideal hosts reside, viz $`\pi _{S_0}()^{\prime \prime }(_{S_0}),`$ and we can make this even more precise: Since $``$ is an ideal host, $`S_0`$ is a folium, and so for the projection associated with $`S_0`$ we have $`P^{}^{\prime \prime }`$ (cf. Remark 1.3) and $`S_0=\left\{\omega S()\right\|\stackrel{~}{\omega }(P)=1\}`$ where $`\stackrel{~}{\omega }`$ denotes the normal extension of $`\omega `$ from $``$ to $`^{\prime \prime }.`$ Thus in $`\pi _{}`$ we have $`P\mathrm{\Omega }_\omega =\mathrm{\Omega }_\omega `$ for all $`\omega S_0,`$ hence by $`P^{}^{\prime \prime }`$ we see $`PF\mathrm{\Omega }_\omega =F\mathrm{\Omega }_\omega `$ for all $`F^{\prime \prime }.`$ Thus $`P`$ is the projector onto $`\underset{\omega S_0}{}_\omega =_{S_0}_{},`$ hence $`P\pi _{}()^{\prime \prime }=\pi _{S_0}()^{\prime \prime }`$ and so we conclude from $`P^{\prime \prime }`$ that $$\pi _{S_0}()^{\prime \prime }\pi _{}()^{\prime \prime }=^{\prime \prime }.$$ $`(11)`$ (3) Since for an ideal host $``$ for a pair $`(,S_0)`$ we know by Theorem 1.5 that $`\xi ()M()^{\prime \prime }=\pi _{}()^{\prime \prime },`$ we conclude that the embedding $`\xi :M()`$ is precisely $`\pi _{S_0}.`$ So $`\xi `$ is faithful iff $`\pi _{S_0}`$ is faithful iff for each $`F`$ we have $`\omega (F)0`$ for some $`\omega S_0.`$ (4) We can easily adapt the proofs of 1.4 and 1.5 to prove a Stone–Weierstrass theorem for von Neumann algebras, i.e. if a von Neumann algebra contains a sub–von Neumann algebra which separates its normal states, then they are equal. However, this fact has also a very short proof via the bipolar theorem (Private communication with R.Longo). (5) By Corollary 1.7 we can also find topological goups which have no group algebras (i.e. ideal hosts for the the pair $`(C^{}(G_d),S_0)`$ where $`S_0`$ are states $`\omega `$ for which $`g\omega (\delta _g)`$ is continuous). For example, let $`G=L^{\mathrm{}}(R)`$ with the group operation being addition, and with the strong operator topology w.r.t. its representation as multiplication operators on $`L^2(R).`$ Then there is a topological isomorphism $`L^{\mathrm{}}(R)C(R_s)`$ where $`R_s`$ is $`R`$ compactified and endowed with a suitable hyperstonean topology (cf. proof of Theorem III.1.18 \[Tak\], or Theorem 2.1 below). Any irreducible representation of $`G`$ must be a character, hence point evaluation on $`C(R_s),`$ and this cannot be continuous w.r.t. the strong operator topology, because points are still of measure zero w.r.t. the extension of the Lebesgue measure to $`R_s.`$ Thus the folium $`S_0`$ of states of $`C^{}(G_d)`$ which are continuous w.r.t. the topology of $`G`$ contains no pure states, hence by Corr. 1.7 we conclude that G has no group algebra. To conclude this section, we would like to make precise the relation between the representations of a host $``$ for a pair $`(,S_0)`$ and the representations of $`.`$ Denote the normal representations of $`\pi _{S_0}()`$ by $`\mathrm{Rep}_{S_0}.`$ By Theorem 1.5 we know that $`^{\prime \prime }\pi _{S_0}()^{\prime \prime },`$ hence for each $`\pi \mathrm{Rep}_{S_0}`$ we can construct a representation $`\mathrm{\Lambda }(\pi )\mathrm{Rep}`$ by first extending $`\pi `$ via strong operator continuity to a representation $`\stackrel{~}{\pi }\mathrm{Rep}\left(\pi _{S_0}()^{\prime \prime }\right),`$ and then defining $`\mathrm{\Lambda }(\pi )`$ as $`\stackrel{~}{\pi }`$ restricted to its essential subspace. This produces a map $`\mathrm{\Lambda }:\mathrm{Rep}_{S_0}\mathrm{Rep}.`$ Theorem 1.9. If $``$ is an ideal host for the pair $`(,S_0),`$ then the map $`\mathrm{\Lambda }:\mathrm{Rep}_{S_0}\mathrm{Rep}`$ is a bijection which takes irreducible representations to irreducible representations. In the case that $``$ is merely a host, $`\mathrm{\Lambda }`$ is a bijection modulo unitary equivalence. Its inverse is via inducing of representations;- $`\mathrm{\Lambda }^1(\{\pi \})=\left\{\mathrm{Ind}_{}^{}(\pi )\right\}`$ where $`\{\}`$ denotes unitary equivalence classes, and the induction is done via the right$`\text{–rigged}`$left$`\text{–module}`$ $`:=[]`$ with rigging map $`u,v:=u^{}v`$ for all $`u,v.`$ Proof: If $``$ is an ideal host, $`^{\prime \prime }=\pi _{S_0}()^{\prime \prime },`$ and now as both $``$ and $`\pi _{S_0}()`$ are strong operator dense in $`^{\prime \prime },`$ it is obvious that each uniquely determines a normal representation, and so the proof for this case ends here. For the case of $``$ just a host, let $`(\pi ,)\mathrm{Rep}_{S_0},`$ so $`\mathrm{\Lambda }(\pi )`$ is $`\stackrel{~}{\pi }`$ restricted to its essential subspace $`[\stackrel{~}{\pi }()].`$ Now for our proof, we will first construct $`\mathrm{Ind}_{}^{}(\mathrm{\Lambda }(\pi )),`$ show it is unitary equivalent to $`\pi ,`$ and then show that every $`\pi \mathrm{Rep}_{S_0}`$ is unitarily equivalent to some $`\mathrm{Ind}_{}^{}(\gamma ),`$ $`\gamma \mathrm{Rep}.`$ Following Fell and Doran XI.4.12 \[FD\] or Rieffel \[Ri\], consider the right$`\text{–rigged}`$left$`\text{–module}`$ $$:=[],u,v:=u^{}vu,v$$ where we used the fact that $``$ is hereditary in $``$ to conclude $`u^{}v.`$ Since $``$ is a C\–algebra, every representation of it is inducible via $``$ (cf. XI.4.12 \[FD\]). Now construct $`\rho =\mathrm{Ind}_{}^{}(\gamma )\mathrm{Rep}()`$ as follows. On $``$ define a pre–inner product $`(,)_0`$ by $$(s\xi ,t\eta )_0:=(\gamma \left(t,s\right)\xi ,\eta )=(\gamma \left(t^{}s\right)\xi ,\eta )$$ $`(12)`$ and define from $``$ the Hilbert space $$𝒦:=\overline{/\mathrm{Ker}(,)_0}$$ where closure is obviously w.r.t. $`(,)_0.`$ Denote the image of an elementary tensor $`s\xi `$ in $`𝒦`$ by $`s\stackrel{~}{\mathrm{}}\xi ,`$ and define the representation $`\rho :(𝒦)`$ by $$\rho (F)(s\stackrel{~}{\mathrm{}}\xi ):=Fs\stackrel{~}{\mathrm{}}\xi s,\xi ,F.$$ If we let $`\gamma =\mathrm{\Lambda }(\pi ),`$ then Equation (12) becomes $$(s\xi ,t\eta )_0:=(\stackrel{~}{\pi }(s)\xi ,\stackrel{~}{\pi }(t)\eta )$$ and so we can identify $`𝒦`$ with the subspace $`[\stackrel{~}{\pi }()]=[\stackrel{~}{\pi }()],`$ via the unitary $`U(s\xi ):=\stackrel{~}{\pi }(s)\xi .`$ Moreover we have for all $`s\stackrel{~}{\mathrm{}}\xi ,`$ $`t\stackrel{~}{\mathrm{}}\eta `$ that $$\begin{array}{ccc}\hfill \left(\rho (F)(s\stackrel{~}{\mathrm{}}\xi ),t\stackrel{~}{\mathrm{}}\eta \right)& =\left(Fs\stackrel{~}{\mathrm{}}\xi ,t\stackrel{~}{\mathrm{}}\eta \right)=(\stackrel{~}{\pi }(Fs)\xi ,\stackrel{~}{\pi }(t)\eta )\hfill & \\ & =(\pi (F)\stackrel{~}{\pi }(s)\xi ,\stackrel{~}{\pi }(t)\eta )\hfill & \end{array}$$ and so $`(\rho ,𝒦)`$ is unitarily equivalent to $`\pi `$ on $`[\stackrel{~}{\pi }()].`$ Now application of Lemma 1.10 (proven below) to $`\stackrel{~}{\pi }\mathrm{Rep}`$ implies that we have $`=[\stackrel{~}{\pi }()]`$ iff $`\pi \mathrm{Rep}_{S_0},`$ and the latter is what we assumed at the start. Thus $`\rho `$ is unitarily equivalent to $`\pi .`$ Furthermore, we see above that the induction process produce representations $`\stackrel{~}{\pi }`$ such that $`=[\stackrel{~}{\pi }()],`$ hence the image under induction via $``$ of $`\mathrm{Rep}`$ is $`\mathrm{Rep}_{S_0},`$ using Lemma 1.10 again. . Lemma 1.10. A representation $`(\pi ,)\mathrm{Rep}`$ satisfies $`=[\pi ()]`$ iff $`\pi `$ is normal with respect to $`\pi _{S_0}().`$ Proof: Let $`=[\pi ()]`$ and choose a normalised vector $`\psi \pi (),`$ i.e. $`\psi =\pi (L)\xi ,`$ $`\psi =1.`$ Let $`\omega _\psi `$ denote the associated vector state on $`,`$ then since $`\omega _\psi =1`$ and $``$ is hereditary, it is everywhere determined on $``$ by its values on $`.`$ Thus $`\omega _\psi =\theta (\omega _\psi )S_0\left(\pi _{S_0}()\right)_{}.`$ Since the normal functionals of any representation is a folium, hence invariant under conjugation, we conclude that also $`\omega _{\pi (F)\psi }=\omega _{\pi \left(FL\right)\xi }\left(\pi _{S_0}()\right)_{}`$ for all $`L,`$ $`F`$ and $`\xi ,`$ i.e. $`\omega _{\pi \left(\right)}\left(\pi _{S_0}()\right)_{}.`$ Since by assumption $`\pi ()`$ spans a dense subspace of $`,`$ it follows from Kadison and Ringrose 7.1.15 \[KR\] that $`\pi `$ is normal with respect to $`\pi _{S_0}().`$ Conversely, let $`\pi `$ be normal with respect to $`\pi _{S_0}(),`$ and assume that $`[\pi ()],`$ i.e. there is some $`\psi [\pi ()].`$ First, we show that $`\omega _\psi \left(\pi _{S_0}()\right)_{}.`$ If not, then $`\theta ^1(\omega _\psi )S(),`$ i.e. the normal extension $`\stackrel{~}{\omega }_\psi `$ to $`^{\prime \prime }`$ restricts to a state on $`.`$ Since this normal extension on $``$ is just $`\stackrel{~}{\omega }_\psi (A)=(\psi ,\pi (A)\psi ),`$ we conclude from the given $`\psi [\pi ()]\pi ()\psi `$ that $`\stackrel{~}{\omega }_\psi ()=0,`$ which contradicts the fact that it must be a state on $`.`$ Thus $`\omega _\psi \left(\pi _{S_0}()\right)_{}`$. Now we know by the first part that $`\omega _{\pi \left(\right)}\left(\pi _{S_0}()\right)_{}`$, and in fact since the normal functionals is a norm–closed folium, we have $`\omega _{\left[\pi \left(\right)\right]}\left(\pi _{S_0}()\right)_{}`$, so since $`=C\psi [\pi ()]`$ it is impossible to find a dense subspace $`𝒮`$ such that $`\omega _\phi \left(\pi _{S_0}()\right)_{}`$ for all $`\phi 𝒮,`$ hence by Kadison and Ringrose 7.1.15 \[KR\] $`\pi `$ cannot be normal w.r.t. $`\pi _{S_0}().`$ This contradicts our hypothesis, hence $`=[\pi ()].`$ . ## 2. Ideal hosts up to a central algebra. Above we saw that a pair $`(,S_0)`$ with $`S_0`$ a folium without pure states, has no ideal host. A particularly bad case of this, is the pair $`(L^{\mathrm{}}(X,\mu ),S_N)`$ where $`\mu `$ has no discrete part, and where $`S_N`$ denotes the set of normal states, i.e. the measures absolutely continuous w.r.t. $`\mu .`$ In this case $`S_N`$ does not even have extreme points, because for any measure $`\nu `$ absolutely continuous w.r.t. $`\mu ,`$ we only need to subdivide its support to write it as a convex combination of other probability measures in this class. Only when $`\mathrm{supp}(\nu )`$ is a point can we not do this, and this case does not occur since $`\mu `$ has no discrete part. In this section we want to argue that this example is symptomatic of the general case, in that if a pair $`(,S_0)`$ has no ideal host, it is because $`(L^{\mathrm{}}(X,\mu ),S_N)`$ is embedded in it, and it acts as an obstruction. To be precise about the embedding, we will show that for any pair $`(,S_0)`$ with $`S_0`$ a folium, we can always find a quasi–host $`\stackrel{~}{}`$ in the sense of the next definition: Def. Given a pair $`(,S_0)`$ where $`S_0`$ is a folium, a quasi–host for it, is a C\–algebra $`\stackrel{~}{}`$ and two embeddings $`M(\stackrel{~}{})`$ and $`L^{\mathrm{}}(X,\mu )ZM(\stackrel{~}{})`$ for some measure space $`(X,\mu )`$ such that $`\stackrel{~}{S}_\mu =S_0`$ where $`S_\mu :=\left\{\omega S(\stackrel{~}{})\right\|\stackrel{~}{\omega }L^{\mathrm{}}(X,\mu )\text{is normal}\}`$ and moreover, $`\stackrel{~}{S}_\mu C^{}(L^{\mathrm{}}(X,\mu ))`$ defines an injection for $`S_\mu .`$ We will show that the given pair $`(,S_0)`$ has no ideal host if the measure $`\mu `$ is purely continuous. This is what we mean by saying $`L^{\mathrm{}}(X,\mu )`$ acts as an obstruction. Example. Consider the von Neumann algebra $`=L^{\mathrm{}}(X,\mu )\overline{\mathrm{}}()`$ acting on the Hilbert space $`L^2(X,\mu ),`$ where $`\overline{\mathrm{}}`$ denotes the W\–tensor product (cf. 11.2 \[KR\]). Its pure states consists of product states $`\omega _1\omega _2`$ such that $`\omega _i`$ are both pure. Thus, if we take the pair $`(,S_0)`$ where $`S_0`$ are the normal states of $`,`$ and assume that $`\mu `$ has no discrete part, then $`S_0`$ has no pure states, hence this pair has no ideal host. Nevertheless, the algebra $`\stackrel{~}{}=L^{\mathrm{}}(X,\mu )𝒦()`$ is a quasi–host for $`(,S_0).`$ To start the analysis, let $`(,S_0)`$ be a pair with $`S_0`$ a folium, then by Theorem 1.2 we construct the representation $`\pi _{S_0}=\underset{\omega S_0}{}\pi _\omega `$ and identify $`S_0`$ with the normal states of the concrete algebra $`\pi _{S_0}(),`$ hence with the normal states of the von Neumann algebra $`\pi _{S_0}()^{\prime \prime }.`$ We first analyze a single cyclic component $`\pi _\omega ,`$ $`\omega S_0`$ of the direct sum. We denote the von Neumann algebra $`\pi _\omega ()^{\prime \prime }`$ by $`𝒩,`$ and its set of normal states by $`S_N.`$ Central to the following constructions, is the usual decomposition theory with respect to some commutative subalgebra $`𝒞𝒩^{}`$ (cf. Takesaki \[Tak\]), however, since the maps involved are only defined up to $`\mu \text{–negligible}`$ sets on some measure space $`(X,\mu ),`$ and we will actually need everywhere defined maps, we now redo some of the basic constructions in order to remedy this. Theorem 2.1. Let $``$ be a C\–algebra with a fixed state $`\omega S(),`$ and a commutative unital C\–algebra $`𝒞\pi _\omega ()^{},`$ then the Gel’fand isomorphism $`\mathrm{\Phi }:𝒞^{\prime \prime }L^{\mathrm{}}(X,\mu )`$ equips the spectrum $`X`$ of $`𝒞`$ with a probability measure $`\mu ,`$ $`\mathrm{supp}\mu =X,`$ and there is a map $`\psi :XS(𝒞^{})`$ such that (i) the map $`x\psi _x(A)`$ is in $`L^{\mathrm{}}(X,\mu )`$ for all $`A𝒞^{},`$ and if $`𝒞`$ is a von Neumann algebra, $`x\psi _x(A)`$ is in $`C(X)`$ for all $`A𝒞^{},`$ (ii) $`\stackrel{~}{\omega }(A):=(\mathrm{\Omega }_\omega ,A\mathrm{\Omega }_\omega )=_X\psi _x(A)𝑑\mu (x)`$ $`A𝒞^{},`$ (iii) $`\psi _x(CA)=\mathrm{\Phi }(C)(x)\psi _x(A)`$ $`C𝒞^{\prime \prime }`$, $`A𝒞^{}.`$ Proof: (This proof is based on Takesaki 6.23, p241 \[Tak\]) Since $`\mathrm{\Omega }_\omega `$ is cyclic for $`\pi _\omega (),`$ it is separating for $`𝒞`$ and $`𝒞^{\prime \prime },`$ so $`\stackrel{~}{\omega }\mathrm{}𝒞`$ is a faithful state of $`𝒞.`$ Thus by the Riesz representation theorem there is a Borel measure $`\mu `$ on the compact set $`X`$ such that $$\stackrel{~}{\omega }(C)=_X\mathrm{\Phi }(C)(x)𝑑\mu (x),C𝒞$$ with $`\mathrm{\Phi }:𝒞C(X)`$ the Gel’fand isomorphism. Moreover $`\mu `$ is a probability measure since $`\stackrel{~}{\omega }(I)=1`$ and $`\mathrm{supp}\mu =X`$ because $`\stackrel{~}{\omega }`$ is faithful. Since $`\mathrm{\Omega }_\omega `$ separates $`𝒞,`$ we can consistently define a unitary $`U:[𝒞\mathrm{\Omega }_\omega ]L^2(X,\mu )`$ by $`U(C\mathrm{\Omega }_\omega ):=\mathrm{\Phi }(C)`$ which produces the representation $`\stackrel{~}{\mathrm{\Phi }}:𝒞(L^2(X,\mu ))`$ by $`\stackrel{~}{\mathrm{\Phi }}(C):=UCU^1=T_{\mathrm{\Phi }\left(C\right)}`$ in terms of multiplication operators $`\left\{T_f\right\|fC(X)\}.`$ Since unitary conjugation is a normal map, $`\stackrel{~}{\mathrm{\Phi }}`$ extends to $`𝒞^{\prime \prime }\pi _\omega ()^{}`$ and $`\stackrel{~}{\mathrm{\Phi }}(𝒞)^{\prime \prime }=\stackrel{~}{\mathrm{\Phi }}(𝒞^{\prime \prime }),`$ i.e. $`\stackrel{~}{\mathrm{\Phi }}(𝒞^{\prime \prime })=\left\{T_f\right\|fL^{\mathrm{}}(X,\mu )\}.`$ If $`𝒞`$ is already a von Neumann algebra, this simplifies to $$\left\{T_f\right\|fC(X)\}=\stackrel{~}{\mathrm{\Phi }}(𝒞)=\stackrel{~}{\mathrm{\Phi }}(𝒞)^{\prime \prime }=\left\{T_f\right\|fL^{\mathrm{}}(X,\mu )\},$$ so for each $`fL^{\mathrm{}}(X,\mu )`$ we can find a $`\stackrel{~}{f}C(X)`$ such that $`T_f=T_{\stackrel{~}{f}}`$ on $`L^2(X,\mu ),`$ producing the isomorphism $`C(X)L^{\mathrm{}}(X,\mu )`$ (this isomorphism also occurs in the proof of Theorem III.1.18 \[Tak\]). Henceforth we will blur the distinction between $`\mathrm{\Phi }`$ and $`\stackrel{~}{\mathrm{\Phi }}`$ and always make the identification with $`C(X)`$ if $`𝒞`$ is a von Neumann algebra. Let $`e`$ be the projection of $`_\omega `$ onto $`[𝒞\mathrm{\Omega }_\omega ]=[𝒞^{\prime \prime }\mathrm{\Omega }_\omega ]`$, then $`e𝒞^{}`$ because $`𝒞`$ preserves $`[𝒞\mathrm{\Omega }_\omega ]`$. Define $`\mathrm{{\rm Y}}:𝒞^{\prime \prime }\left([𝒞\mathrm{\Omega }_\omega ]\right)`$ by $`\mathrm{{\rm Y}}(C):=eC`$, then $`\mathrm{{\rm Y}}(𝒞^{\prime \prime })=e𝒞^{\prime \prime }`$ is maximally commutative in $`\left([𝒞\mathrm{\Omega }_\omega ]\right)`$ because it has a cyclic vector $`\mathrm{\Omega }_\omega `$, cf. Takesaki Corr. 1.3, p104 \[Tak\]. Moreover $`\mathrm{{\rm Y}}`$ is injective because $`\mathrm{\Omega }_\omega `$ is separating for $`𝒞^{\prime \prime }`$. Now $`(e𝒞)^{}\mathrm{}[𝒞\mathrm{\Omega }_\omega ]=e𝒞^{}e([𝒞\mathrm{\Omega }_\omega ])`$ (easily verified, but also in Takesaki 3.10 \[Tak\]), hence $`e𝒞^{}ee𝒞^{\prime \prime }`$ because the latter is maximally commutative in $`([𝒞\mathrm{\Omega }_\omega ])`$ (also by Takesaki 3.10). Define $`\delta :𝒞^{}\mathrm{{\rm Y}}(𝒞^{\prime \prime })`$ by $`\delta (A):=eAee𝒞^{\prime \prime }`$. Clearly $`\delta (I)=e`$ and $`\delta `$ is positive and extends $`\mathrm{{\rm Y}}`$. Next define $`\psi :XS(𝒞^{})`$ by $$\psi _x(A):=\mathrm{\Phi }(\mathrm{{\rm Y}}^1(\delta (A)))(x)$$ and note $`\mathrm{\Phi }(\mathrm{{\rm Y}}^1(\delta (A)))`$ is in $`L^{\mathrm{}}(X,\mu )`$, so $`\psi `$ is defined $`\mu \text{–almost}`$ everywhere, and if $`𝒞`$ is a von Neumann algebra we can identify $`\mathrm{\Phi }(\mathrm{{\rm Y}}^1(\delta (A)))`$ with an element of $`C(X)`$ in which case $`\psi `$ is everywhere defined. By positivity and normalisation of the various maps each $`\psi _x`$ is a state. Now $$\begin{array}{ccc}\hfill \stackrel{~}{\omega }(\mathrm{{\rm Y}}^1(\delta (A)))& =_X\mathrm{\Phi }(\mathrm{{\rm Y}}^1(\delta (A)))(x)𝑑\mu (x)\hfill & \\ & =_X\psi _x(A)𝑑\mu (x)\hfill & \\ \hfill \stackrel{~}{\omega }(A)=(\mathrm{\Omega }_\omega ,A\mathrm{\Omega }_\omega )& =(\mathrm{\Omega }_\omega ,eAe\mathrm{\Omega }_\omega )\text{using }e\mathrm{\Omega }_\omega =\mathrm{\Omega }_\omega \hfill & \\ & =(\mathrm{\Omega }_\omega ,\delta (A)\mathrm{\Omega }_\omega )=(\mathrm{\Omega }_\omega ,\mathrm{{\rm Y}}^1(\delta (A))\mathrm{\Omega }_\omega )\hfill & \\ & =\stackrel{~}{\omega }(\mathrm{{\rm Y}}^1(\delta (A)))=_X\psi _x(A)𝑑\mu (x).\hfill & \end{array}$$ For (ii), observe that for $`C𝒞^{\prime \prime }`$, $`A𝒞^{}`$ we have $$\psi _x(CA)=\mathrm{\Phi }(\mathrm{{\rm Y}}^1(\delta (CA)))(x)=\mathrm{\Phi }(C)(x)\psi _x(A)$$ simply using the fact that $`\mathrm{\Phi },\mathrm{{\rm Y}}^1`$ and $`\delta `$ are homomorphisms and by definition $`\mathrm{{\rm Y}}^1(\delta (C)=C`$ for $`C𝒞^{\prime \prime }`$. . Remarks. (1) A very important point here is that for $`𝒞`$ a von Neumann algebra, $`\psi :XS(^{})`$ is everywhere defined. This comes from the identification which $`\mathrm{\Phi }`$ makes between $`C(X)`$ and $`L^{\mathrm{}}(X,\mu ).`$ The spectrum $`X`$ of $`𝒞`$ in this case is of course hyperstonean hence extremely disconnected (cf. \[Tak\]). (2) Note that no separability assumption was needed here. (3) The map $`\psi `$ and measure $`\mu `$ produces precisely the orthogonal barycentric decomposition of $`\stackrel{~}{\omega }`$ used in decomposition theory in e.g. Takesaki \[Tak\] or Bratteli and Robinson \[BR\]. One uses $`\psi `$ to write $`\mu `$ as a measure on the state space $`S().`$ (4) By restricting $`\stackrel{~}{\omega }`$ and $`\psi _x`$ to any subalgebra of $`𝒞^{}`$ (e.g. $`\pi _\omega ()`$ or $`\pi _\omega ()^{\prime \prime }`$), we obtain a decomposition theorem for states on these. Note that the usual theory assumes that $`𝒞`$ is a von Neumann algebra. Def. Given the data $`,\omega ,𝒞`$ of theorem 2.1 assume that $`𝒞`$ is a von Neumann algebra with subsequent map $`\psi :XS(𝒞^{})`$ and define the following two bundles on $`X`$: $``$ $`(\psi )`$ is the bundle with projection $`p:(\psi )X`$ by $`p^1(x)=_{\psi _x}`$, $``$ $`(\psi )`$ the bundle with $`q:(\psi )X`$ by $`q^1(x)=(_{\psi _x})`$. At this point there is no topology, so the total spaces of these bundles are just the unions of their fibres. Clearly the sections $`\mathrm{\Gamma }((\psi ))`$ act pointwise on the sections $`\mathrm{\Gamma }((\psi ))`$. Now there are some canonical families of sections: $``$$`\mathrm{\Omega }\mathrm{\Gamma }((\psi ))`$ is the section $`\mathrm{\Omega }(x):=\mathrm{\Omega }_{\psi _x}`$, $``$$`\mathrm{\Pi }(A)\mathrm{\Gamma }((\psi ))`$ is the section $`\mathrm{\Pi }(A)(x):=\pi _{\psi _x}(A)`$, $`A𝒞^{}`$. Clearly the latter specialises on $`𝒞`$ to $``$ $`\mathrm{\Pi }(C)(x)\phi =\mathrm{\Phi }(C)(x)\phi \phi _{\psi _x},C𝒞`$. Let $`c=\mathrm{\Pi }(C)\mathrm{\Omega }`$, $`d=\mathrm{\Pi }(D)\mathrm{\Omega }`$ with $`C,D𝒞^{}`$, then $$(c(x),d(x))_{_{\psi _x}}=(\mathrm{\Pi }(C)\mathrm{\Omega }(x),\mathrm{\Pi }(D)\mathrm{\Omega }(x))_{_{\psi _x}}=\psi _x(C^{}D)$$ which is a continuous function in $`x`$ by theorem 2.1. Obviously $`\mathrm{\Pi }(𝒞^{})\mathrm{\Omega }`$ is a linear space, and we have that $`x(c(x),d(x))_{_{\psi _x}}`$ is continuous, hence integrable (since $`X`$ compact, $`\mu `$ a probability measure). Thus we can equip $`\mathrm{\Pi }(𝒞^{})\mathrm{\Omega }`$ with the inner product $$(c,d):=_X(c(x),d(x))_{_{\psi _x}}𝑑\mu (x).$$ and hence form the Hilbert space $$_\mathrm{\Gamma }:=\overline{\mathrm{\Pi }(𝒞^{})\mathrm{\Omega }/\mathrm{Ker}(,)}$$ with factorisation map $`\kappa :\mathrm{\Pi }(𝒞^{})\mathrm{\Omega }_\mathrm{\Gamma }`$. Theorem 2.2. $`\mathrm{\Pi }(𝒞^{})`$ lifts through $`\kappa `$ to define a representation $`\pi :𝒞^{}(_\mathrm{\Gamma })`$. Then $`\kappa (\mathrm{\Omega })`$ is cyclic for $`\pi ()`$, and there is a unitary $`U:_\omega _\mathrm{\Gamma }`$ which intertwines the representation $`(\pi _{\stackrel{~}{\omega }},_{\stackrel{~}{\omega }},\mathrm{\Omega }_{\stackrel{~}{\omega }})`$ of $`𝒞^{}`$ with $`(\pi ,_\mathrm{\Gamma },\kappa (\mathrm{\Omega }))`$. Proof: Obviously $`\mathrm{\Pi }(𝒞^{})`$ preserves $`\mathrm{\Pi }(𝒞^{})\mathrm{\Omega }`$;- we show that it preserves $`\mathrm{Ker}(,)`$. Let $`\phi =\mathrm{\Pi }(A)\mathrm{\Omega }\mathrm{Ker}(,)`$, i.e. $$\begin{array}{ccc}\hfill 0=\phi ^2=_X(\mathrm{\Omega }(x),\mathrm{\Pi }(A^{}A)(x)\mathrm{\Omega }(x))𝑑\mu (x)=_X\psi _x(A^{}A)𝑑\mu (x)& & \end{array}$$ Thus $`\psi _x(A^{}A)=0`$ by positivity and continuity of this map. Let $`B𝒞^{}`$, then $$\mathrm{\Pi }(B)\phi ^2=_X\psi _x(A^{}B^{}BA)𝑑\mu (x).$$ By the Cauchy–Schwartz inequality $`\psi _x(A^{}B^{}BA)=0`$ using $`\psi _x(A^{}A)=0`$. Thus $`\mathrm{\Pi }(B)`$ preserves $`\mathrm{Ker}(,)`$ and hence lifts through $`\kappa `$. To show that $`\kappa (\mathrm{\Omega })`$ is cyclic for $`\mathrm{\Pi }()`$, assume the converse, i.e. there is a sequence $`\{\phi _n\}\mathrm{\Pi }(𝒞^{})\mathrm{\Omega }`$ converging with respect to the seminorm $`\phi :=\left[_X\phi (x)_{_{\psi _x}}𝑑\mu (x)\right]^{1/2}`$ and $`(\phi _n,\mathrm{\Pi }()\mathrm{\Omega })0`$ as $`n\mathrm{}`$, but $`\phi _n`$ does not converge to zero. (This means there is some nonzero $`\stackrel{~}{\phi }_\mathrm{\Gamma }`$ which is orthogonal to $`\mathrm{\Pi }()\mathrm{\Omega }`$). Let $`\phi _n=\mathrm{\Pi }(A_n)\mathrm{\Omega },`$ $`A_n𝒞^{},`$ then $$\begin{array}{ccc}\hfill (\phi _n,\mathrm{\Pi }(B)\mathrm{\Omega })& =_X\psi _x(A_n^{}B)𝑑\mu (x)=(\mathrm{\Omega }_\omega ,A_n^{}B\mathrm{\Omega }_\omega )\hfill & \\ & =(A_n\mathrm{\Omega }_\omega ,B\mathrm{\Omega }_\omega )0B\pi _\omega ().\hfill & \end{array}$$ However $`\mathrm{\Omega }_\omega `$ is cyclic for $`\pi _\omega ()`$ in $`_\omega ,`$ hence $`A_n\mathrm{\Omega }_\omega 0`$ and so since $`\phi _n^2=A_n\mathrm{\Omega }_\omega ^20`$, we have contradicted the hypothesis, so $`\kappa (\mathrm{\Omega })`$ must be cyclic for $`\mathrm{\Pi }()`$. Now since $$(\mathrm{\Omega }_\omega ,A\mathrm{\Omega }_\omega )=\stackrel{~}{\omega }(A)=_X\psi _x(A)𝑑\mu (x)=(\mathrm{\Omega },\mathrm{\Pi }(A)\mathrm{\Omega })$$ the unitary equivalence follows from the GNS–theorem. . Remarks. (i) Since $`(\pi _{\stackrel{~}{\omega }},_{\stackrel{~}{\omega }},\mathrm{\Omega }_{\stackrel{~}{\omega }})`$ is the representation in which $`𝒞^{}`$ is defined, $`\mathrm{\Pi }`$ is an isomorphism. (ii) We can realise vectors $`\xi _\mathrm{\Gamma }`$ as sections $`x\xi (x)_{\psi _x}`$ (almost everywhere) such that $$(\xi ,\varphi )=_X(\xi (x),\varphi (x))_{_{\psi _x}}𝑑\mu (x),$$ but we will not need it here, so omit it. Def. Let $`\phi ,\zeta \mathrm{\Pi }(𝒞^{})\mathrm{\Omega }\mathrm{\Gamma }((\psi ))`$ and define a section $`\stackrel{~}{A}_{\phi ,\zeta }\mathrm{\Gamma }((\psi ))`$ by $$\left(\stackrel{~}{A}_{\phi ,\zeta }c\right)(x):=\phi (x)(\zeta (x),c(x))_{_{\psi _x}}c\mathrm{\Gamma }((\psi )).$$ Lemma 2.3. $`\stackrel{~}{A}_{\phi ,\zeta }`$ preserves $`\mathrm{\Pi }(𝒞^{})\mathrm{\Omega }`$ and $`\mathrm{Ker}(,)`$, and is bounded on these spaces, hence defines an operator $`A_{\phi ,\zeta }(_\mathrm{\Gamma })`$ by $`A_{\phi ,\zeta }\kappa (\eta ):=\kappa (\stackrel{~}{A}_{\phi ,\zeta }\eta )`$ for $`\eta \mathrm{\Pi }(𝒞^{})\mathrm{\Omega }`$. Proof: Let $`\phi =\mathrm{\Pi }(A)\mathrm{\Omega }`$, $`\zeta =\mathrm{\Pi }(B)\mathrm{\Omega }`$, $`\eta =\mathrm{\Pi }(C)\mathrm{\Omega }`$, $`A,B,C𝒞^{}`$, then $$\begin{array}{ccc}\hfill \left(\stackrel{~}{A}_{\phi ,\zeta }\eta \right)(x)& =\pi _{\psi _x}(A)\mathrm{\Omega }_{\psi _x}(\pi _{\psi _x}(B)\mathrm{\Omega }_{\psi _x},\pi _{\psi _x}(C)\mathrm{\Omega }_{\psi _x})\hfill & \\ & =\psi _x(B^{}C)\pi _{\psi _x}(A)\mathrm{\Omega }_{\psi _x}=\mathrm{\Phi }(\mathrm{{\rm Y}}^1(\delta (B^{}C)))(x)\mathrm{\Pi }(A)\mathrm{\Omega }(x)\hfill & \\ & =\mathrm{\Pi }(\mathrm{{\rm Y}}^1(\delta (B^{}C))A)\mathrm{\Omega }(x)\hfill & \end{array}$$ where we used the fact that $`\mathrm{{\rm Y}}^1(\delta (B^{}C))𝒞^{\prime \prime }=𝒞`$ and theorem 2.1. Thus $`\stackrel{~}{A}_{\phi ,\zeta }\eta \mathrm{\Pi }(𝒞^{})\mathrm{\Omega }`$. That $`\stackrel{~}{A}_{\phi ,\zeta }`$ preserves $`\mathrm{Ker}(,)`$ will follow from the next calculation as well as the remaining claims. Since by theorem 2.2, $`\kappa (\mathrm{\Omega })`$ is cyclic for $`\pi ()`$, it suffices to do the calculation on $`\mathrm{\Pi }()\mathrm{\Omega }`$, so now let $`\eta =\mathrm{\Pi }(C)\mathrm{\Omega }`$, $`C\pi _\omega ()`$. Then $$\begin{array}{ccc}\hfill \stackrel{~}{A}_{\phi ,\zeta }\eta ^2& =\left\|\psi _x(B^{}C)\right\|^2\psi _x(A^{}A)𝑑\mu (x)\hfill & \\ & _X\psi _x(B^{}B)\psi _x(C^{}C)\psi _x(A^{}A)𝑑\mu (x)\hfill & \\ & A^2B^2_X\psi _x(C^{}C)𝑑\mu (x)\hfill & \\ & =A^2B^2\eta ^2.\hfill & \end{array}$$ . Define the norm $`a=\underset{xX}{sup}a(x)_{\left(_{\psi _x}\right)}`$ on the space of bounded sections: $`\mathrm{\Gamma }_0((\psi )):=\left\{a\mathrm{\Gamma }((\psi ))\right\|a<\mathrm{}\}`$ which makes it into a C\–algebra. Lemma 2.4. $`\left\{\stackrel{~}{A}_{\phi ,\zeta }\right\|\phi ,\zeta \mathrm{\Pi }(𝒞^{})\mathrm{\Omega }\}\mathrm{\Gamma }_0((\psi ))`$. Proof: $`(\stackrel{~}{A}_{\phi ,\zeta }c)(x)=\phi (x)(\zeta (x),c(x))`$ for $`c\mathrm{\Gamma }((\psi ))`$, hence $$\begin{array}{ccc}\hfill (\stackrel{~}{A}_{\phi ,\zeta }c)(x)\|& \phi (x)\zeta (x)c(x),\text{and}\hfill & \\ \hfill \stackrel{~}{A}_{\phi ,\zeta }(x)_{\left(_{\psi _x}\right)}& =\phi (x)\zeta (x),\text{so}\hfill & \\ \hfill \stackrel{~}{A}_{\phi ,\zeta }& =\underset{xX}{sup}\left(\phi (x)\zeta (x)\right)=\underset{xX}{sup}\left(\psi _x(A^{}A)\psi _x(B^{}B)\right)^{1/2}\hfill & \\ & AB<\mathrm{}\hfill & \end{array}$$ where we assumed $`\phi =\mathrm{\Pi }(A)\mathrm{\Omega }`$, $`\zeta =\mathrm{\Pi }(B)\mathrm{\Omega }`$. . Using the C\–operations of $`\mathrm{\Gamma }_0((\psi ))`$, we now define $$\stackrel{~}{}:=C^{}\left\{\stackrel{~}{A}_{\phi ,\zeta }\right\|\phi ,\zeta \mathrm{\Pi }(𝒞^{})\mathrm{\Omega }\}\mathrm{\Gamma }_0((\psi )).$$ Theorem 2.5. With the data $`,\omega ,𝒞=𝒞^{\prime \prime }`$ above;– $`(i)\stackrel{~}{}(x):=\left\{a(x)\right\|a\}=𝒦(_{\psi _x})`$ $`(ii)`$ $`\mathrm{\Pi }(𝒞^{})\stackrel{~}{}\stackrel{~}{}\stackrel{~}{}\mathrm{\Pi }(𝒞^{})`$ and $`[\stackrel{~}{},\mathrm{\Pi }(𝒞)]=0`$, $`(iii)`$ $`\stackrel{~}{}`$ lifts through $`\kappa `$ to define a representation $`\rho :\stackrel{~}{}(_\mathrm{\Gamma })`$. Thus by $`(ii)`$ $`\rho (\stackrel{~}{})\pi (𝒞^{})`$. Proof: $`(i)`$ In a fibre we have $`\stackrel{~}{A}_{\phi ,\zeta }(x)_{\left(_{\psi _x}\right)}=\phi (x)\zeta (x)`$ so $`\stackrel{~}{A}_{\phi ,\zeta }(x)`$ is continuous in $`\phi (x)`$, $`\zeta (x)`$, and thus the norm closure of $`\left\{\stackrel{~}{A}_{\phi ,\zeta }(x)\right\|\phi ,\zeta \mathrm{\Pi }(𝒞^{})\mathrm{\Omega }\}`$ contains all rank one operators, using the fact that $`\left(\mathrm{\Pi }(𝒞^{})\mathrm{\Omega }\right)(x)`$ is dense in $`_{\psi _x}`$. Since $`𝒦(_{\psi _x})`$ is spanned by its rank one operators and closure in a supremum norm produces pointwise closure, we get that $`\stackrel{~}{}(x)=𝒦(_{\psi _x})`$. $`(ii)`$ Let $`\phi =\mathrm{\Pi }(A)\mathrm{\Omega }`$, $`\zeta =\mathrm{\Pi }(B)\mathrm{\Omega }`$, then for $`E𝒞^{}`$: $$\begin{array}{ccc}\hfill \left(\pi (E)\stackrel{~}{A}_{\phi ,\zeta }c\right)(x)& =\pi _{\psi _x}(EA)\mathrm{\Omega }_{\psi _x}(\pi _{\psi _x}(B)\mathrm{\Omega }_{\psi _x},c(x))\hfill & \\ & =\left(\stackrel{~}{A}_{\xi ,\zeta }c\right)(x)\text{where}\xi :=\mathrm{\Pi }(EA)\mathrm{\Omega }.\hfill & \end{array}$$ So $`\mathrm{\Pi }(E)\stackrel{~}{A}_{\phi ,\zeta }\stackrel{~}{}`$. Similarly $`\stackrel{~}{A}_{\phi ,\zeta }\mathrm{\Pi }(E)\stackrel{~}{}`$, and as $`\stackrel{~}{A}_{\phi ,\zeta }(x)^{}=\stackrel{~}{A}_{\zeta ,\phi }(x)`$ it follows that $`\mathrm{\Pi }(𝒞^{})\stackrel{~}{}\stackrel{~}{}\stackrel{~}{}\mathrm{\Pi }(𝒞^{})`$. To see that $`[\stackrel{~}{},\mathrm{\Pi }(𝒞)]=0`$, let $`F𝒞`$, so $$\begin{array}{ccc}\hfill \left(\stackrel{~}{A}_{\phi ,\zeta }\mathrm{\Pi }(F)c\right)(x)& =\phi (x)(\zeta (x),\mathrm{\Pi }(F)c(x))=\phi (x)(\zeta (x),\mathrm{\Phi }(F)(x)c(x))\hfill & \\ & =\mathrm{\Phi }(F)(x)\left(\stackrel{~}{A}_{\phi ,\zeta }c\right)(x)=\left(\mathrm{\Pi }(F)\stackrel{~}{A}_{\phi ,\zeta }c\right)(x).\hfill & \end{array}$$ $`(iii)`$ By lemma 2.3 we already know that $`\stackrel{~}{A}_{\phi ,\zeta }`$ lifts through $`\kappa `$, so since lifting is a homomorphism, we define $`\rho (\stackrel{~}{A}_{\phi ,\zeta }):=A_{\phi ,\zeta }(_\mathrm{\Gamma })`$ and check continuity. Let $`\phi =\mathrm{\Pi }(A)\mathrm{\Omega }`$, $`\zeta =\mathrm{\Pi }(B)\mathrm{\Omega }`$, $`\eta =\mathrm{\Pi }(C)\mathrm{\Omega }`$ then $$\begin{array}{ccc}\hfill \rho (\stackrel{~}{A}_{\phi ,\zeta })& \kappa (\eta )^2=\kappa (\stackrel{~}{A}_{\phi ,\zeta }\eta )^2=_X(\stackrel{~}{A}_{\phi ,\zeta }\eta )(x)^2d\mu (x)\hfill & \\ & =_X\left\|\psi _x(B^{}C)\right\|^2\psi _x(A^{}A)𝑑\mu (x)\hfill & \\ & _X\psi _x(B^{}B)\psi _x(C^{}C)\psi _x(A^{}A)𝑑\mu (x)\hfill & \\ & \underset{xX}{sup}\left(\psi _x(A^{}A)\psi _x(B^{}B)\right)_X\psi _x(C^{}C)𝑑\mu (x)\hfill & \\ & =\stackrel{~}{A}_{\phi ,\zeta }^2\kappa (\eta )^2\hfill & \end{array}$$ . Remarks. (1) Note that theorem 2.5(ii) expresses that $`\stackrel{~}{}`$ is a $`C(X)\text{–algebra}`$, given that $`\mathrm{\Pi }(𝒞)=C(X)`$. Thus the family $`\left\{\stackrel{~}{}(x)\right\|xX\}`$ can be topologised as a Fell bundle (cf. M. Nilsen \[Ni\]). In fact, since $`𝒞^{}`$ is also obviously a $`C(X)\text{–algebra}`$, it is also the set of continuous sections of a Fell bundle. Note also that $`\stackrel{~}{}`$ has ideals corresponding to closed subsets $`Y`$ of $`X`$ by $`_Y:=\left\{L\right\|\mathrm{\Pi }(L)\mathrm{}Y=0\}.`$ Obviously, since $`\stackrel{~}{}`$ is fibrewise the compacts, it is a nonunital algebra. (2) Another, possibly smaller choice for $`\stackrel{~}{}`$ is $$_{}:=C^{}\left\{\stackrel{~}{A}_{\phi ,\zeta }\right\|\phi ,\zeta \mathrm{Span}\{\mathrm{\Pi }(𝒞\pi _\omega ()^{\prime \prime })\mathrm{\Omega }\}\}$$ in which case we still have fibrewise the compacts, $`_{}(x)=𝒦(_{\psi _x})`$, but instead of 2.4(ii) we now have only the weaker property that $`\mathrm{\Pi }(𝒞)`$ and $`\mathrm{\Pi }(\pi _\omega ()^{\prime \prime })`$ are in the relative multiplier algebra of $`_{}`$ in $`\mathrm{\Gamma }_0((\psi ))`$. Lemma 2.6. $`(i)`$ Let $`a\stackrel{~}{}`$, then the map $`xa(x)`$ is continuous. $`(ii)`$ Fix an $`xX`$, then $`\left\{a\stackrel{~}{}\right\|a(x)=0\}`$ is dense in the set $`\left\{fa\right\|a\stackrel{~}{},fC(X),f(x)=0\}`$. Proof: $`(i)`$ First we show that $`xa(x)`$ is continuous for the generating set of $`\stackrel{~}{}`$. Let $`a=\stackrel{~}{A}_{\phi ,\zeta }`$ with $`\phi =\mathrm{\Pi }(A)\mathrm{\Omega }`$, $`\zeta =\mathrm{\Pi }(B)\mathrm{\Omega }`$, $`A,B𝒞^{}`$, then $$\begin{array}{ccc}\hfill a(x)& =\stackrel{~}{A}_{\phi ,\zeta }(x)=\phi (x)\zeta (x)\hfill & \\ & =\pi _{\psi _x}(A)\mathrm{\Omega }_{\psi _x}\pi _{\psi _x}(B)\mathrm{\Omega }_{\psi _x}=\left[\psi _x(A^{}A)\psi _x(B^{}B)\right]^{1/2}\hfill & \end{array}$$ which is continuous in $`x`$ because $`\psi _x(A^{}A)=\mathrm{\Phi }(\mathrm{{\rm Y}}^1(\delta (A^{}A)))(x)`$ and this is continuous by theorem 2.1(i). It is easy to see that $`xa(x)`$ is continuous for linear combinations of the $`\stackrel{~}{A}_{\phi ,\zeta }`$: let $`a,b\mathrm{\Gamma }((\psi ))`$ be such that $`xa(x)`$ and $`xb(x)`$ is continuous, then for a convergent net $`x_\nu x`$ in $`X`$, $$\begin{array}{ccc}\hfill \|a(x_\nu )+b(x_\nu )& a(x)+b(x)\|a(x_\nu )+b(x_\nu )a(x)b(x)\hfill & \\ & a(x_\nu )a(x)+b(x_\nu )b(x)\mathrm{\hspace{0.33em}0}.\hfill & \end{array}$$ Next we show that $`𝒴:=\mathrm{Span}\left\{\stackrel{~}{A}_{\phi ,\zeta }\right\|\phi ,\zeta \mathrm{\Pi }()\mathrm{\Omega }\}`$ is in fact a dense \–subalgebra of $`\stackrel{~}{}`$, not just a generating set. That it is involutive follows from $`\stackrel{~}{A}_{\phi ,\zeta }^{}=\stackrel{~}{A}_{\zeta ,\phi }`$. We only need to show that $`\stackrel{~}{A}_{\phi ,\zeta }\stackrel{~}{A}_{\xi ,\eta }𝒴`$. Let $`c\mathrm{\Gamma }((\psi ))`$, then $$\begin{array}{ccc}\hfill (\stackrel{~}{A}_{\phi ,\zeta }& \stackrel{~}{A}_{\xi ,\eta }c)(x)=\phi (x)(\zeta (x),\left(\stackrel{~}{A}_{\xi ,\eta }c\right)(x))\hfill & \\ & =\phi (x)(\zeta (x),\xi (x))(\eta (x),c(x))=\left(\stackrel{~}{A}_{\overline{\phi },\eta }c\right)(x)\hfill & \end{array}$$ where $`\overline{\phi }:=\mathrm{\Pi }(\stackrel{~}{(\zeta ,\xi )})\phi \mathrm{\Pi }(𝒞^{})\mathrm{\Omega }`$, since $`x\stackrel{~}{(\zeta ,\xi )}(x):=(\zeta (x),\xi (x))`$ is continuous so corresponds to an element $`\stackrel{~}{(\zeta ,\xi )}`$ of $`𝒞`$. Thus $`\stackrel{~}{A}_{\phi ,\zeta }\stackrel{~}{A}_{\xi ,\eta }𝒴`$. Finally, we need to show that if $`\{a_n\}𝒴`$ is a sequence converging to an $`a\stackrel{~}{}`$ in norm, then $`xa(x)`$ is continuous. Given a converging net $`x_\nu x`$ in $`X`$, $$\begin{array}{ccc}\hfill \left\|a(x_\nu )a(x)\right\|& a(x_\nu )a(x)\hfill & \\ & =a(x_\nu )a_n(x_\nu )+a_n(x_\nu )a_n(x)+a_n(x)a(x)\hfill & \\ & a(x_\nu )a_n(x_\nu )+a_n(x_\nu )a_n(x)+a_n(x)a(x)\hfill & \\ & 2aa_n+a_n(x_\nu )a_n(x)\hfill & \end{array}$$ and this can be made arbitrary small for suitable choices of $`n`$ and $`\nu `$. $`(ii)`$ Denote $`_x:=\left\{a\stackrel{~}{}\right\|a(x)=0\}`$ and $`𝒥_x:=\overline{\left\{f\stackrel{~}{}\right\|fC(X),f(x)=0\}}`$. If $`g=fa𝒥_x`$, we clearly have that $`g(x)=f(x)a(x)=0`$, i.e. $`𝒥_x_x`$. Conversely, let $`a_x`$, then by part (i) of this lemma, $`U_\epsilon :=\left\{yX\right\|a(y)<\epsilon \}`$ is open for any $`\epsilon >0`$. Since $`X`$ is the spectrum of a von Neumann algebra, it is completely regular, so there is a continuous function $`f:X[0,\mathrm{\hspace{0.17em}1}]`$ such that $`f(x)=0`$, $`f(y)=1`$ for all $`yU_\epsilon `$. So $$(afa)(x)=a(x)(1f(x))=a(x)\|1f(x)\|=0$$ as $`a(x)=0`$. Since $`(afa)(y)<\epsilon `$ for all $`yU_\epsilon `$ and $`(afa)(y)=0`$ when $`yU_\epsilon `$, it is clear that $`_x=𝒥_x`$. . Lemma 2.7. Any pure state $`\gamma `$ of $`\stackrel{~}{}`$ is of the form $`\gamma (a)=\gamma _x(a(x))`$ for all $`a\stackrel{~}{}`$ where $`x`$ is a distinguished point $`xX`$, and $`\gamma _x`$ is a pure state of $`\stackrel{~}{}(x)=𝒦(_{\psi _x})`$. Proof: From 2.5 we have $`C(X)\stackrel{~}{}\stackrel{~}{}\stackrel{~}{}C(X)`$ using $`\mathrm{\Pi }(𝒞)=C(X)`$ in the explicit action, cf. the definition below 2.1. Hence by Dixmier 2.11.7 \[Di\] there is a unique extension $`\stackrel{~}{\gamma }`$ of $`\gamma `$ to a pure state of $`C^{}(\stackrel{~}{}\mathrm{\Pi }(𝒞))M(\stackrel{~}{})`$. Since $`C(X)=\mathrm{\Pi }(𝒞)Z(C^{}(\stackrel{~}{}\mathrm{\Pi }(𝒞))`$, $`\stackrel{~}{\gamma }\mathrm{}C(X)`$ is pure (since the image of the centre of a C\–algebra under an irreducible representation is one–dimensional). Thus $`\stackrel{~}{\gamma }\mathrm{}C(X)`$ is evaluation at some distinguished point $`xX`$, and obviously the left kernel $$N_{\stackrel{~}{\gamma }\mathrm{}C\left(X\right)}=\mathrm{Ker}\stackrel{~}{\gamma }\mathrm{}C(X)=\left\{fC(X)\right\|f(x)=0\}.$$ Since $`C(X)`$ commutes with $`\stackrel{~}{}`$, this means $$\left\|\gamma (Lf)\right\|^2\gamma (L^{}L)\gamma (f^{}f)=0f\mathrm{Ker}\stackrel{~}{\gamma }\mathrm{}C(X),L\stackrel{~}{}$$ i.e. $`\left\{\stackrel{~}{}f\right\|f(x)=0\}N_\gamma `$. Since $`\gamma `$ is pure, $$\mathrm{Ker}\gamma =N_\gamma +N_\gamma ^{}\overline{\left\{\stackrel{~}{}f\right\|f(x)=0\}}$$ using Dixmier 2.9.1 \[Di\]. The last set is a two–sided ideal, hence in $`\mathrm{Ker}\pi _\gamma `$, and moreover by lemma 2.6(ii), $`\stackrel{~}{}(x)=\stackrel{~}{}/\overline{\left\{\stackrel{~}{}f\right\|f(x)=0\}}`$ and so $$\gamma (a)=\gamma \left(a+\overline{\left\{\stackrel{~}{}f\right\|f(x)=0\}}\right)=\widehat{\gamma }(a(x))$$ defines a state $`\widehat{\gamma }`$ on $`\stackrel{~}{}(x)`$, henceforth denoted by $`\gamma _x`$, i.e. $`\gamma (a)=\gamma _x(a(x))`$. Clearly if $`\gamma _x`$ is not pure, we can write it as a convex combination of states, which through the last expression produces a convex combination for $`\gamma `$, contradicting the fact that $`\gamma `$ is pure. Thus $`\gamma _x`$ must also be pure. . Remarks. (1) Given that $`\stackrel{~}{}`$ can be realised as the continuous sections of a Fell bundle, lemma 2.7 is well–known, cf. e.g. Fell and Doran Prop. 8.8 p582 \[FD\]. (2) C\–algebras of the form of $`\stackrel{~}{}`$ are well–studied in the literature as fields of elementary algebras, cf. Dixmier \[Di\]. Theorem 2.8. Let $`\gamma `$ be a state on $`\stackrel{~}{}`$, then there is a probability measure $`\nu `$ on $`X`$ and a $`\nu \text{–almost}`$ everywhere defined map $`\rho :\mathrm{supp}\nu \underset{xX}{}S(\stackrel{~}{}(x))`$ such that $`\rho _xS(\stackrel{~}{}(x))`$, and $$\gamma (a)=_X\rho _x(a(x))𝑑\nu (x)a\stackrel{~}{}.$$ Proof: Since by 2.5 we have $`\mathrm{\Pi }(𝒞^{})`$ in $`M(\stackrel{~}{})`$, both $`\gamma `$ and $`\pi _\gamma `$ extend uniquely (on the same space) to it. Consider the unital C\–algebra $`\stackrel{~}{𝒞}:=\pi _\gamma (\mathrm{\Pi }(𝒞))\pi _\gamma (\stackrel{~}{})^{}`$. It will only be a von Neumann algebra if $`\gamma `$ is normal on $`\mathrm{\Pi }(𝒞)`$, which we cannot assume. Nevertheless, we can apply theorem 2.1 (which works for also for a commutative C\–algebra in the commutant) to the triple $`,\gamma ,\stackrel{~}{𝒞}`$ to obtain a probability measure $`\stackrel{~}{\nu }`$ on the spectrum $`Y`$ of $`\stackrel{~}{𝒞}`$ with support $`Y`$, and the Gel’fand isomorphism $`\stackrel{~}{\mathrm{\Phi }}:\stackrel{~}{𝒞}L^{\mathrm{}}(Y,\stackrel{~}{\nu })`$ and a $`\stackrel{~}{\nu }\text{–a.e.}`$ defined map $`\stackrel{~}{\psi }:YS(\stackrel{~}{})`$ such that $$\gamma (a)=_Y\stackrel{~}{\psi }_x(a)d\stackrel{~}{\nu }(x)\text{and}\stackrel{~}{\psi }_x(C.a)=\stackrel{~}{\mathrm{\Phi }}(C)(x)\stackrel{~}{\psi }_x(a)$$ for all $`a\stackrel{~}{}`$, $`C\stackrel{~}{𝒞}`$. Now $`\stackrel{~}{𝒞}`$ is a homomorphic image of $`𝒞C(X)`$, and as all ideals in $`C(X)`$ are of the form $`\left\{fC(X)\right\|f(Z)=0\}`$ for some closed set $`ZX`$, the homomorphic images of $`C(X)`$ are all isomorphic to the algebras $`C(\overline{X\backslash Z})=C(V)`$ for $`VX`$ the closure of an open set. Thus there is a homeomorphism $`\beta :YVX`$, with $`V`$ the closure of an open set and we obtain it as follows. Since $`\pi _\gamma `$ must map maximal ideals of $`𝒞`$ to maximal ideals of $`\stackrel{~}{𝒞}`$, define $`\beta :YX`$ by $`\pi _\gamma (_{\beta \left(x\right)})=\stackrel{~}{}_x`$ where $`_x:=\left\{C𝒞\right\|\mathrm{\Phi }(C)(x)=0\}`$ and $`\stackrel{~}{}_x:=\left\{D\stackrel{~}{𝒞}\right\|\stackrel{~}{\mathrm{\Phi }}(D)(x)=0\}`$. Observe that this implies $$\stackrel{~}{\mathrm{\Phi }}(\pi _\gamma (C))(x)=\mathrm{\Phi }(C)(\beta (x)).$$ We can thus write the integral of $`\gamma `$ over $`X`$ instead of $`Y`$;- define a measure $`\nu `$ on $`X`$ by setting it equal to $`\nu =\stackrel{~}{\nu }\beta ^1`$ on $`\beta (Y)X`$, and zero outside of this set. Then $`\gamma (a)=_X\stackrel{~}{\psi }_{\beta ^1(x)}(a)𝑑\nu (x)`$ for $`a\stackrel{~}{}`$. In order to define the map $`\rho _x`$ from this, we need to show for $`\nu \text{–almost}`$ all $`x`$ that $`\stackrel{~}{\psi }_{\beta ^1(x)}(a)`$ only depends on the value $`a(x)`$ of $`a\stackrel{~}{}`$. Now $$\stackrel{~}{\psi }_{\beta ^1(x)}(C.a)=\stackrel{~}{\mathrm{\Phi }}(C)(\beta ^1(x))\stackrel{~}{\psi }_{\beta ^1(x)}(a)$$ so if $`C\stackrel{~}{}_{\beta ^1(x)}=\pi _\gamma (_x)`$, then $`\stackrel{~}{\psi }_{\beta ^1(x)}(C.a)=0`$ for all $`a\stackrel{~}{}`$. But by lemma 2.6(ii), $$\overline{_x\stackrel{~}{}}=\{a\stackrel{~}{}\|a(x)=0\}=:K_x$$ so $`\stackrel{~}{\psi }_{\beta ^1(x)}(K_x)=0`$, and thus using $`\stackrel{~}{}(x)=\stackrel{~}{}/K_x`$, we have $$\stackrel{~}{\psi }_{\beta ^1(x)}(a)=\stackrel{~}{\psi }_{\beta ^1(x)}(a+K_x)=:\rho _x(a(x)),a\stackrel{~}{}$$ and obviously $`\rho _xS(\stackrel{~}{}(x))`$. Thus finally, $`\gamma (a)=_X\rho _x(a(x))𝑑\nu (x)`$ for $`a\stackrel{~}{}`$. . Corollary 2.9. With the data $`,\omega ,𝒞`$ above, $`(i)`$ for any state $`\gamma `$ of $`\stackrel{~}{}`$ there is a probability measure $`\nu `$ on $`X`$ and a $`\nu \text{–almost}`$ everywhere defined map $`T:\mathrm{supp}\nu \underset{xX}{}𝒦(_{\psi _x})`$ such that $`T_x𝒦(_{\psi _x})`$ is positive, trace class, normalised, and satisfies $$\gamma (a)=_X\mathrm{Tr}\left(a(x)T_x\right)𝑑\nu (x),a\stackrel{~}{}.$$ $`(ii)`$ Let $`\theta (\gamma )`$ denote the unique extension of a state $`\gamma S(\stackrel{~}{})`$ to $`𝒞^{}`$, then $$\theta (\gamma )(A)=_X\mathrm{Tr}\left(\pi _{\psi _x}(A)T_x\right)𝑑\nu (x),A𝒞^{}$$ where $`\nu `$ and $`T`$ are as above. Proof: $`(i)`$ This follows directly from theorem 2.8 and the fact that $`\stackrel{~}{}(x)=𝒦(_{\psi _x})`$ for all $`xX`$. So, since all states on the compacts are of the form $`\varphi (A)=\mathrm{Tr}(AT)`$ with $`T`$ positive trace–class, it follows that $`\rho _x(a(x))=\mathrm{Tr}\left(a(x)T_x\right)`$. $`(ii)`$ Let $`\gamma S(\stackrel{~}{})`$ be as above in $`(i)`$, i.e. $$\gamma (a):=_X\mathrm{Tr}\left(a(x)T_x\right)𝑑\nu (x)$$ where $`a\stackrel{~}{}`$. Now we know the unique extension of $`\gamma `$ to $`𝒞^{}`$ is given by $$\stackrel{~}{\gamma }(a)=\underset{\alpha }{lim}\gamma (E_\alpha a),a𝒞^{}.$$ In particular, for the states $`\gamma _xS(\stackrel{~}{})`$ given by $`\gamma _x(a):=\mathrm{Tr}(a(x)T_x),`$ we get $$\stackrel{~}{\gamma }_x(a)=\underset{\alpha }{lim}\mathrm{Tr}\left(E_\alpha (x)a(x)T_x\right)=\mathrm{Tr}(a(x)T_x)$$ since $`\gamma _x`$ is normal, so has a unique extension by the last formula to $`(_{\gamma _x}).`$ Furthermore, $`\stackrel{~}{\gamma }_x=1`$ for all $`x,`$ and as $`\nu `$ is a probability measure on $`X,`$ the function $`1`$ is in $`L^1(X,\nu ).`$ Thus we may apply the Lebesgue dominated convergence theorem: $$\begin{array}{ccc}\hfill \stackrel{~}{\gamma }(a)& =\underset{\alpha }{lim}_X\mathrm{Tr}\left(E_\alpha (x)a(x)T_x\right)𝑑\nu (x)\hfill & \\ & =_X\underset{\alpha }{lim}\mathrm{Tr}\left(E_\alpha (x)a(x)T_x\right)d\nu (x)\hfill & \\ & =_X\mathrm{Tr}(a(x)T_x)𝑑\nu (x).\hfill & \end{array}$$ . Theorem 2.10. $`\gamma `$ is a normal state of $`𝒞^{}`$ or of $`\pi _\omega ()^{\prime \prime }𝒞^{}`$ iff it can be written $$\gamma (A)=_X\mathrm{Tr}\left(\pi _{\psi _x}(A)T_x\right)f(x)𝑑\mu (x),A𝒞^{}$$ where $`T_x𝒦(_{\psi _x})`$ is a.e. trace–class, positive and normalised, and $`fL_+^1(X,\mu )`$ with $`\mu `$ the measure associated to the initial choice $`\omega `$ and $`𝒞`$. Proof: Let $`\gamma `$ be a normal state on $`𝒞^{}`$ or $`\pi _\omega ()^{\prime \prime }`$, then by theorem 2.2 there is a unitary $`U:_\omega _\mathrm{\Gamma }`$ intertwining $`\pi _\omega `$ with $`\mathrm{\Pi }`$. Thus by Kadison and Ringrose 7.1.12 \[KR\], there is a countable set of vectors $`\{\phi _n\}_\mathrm{\Gamma }`$ such that $`1=\underset{n=1}{\overset{\mathrm{}}{}}\phi ^2`$ and $`\gamma (A)=\underset{n}{}(\phi _n,\mathrm{\Pi }(A)\phi _n)`$. Now let $`\{\zeta _k^n\}\mathrm{\Pi }(𝒞^{})\mathrm{\Omega }`$ be sequences such that $`\kappa (\zeta _k^n)\phi _n_\mathrm{\Gamma }`$ where the convergence is in $`k`$. We can in fact choose such sequences with $`\kappa (\zeta _k^n)=\phi _n`$ because if $`\zeta _k^n`$ is a nonzero sequence converging to $`\phi _n`$, so is $`\zeta _k^n\phi _n/\zeta _k^n`$. Below we will blur the distinction between $`\kappa (\zeta _k^n)`$ and $`\zeta _k^n`$. Thus $$\begin{array}{ccc}\hfill 1=\underset{n=1}{\overset{\mathrm{}}{}}\phi _n^2& =\underset{n}{}\zeta _k^n^2=\underset{n}{}_X\zeta _k^n(x)_{_{\psi _x}}^2𝑑\mu (x)\hfill & \\ & =_X\underset{n}{}\zeta _k^n(x)_{_{\psi _x}}^2d\mu (x)\hfill & \end{array}$$ by Fubini and absolute convergence. Thus $$\begin{array}{ccc}\hfill \underset{n}{}\zeta _k^n(x)_{_{\psi _x}}^2& L^1(X,\mu )_+.\text{Now}\hfill & \\ \hfill \left\|(\zeta _k^n(x),B\zeta _k^n(x))_{_{\psi _x}}\right\|& B\zeta _k^n(x)_{_{\psi _x}}^2\text{hence}\hfill & \\ \hfill \gamma _x^k(B):=\underset{n}{}& (\zeta _k^n(x),B\zeta _k^n(x))_{_{\psi _x}}\hfill & \end{array}$$ defines a positive functional on $`(_{\psi _x})`$ for $`\mu \text{–almost}`$ all $`x`$. Moreover, it is normal by Kadison and Ringrose 7.1.12 \[KR\]. Now for $`A`$ a positive element of $`𝒞^{}`$ $$\begin{array}{ccc}\hfill \gamma (A)& =\underset{n}{}(\phi _n,A\phi _n)=\underset{n}{}\underset{k}{lim}(\zeta _k^n,A\zeta _k^n)=\underset{k}{lim}\underset{n}{}(\zeta _k^n,A\zeta _k^n)\hfill & \end{array}$$ by dominated convergence, since $`\left\|(\zeta _k^n,A\zeta _k^n)\right\|A\zeta _k^n^2=A\phi _n^2`$. So $$\begin{array}{ccc}\hfill \gamma (A)& =\underset{k}{lim}\underset{n}{}_X(\zeta _k^n(x),A(x)\zeta _k^n(x))_{_{\psi _x}}𝑑\mu (x)\hfill & \\ & =\underset{k}{lim}_X\underset{n}{}(\zeta _k^n(x),A(x)\zeta _k^n(x))_{_{\psi _x}}d\mu (x)\hfill & \\ & =\underset{k}{lim}_X\gamma _x^k(A(x))𝑑\mu (x)\hfill & \end{array}$$ where we used Fubini’s theorem and absolute convergence. Moreover, since $$\left\|\gamma _x^k(A(x))\right\|A\underset{n}{}\zeta _k^n(x)_{_{\psi _x}}^2L^1(X,\mu )_+$$ we can use the dominated convergence theorem to conclude $$\gamma (A)=_X\underset{k}{lim}\gamma _x^k(A(x))d\mu (x)$$ $`(1)`$ providing we can show that the pointwise limits $`\underset{k}{lim}\gamma _x^k(A(x))`$ exist a.e. which is what we prove now. Since $`\zeta _k^n\mathrm{\Pi }(𝒞^{})\mathrm{\Omega }`$, let $`\zeta _k^n=\mathrm{\Pi }(A_k^n)\mathrm{\Omega }`$, $`A_k^n𝒞^{}`$, so for $`B𝒞^{}`$ $$\begin{array}{ccc}\hfill \gamma _x^k(B)& =\underset{n}{}(\zeta _k^n(x),B(x)\zeta _k^n(x))=\underset{n}{}(\mathrm{\Pi }(A_k^n)\mathrm{\Omega }(x),\mathrm{\Pi }(BA_k^n)\mathrm{\Omega }(x))\hfill & \\ & =\underset{n}{}\psi _x(A_{k}^{n}{}_{}{}^{}BA_k^n)\text{so}\hfill & \\ \hfill \|\gamma _x^k(B)& \gamma _x^{\mathrm{}}(B)\|=\|\underset{n}{}\psi _x(A_{k}^{n}{}_{}{}^{}BA_k^nA_{\mathrm{}}^{n}{}_{}{}^{}BA_{\mathrm{}}^n)\|\hfill & \\ & \|\underset{n}{}(\psi _x((A_k^nA_{\mathrm{}}^n)^{}B(A_k^nA_{\mathrm{}}^n))+\psi _x(A_{\mathrm{}}^{n}{}_{}{}^{}B(A_k^nA_{\mathrm{}}^n))\hfill & \\ & +\psi _x((A_k^nA_{\mathrm{}}^n)^{}BA_{\mathrm{}}^n)\left)\right\|\hfill & \\ & B\underset{n}{}(\psi _x((A_k^nA_{\mathrm{}}^n)^{}(A_k^nA_{\mathrm{}}^n))\hfill & \\ & +2\psi _x(A_{\mathrm{}}^{n}{}_{}{}^{}A_{\mathrm{}}^n)^{1/2}\psi _x((A_k^nA_{\mathrm{}}^n)^{}(A_k^nA_{\mathrm{}}^n))^{1/2})\hfill & (2)\hfill \end{array}$$ using the Cauchy–Schwartz inequality. Now $$\underset{n}{}\psi _x((A_k^nA_{\mathrm{}}^n)^{}(A_k^nA_{\mathrm{}}^n))=\zeta _k^n(x)\zeta _{\mathrm{}}^n(x)^2$$ which must converge to zero a.e. as $`k`$ and $`\mathrm{}`$ approach infinity because $$\zeta _k^n\zeta _{\mathrm{}}^n=\left(_X\zeta _k^n(x)\zeta _{\mathrm{}}^n(x)^2𝑑\mu (x)\right)^{1/2}$$ and this approaches zero with $`k`$, $`\mathrm{}`$ since $`\{\zeta _k^n\}`$ is a convergent sequence. Applying the Cauchy–Schwartz inequality to the sum: $$\begin{array}{ccc}\hfill \|\underset{n}{}\psi _x(A_{\mathrm{}}^{n}{}_{}{}^{}A_{\mathrm{}}^n)^{1/2}& \psi _x((A_k^nA_{\mathrm{}}^n)^{}(A_k^nA_{\mathrm{}}^n))^{1/2}\|^2\hfill & \\ & \underset{n}{}\psi _x(A_{\mathrm{}}^{n}{}_{}{}^{}A_{\mathrm{}}^n)\underset{m}{}\psi _x((A_k^mA_{\mathrm{}}^m)^{}(A_k^nA_{\mathrm{}}^n))\hfill & \end{array}$$ which as we saw approach zero a.e. Thus (2) converges to zero a.e. and so the pointwise limit $`lim_k\gamma _x^k(B(x))`$ exists a.e. and (1) is justified. Since the normal states are sequentially weak\–closed (cf. Takesaki 5.2 p148 \[Tak\]), we conclude that $`\underset{k}{lim}\gamma _x^k`$ is a normal state on $`(_{\psi _x})`$. Thus there is a positive normalised trace–class operator $`T_x𝒦(_{\psi _x})`$ such that $`\gamma _x(A(x))=f(x)\mathrm{Tr}\left(A(x)T_x\right)`$ where $`f(x)=\gamma _x(I)>0`$. That $`fL^1(X,\mu )`$ follows from $$\begin{array}{c}\text{ }1=\gamma (I)=_X\gamma _x(I)𝑑\mu (x)=_Xf(x)𝑑\mu (x).\text{ }\hfill \\ \text{ }\text{Thus}\gamma (A)=_X\mathrm{Tr}\left(\pi _{\psi _x}(A)T_x\right)f(x)𝑑\mu (x)A𝒞^{}.\text{ }\hfill \end{array}$$ Conversely, assume $`\gamma (A)`$ to have this form. Then by Kadison and Ringrose 7.1.12 \[KR\], it suffices to show $`\gamma `$ is strong operator continuous on the unit ball of $`𝒞^{}`$. Let $`\{A_\alpha \}𝒞^{}`$ be a net converging to zero in strong operator topology and with $`A_\alpha 1`$. That is, for all $`\phi _\omega `$ $$A_\alpha \phi ^2=_XA_\alpha (x)\phi (x)_{_{\psi _x}}^2𝑑\mu (x)0$$ which implies that $`A_\alpha (x)\phi (x)0`$ almost everywhere. Since $`A_\alpha (x)=\pi _{\psi _x}(A_\alpha )`$, $$\begin{array}{ccc}\hfill \left\|\gamma (A_\alpha )\right\|=\left\|_X\mathrm{Tr}\left(\pi _{\psi _x}(A_\alpha )T_x\right)f(x)𝑑\mu (x)\right\|_X\left\|\mathrm{Tr}\left(\pi _{\psi _x}(A_\alpha )T_x\right)\right\|f(x)𝑑\mu (x).& & \end{array}$$ Moreover $`\left\|\mathrm{Tr}\left(\pi _{\psi _x}(A_\alpha )T_x\right)\right\|f(x)A_\alpha f(x)`$ which is of course an $`L^1\text{–function}`$ and $`\underset{\alpha }{lim}\left\|\mathrm{Tr}\left(\pi _{\psi _x}(A_\alpha )T_x\right)\right\|0`$ because $`\pi _{\psi _x}(A_\alpha )1`$ and $`\mathrm{Tr}(T_x)`$ is a normal functional on $`(_{\psi _x})`$ (recall that $`A_\alpha (x)`$ converges almost everywhere to $`0`$ in the strong operator topology). Thus by the dominated convergence theorem, $$\begin{array}{ccc}\hfill \underset{\alpha }{lim}\left\|\gamma (A_\alpha )\right\|& \underset{\alpha }{lim}_X\left\|\mathrm{Tr}\left(\pi _{\psi _x}(A_\alpha )T_x\right)\right\|f(x)𝑑\mu (x)\hfill & \\ & =_X\underset{\alpha }{lim}\left\|\mathrm{Tr}\left(\pi _{\psi _x}(A_\alpha )T_x\right)\right\|f(x)d\mu (x)=0.\hfill & \end{array}$$ The argument generalises to nets $`A_\alpha A0`$ in strong operator topology for $`A_\alpha 1`$, hence $`\gamma `$ is normal. . Remark. Observe that the theorem automatically constructs a state on $`𝒞^{}`$ even if one starts from a state on $`\pi _\omega ()^{\prime \prime }`$, though due to the choices involved this need not be unique. Def. Given $`,\omega ,𝒞`$ as above, define the set $`S_\omega S(\stackrel{~}{})`$ as the set of those states $`\gamma `$ with $`\gamma (a)=_X\mathrm{Tr}\left(a(x)T_x\right)𝑑\nu (x)`$ as in 2.9, where $`\nu `$ is absolutely continuous with respect to $`\mu `$. As before, we denote the set of normal states of a concrete von Neumann algebra $`𝒩`$ by $`S_N(𝒩)`$. Theorem 2.11. Assume the data and notation of Corr. 2.9, then $`(i)`$ For a state $`\gamma S(\stackrel{~}{})`$ the extension $`\theta (\gamma )`$ is normal on $`𝒞^{}`$ or on $`\pi _\omega ()^{\prime \prime }`$ iff its measure $`\nu `$ on $`X`$ is absolutely continuous with the measure $`\mu `$ associated with $`\omega `$. $`(ii)`$ For each normal state $`\eta `$ of $`𝒞^{}`$ or $`\pi _\omega ()^{\prime \prime }`$ there is a $`\gamma S(\stackrel{~}{})`$ such that $`\eta =\theta (\gamma )`$. $`(iii)`$ We have: $`\theta :S_\omega S_N(\pi _\omega ()^{\prime \prime })`$ is a surjection. Since $`\stackrel{~}{}𝒞^{}`$, we see that $`\theta :S_\omega S_N(𝒞^{})`$ is a bijection. Proof: $`(i)`$ By Corr. 2.9(ii) and theorem 2.10, this follows immediately. $`(ii)`$ Now $`𝒞^{}\stackrel{~}{}^{\prime \prime },`$ so normal states of $`𝒞^{}`$ are restrictions of normal states of $`\stackrel{~}{}^{\prime \prime },`$ and these in turn are the unique extensions of states from $`\stackrel{~}{}`$ by strong operator continuity. Since $`𝒞^{}M(\stackrel{~}{})\stackrel{~}{}^{\prime \prime },`$ these strong operator extensions are just the unique extensions by $`\theta .`$ Thus we have surjectivity as claimed. $`(iii)`$ This is just a restatement of the preceding parts. . Remark. Since $`S_\omega `$ is a proper subset of $`S(\stackrel{~}{}),`$ this means that $`\stackrel{~}{}`$ cannot in general be an ideal host for $`𝒞^{}`$ or for $`\pi _\omega ()^{\prime \prime },`$ a fact which is also obvious from the commutative situation: $`\pi _\omega ()^{\prime \prime }=L^{\mathrm{}}(X,\mu )=\stackrel{~}{}.`$ But we have almost finished showing that it is a quasi–host. Theorem 2.12. Let $`𝒞L^{\mathrm{}}(X,\mu )`$ be maximally commutative in $`\pi _\omega ()^{}.`$ Then $`𝒞^{}=(𝒩𝒞)^{\prime \prime }\stackrel{~}{},`$ $`\pi _\omega ()M(\stackrel{~}{}),`$ $$S_\omega =\left\{\phi S(\stackrel{~}{})\right\|\stackrel{~}{\phi }L^{\mathrm{}}(X,\mu )\text{is normal}\},$$ $`\stackrel{~}{S}_\omega \pi _\omega ()=S_N(𝒩)=S_N(\pi _\omega ())`$ and $`\stackrel{~}{S}_\omega C^{}(𝒩𝒞)`$ is in bijection with $`S_\omega .`$ In short, $`\stackrel{~}{}`$ is a quasi–host for $`(𝒩,S_N(𝒩))`$ (hence for $`(\pi _\omega (),S_N(\pi _\omega ())).`$ Proof: Since $`𝒞`$ is maximally commutative in $`𝒩^{},`$ we have $`𝒞=𝒩^{}𝒞^{}=(𝒩𝒞)^{}`$ and thus $`𝒞^{}=(𝒩𝒞)^{\prime \prime }`$ hence a normal state on $`𝒞^{}`$ is uniquely determined by its values on $`𝒞`$ and on $`𝒩.`$ Thus by Theorem 2.11(iii) we have that $`\stackrel{~}{S}_\omega C^{}(𝒩𝒞)`$ is in bijection with $`S_\omega .`$ We already have the embeddings stated, so to check the claimed characterisation of $`S_\omega ,`$ recall that it consists of states $`\gamma S(\stackrel{~}{})`$ such that $$\gamma (A)=_X\mathrm{Tr}\left(A(x)T_x\right)𝑑\nu (x),A\stackrel{~}{}$$ where $`\nu `$ is absolutely continuous w.r.t. $`\mu .`$ Then for $`f𝒞=L^{\mathrm{}}(X,\mu )`$ we have for any approximate identity $`\{E_\alpha \}`$ of $`\stackrel{~}{}:`$ $$\stackrel{~}{\gamma }(f)=\underset{\alpha }{lim}\gamma (fE_\alpha )=\underset{\alpha }{lim}_Xf(x)\mathrm{Tr}\left(E_\alpha (x)T_x\right)𝑑\nu (x)=_Xf(x)𝑑\nu (x)$$ using the argument in the proof of Corollary 2.9(ii). These are precisely the normal states of $`L^{\mathrm{}}(X,\mu ).`$ This completes the proof. . Return now to the original pair $`(,S_0)`$ at the start of the investigation, which was equivalent to the examination of $`(\pi _{S_0}(),S_N(\pi _{S_0}())),`$ where $`\pi _{S_0}=\underset{\omega S_0}{}\pi _\omega .`$ Take the direct sum $`\stackrel{~}{}_0:=\underset{\omega S_0}{}\stackrel{~}{}_\omega (_{S_0})`$ where $`\stackrel{~}{}_\omega `$ is the quasi–host constructed above for the pair $`(\pi _\omega (),S_N(\pi _\omega ())).`$ Then $$\pi _{S_0}()M(\stackrel{~}{}_0),\text{and}\underset{\omega S_0}{}L^{\mathrm{}}(X_\omega ,\mu _\omega )ZM(\stackrel{~}{}_0).$$ Let $`(X,\mu )`$ be the disjoint union of the measure spaces $`(X_\omega ,\mu _\omega )`$ (hence $`\mu `$ is not a probability measure), so we can write $`L^{\mathrm{}}(X,\mu )=\underset{\omega S_0}{}L^{\mathrm{}}(X_\omega ,\mu _\omega ),`$ and thus $`L^{\mathrm{}}(X,\mu )ZM(\stackrel{~}{}_0).`$ Moreover with notation $$S_\mu :=\left\{\phi S(\stackrel{~}{}_0)\right\|\stackrel{~}{\phi }L^{\mathrm{}}(X,\mu )\text{is normal}\}$$ we see that for $`\phi S_\mu `$ that $`\stackrel{~}{\phi }L^{\mathrm{}}(X_\omega ,\mu _\omega )`$ is normal for all $`\omega S_0,`$ hence $`S_\mu `$ is the norm closed convex hull of $`\underset{\omega S_0}{}S_\omega `$ and by Theorem 2.12 $`\stackrel{~}{S}_\mu \pi _{S_0}()=S_0`$ and $`\stackrel{~}{S}_\mu C^{}(\pi _{S_0}()L^{\mathrm{}}(X,\mu ))`$ is in bijection with $`S_\mu .`$ In other words, $`\stackrel{~}{}_0`$ is a quasi–host for $`(,S_0).`$ Note that whilst $`\stackrel{~}{}_0`$ is in $`(_{S_0}),`$ it need not be in $`\pi _{S_0}()^{\prime \prime },`$ (unlike ideal hosts) because $`𝒞`$ may have a part outside $`\pi _{S_0}()^{\prime \prime }.`$ Theorem 2.13. Given the preceding notation, if $`(,S_0)`$ has an ideal host, then $`\mu `$ must have discrete points. Proof: If $`(,S_0)`$ has an ideal host, then by Corollary 1.7, $`S_0`$ has pure states. So in the direct sum $`\pi _{S_0}()=\underset{\omega S_0}{}\pi _\omega ()`$ we have that for some $`\omega S_0,`$ that $`\pi _\omega ()^{}=CI,`$ and so $`𝒞_\omega =CI`$ and thus $`(X_\omega ,\mu _\omega )`$ is the discrete trivial measure space $`(\{x\},\delta ),`$ $`\delta (\{x\})=1,`$ $`\delta (\mathrm{})=0.`$ Since this is a summand of $`L^{\mathrm{}}(X,\mu ),`$ we conclude that $`\mu `$ has discrete points. . This makes now explicit the sense in which we meant that $`L^{\mathrm{}}(X,\mu )`$ with $`\mu `$ continuous, is an obstruction to the existence of an ideal host. It would be nice to get a converse, i.e. to argue that an ideal host exists iff $`\mu `$ has some specific structure, but we do not have this yet. ## 3. Applications. Here we will do a couple of applications of ideal hosts, mainly following the well–trodden path of group algebras. First, there is the question of useful decompositions of states in $`S_0`$ into other states in $`S_0:`$ Theorem 3.1. Let $``$ be an ideal host for a pair $`(,S_0),`$ and let $`\phi S_0.`$ Then there is an integral decomposition $$\phi (F)=_{S_0}\omega (F)𝑑\mu (\omega ),F,$$ such that $`\mu `$ is pseudosupported on the pure states of $`S_0.`$ If $``$ is separable, $`\mu `$ is actually supported on the pure states of $`S_0.`$ Proof: Most of the work for the proof has already been done in the last section. First observe that from Theorem 1.5 we know that $`^{\prime \prime }=\pi _{S_0}()^{\prime \prime },`$ and hence for an $`\omega S()`$ we have that $`\pi _\omega ()^{\prime \prime }=\pi _{\theta (\omega )}()^{\prime \prime }.`$ For a given $`\phi S_0,`$ choose $`\omega =\theta ^1(\phi ).`$ From Theorem 2.1, for a choice of commutative algebra $`𝒞\pi _\omega ()^{},`$ we have the decomposition $$\stackrel{~}{\omega }(A):=(\mathrm{\Omega }_\omega ,A\mathrm{\Omega }_\omega )=_X\psi _x(A)𝑑\mu (x)$$ for all $`A𝒞^{}\pi _{\stackrel{~}{\omega }}()=\pi _\phi (),`$ and hence by restriction to $`\pi _\phi ()`$ this is also a decomposition for $`\phi =\stackrel{~}{\omega }S_0.`$ Make the usual identification of $`X`$ with a subset of $`S()`$ by $`x\psi _x,`$ and maintain the same symbol for the measure $`\mu `$ carried by this identification to $`S().`$ Of course via $`\theta `$ we can now carry this measure to $`S_0`$ too. Thus we have $$\begin{array}{ccc}\hfill \omega (L)& =_{S()}\gamma (L)𝑑\mu (\gamma ),L,\text{hence:}\hfill & \\ \hfill \phi (F)& =\theta (\omega )(F)=\underset{\alpha }{lim}_{S()}\gamma (FE_\alpha )𝑑\mu (\gamma )\hfill & \\ & =_{S()}\theta (\gamma )(F)𝑑\mu (\gamma )=_{\theta (S())}\psi (F)𝑑\mu (\psi )\hfill & \\ & =_{S_0}\psi (F)𝑑\mu (\psi )\hfill & \end{array}$$ for all $`F,`$ any approximate identity $`\{E_\alpha \}`$ of $`,`$ and where we used the argument in the proof of Corollary 2.9(ii) to bring the limit into the integral. This then provides the basic decomposition formula. Now, since the measure $`\mu `$ on $`S()`$ is the same one constructed in the decomposition theory in Takesaki \[Tak\], we can use his Theorem 6.28, p246, to conclude that if $`𝒞`$ is maximally commutative in $`\pi _\omega ()^{},`$ then $`\mu `$ is pseudosupported by the pure states of $`,`$ and supported by them if $``$ is separable. Because $`\theta `$ restricts to a bijection between the pure states on $``$ and the pure states of $`S_0,`$ this proves the assertion. . Remark. Since by Theorem 1.5 we deduce that $`\theta `$ also restricts to a bijection between the factor states of $``$ and the factor states of $`S_0,`$ we can use the central decomposition of a state on $``$ to obtain a similar decomposition of a state in $`S_0,`$ in terms of factor states. Again, we will have the two claims for the measure;– in general it is pseudosupported on the factor states, but if $``$ is separable, it is actually supported by the factor states. As a second application for host algebras, we mention that of inducing representations. For group algebras, this is of course one of the main applications (cf. Rieffel \[Ri\]). What allows one to do the same here, is the relationship given in Theorem 1.9 between the representations of $``$ and those of a host algebra. To be more specific, given two pairs $`(_i,S_{0i}),`$ $`i=1,\mathrm{\hspace{0.33em}2}`$ together with respective hosts $`_i`$ then providing one has constructed a right$`_1\text{–rigged}`$left$`_2\text{–module}`$ $`,`$ then we can induce a representation $`\pi \mathrm{Rep}_{S_{01}}_1`$ to a representation $`\rho \mathrm{Rep}_{S_{02}}_2`$ by using the map $`\mathrm{\Lambda }`$ in Theorem 1.9 to identify $`\pi `$ with a representation of $`_1,`$ inducing this representation via $``$ to $`_2,`$ and then identifying the result with a representation $`\rho \mathrm{Rep}_{S_{02}}_2`$ with the map $`\mathrm{\Lambda }.`$ The benefit of doing an induction via host algebras is that we remain within the class of representations normal w.r.t. the representations $`\pi _{S_{0i}}`$ whereas induction between the $`_i\text{’s}`$ directly, can move us out of these classes. The crux comes of course with the construction of the module $`.`$ For concrete examples of this, we refer to any example of induction via group algebras (cf. Rieffel \[Ri\]). ## Acknowledgements. I am deeply grateful to Detlev Buchholz, who has been as much responsible for the development of host algebras as I have. Not only did the current line of thought develop out of his question at a seminar I gave in 1997 and many subsequent lively discussions, but he also read several previous misguided attempts, discovering deeply buried and serious errors. Moreover, over the years he prodded me to return to this project, and most recently at Göttingen provided the opportunity and support for me to pursue it, of which this paper is the result. Thank you Detlev! I am also grateful to the Erwin Schrödinger Institute in Vienna (and in particular to Jakob Yngvason and Heide Narnhofer) who provided me with the opportunity of completing the last part of the paper. ## Bibliography. \[AP\] C. 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Kastler, Nets of C\–algebras and Classification of States. Commun. Math. Phys. 16, 81–104 (1970) \[KR\] Kadison, R.V., Ringrose, J.R.: Fundamentals of the Theory of Operator Algebras II. Orlando: Academic Press 1986 \[Ku\] Kusuda, M.: Unique state extension and hereditary C\–subalgebras. Math. Ann. 288, 201–209 (1990) \[Ma\] Manuceau, J.: C\–algèbre de relations de commutation. Ann. Inst. H. Poincaré 8, 139–161 (1968) \[Mu\] Murphy, G.J.: C\–Algebras and Operator Theory. Boston: Academic Press 1990 \[Ni\] Nilsen, M.: C\–bundles and $`C_0(X)\text{–algebras.}`$ Indiana Univ. Math. J., 45, 463–477 (1996) \[Pe\] Pedersen, G.K.: C\–Algebras and their Automorphism Groups. London: Academic Press 1989 \[Ri\] M.A. Rieffel: Induced representations of C\–algebras. Adv. Math. 13, 176–257 (1974) \[Se\] I.E. Segal: Representations of the canonical commutation relations, Cargése lectures in theoretical physics, Gordon and Breach, 1967. \[Ta\] A. Taylor: Introduction to Functional Analysis. John Wiley & Sons, New York 1958. \[Tak\] Takesaki, M.: Theory of operator algebras I. New York: Springer-Verlag 1979
warning/0007/hep-ph0007239.html	ar5iv	text	# The Angular Momentum and 𝑔₁^𝑝 Sum Rules for the ProtonResearch supported in part by PPARC grant GR/L56374 and EC TMR grant FMRX-CT96-0008. Invited talk at QCD00, Montpellier, July 2000. CERN-TH/2000-214, SWAT/00-263, UGVA-DPT-00-7-1087. ## 1 Introduction In this talk, based on work in collaboration with B. White, I review our recent formulation of the angular momentum sum rule for the proton and discuss its relation to the sum rule for the first moment of the polarised structure function $`g_1^p`$. In particular, the role of the axial charge $`a^0`$, which is measured in the $`g_1^p`$ sum rule, is highlighted and it is shown how this decouples from the angular momentum sum rule. This emphasises the limitations of attempting to identify $`a^0`$ with quark (and gluon) spin and an alternative interpretation in terms of topological charge is briefly reviewed. The angular momentum sum rule is shown to take the simple form $`\frac{1}{2}=J_q+J_g`$, where $`J_q`$ and $`J_g`$ are gauge and Lorentz invariant form factors of the forward matrix elements of local operators, which may reasonably be interpreted as quark and gluon components of the total angular momentum of the proton. Experimentally, they may be measured in, for example, deeply virtual Compton scattering. Their RG evolution properties are derived from a careful analysis of operator mixing. We also discuss critically whether there is indeed a QCD field-theoretic basis for a decomposition of the proton spin into separate quark and gluon spin and orbital angular momentum components, as in the frequently-quoted sum rule $`\frac{1}{2}=\frac{1}{2}\mathrm{\Delta }q+\mathrm{\Delta }g+L_q+L_g`$. ## 2 The $`g_1^p`$ sum rule and topological charge The sum rule for the first moment of $`g_1^p`$ (see e.g. ref. for reviews of our earlier work and references) is $$_0^1𝑑xg_1^p(x,Q^2)=\frac{1}{12}C_1^{\mathrm{NS}}\left(a^3+\frac{1}{3}a^8\right)+\frac{1}{9}C_1^\mathrm{S}a^0(Q^2)$$ (1) where the $`C_1`$ are Wilson coefficients and the flavour singlet axial charge $`a^0`$ is defined as the form factor in the forward matrix element of the corresponding axial current, $$p,s\|A_\mu ^0\|p,s=a^0s_\mu $$ (2) where $`s_\mu `$ is the covariant spin vector. Since $`A_\mu ^0`$ is not a conserved current, due to the $`U_A(1)`$ anomaly $$^\mu A_\mu ^02n_fQ0$$ (3) where $`Q=\frac{\alpha _s}{8\pi }ϵ^{\mu \nu \rho \sigma }\mathrm{tr}F_{\mu \nu }F_{\rho \sigma }`$ is the gluon topological charge density, the form factor $`a^0`$ is scale dependent and satisfies the non-trivial RG evolution equation $$\frac{d}{dt}a^0(Q^2)=\gamma a^0(Q^2)$$ (4) where $`t=\mathrm{ln}Q^2/\mathrm{\Lambda }^2`$ and the anomalous dimension is $`\gamma =n_f\frac{\alpha _s}{2\pi ^2}`$. Given $`a^3`$ and $`a^8`$ from low-energy neutron and hyperon decays, $`a^0`$ can be extracted from polarised inclusive DIS processes $`e(\mu )pe(\mu )X`$. The simplest prediction, $`a^0a^8`$, is an immediate consequence of the OZI rule and transcribed into the sum rule for $`g_1^p`$ gives the Ellis-Jaffe sum rule. However, the flavour singlet pseudovector or pseudoscalar channel is precisely where we would expect to find strong OZI violations related to the $`U_A(1)`$ anomaly and indeed it is found experimentally that $`a^0a^8`$. Since we can rewrite $`a^0`$ using eq.(3) in terms of the matrix element of the topological charge density, $$a^0=p,s\|Q\|p,s$$ (5) we see immediately that the observed suppression in $`a^0`$ is a manifestation of topological charge screening. This interpretation has been developed in a series of papers written in collaboration with Veneziano, Narison and De Florian. Our proposal is that this screening is universal, i.e. target-independent, being an intrinsic property of the QCD vacuum itself. Specifically, we showed that (in the chiral limit) $$a^0=\frac{1}{2M}2n_f\left[\chi (0)\mathrm{\Gamma }_{Qpp}+\sqrt{\chi ^{}(0)}\mathrm{\Gamma }_{\mathrm{\Phi }_5pp}\right]$$ (6) where the $`\mathrm{\Gamma }`$ are suitably defined 1PI vertex functions and $`\chi (k^2)=id^4xe^{ik.x}0\|TQ(x)Q(0)\|`$ is the topological susceptibility, $`\chi ^{}(0)`$ being its slope at $`k=0`$. Since the anomalous chiral Ward identity implies $`\chi (0)=0`$ in the chiral limit, the sole contribution to $`a^0`$ comes from the second term in eq.(6). Making the motivated assumption that the RG invariant vertex $`\mathrm{\Gamma }_{\mathrm{\Phi }_5pp}`$ obeys the OZI rule to a good approximation, we conjecture that the principal origin of the suppression in $`a^0`$ is an anomalously small value of $`\chi ^{}(0)`$ due to universal topological charge screening by the QCD vacuum.<sup>1</sup><sup>1</sup>1 Explanations which favour OZI violations due to a large polarised strange quark component of the proton or the implications of the Skyrme model of proton structure have been recently reviewed in, for example, ref. This is anticipated in certain instanton-based models of the QCD vacuum, has been confirmed quantitatively by QCD spectral sum rule calculations, and is currently being investigated in lattice gauge theory. The QCD parton model gives an alternative interpretation of $`a^0`$ through the identification $$a^0(Q^2)=\mathrm{\Delta }q2n_f\frac{\alpha _s}{4\pi }\mathrm{\Delta }g(Q^2)$$ (7) where $`\mathrm{\Delta }q=\mathrm{\Delta }u+\mathrm{\Delta }d+\mathrm{\Delta }s`$ and $`\mathrm{\Delta }g(Q^2)`$ are the first moments of the polarised flavour singlet quark and gluon distributions and we have used the AB class of renormalisation schemes where $`\mathrm{\Delta }q`$ is defined to be $`Q^2`$ independent. The OZI/Ellis-Jaffe relation $`a^0=a^8`$ follows immediately from the assumption that in the proton, $`\mathrm{\Delta }s=0`$ and $`\mathrm{\Delta }g=0`$. In this model, $`\frac{1}{2}\mathrm{\Delta }q`$ and $`\mathrm{\Delta }g`$ are interpreted as the quark and gluon spins, which led to the initial interpretation of the experimental observation $`a^0a^8`$ as indicating that the quarks carry only a small fraction of the spin of the proton – the so-called ‘proton spin crisis’. In the rest of this talk, we derive the actual angular momentum sum rule for the proton in terms of gauge invariant operator matrix elements, with particular emphasis on whether and how the axial charge $`a^0`$ appears and whether the interpretation of $`\mathrm{\Delta }q`$ and $`\mathrm{\Delta }g`$ as spin components can be complemented by corresponding definitions of orbital angular momentum components $`L_q`$ and $`L_g`$. ## 3 The angular momentum sum rule The angular momentum sum rule is derived by taking the forward matrix element of the conserved angular momentum current $`M^{\mu \nu \lambda }`$, defined from the energy-momentum tensor as $$M^{\mu \nu \lambda }=x^{[\nu }T^{\lambda ]\mu }+_\rho X^{\rho \mu \nu \lambda }$$ (8) The inclusion of the arbitrary tensor $`X^{\rho \mu \nu \lambda }`$ (antisymmetric under $`\rho \mu `$ and $`\nu \lambda `$) just reflects the usual freedom in QFT in defining conserved currents. However, this arbitrariness allows us to write different equivalent expressions for $`M^{\mu \nu \lambda }`$ as a sum of local operators, suggesting corresponding interpretations of the total angular momentum as a sum of ‘components’, which we can try to identify as reasonable definitions of quark and gluon spin and orbital angular momentum. The best decomposition is the following: $`M^{\mu \nu \lambda }=O_1^{\mu \nu \lambda }+O_2^{\mu [\lambda }x^{\nu ]}+O_3^{\mu [\lambda }x^{\nu ]}+g^{\mu [\lambda }x^{\nu ]}_{\mathrm{gi}}`$ $`+\left\{i^{\{\mu }\overline{c}D^{[\lambda \}}c+^{\{\mu }BA^{[\lambda \}}+g^{\mu [\lambda }_{\mathrm{gf}}\right\}x^{\nu ]}`$ $`{\displaystyle \frac{1}{4}}_\rho \left[x^{[\nu }ϵ^{\lambda ]\mu \rho \sigma }\overline{\psi }\gamma _\sigma \gamma _5\psi \right\}+\mathrm{EOM}+_\rho X^{\rho \mu \nu \lambda }`$ (9) The tensor $`X^{\rho \mu \nu \lambda }`$ is chosen to cancel the divergence and equation of motion (EOM) terms, while the forward matrix element of the operator $`g^{\mu [\lambda }x^{\nu ]}_{\mathrm{gi}}`$ vanishes. The term in $`\{\mathrm{}\}`$ is the contribution from the covariant gauge-fixing and ghost terms in the Lagrangian, but this turns out to be a BRS variation so its matrix element between physical states vanishes. The remaining (bare) operators are gauge invariant: $`O_1^{\mu \nu \lambda }={\displaystyle \frac{1}{2}}ϵ^{\mu \nu \lambda \sigma }\overline{\psi }\gamma _\sigma \gamma _5\psi {\displaystyle \frac{1}{2}}ϵ^{\mu \nu \lambda \sigma }A_\sigma ^0`$ $`O_2^{\mu \lambda }=i\overline{\psi }\gamma ^\mu \stackrel{}{D^\lambda }\psi `$ $`O_3^{\mu \lambda }=F^{\mu \rho }F_\rho ^\lambda `$ (10) We also find it convenient for later use to define $`O_4^{\mu \lambda }=\frac{1}{2}O_2^{\{\mu \lambda \}}`$. At first sight, $`O_1^{\mu \nu \lambda }`$, which is just the flavour singlet axial current considered in section 2, looks as if it may be associated with ‘quark spin’, with $`O_2^{\mu [\lambda }x^{\nu ]}`$ corresponding to a gauge invariant definition of ‘quark orbital angular momentum’. This leaves $`O_3^{\mu [\lambda }x^{\nu ]}`$ to be associated with the ‘gluon total angular momentum’. In fact, there is no further decomposition of the gluon contribution as long as we restrict to gauge invariant operators. We can certainly make the alternative decomposition: $`M^{\mu \nu \lambda }=`$ $`\stackrel{~}{O}_1^{\mu \nu \lambda }+\stackrel{~}{O}_2^{\mu [\lambda }x^{\nu ]}+\stackrel{~}{O}_3^{\mu \nu \lambda }+\stackrel{~}{O}_4^{\mu [\lambda }x^{\nu ]}+g^{\mu [\lambda }x^{\nu ]}_{\mathrm{gi}}`$ $`+\left\{A^\mu ^{[\lambda }B+i^\mu \overline{c}^{[\lambda }c+i^{[\lambda }\overline{c}D^\mu c+g^{\mu [\lambda }_{\mathrm{gf}}\right\}x^{\nu ]}`$ $`+_\rho \left[x^{[\nu }A^{\lambda ]}F^{\mu \rho }\right]+\mathrm{EOM}+_\rho X^{\rho \mu \nu \lambda }`$ (11) where $`\stackrel{~}{O}_1^{\mu \nu \lambda }={\displaystyle \frac{1}{2}}ϵ^{\mu \nu \lambda \sigma }\overline{\psi }\gamma _\sigma \gamma _5\psi `$ $`\stackrel{~}{O}_2^{\mu \lambda }=i\overline{\psi }\gamma ^\mu ^\lambda \psi `$ $`\stackrel{~}{O}_3^{\mu \nu \lambda }=F^{\mu [\nu }A^{\lambda ]}`$ $`\stackrel{~}{O}_4^{\mu \nu }=F^{\mu \rho }^\lambda A_\rho `$ (12) and try to identify the first four operators with, respectively, quark spin and orbital and gluon spin and orbital angular momentum. However, even the forward matrix elements of these operators turn out not to be gauge invariant. Moreover, the gauge-fixing and ghost term in $`\{\mathrm{}\}`$ is no longer a BRS variation so would contribute a non-vanishing ‘ghost orbital angular momentum’. Further discussion of the problems with such gauge non-invariant decompositions is given in ref., and from now on we restrict attention to the gauge invariant formulation based on eq.(9). The next step is to express the matrix elements of the operators $`O_1^{\mu \nu \lambda }`$, $`O_2^{\mu [\lambda }x^{\nu ]}`$ and $`O_3^{\mu [\lambda }x^{\nu ]}`$ in terms of form factors. There are technical subtleties connected with defining operators of the form $`Ox`$ and their renormalisation mixing which are explained precisely in ref.. The prescription is essentially to define the forward matrix elements of $`Ox`$ in terms of the limit of an off-forward matrix element. We therefore write (a little loosely) $$p\|O^{\mu \lambda }x^\nu \|p=i\frac{}{\mathrm{\Delta }_\nu }p\|O^{\mu \lambda }\|p^{}\|_{p^{}=p}$$ (13) where $`\mathrm{\Delta }=pp^{}`$. We then have $`p,s\|O_1^{\mu \nu \lambda }\|p,s=a^0Mϵ^{\mu \nu \lambda \sigma }s_\sigma `$ $`p,s\|O_2^{\mu [\lambda }x^{\nu ]}\|p,s=B_q(0){\displaystyle \frac{1}{2M}}p_\rho p^{\{\mu }ϵ^{[\lambda \}\nu ]\rho \sigma }s_\sigma `$ $`+\stackrel{~}{B}_q(0){\displaystyle \frac{1}{2M}}p_\rho p^{[\mu }ϵ^{[\lambda ]\nu ]\rho \sigma }s_\sigma 2D_q(0)Mϵ^{\mu \nu \lambda \sigma }s_\sigma `$ $`p,s\|O_3^{\mu [\lambda }x^{\nu ]}\|p,s=B_g(0){\displaystyle \frac{1}{2M}}p_\rho p^{\{\mu }ϵ^{[\lambda \}\nu ]\rho \sigma }s_\sigma `$ (14) The crucial observation now follows from the identity $$O_1^{\mu \nu \lambda }+O_2^{\mu [\lambda }x^{\nu ]}=O_4^{\mu [\lambda }x^{\nu ]}+\mathrm{divergence}+\mathrm{EOM}$$ (15) Since the matrix elements of the divergence and EOM terms vanish, and recalling that the operator $`O_4^{\mu \lambda }`$ is symmetric in $`\mu ,\lambda `$, we see that $$\stackrel{~}{B}_q(0)=02D_q(0)=a^0$$ (16) Thus $`a^0`$ appears not only as the unique form factor in the matrix element $`O_1`$ of the axial current, but also as a contribution to $`O_2x`$. It therefore cancels from the angular momentum sum rule. Introducing the notation $`J_q=B_q(0)`$, $`J_g=B_g(0)`$, then from eq.(14) we may write the sum rule as $$\frac{1}{2}=J_q+J_g$$ (17) where $`J_q`$ and $`J_g`$ are gauge and Lorentz invariant quantities which may reasonably be identified as total ‘quark’ and ‘gluon’ angular momenta respectively. Of course, this is a rather non-rigorous terminology since the corresponding operators $`O_2x`$ and $`O_3x`$ mix, and indeed $`O_2x`$ itself contains an explicit gluon field component in the covariant derivative, but it is convenient and the closest approximation to a quark-gluon decomposition that can be given in an interacting QFT such as QCD. Moreover, as we discuss below, $`J_q`$ and $`J_g`$ are still not RG scheme/scale dependent quantities. The point we wish to stress is that provided we restrict to gauge and Lorentz invariant quantities, the true angular momentum sum rule involves only the two form factors $`J_q`$ and $`J_g`$. The axial charge $`a^0`$ is simply not present in the sum rule (17). We return to this point in section 5. Just as the axial charge form factor $`a^0`$ can be measured in polarised inclusive DIS, the angular momentum form factors $`J_q`$ and $`J_g`$ can in principle be extracted from measurements of unpolarised off-forward parton distribution functions in processes such as deeply virtual Compton scattering $`\gamma ^{}p\gamma p`$. The required identifications are $`iP^+{\displaystyle \frac{}{\mathrm{\Delta }_\mu }}{\displaystyle _1^1}dxxf_{q(g)/p}(x,\xi ,\mathrm{\Delta })\|_{\mathrm{\Delta }=0}`$ $`=J_{q(g)}{\displaystyle \frac{1}{M}}ϵ^{+\mu \rho \sigma }P_\rho s_\sigma `$ (18) where $`\xi =\frac{q.\mathrm{\Delta }}{2q.P}`$ and the incoming(outgoing) proton momenta are $`P(+)\mathrm{\Delta }`$. These form factors may also be calculated non-perturbatively in lattice gauge theory and some initial results for $`J_q`$ in the quenched approximation have recently been obtained in ref.. ## 4 Operator mixing and RG evolution The operators $`O_1`$, $`O_2x`$ and $`O_3x`$ in the angular momentum sum rule renormalise and mix in a non-trivial way. The analysis is made more subtle by the explicit factors of the coordinate $`x`$ which have to be carefully treated. A detailed discussion is presented in ref. and here we only sketch the main features. First, note that when inserted into forward matrix elements, operators of the form $`O_a`$ and $`O_ix`$ mix with a block triangular structure: $$\left(\begin{array}{c}O_a\\ O_ix\end{array}\right)_R=\left(\begin{array}{cc}Z_{ab}^1& 0\\ Z_{ib}^1& Z_{ij}^1\end{array}\right)\left(\begin{array}{c}O_b\\ O_jx\end{array}\right)_B$$ (19) since gauge-invariant operators with no factors of $`x`$ only mix with other similar operators. Then since $`O_3^{\mu \lambda }`$ is symmetric, it can only mix with the symmetric operators $`O_3^{\mu \lambda }`$ and $`O_4^{\mu \lambda }`$, which for forward matrix elements implies that $`O_3^{\mu [\lambda }x^{\nu ]}`$ only mixes with itself and $`O_1^{\mu \nu \lambda }+O_2^{\mu [\lambda }x^{\nu ]}`$. Finally, since the full angular momentum current is conserved and therefore not renormalised, the columns of the mixing matrix must all add to one. This implies the following form for the mixing matrix for forward matrix elements: $`\left(\begin{array}{c}O_1^{\mu \nu \lambda }\\ O_2^{\mu [\lambda }x^{\nu ]}\\ O_3^{\mu [\lambda }x^{\nu ]}\end{array}\right)_B=`$ (23) $`\left(\begin{array}{ccc}1+X& 0& 0\\ ZX& 1+Z& Y\\ Z& Z& 1+Y\end{array}\right)\left(\begin{array}{c}O_1^{\mu \nu \lambda }\\ O_2^{\mu [\lambda }x^{\nu ]}\\ O_3^{\mu [\lambda }x^{\nu ]}\end{array}\right)_R`$ (30) and one-loop calculations show $`Y=\frac{2}{3}n_f\frac{\alpha _s}{4\pi }\frac{1}{ϵ}`$ and $`Z=\frac{8}{3}C_F\frac{\alpha _s}{4\pi }\frac{1}{ϵ}`$. $`X`$ is due to the anomaly and is $`O(\alpha _s^2)`$. For the form factors, this gives: $`\left(\begin{array}{c}a^0\\ B_q\\ B_g\end{array}\right)_B=`$ (34) $`\left(\begin{array}{ccc}1+X& 0& 0\\ 0& 1+Z& Y\\ 0& Z& 1+Y\end{array}\right)\left(\begin{array}{c}a^0\\ B_q\\ B_g\end{array}\right)_R`$ (41) We therefore find the evolution equations for the quark and gluon components of the proton angular momentum: $$\frac{d}{dt}\left(\begin{array}{c}J_q\\ J_g\end{array}\right)=\frac{\alpha _s}{4\pi }\left(\begin{array}{cc}\frac{8}{3}C_F& \frac{2}{3}n_f\\ \frac{8}{3}C_F& \frac{2}{3}n_f\end{array}\right)\left(\begin{array}{c}J_q\\ J_g\end{array}\right)$$ (42) together with eq.(4) for $`a^0`$. It follows that in the asymptotic limit $`Q^2\mathrm{}`$, the partitioning of quark:gluon angular momenta is $`3n_f:16`$ . Interestingly, this is the same result as for the partitioning of momentum obtained from the first moment of the unpolarised pdfs. ## 5 Partons and orbital angular momentum We have seen how the $`g_1^p`$ and angular momentum sum rules involve three Lorentz invariant form factors $`a^0`$, $`J_q`$ and $`J_g`$, where $`J_{q(g)}`$ may be reasonably identified as total quark (gluon) angular momentum components in the sum rule $`\frac{1}{2}=J_q+J_g`$ while the axial charge $`a^0`$, which enters the $`g_1^p`$ sum rule, decouples and may instead be interpreted in terms of topological charge density. There is no gauge and Lorentz invariant operator identification of ‘orbital angular momenta’ $`L_q`$ and $`L_g`$. In the parton model, the axial charge $`a^0`$ is interpreted as a sum of polarised quark $`\mathrm{\Delta }q`$ and gluon $`\mathrm{\Delta }g`$ distributions as in eq.(7). In the AB scheme, the RG evolution equations are $$\frac{d\mathrm{\Delta }q}{dt}=0\frac{d\mathrm{\Delta }g}{dt}=\frac{\alpha _s}{4\pi }\left(3C_F\mathrm{\Delta }q+\beta _0\mathrm{\Delta }g\right)$$ (43) where $`\beta _0=11\frac{2}{3}n_f`$, compatible with eq.(4). This evolution can also be directly obtained from the splitting functions. Now if as the parton model suggests, $`\frac{1}{2}\mathrm{\Delta }q`$ and $`\mathrm{\Delta }g`$ are to be interpreted as quark and gluon spins, then to complete the angular momentum sum rule we are forced to write $$\frac{1}{2}=\frac{1}{2}\mathrm{\Delta }q+\mathrm{\Delta }g+L_q+L_g$$ (44) where $`L_{q(g)}`$ are orbital angular momenta. Since there is no intrinsic operator definition of these partonic quantities, the best we can do is to define them to be consistent with the sum rule (17), i.e. $$L_q=J_q\frac{1}{2}\mathrm{\Delta }qL_g=J_g\mathrm{\Delta }g$$ (45) Of course this means the sum rule (24) is not really predictive, since there is no way to independently measure $`L_{q(g)}`$ – only the form factors $`J_{q(g)}`$ can be extracted from experiment. Moreover, the identifications (25) are not very natural, since they involve subtracting quantities belonging to form factors for different Lorentz structures. They are therefore frame-dependent, not surprisingly given that spin and orbital angular momentum are associated with different representations of the Lorentz group. Nevertheless, if we adopt (7),(25),(24) as the best possible gauge-invariant definition of a quark/gluon spin/orbital angular momentum sum rule, we may determine the RG evolution for the components $`\mathrm{\Delta }q`$, $`\mathrm{\Delta }g`$, $`L_q`$ and $`L_g`$ from eqs. (22),(23). We find $`{\displaystyle \frac{d}{dt}}\left(\begin{array}{c}\mathrm{\Delta }q\\ \mathrm{\Delta }g\\ L_q\\ L_g\end{array}\right)=`$ (50) $`{\displaystyle \frac{\alpha _s}{4\pi }}\left(\begin{array}{cccc}0& 0& 0& 0\\ 3C_F& \beta _0& 0& 0\\ \frac{4}{3}C_F& \frac{2}{3}n_f& \frac{8}{3}C_F& \frac{2}{3}n_f\\ \frac{5}{3}C_F& 11& \frac{8}{3}C_F& \frac{2}{3}n_f\end{array}\right)\left(\begin{array}{c}\mathrm{\Delta }q\\ \mathrm{\Delta }g\\ L_q\\ L_g\end{array}\right)`$ (59) which is consistent with the non-renormalisation of the full angular momentum current and may, at least in part, also be derived in a splitting function approach .
warning/0007/math0007177.html	ar5iv	text	# Transitive and Co–Transitive caps ## 1 Introduction Let $`PG(r,q)`$ be the projective space of dimension $`r`$ over $`GF(q)`$. A $`k`$cap $`\overline{K}`$ in $`PG(r,q)`$ is a set of $`k`$ points, no three of which are collinear , and a $`k`$–cap is said to be complete if it is maximal with respect to set–theoretic inclusion. The maximum value of $`k`$ for which there is known to exist a $`k`$–cap in $`PG(r,q)`$ is denoted by $`m_2(r,q)`$. Some known bounds for $`m_2(r,q)`$ are given below. Suppose that $`\overline{K}`$ is a cap in $`PG(r,q)`$ with automorphism group $`\overline{G}_0P\mathrm{\Gamma }L(r+1,q)`$. Then $`\overline{K}`$ is said to be transitive if $`\overline{G}_0`$ acts transitively on $`\overline{K}`$, and co-transitive if $`\overline{G}_0`$ acts transitively on $`PG(r,q)\overline{K}`$. Our main result is the following theorem. ###### Theorem 1 Suppose $`\overline{K}`$ is a transitive, co–transitive cap in $`PG(r,q)`$. Then one of the following occurs: 1. $`\overline{K}`$ is an elliptic quadric in $`PG(3,q)`$ and $`q`$ is a square when $`q`$ is odd; 2. $`\overline{K}`$ is the Suzuki–Tits ovoid in $`PG(3,q)`$ and $`q=2^h`$, with $`h`$ odd and $`3`$; 3. $`\overline{K}`$ is a hyperoval in $`PG(2,4)`$; 4. $`\overline{K}`$ is an $`11`$–cap in $`PG(4,3)`$ and $`\overline{G_0}M_{11}`$; 5. $`\overline{K}`$ is the complement of a hyperplane in $`PG(r,2)`$; 6. $`\overline{K}`$ is a union of Singer orbits in $`PG(r,q)`$ and $`G_0\mathrm{\Gamma }L(1,p^d)GL(d,p)`$. In each of $`1`$$`5`$ $`\overline{K}`$ is indeed a transitive co–transitive cap. Our conclusion is that transitive, co–transitive caps are rare with the possible exception of unions of Singer cyclic orbits. The origin of this problem are papers by Hill , , in which he studies such caps whose automorphism group acts $`2`$–transitively on the cap. \[As he notes \[8, Theorem 1\], it is trivial to show that if $`\overline{K}`$ is a subset of $`PG(r,q)`$ lying in no proper subspace and admitting a $`3`$–transitive group then $`\overline{K}`$ must be a cap.\] Hill gives a short list of possibilities (omitting Suzuki–Tits ovoids) but excludes caps in $`PG(r,q)`$ for $`q>2`$ and $`r13`$. We find no new caps but show that any other transitive, co–transitive cap is a union of Singer cyclic orbits. The known upper bounds on cap sizes are summarised in the following Result. ###### Result 2 \[10, Theorem 27.3.1\] $`m_2(2,q)=q+1`$ (for $`q`$ odd); $`m_2(2,q)=q+2`$ (for $`q`$ even); $`m_2(3,q)=q^2+1`$ for $`q>2`$; $`m_2(r,2)=2^r`$; and $`m_2(r,q)q^{r1}`$ for $`q>2`$ and $`r4`$. We begin by showing that as a consequence of Result 2, a cap must be smaller than its complement (with one exception). It then follows that in considering subgroups of $`P\mathrm{\Gamma }L(r+1,q)`$ having two orbits, we need only consider the smaller orbit when looking for transitive, co-transitive caps. ###### Lemma 3 Suppose that $`\overline{K}`$ is a cap in $`PG(r,q)`$. Then either $`\|\overline{K}\|<(q^{r+1}1)/2(q1)`$, or $`q=2`$ and $`\overline{K}`$ is the complement of a hyperplane. Proof. It is easy to deduce from Result 2, that the result holds when $`q2`$. Thus suppose now that $`q=2`$ and that $`\|\overline{K}\|(2^{r+1}1)/2`$. The only possiblity is that $`\|\overline{K}\|=2^2`$. Let $`x\overline{K}`$. For each $`y\overline{K}\{x\}`$ there is a line through $`x`$ and $`y`$ and the $`2^r1`$ such lines must be distinct since $`\overline{K}`$ is a cap. However $`x`$ lies on exactly $`2^r1`$ lines in $`PG(r,2)`$ and so every line in $`PG(r,2)`$ through $`x`$ meets $`\overline{K}`$ in two points and $`PG(r,q)\overline{K}`$ in one point. Therefore any line meeting $`PG(r,q)\overline{K}`$ in at least two points is contained in $`PG(r,q)\overline{K}`$. This shows that $`PG(r,q)\overline{K}`$ is a subspace of $`PG(r,2)`$ and its size shows that it is a hyperplane. $`\mathrm{}`$ Using Result 6 we shall know orbit lengths when looking at candidates for transitive, co-transitive caps. Lemma 5 below helps in eliminating a number of possibilities. ###### Definition 4 Suppose that $`\overline{K}`$ is a cap in $`PG(r,q)`$. For any $`xPG(r,q)`$, the chord-number of $`x`$ is the number of chords (2-secants) of $`\overline{K}`$ passing through $`x`$. ###### Lemma 5 Suppose that $`\overline{K}`$ is a tranistive, co-transitive cap in $`PG(r,q)`$ and suppose that $`xPG(r,q)\overline{K}`$. Let $`k=\|\overline{K}\|`$ and $`m=\|PG(r,q)\overline{K}\|`$. Then the chord-number, $`c`$, of $`x`$ is given by $`c=\frac{k(k1)(q1)}{2m}`$. In particular the expression for $`c`$ always yields an integer. Proof. We count combinations of chords and points of $`PG(r,q)\overline{K}`$ in two ways. Firstly there are $`k(k1)/2`$ chords of $`\overline{K}`$ and each has $`q1`$ points not in $`\overline{K}`$. There is a subgroup $`\overline{G_0}`$ of $`\mathrm{\Gamma }L(r+1,q)`$ acting transitively on $`PG(r,q)\overline{K}`$, so each of these $`m`$ points has the same chord-number $`c`$ and a second count gives $`mc`$ chord-point combinations. Thus $`mc=k(k1)(q1)/2`$ leading to the required expression for $`c`$. $`\mathrm{}`$ The main tool in our investigation is the substantial result by M.W. Liebeck , where the affine permutation groups of rank three are classified. ###### Result 6 Let $`G`$ be a finite primitive affine permutation group of rank three and of degree $`n=p^d`$, with socle $`V`$, where $`V(Z_p^d)`$ for some prime $`p`$, and let $`G_0`$ be the stabilizer of the zero vector in $`V`$. Then $`G_0`$ belongs to one of the following families: (A) 11 Infinite classes; (B) Extraspecial classes with $`G_0N_{\mathrm{\Gamma }L(d,p)}(R)`$, where $`R`$ is a $`2`$–group or $`3`$–group irreducible on $`V`$; (C) Exceptional classes. Here the socle $`L`$ of $`G_0/Z(G_0)`$ is simple (where $`Z(G_0)`$ denotes the centre of $`G_0`$). We shall recall the details of the groups belonging to the classes in (A), (B) and (C) as we need them. Suppose $`\overline{K}`$ is a cap in $`PG(r,q)`$ such that a subgroup $`\overline{G_0}`$ of $`P\mathrm{\Gamma }L(r+1,q)`$ acts transitively on each of $`\overline{K}`$ and its complement. Then $`\overline{G_0}`$ corresponds to a subgroup $`G_0`$ of $`GL(d,p)`$ having three orbits on the vectors of $`V(d,p)`$, where $`p`$ is prime and $`p^d=q^{r+1}`$. Moreover $`G_0`$ will contain matrices corresponding to scalar multiplication by elements of $`GF(q)^{}`$. As we demonstrate shortly, with one exception, $`V(d,p)G_0`$ is primitive as a permutation group, so Liebeck’s theorem may be applied. Notice that since we are interested in groups $`G_0`$ containing $`GF(q)^{}`$ we avoid the possibility of two orbits of vectors in $`V(d,p)`$ giving rise to a single orbit of points in $`PG(r,q)`$. Clearly $`G_0`$ may be embedded in $`\mathrm{\Gamma }L(r+1,q)`$. At the beginning of Section 1 of , Liebeck notes that in his result $`G_0GL(d,p)`$ is embedded in $`\mathrm{\Gamma }L(a,p^{d/a})`$ with $`a`$ minimal. Thus $`r+1a`$ i.e. $`qp^{d/a}`$. Moreover in almost all cases it is clear that the groups he identifies have orbits that are unions of $`1`$–dimensional subspaces of $`V(a,p^{d/a})`$ (excluding the zero vector). If a $`1`$-dimensional subspace over $`GF(p^{d/a})`$ does contains vectors $`u,v`$ that are linearly independent over $`GF(q)`$, then $`u,v`$ and $`u+v`$ correspond to three collinear points in $`PG(r,q)`$ and the orbit in $`PG(r,q)`$ cannot be a cap. Thus in our setting we usually have $`q=p^{d/a}`$: there is just one exception, the class A1, although we have to justify $`q=p^{d/a}`$ for the class A2. ###### Lemma 7 Suppose $`\overline{K}`$ is a transitive, co-transitive cap in $`PG(r,q)`$ with $`\overline{G_0}P\mathrm{\Gamma }L(r+1,q)`$ acting transitively on each of $`\overline{K}`$ and $`PG(r,q)\overline{K}`$ and suppose that $`G_0`$ is the pre-image of $`\overline{G}_0`$ in $`GL(d,p)`$. Let $`H=V(d,p)G_0`$. Then $`H`$ is imprimitive on $`V=V(d,p)`$ if and only if $`q=2`$ and $`\overline{K}`$ is the complement of a hyperplane. Proof. Suppose that $`H`$ is imprimitive on $`V`$. Let $`\mathrm{\Omega }`$ be a block containing $`0`$. Then the two orbits of non–zero vectors of $`G_0`$ are $`\mathrm{\Omega }0`$ and $`V\mathrm{\Omega }`$. Let $`u`$ and $`v`$ be any two vectors in $`\mathrm{\Omega }`$, then $`\mathrm{\Omega }+v`$ is a block containing $`0+v`$ and $`u+v`$ so $`\mathrm{\Omega }+v=\mathrm{\Omega }`$. In other words $`u+v`$ is in $`\mathrm{\Omega }`$ and so $`\mathrm{\Omega }`$ is a $`GF(p)`$–subspace of $`V`$. More than this $`G_0`$ contains the scalars in $`GF(q)^{}`$ and so $`\mathrm{\Omega }`$ is actually a $`GF(q)`$–subspace. Thus $`\mathrm{\Omega }`$ cannot correspond to a cap. In $`PG(r,q)`$ our two orbits consist of points in a subspace and the complement. A line not in the subspace meets the subspace in at most one point so the complement cannot form a cap except perhaps when $`p=q=2`$ and the subspace has projective dimension $`r1`$. Conversely, as is well known, the complement of a hyperplane is indeed a cap in $`PG(r,2)`$ and it is the only way in which the complement of a subspace is a cap. It is easy to see that this cap is transitive and co–transitive. $`\mathrm{}`$ We recall for the reader that the socle of a finite group is the product of its minimal normal subgroups. In our setting $`V(d,p)G_0`$ has $`V`$ as its unique minimal normal subgroup. Liebeck’s theorem tells us the possibilities for $`G_0`$ and gives two orbits of $`G_0`$ on the non–zero vectors of $`V(d,p)`$. We denote these by $`K_1`$ and $`K_2`$, and the corresponding sets of points in $`PG(r,q)`$ by $`\overline{K_1}`$ and $`\overline{K_2}`$. We assume that neither $`K_1`$ nor $`K_2`$ lies in a subspace of $`V(r+1,q)`$; given $`GF(q)^{}G_0`$ this means that neither $`K_1`$ nor $`K_2`$ lies in a subspace of $`V(d,p)`$. We may henceforth assume that $`V(d,p)G_0`$ is a finite primitive affine permutation group of rank $`3`$ and degree $`p^d`$, so we may apply Result 6. We begin with the case by case analysis. In many cases we use data from Result 6 and apply Lemmas 2, 5, but there are occasions when we need to look at the structure of orbits in detail; there are also occasions when using the structure of the orbits is more illuminating and yet no less efficient than the bound and chord-number arguments. ## 2 The infinite classes A ### 2.1 The class A1 In this case $`G_0`$ is a subgroup of $`\mathrm{\Gamma }L(1,p^d)`$ containing $`GF(q)^{}`$. Such a subgroup is generated by $`\omega ^N`$ and $`\omega ^e\alpha ^s`$, for some $`N,e,s`$ where $`\omega `$ is a primitive element of $`GF(p^d)`$ and $`\alpha `$ is the generating automorphism $`xx^p`$ of $`GF(p^d)`$; if we write $`p^d=q^a`$, then $`N`$ divides $`(q^a1)/(q1)`$. Let $`H_0`$ be the subgroup of $`\mathrm{\Gamma }L(1,p^d)`$ generated by $`\omega ^N`$. Then $`H_0`$ is a Singer subgroup of $`GL(1,p^d)`$ and the orbits of $`H_0`$ in $`PG(r,q)`$ are called Singer orbits. Clearly if $`G_0`$ has two orbits in $`PG(r,q)`$, then each orbit is the union of Singer orbits. If the smaller orbit is to be a cap, then each Singer orbit must itself be a cap. A precise criterion for deciding when Singer orbits are caps in $`PG(r,q)`$ is given by Szőnyi \[14, Proposition 1\]. Precise criteria for there to be two orbits for $`G_0`$ on non–zero vectors of $`V(d,p)`$ are given by Foulser and Kallaher . These involve numbers $`m_1`$ and $`v`$ such that the primes of $`m_1`$ divide $`p^s1`$, $`v`$ is a prime $`2`$ and ord$`{}_{v}{}^{}p_{}^{sm_1}=v1`$, $`(e,m_1)=1`$, $`m_1s(v1)\|d`$, $`N=vm_1`$. The orbit lengths are $`m_1(p^d1)/N`$ and $`(v1)m_1(p^d1)/N`$. Notice that when $`p=2`$ the smaller orbit has odd size. Hill suggests the possibility of transitive, co–transitive caps of size $`78`$ in $`PG(5,4)`$ and $`430`$ in $`PG(6,4)`$. It is now clear that these cannot be caps from class $`A1`$ and our main theorem then shows that they cannot be caps at all. ### 2.2 The class A2 $`G_0`$ preserves a direct sum $`V_1V_2`$, where $`V_1,V_2`$ are subspaces of $`V(d,p)`$. One orbit must be $`K_1=(V_1V_2)\{0\}`$ and the other $`K_2=\{v_1+v_2:0v_1V_1,0v_2V_2\}`$. We first show that $`V_1,V_2`$ are subspaces over $`GF(q)`$. Observe that for any $`\lambda GF(q)^{}G_0`$, $`\lambda V_1=V_1`$ or $`V_2`$ and let $`F=\{\lambda GF(q)^{}:\lambda V_1=V_1\}\{0\}`$. Then $`F`$ is a subfield of $`GF(q)`$ having order greater than $`q/2`$ so must be $`GF(q)`$. It is now clear that $`V_1,V_2`$ are subspaces of $`V(r+1,q)`$ of dimension $`t=(r+1)/2`$. Given that $`r2`$, we must have $`t2`$, so $`\overline{K_1}`$ contains lines of $`PG(r,q)`$ and cannot be a cap. Moreover $`\|\overline{K_1}\|=2(q^t1)/(q1)<(q^{r+1}1)/2`$ so $`\overline{K_1}`$ is the smaller orbit and therefore $`\overline{K_2}`$ cannot be a cap. ### 2.3 The class A3 $`G_0`$ preserves a tensor product $`V_1V_2`$ over $`GF(q)`$, with $`V_1`$ having dimension 2 over $`GF(q)`$. One orbit must be $`K_1=\{v_1v_2:0v_1V_1,0v_2V_2\}`$ and the other $`K_2=V(K_1\{0\})`$. Consider the $`GF(q)`$–subspace $`V_1v_2`$ of V for some $`0v_2V_2`$. It has dimension $`2`$ in $`V(r+1,q)`$ so corresponds to a line in $`PG(r,q)`$ inside $`\overline{K_1}`$. Hence $`\overline{K_1}`$ is not a cap. Let $`b`$ be the dimension of $`V_2`$ over $`GF(q)`$. Then $`r+1=2b`$ and $`\|\overline{K_1}\|=(q+1)(q^b1)/(q1)`$ (\[12, Table12\]) so $`\|\overline{K_2}\|=q(q^b1)(q^{b1}1)/(q1)>\|\overline{K_1}\|`$ except when $`q=2,b=2`$ (i.e., $`r+1=d=4`$). Thus there is only one case in which $`\overline{K_2}`$ can possibly be a cap. Suppose that $`q=p=2`$ and $`d=4`$, i.e. we are reduced to considering caps in $`PG(3,2)`$. In $`PG(3,2)`$, we see that $`\|\overline{K_1}\|=9`$ and $`\|\overline{K_2}\|=6`$ . Thus here $`\overline{K_1}`$ is too big and for $`\overline{K_2}`$ it is simplest to note that $`(\mathrm{6.5.1})/(2.9)`$, so neither is a cap (by Lemmas 2 and 5). ### 2.4 The class A4 $`G_0SL(a,s)`$ and $`p^d=s^{2a}`$. Here $`q=s^2`$ and $`p^d=q^a`$ with $`SL(a,s)`$ embedded in $`GL(d,p)`$ as a subgroup of $`SL(a,q)`$: let $`e_1,e_2,\mathrm{},e_a`$ be a basis for $`V`$ over $`GF(q)`$ then with respect to this basis $`SL(a,s)`$ consists of the matrices in $`SL(a,q)`$ having every entry in $`GF(s)`$. If $`G_0`$ has two orbits on non-zero vectors of $`V`$ then those orbits must be $`K_1=\{\gamma \lambda _ie_i`$ ($`\lambda _iGF(s)`$, not all $`0`$),$`0\gamma GF(q)\}`$ and $`K_2`$ the set of all remaining non-zero vectors. In other words $`\overline{G_0}`$ preserves a subgeometry of $`PG(r,q)`$. We have $`r>1`$ so that $`a3`$. Thus three collinear points of $`PG(r,s)`$ are still three collinear points in $`PG(r,q)`$ and so $`\overline{K_1}`$ is not a cap. Let us turn to $`\overline{K_2}`$. As noted above, $`r>1`$ so $`a3`$. Let $`u=e_1+\sigma e_2,v=e_2+\sigma e_3`$, where $`\sigma GF(q)GF(s)`$. Then $`u,v`$ and $`u+v=e_1+(\sigma +1)e_2+\sigma e_3K_2`$ correspond to collinear points of $`PG(r,q)`$, all in $`\overline{K_2}`$. Hence $`\overline{K_2}`$ is not a cap. ### 2.5 The class A5 $`G_0SL(2,s)`$ and $`p^d=s^6`$. Here $`q=s^3`$ and $`p^d=q^2`$ with $`SL(2,s)`$ embedded in $`GL(d,p)`$ as a subgroup of $`SL(2,q)`$. However $`r=1`$ in this case so it does not concern us. ### 2.6 The class A6 $`G_0SU(a,q^{})`$ and $`p^d=((q^{})^2)^a`$. In this case $`q=(q^{})^2`$. Here one orbit $`K_1`$ consists of the non-zero isotropic vectors and the other orbit $`K_2`$ consists of the non-isotropic vectors with respect to an appropriate non-degenerate hermitian form. Each orbit is a union of $`1`$–dimensional subspaces of $`V(a,q)`$ (excluding the zero vector). To begin with, a non–isotropic line of $`PG(r,q)`$ contains at least three isotropic points, i.e., three points of $`\overline{K_1}`$. Therefore $`\overline{K_1}`$ cannot be a cap. Now consider $`\overline{K_2}`$. Given $`a3`$, consider a line of $`PG(r,q)`$ that is isotropic but not totally isotropic, then it contains one point of $`\overline{K_1}`$ and $`q4`$ points of $`\overline{K_2}`$. Hence $`\overline{K_2}`$ is not a cap. ### 2.7 The class A7 $`G_0\mathrm{\Omega }^\pm (a,q)`$ and $`p^d=(q)^a`$ with $`a`$ even (and if $`q`$ is odd , $`G_0`$ contains an automorphism interchanging the two orbits of $`\mathrm{\Omega }^\pm (a,q)`$ on non-singular 1-spaces). The arguments here are similar to the Unitary case. $`K_1`$ consists of the non-zero singular vectors and $`K_2`$ consists of the non-singular vectors. Let $`b`$ be the Witt index of the approppriate quadratic form on $`V(a,q)`$ i.e., the dimension of a maximal totally singular subspace. Then $`a`$ is one of $`2b,2b+2`$. Any totally singular line would be a line of $`PG(r,q)`$ lying inside $`\overline{K_1}`$. Given that $`a3`$, it follows that the only possibility for $`\overline{K_1}`$ being a cap is when $`\overline{K_1}`$ is an elliptic quadric in $`PG(3,q)`$. In passing we note that for odd $`q`$, the necessary automorphism is contained in $`G_0`$ only when $`q`$ is square; in this case and in the case $`q`$ even, the elliptic quadric gives a well known cap. Let us turn to $`\overline{K_2}`$. Any anisotropic line of $`PG(r,q)`$ lies inside $`\overline{K_2}`$ so $`\overline{K_2}`$ can never be a cap. ### 2.8 The class A8 $`G_0SL(5,q)`$ and $`p^d=(q)^{10}`$ (from the action of $`SL(5,q)`$ on the skew square $`^2(V(5,q))`$. Here one orbit of non-zero vectors must be $`K_1=\{0uv:u,vV(5,q)\}`$ with the other non-zero vectors belonging to $`K_2`$. One can argue in a similar manner to the case of the tensor product. However it is quicker here to note that the orbits of $`\overline{G_0}`$ on $`PG(r,q)`$ have sizes $`k=(q^51)(q^2+1)/(q1)`$ and $`m=q^2(q^51)(q^31)/(q1)`$ (\[12, Table12\]) with $`k<m`$ for all values of $`q`$. The chord-number is then given by $`c=k(k1)(q1)/2m`$ by Lemma 5 i.e., $`c=(q^2+1)(q^3+q+1)/2q`$. Hence neither $`\overline{K_1}`$ nor $`\overline{K_2}`$ is a cap. ### 2.9 The class A9 $`G_0/Z(G_0)\mathrm{\Omega }(7,q)Z_{(2,q1)}`$ and $`p^d=q^8`$ (from the action of $`B_3(q)`$ on a spin module) , . The study of Clifford algebras leads to the construction of ”spin modules” for $`P\mathrm{\Omega }(m,q)`$. When $`m=8`$ this leads to the triality automorphism of $`P\mathrm{\Omega }^+(8,q)`$. One finds that it is possible (via this automorphism) to embed $`\mathrm{\Omega }(7,q)P\mathrm{\Omega }(7,q)`$ inside $`P\mathrm{\Omega }^+(8,q)`$ as an irrdeucible subgroup. The important thing from our point of view is that two non–trivial orbits of $`G_0`$ must be the set of all non–zero singular vectors and the set of all non–singular vectors with respect to a non–degenerate quadratic form on $`V(8,q)`$. In this setting the arguments employed for class $`A7`$ apply: neither orbit can be a cap. ### 2.10 The class A10 $`G_0/Z(G_0)P\mathrm{\Omega }^+(10,q)`$ and $`p^d=q^{16}`$ (from the action of $`D_5(q)`$ on a spin module) , . Once again we have a spin representation, this time of $`P\mathrm{\Omega }^+(10,q)`$ on $`PG(15,q)`$. On this occasion it is quickest to work from the orbit lengths. The orbits of $`\overline{G_0}`$ on $`PG(r,q)`$ have sizes $`k=(q^81)(q^3+1)/(q1)`$ and $`m=q^3(q^81)(q^51)/(q1)`$ (\[12, Table12\]) with $`k<m`$ for all values of $`q`$. The chord-number is then given by $`c=k(k1)(q1)/2m`$ by Lemma 5 i.e., $`c=(q^3+1)(q^5+q^2+1)/2q^2`$. Hence neither $`\overline{K_1}`$ nor $`\overline{K_2}`$ is a cap. ### 2.11 The class A11 $`G_0Sz(q)`$ and $`p^d=(q)^4`$, with $`q8`$ an odd power of $`2`$ (from the embedding $`Sz(q)Sp(4,q)`$). Here the smaller orbit $`\overline{K_1}`$ on $`PG(3,q)`$ is a Suzuki–Tits ovoid containing $`q^2+1`$ points and this is indeed a cap , \[9, 16.4.5\]. ## 3 The Extraspecial classes In most cases here $`G_0M`$ where $`M`$ is the normalizer in $`\mathrm{\Gamma }L(a,q)`$ of a $`2`$–group $`R`$, where $`p^d=(q)^a`$ and $`a=2^m`$ for some $`m1`$; either $`R`$ is an extraspecial group $`2^{1+2m}`$ or $`R`$ is isomorphic to $`Z_42^{1+2m}`$. In all cases here $`p`$ is odd. There are two types of extraspecial group $`2^{1+2m}`$, denoted $`R_1^m`$ and $`R_2^m`$; the first of these has the structure $`D_8D_8\mathrm{}D_8`$ ($`m`$ copies) and the second $`D_8D_8\mathrm{}D_8Q_8`$ ($`m1`$ copies of $`D_8`$), where $`D_8`$ and $`Q_8`$ are respectively the dihedral and quaternion groups of order $`8`$, and ’$``$’ indicates a central product. The group $`Z_42^{1+2m}`$ is again a central product, this time $`Z_4D_8D_8\mathrm{}D_8`$ ($`m`$ copies of $`D_8`$) and is denoted by $`R_3^m`$. Notice that $`R`$ modulo its centre is an elementary abelian $`2`$–group, i.e. a $`2m`$–dimensional vector space over $`GF(2)`$ and in fact $`M/RZ`$ ($`Z`$ being the centre of $`\mathrm{\Gamma }L(a,q)`$) may be embedded in $`GSp(2m,2)`$. In just one case $`G_0M`$ with $`M`$ the normalizer in $`\mathrm{\Gamma }L(3,4)`$ of a $`3`$–group of order $`27`$. We record from \[12, Table 13\] that in this case the non–trivial orbit sizes of $`G_0`$ on $`V(3,4)`$ are $`27`$ and $`36`$, i.e. the point orbit sizes in $`PG(2,4)`$ are $`9`$ and $`12`$, but the largest possible size of a cap (here better termed an arc) in $`PG(2,4)`$ is $`6`$. Hence there are no caps here and we may henceforth assume that $`R`$ is a $`2`$–group, with $`p`$ odd. There are sixteen instances where $`G_0`$ has two non–trivial orbits on $`V(d,p)V(a,q)`$, but ten of these have $`a=2`$ (i.e. $`m=1`$) and so refer to action on a projective line, i.e. $`r<2`$; note that two of these cases have $`q>p`$. Thus we concentrate on the remaining six cases. In each of these cases $`q=p`$ and in all but the last case the vector space is $`V(4,p)`$. In the last case the vector space is $`V(8,3)`$. Four cases follw immediately from known bounds - they are listed in the table below. \| p=q \| r \| R \| smaller orbit size \| max. cap size \| \| --- \| --- \| --- \| --- \| --- \| \| 3 \| 3 \| $`R_1^2`$ \| 16 \| 10 \| \| 5 \| 3 \| $`R_2^2`$ \| 60 \| 26 \| \| 5 \| 3 \| $`R_3^2`$ \| 60 \| 26 \| \| 7 \| 3 \| $`R_2^2`$ \| 80 \| 50 \| The case $`𝐩=𝐪=\mathrm{𝟑}`$, $`𝐫=\mathrm{𝟕}`$, $`𝐑=𝐑_\mathrm{𝟐}^\mathrm{𝟑}`$. In this case smaller orbit of $`\overline{G_0}`$ on $`PG(7,3)`$ has size $`720`$, while the maximum size for a cap in $`PG(7,3)`$ is only known to be $`729`$. Instead we use Lemma 5: the larger orbit has size $`2560`$ and $`(\mathrm{720.719.2})/(2.2560)`$. The case $`𝐩=𝐪=\mathrm{𝟑}`$, $`𝐫=\mathrm{𝟑}`$, $`𝐑=𝐑_\mathrm{𝟐}^\mathrm{𝟐}`$. Here Liebeck notes that $`R`$ has five orbits of size $`16`$ on $`V(4,3)`$ and $`M`$ permutes these orbits acting as $`S_5`$, the symmetric group of degree $`5`$. Thus there are a number of possibilities for $`G_0`$ having two non–trivial orbits on $`V(4,3)`$. However it is straightforward to construct generating matrices for $`R`$ and we see immediately that one orbit of size $`16`$ on $`V(4,3)`$ cannot correspond to a cap in $`PG(3,3)`$. Therefore none of the orbits of size $`16`$ can correspond to a cap and hence no possible choices of $`G_0`$ can give rise to a cap. ## 4 The Exceptional classes Finally we turn to the exceptional classes where the socle $`L`$ of $`G_0/Z(G_0)`$ is simple. There are just thirteen different possibilities for $`L`$, although on occasion more than one possibility for $`G_0`$ corresponds to a given $`L`$. For example for $`L=A_5`$ there are seven different possibilities for $`G_0`$ (one of which leads to a single orbit in $`PG(d1,p)`$); however all of these lead to $`r<2`$ and so do not concern us. We employ a variety of techniques to tackle these cases. Liebeck \[12, Table 14\] gives the orbit sizes in $`V(d,p)`$ and sometimes we can use these to rule out the possibility of caps. On other occasions we can use the fact that the chord-number is an integer. On two occasions, neither of these appraoches works and we have to investigate the known structure of the smaller orbit. There remain two cases where a cap does occur. The cases where caps occur. When $`L=A_6`$ and $`(d,p)=(6,2)`$, $`L`$ admits an embedding in $`PSL(3,4)`$ (so here $`q=4`$) and $`G_0`$ has an orbit of size $`6`$. In fact this in a hyperoval in $`PG(2,4)`$ , so we do have a cap. When $`L=M_{11}`$ and $`(d,p)=(5,3)`$ there is a representation of $`L`$ in which one orbit has size $`11`$ and in fact this is a cap. In passing we note that this cap arises as an orbit of a Singer cyclic subgroup of $`PG(4,3)`$ ; moreover $`PG(4,3)`$ is partitioned into eleven $`11`$–caps (the eleven orbits of the Singer cyclic subgroup). Note also that there is a second representation of $`L=M_{11}`$ on $`PG(4,3)`$ (see below). In fact both representations appear in the context of the ternary Golay code \[1, Ch. 6\]. Cases where known bounds rule out caps. In each of the following cases the smaller orbit is larger than the known upper bound for a cap size, so cannot be a cap. In the table $`k`$ is the smaller orbit size. \| $`L`$ \| $`(d,p)`$ \| $`r`$ \| $`q`$ \| $`k`$ \| max. cap size \| \| --- \| --- \| --- \| --- \| --- \| --- \| \| $`A_6`$ \| $`(4,5)`$ \| 3 \| 5 \| 36 \| 26 \| \| $`A_7`$ \| $`(4,7)`$ \| 3 \| 7 \| 120 \| 50 \| \| $`M_{11}`$ \| $`(5,3)`$ \| 4 \| 3 \| 55 \| $`27`$ \| \| $`J_2`$ \| $`(6,5)`$ \| 5 \| 5 \| 1890 \| $`625`$ \| \| $`J_2`$ \| $`(12,2)`$ \| 5 \| 4 \| 525 \| $`256`$ \| Cases where $`c`$ an integer rules out caps. In each of the following cases a calculation $`c=k(k1)(q1)/2m`$ yields a non-integer and so by Lemma 5, the smaller orbit does not correspond to a cap. In the table $`k`$ is the smaller orbit size and $`m`$ the larger orbit size. \| $`L`$ \| $`(d,p)`$ \| $`r`$ \| $`q`$ \| $`k`$ \| $`m`$ \| \| --- \| --- \| --- \| --- \| --- \| --- \| \| $`A_9`$ \| $`(8,2)`$ \| 7 \| 2 \| 120 \| 135 \| \| $`A_{10}`$ \| $`(8,2)`$ \| 7 \| 2 \| 45 \| 210 \| \| $`L_2(17)`$ \| $`(8,2)`$ \| 7 \| 2 \| 102 \| 153 \| \| $`M_{24}`$ \| $`(11,2)`$ \| 10 \| 2 \| 276 \| 1771 \| \| $`M_{24}`$ \| $`(11,2)`$ \| 10 \| 2 \| 759 \| 1288 \| \| $`Suz`$ or $`J_4`$ \| $`(12,3)`$ \| 11 \| 2 \| 65520 \| 465920 \| The case $`𝐋=𝐀_\mathrm{𝟕}`$ and $`(𝐝,𝐩)=(\mathrm{𝟖},\mathrm{𝟐})`$. Here $`L`$ is embedded in $`PSL(4,4)`$ (so $`q=4`$). In fact $`L`$ may actually embedded in $`A_8PSL(4,2)PSL(4,4)`$. The group $`A_8`$ and therefore $`A_7`$ preserve a subgeometry whose $`15`$ points form the smaller orbit. There are numerous examples of three points on a line in the subgeometry. Thus we have no caps. The case $`𝐋=\mathrm{𝐏𝐒𝐔}(\mathrm{𝟒},\mathrm{𝟐})`$ and $`(𝐝,𝐩)=(\mathrm{𝟒},\mathrm{𝟕})`$. The vectors in the smaller orbit are given by Liebeck \[12, Lemma 3.4\]: $$(\theta ;0,0,0),(0;\theta ,0,0),(0;\omega ^a,\omega ^b,\omega ^c),(\omega ^a;0,\omega ^b,\omega ^c),$$ (together with all scalar multiples) where $`\theta =\omega =2`$; $`a,b,c`$ take any integral values; and the last three coordinates may be permuted cyclically. It suffices here to observe that $`(1;0,0,0)`$, $`(1;0,1,6)`$ and $`(2;0,1,6)`$ all lie in this orbit and give three collinear points in $`PG(3,7)`$. So no cap arises here. AMS Mathematics Subject Classification: Primary 51E22; Secondary 20B15, 20B25 Keywords and Phrases: Caps, Rank 3 permutation groups A. Cossidente, Dipartimento di Matematica, Università della Basilicata, via N.Sauro 85, 85100 Potenza, Italy. e-mail: cossidente@unibas.it O.H. King, Department of Mathematics, The University of Newcastle, Newcastle Upon Tyne, NE1 7RU, United Kingdom. e-mail: o.h.kink@ncl.ac.uk
warning/0007/cond-mat0007177.html	ar5iv	text	# Fermi-Liquid Interactions in 𝑑-Wave Superconductors \[ ## Abstract This article develops a quantitative quasiparticle model of the low-temperature properties of $`d`$-wave superconductors which incorporates both Fermi-liquid effects and band-structure effects. The Fermi-liquid interaction effects are found to be classifiable into strong and negligible renormalizaton effects, for symmetric and antisymmetric combinations of the energies of $`k`$ and $`k`$ quasiparticles, respectively. A particularly important conclusion is that the leading clean-limit temperature-dependent correction to the superfluid density is not renormalized by Fermi-liquid interactions, but is subject to a Fermi velocity (or mass) renormalization effect. This leads to difficulties in accounting for the penetration depth measurements with physically acceptable parameters, and hence reopens the question of the quantitative validity of the quasiparticle picture. \] There is now considerable experimental evidence that the cuprate high $`T_c`$ superconductors exhibit the simple power law temperature dependences predicted by the quasiparticle picture for their thermodynamic and transport properties at temperatures well below $`T_c`$. For example, penetration depth measurements find that the superfluid density exhibits a low-temperature clean-limit linear in T temperature dependence , in agreement with theory . The NMR relaxation rate exhibits the expected $`T^3`$ temperature dependence . The predicted effect of impurities in giving rise to a universal thermal conductivity has been confirmed . The clean limit specific heat varying as $`T^2`$ appears to have been observed . Even the electrical transport relaxation rate observed in microwave conductivity experiments , which had resisted explanation for some time, has now been explained in terms of a quasiparticle picture . Whether or not the magnitudes of the coefficients of the above power law temperature dependences are accurately given by a quasiparticle description is at present an open question. A recent study correlating these different coefficients concludes that the quasiparticle model may be successful here also provided a Fermi-liquid interaction factor multiplying the superfluid density is treated as an adjustable parameter. The contention of this article is that there is however no Fermi-liquid interaction renormalization of the linear in T contribution to the inverse square penetration depth. There is instead a renormalization by a factor involving the ratio of a band Fermi velocity to a Landau quasiparticle Fermi velocity. The difficulty now is that a physically unreasonable value of this renormalization parameter is obtained from experiment. This reopens the question of to what extent the quasiparticle picture can provide an accurate quantitative picture of the low temperature behavior of the high $`T_c`$ superconductors. Recent debate on correctness of the quasiparticle picture is also occurring in connection with ARPES experiments , and in connection with the role of phase fluctuations of the complex order parameter in the determination of the temperature dependence of the superfluid density . The potential importance of Fermi-liquid interactions in renormalizing the superfluid density has been emphasized in Refs. and . These papers note that Fermi-liquid renormalization effects in $`d`$-wave superconductors can be either strong or weak according as the contributing quasiparticles are from the entire Fermi surface, or confined the the nodal points where the $`d`$-wave gap goes to zero. Both expect a strong renormalization effect for the superfluid density, whereas this article does not find such an effect. This article shows that the physics of Fermi liquid effects in $`d`$-wave superconductors has an interesting symmetry property. This manifests itself when the quasiparticle energies are separated into parts that are $`𝒮`$ymmetric and $`𝒜`$ntisymmetric combinations of the energies of the $`+k`$ and $`k`$ states. (The calligraphic letters $`𝒮`$ and $`𝒜`$ are used here to emphasize the difference with the more usual definition of the symmetric and antisymmetric combinations with respect to $`+k`$ and $`+k`$ states common in normal state analyses, e.g. see Eq. 1.32 of Ref. .) In the presence of Fermi-liquid interactions, the $`𝒮`$ymmetric and $`𝒜`$ntisymmetric corrections to the quasiparticle energies obey integral equations that are independent of each other, and they are renormalized differently. The $`𝒮`$ymmetric energy corrections exhibit strong Fermi-liquid renormalization effects, while the the $`𝒜`$ntisymmetric energy corrections exhibit relatively weak temperature-dependent renormalizations that can often be neglected. Temperature gives a $`𝒮`$ymmetric correction to the quasiparticle energy because $`+k`$ and $`k`$ states are affected in the same way by temperature. A superfluid flow generates an $`𝒜`$ntisymmetric correction since the components of $`+k`$ and $`k`$ along the superfluid velocity have opposite signs. Also the Zeeman interaction generates an $`𝒜`$ntisymmetric correction because the spin $``$ and spin $``$ contributions to the energy have opposite signs. Thus the superfluid density and the magnetic susceptibility are negligibly renormalized by Fermi-liquid interactions, while the effects of temperature (although relatively small) are strongly renormalized by Fermi-liquid interactions. The approach of this article to the inclusion of Fermi-liquid interactions in the study of the superconducting state follows the intuitively appealing approach of Ref. , which is consistent with more formal correlation function approaches. Rather than starting with a band energy in the absence of electron-electron interactions of $`ϵ_k^b=\mathrm{}^2k^2/(2m)`$ as in Ref. , however, this article allows $`ϵ_k^b`$ to be an arbitrary function of $`𝐤`$ so as to be able to account for anisotropic Fermi surface effects. To form a Hamiltonian from $`ϵ_k^b`$ the substitution $`\mathrm{}𝐤i\mathrm{}e𝐀/c`$ is made. Also, this article studies only equilibrium properties, and does not develop a kinetic equation. Other studies of Fermi liquid interactions in unconventional superconductivity include Refs. . The Hamiltonian describing the excitations of the superconducting state has the following form: $$=\underset{k}{}\begin{array}{cc}[c_k^{}& c_k]\end{array}\left[\begin{array}{cc}\zeta _k+\lambda _k& \mathrm{\Delta }_k\\ \mathrm{\Delta }_k& \zeta _k+\lambda _k\end{array}\right]\left[\begin{array}{c}c_k\\ c_k^{}\end{array}\right].$$ (1) Here $`\zeta _k=\xi _k+\delta \epsilon _k^𝒮+h_k^𝒮`$ where $`\xi _k`$ is the Landau quasiparticle energy relative to the chemical potential (neglecting quasiparticle interactions), $`\lambda _k=\delta \epsilon _k^𝒜+h_k^𝒜`$, and $`\mathrm{\Delta }_k`$ is the momentum-dependent gap function appropriate for $`d`$-wave symmetry. Fermi liquid interactions give a contribution to the energy of a quasiparticle with momentum k and spin $`\sigma `$ due to other excited quasiparticles which is . $$\delta \epsilon _{k\sigma }=\frac{1}{L^2}\underset{k^{}}{}\left[f_{kk^{}}^{\sigma \sigma }\delta n_{k^{}\sigma }+f_{kk^{}}^{\sigma \overline{\sigma }}\delta n_{k^{}\overline{\sigma }}\right].$$ (2) where $`\overline{\sigma }\sigma `$, $`n_{k\sigma }=c_{k\sigma }^{}c_{c\sigma }`$, and $`\delta `$ indicates a variation due to the excitation of other electrons and holes, either by temperature or by the presence of external fields. (The factor $`L^2`$ occurs in Eq. 2 because the intention is to develop a model applicable to superconductivity in a two-dimensional copper-oxide plane of a high $`T_c`$ superconductor having area $`L^2`$.) An important step in the analysis, as described qualitatively above, is the separation of the Fermi liquid interactions into $`𝒮`$ymmetric and $`𝒜`$ntisymmetric parts defined by $$\delta \epsilon _k^𝒜=\frac{1}{2}\left[\delta \epsilon _k\delta \epsilon _k\right],\delta \epsilon _k^𝒮=\frac{1}{2}\left[\delta \epsilon _k+\delta \epsilon _k\right].$$ (3) The quantities $`h_k^𝒜`$ and $`h_k^𝒮`$ in $``$ represent generalized external fields. For example, in the case of an external magnetic field acting on the orbital motion of the electrons, the gap function will acquire a complex phase. This phase factor can be removed by a gauge transformation, $`c_{k\sigma }c_{k\sigma }exp(i\theta )`$, the end result of which is the addition of the field $`h_k^𝒜=𝐯_k^b𝐩_s,h_k^𝒮=0`$, where $`𝐯_k^bϵ_k^b/𝐤`$, to the Hamiltonian. Here $`𝐩_s=\mathrm{}\theta e𝐀/c`$ is the superfluid momentum, which is assumed to be sufficiently slowly varying spatially that its gradients can be neglected. The velocity $`𝐯_k^b`$ that appears in $`h_k^𝒜`$ is the bare band velocity, unrenormalized by the electron-electron interaction, as noted following Eq. 12 of Ref. , and it is this same velocity that appears in the expression for the quasiparticle contribution to the current density (Eq. 5 below). On the other hand, the electron-electron interaction contributes to the quasiparticle energy $`\xi _k`$ defined in and following Eq. 1, and hence affects the quasiparticle velocity $`v_F`$ that occurs in $`E_k^{(0)}`$ below. The differences in these two velocities have important quantitative consequences for the interpretation of the penetration depth data, as will be seen below. The Hamiltonian of Eq. 1 can be used to find the thermal equilibrium expectation value of the electrical current density operator giving, for the electrical current density $`𝐉=\eta _g𝐩_s+𝐉_{qp}`$, with the gauge contribution determined by $$\eta _g=\frac{e}{2L^2}\underset{k\sigma }{}n_{k\sigma }\frac{1}{\mathrm{}^2}\left(\frac{^2}{k_x^2}+\frac{^2}{k_y^2}\right)ϵ_k^b,$$ (4) and the quasiparticle contribution given by $$𝐉_{qp}=\frac{e}{L^2}\underset{k\sigma }{}𝐯_k^bf(E_{k,\sigma }),$$ (5) $`f(E_{k,\sigma })`$ is the Fermi-Dirac distribution function, and the $`E_{k,\sigma }`$ are the Bogoliubov quasiparticle energies defined below. The case of an external magnetic field H acting on the spin degrees of freedom is described by taking $`h_k^𝒜=\mu _BH,h_k^𝒮=0`$. In both of these cases, the magnetic field acts only on the $`𝒜`$ntisymmetric mode, and has no effect on the $`𝒮`$ymmetric mode of excitation. In addition to causing changes in the energy of a quasiparticle (as in Eq. 2), excited quasiparticles can give rise to changes in the gap function . There are however no changes that are linear in the superfluid momentum , and this effect will therefore be neglected. The diagonalization of the Hamiltonian of Eq. 1 gives $$=\underset{k\sigma }{}E_{k,\sigma }\gamma _{k,\sigma }^{}\gamma _{k,\sigma },E_{k,\sigma }=E_{\sigma k}+\sigma (\delta \epsilon _{\sigma k}^𝒜+h_{\sigma k}^𝒜)$$ (6) where $`\sigma =\pm 1`$, $`E_k=\sqrt{\zeta _k^2+\mathrm{\Delta }_k^2}`$, and the $`\gamma _{k,\sigma }^{}`$ are operators creating Bogoliubov quasiparticles. Later, the energy $`E_k^{(0)}=\sqrt{\xi _k^2+\mathrm{\Delta }_k^2}`$ describing the quasiparticle spectrum in the absence of other excited quasiparticles is also used. For a $`d`$-wave superconductor, the quasiparticle energy can be parameterized in the neighborhood of the Fermi-surface nodal points (see Fig. 1) as $`E_k^{(0)}=\sqrt{(p_1v_F)^2+(p_2v_2)^2}`$, where $`p_1`$ and $`p_2`$ are components of the momentum relative to the nodal point in directions perpendicular and parallel to the Fermi line. At low temperatures, only quasiparticles close to these four points can be thermally excited. Using Eq. 2 in Eq. 3, keeping only terms up to linear order in the $`\delta \epsilon `$’s and $`h`$’s, and dropping some terms of order $`k_BT/(\mathrm{}k_Fv_2)`$ relative to those kept, yields the integral equations $$\delta \epsilon _k^𝒮=\frac{1}{L^2}\underset{k^{}}{}f_{kk^{}}^{(+)}\left[\frac{\xi _k^{}}{E_k^{}^{(0)}}f(E_k^{}^{(0)})\frac{\mathrm{\Delta }_k^{}^2}{E_k^{}^{(0)3}}(\delta \epsilon _k^{}^𝒮+h_k^{}^𝒮)\right]$$ (7) and $$\delta \epsilon _k^𝒜=\frac{1}{L^2}\underset{k^{}}{}f_{kk^{}}^{()}\frac{f}{E_k^{}^{(0)}}\left(\delta \epsilon _k^{}^𝒜+h_k^{}^𝒜\right),$$ (8) where $`f_{kk^{}}^{(\pm )}=f_{kk^{}}^{\sigma \sigma }\pm f_{k,k^{}}^{\sigma \overline{\sigma }}`$. Consider first the $`𝒮`$ymmetrical corrections to the quasiparticle energies (Eq. 7), and assume that there are no $`𝒮`$ymmetrical external fields other than temperature, i.e. $`h_k^𝒮=0`$ (as is the case for the external magnetic fields of most interest in this article, which are purely $`𝒜`$ntisymmetrical). Then the only term driving a nonzero contribution to $`\delta \epsilon _k^𝒮`$ is the term on the right hand side proportional to $`f(E_k^{}^{(0)})`$ and describing the effect of temperature. This term is proportional to $`T^3`$, thus giving a $`\delta \epsilon _k^𝒮T^3`$, and will not be important in contributing to the properties of interest at the temperatures satisfying $`k_BT\mathrm{\Delta }_0`$ ($`\mathrm{\Delta }_0`$ is the maximum gap). Thus $`\delta \epsilon _k^𝒮`$ will be neglected in calculations below. Now from Eq. 8, which determines the $`𝒜`$ntisymmetric corrections to the quasiparticle energies, it is clear that only the values of $`\delta \epsilon _k^𝒜`$ and $`h_k^𝒜`$ at the Fermi surface nodes are relevant to the low energy properties. Also, the solutions of Eq. 8 can be classified according to the irreducible representation of the point group $`C_{4v}`$ (or 4mm) describing a tetragonal copper-oxide plane of a high T<sub>c</sub> superconductor, the independent solutions being $`\delta \epsilon _{A_g}^𝒜`$ $`=`$ $`\left(\delta \epsilon _1^𝒜+\delta \epsilon _2^𝒜+\delta \epsilon _3^𝒜+\delta \epsilon _4^𝒜\right)/4`$ (9) $`\delta \epsilon _{xy}^𝒜`$ $`=`$ $`\left(\delta \epsilon _1^𝒜+\delta \epsilon _2^𝒜\delta \epsilon _3^𝒜\delta \epsilon _4^𝒜\right)/4`$ (10) $`\delta \epsilon _{Ex}^𝒜`$ $`=`$ $`\left(\delta \epsilon _1^𝒜+\delta \epsilon _2^𝒜+\delta \epsilon _3^𝒜+\delta \epsilon _4^𝒜\right)/4`$ (11) $`\delta \epsilon _{Ey}^𝒜`$ $`=`$ $`\left(\delta \epsilon _1^𝒜+\delta \epsilon _2^𝒜+\delta \epsilon _3^𝒜+\delta \epsilon _4^𝒜\right)/4`$ (12) where the indices 1,2,3 and 4 refer to the four nodes in the excitation spectrum, as defined in Fig. 1. The external fields $`h_k^𝒜`$ at the nodes can be similarly classified. The solution of Eq. 8 now yields $$\delta \epsilon _\mathrm{\Gamma }^𝒜(T)=h_\mathrm{\Gamma }^𝒜F_\mathrm{\Gamma }^𝒜(T)/[1+F_\mathrm{\Gamma }^𝒜(T)]$$ (13) with $`F_\mathrm{\Gamma }^𝒜(T)=f_\mathrm{\Gamma }^𝒜ln(2)k_BT/(2\pi \mathrm{}^2v_Fv_2)`$. Here $`\mathrm{\Gamma }`$ represents any of the irreducible representations present in Eqs. 12. The $`f_\mathrm{\Gamma }^𝒜`$’s are defined by $`f_{A_g}^𝒜`$ $`=`$ $`f_{11}^a+f_{13}^a+2f_{12}^a`$ (14) $`f_{xy}^𝒜`$ $`=`$ $`f_{11}^a+f_{13}^a2f_{12}^a`$ (15) $`f_E^𝒜`$ $`=`$ $`f_{11}^sf_{13}^s`$ (16) where $$f_{kk^{}}^{s,a}f_{kk^{}}^{\sigma \sigma }\pm f_{kk^{}}^{\sigma \overline{\sigma }}$$ (17) are the symmetric and antisymmetric combinations of the Fermi-liquid parameters familiar from normal state analyses. , and $`f_{12}^a`$ for example is $`f_{kk^{}}^a`$ for $`k`$ and $`k^{}`$ at nodes 1 and 2, respectively. It is also useful to use Eq. 7 to obtain an idea of how the $`𝒮`$ymmetrical external fields are renormalized by Fermi-liquid interactions. It is clear from Eq. 7 that a knowledge of the Fermi-liquid interaction on the entire Fermi surface is required and that $`\delta \epsilon _k^𝒮`$ must be determined on the entire Fermi surface. To obtain a rough idea of the nature of the solutions, consider a circular Fermi surface of radius $`k_F`$ and look for a solution of $`A_g`$ symmetry by considering a Fermi liquid interaction $`f_{kk}^{(+)}=f_{A_g}^𝒮`$, independent of $`k`$ and $`k^{}`$, and a $`𝒮`$ymmetrical external field $`h_{A_g}^𝒮`$ independent of $`k`$. The solution, which is also independent of $`k`$ on the Fermi surface is $$\delta \epsilon _{A_g}^𝒮=\frac{F_{A_g}^𝒮}{1+F_{A_g}^𝒮}h_{A_g}^𝒮$$ (18) where $`F_{A_g}^𝒮=f_{A_g}^𝒮k_F/(\pi \mathrm{}v_F)`$. In contrast to the $`𝒜`$ntisymmetrical Fermi liquid parameters $`F_\mathrm{\Gamma }^𝒜(T)`$ obtained above, which go to zero linearly with temperature in the superconducting state in the clean limit (and hence have a dependence on temperature $`T`$ explicitly indicated), the $`𝒮`$ymmetrical Fermi liquid parameter $`F_{A_g}^𝒮`$ is temperature independent and of approximately the same magnitude as the corresponding normal state Fermi liquid parameter. The same can be seen to be true of the $`𝒮`$ymmetrical Fermi liquid parameters corresponding to other irreducible representations of $`C_{4v}`$. Note that the ratio of the $`𝒜`$ntisymmetrical to the $`𝒮`$ymmetrical Fermi-liquid $`F`$ parameters is $`F^𝒜/F^𝒮(f^𝒜/f^𝒮)[k_BT/(\mathrm{}k_Fv_2)]`$. As noted above, the presence of a superfluid momentum contributes an $`𝒜`$ntisymmetrical external field to the Hamiltonian of Eq. 1. This external field corresponds to the E irreducible representation of $`C_{4v}`$ with the $`p_{sx}`$ and $`p_{sy}`$ components of $`𝐩_s`$ corresponding to the components $`Ex`$ and $`Ey`$ of Eq. 12. The current density is thus easily evaluated using Eq. 5 with Eqs. 6, 13 and 16. The result is $`𝐉_{qp}=\eta _{qp}𝐩_s`$ where $$\eta _{qp}(T)=\frac{2ln2e(v_F^b)^2k_BT}{[1+F_E^𝒜(T)]\pi \mathrm{}^2v_Fv_2}$$ (19) Note that the Fermi liquid correction does not alter the clean limit linear in T contribution to the $`\eta _{qp}(T)`$, but rather makes a $`T^2`$ contribution (using $`(1+F)^1(1F+\mathrm{})`$). Thus there are no Fermi liquid corrections to the experimentally measured linear in $`T`$ contribution to inverse square penetration depth. The penetration depth $`\lambda `$ is thus given by $$\lambda ^2(T)=\lambda ^2(0)\frac{8ln2e^2}{c^2\mathrm{}^2}\alpha ^2\frac{v_F}{v_2}k_BT+\mathrm{}$$ (20) where $`\alpha =(v_F^b/v_F)`$, and $`\lambda ^2(0)`$ is determined by $`\eta _g`$ given in Eq. 4. This has exactly the same form as Eq. 6 of Ref. , which finds (from a detailed analysis of a number of experiments) $`\alpha ^2=`$ 0.43 for Bi<sub>2</sub>Sr<sub>2</sub>CaCu<sub>2</sub>O<sub>8</sub> and $`\alpha ^2=`$ 0.46 for YBa<sub>2</sub>Cu<sub>3</sub>O<sub>7-δ</sub>. The conclusions here are however completely different from those drawn in Ref. which, based on previous theoretical work, considered $`\alpha ^2`$ to be a Fermi liquid correction with a not unreasonable value. The conclusion of this article is that the experimentally determined value of $`\alpha ^2`$ implies a value of $`v_F^b`$ significantly smaller than $`v_F`$, which is not physically reasonable. One of the essential features of strongly correlated electron systems such as the copper-oxide superconductors is that the strong electron-electron correlations are expected to produce narrow energy bands and large quasiparticle masses, leading to $`v_F<v_F^b`$. The renormalization of the spin susceptibility due to Fermi-liquid interactions can be calculated in a similar way. The Zeeman interaction of the spin of an electron with the magnetic field contributes an $`𝒜`$ntisymmetric external field of $`A_g`$ symmetry to the Hamiltonian. It follows that the magnetic moment per unit area of a copper oxide plane is $$M=\frac{\mu _B}{L^2}\underset{k}{}\left[f(E_{k,1})f(E_{k,1})\right]=\chi H$$ (21) where $$\chi (T)=\frac{\chi _0(T)}{1+F_{A_g}^𝒜(T)},\chi _0(T)=\frac{\mu _B^2ln2k_BT}{\pi \mathrm{}^2v_Fv_2}.$$ (22) Note that here also the low-temperature clean-limit linear in T magnetic susceptibility is not changed by Fermi-liquid interactions. These affect only terms of order $`T^2`$ and higher in the susceptibility. This article has given a detailed description of both Fermi-liquid effects and band structure effects within the framework of a quasiparticle picture of the low-temperature properties of $`d`$-wave superconductors. This opens the way for a detailed quantitative experimental test of the quasiparticle picture of the low-temperature properties. A classification of Fermi-liquid effects is given that separates the strong renormalization and weak renormalization effects according to a symmetry property. The application of the results to the interpretation of penetration depth measurements is of particular interest. This corresponds to the weak (and in fact negligible) Fermi-liquid renormalization case, and the ultimate conclusion is that the experimental results imply a band Fermi velocity which is smaller than the corresponding quasiparticle Fermi velocity, which is an unphysical result. Clearly there are at present problems with the quantitative aspects of the quasiparticle picture of the low-temperature properties of the high T<sub>c</sub> superconductors and further study is desirable. I acknowledge stimulating discussion with L. Taillefer, the hospitality of P. Nozières, and the Theory Group of the Institut Laue Langevin where much of this work was done, and the support of the Canadian Institute for Advanced Research and of the Natural Sciences and Engineering Research Council of Canada.
warning/0007/hep-ph0007181.html	ar5iv	text	# I Introduction ## I Introduction The recent results from the Super-Kamiokande (SuperK) detector confirm the solar and atmospheric neutrino deficits and strongly suggest the existence of neutrino oscillations. The azimuthal angle and energy dependence of the atmospheric data indicates a mass-squared difference scale, $`\delta m_{\mathrm{atm}}^23.5\times 10^3`$ eV<sup>2</sup>. The LSND measurements indicate neutrino oscillations with a different scale, $`\delta m_{\mathrm{LSND}}^20.3`$$`2.0`$ eV<sup>2</sup>; there is also a small region of acceptable parameter space at 6 eV<sup>2</sup>. The evidence for oscillations in solar neutrino data , when taken as a whole, prefer yet a third, lower mass-squared difference scale, $`\delta m_{\mathrm{sun}}^210^4`$ eV<sup>2</sup>. Once oscillation explanations for some or all of the data are accepted, the next step is to attempt to find a neutrino mass and mixing pattern that can provide a unified description of all the relevant neutrino data. In this paper we consider two possible neutrino scenarios which have been proposed to account for the LSND results: (i) a three-neutrino model that also describes the atmospheric neutrino data (in which case the solar data would be explained by a phenomenon other than oscillations ), and (ii) a four-neutrino model that also describes both the solar and atmospheric data . As we show, the three-neutrino model may be considered a sub-case of the four-neutrino model. We examine the ability of short-baseline experiments at a muon storage ring-based neutrino factory to determine the oscillation parameters in each case, and discuss how the three- and four-neutrino scenarios may be distinguished. We concentrate on experiments in which muon and tau neutrinos are detected via charge-current interactions and there is good sign determination of the detected muons and tau leptons. The sign determination allows one to distinguish between $`\overline{\nu }_e\overline{\nu }_\mu `$ and $`\nu _\mu \nu _\mu `$ for stored $`\mu ^{}`$ in the ring, and between $`\nu _e\nu _\mu `$ and $`\overline{\nu }_\mu \overline{\nu }_\mu `$ for stored $`\mu ^+`$. We consider primarily the appearance channels $`\nu _e\nu _\mu `$, $`\nu _e\nu _\tau `$, and $`\nu _\mu \nu _\tau `$ (and the corresponding antineutrino channels), for which the uncertainty on the beam flux is not critical to the sensitivity of the measurements. We find that the $`\nu _\mu \nu _\tau `$ and $`\nu _e\nu _\mu `$ channels provide the best sensitivity to the $`CP`$-violating phase; the $`\nu _e\nu _\tau `$ channel allows measurement of an independent combination of mixing parameters. Short-baseline experiments at a neutrino factory are sufficient to determine the complete three-neutrino parameter space in these scenarios when their results are combined with other accelerator-based experiments that are currently underway or being planned. Measuring $`CP`$-violation in both $`\nu _\mu \nu _\tau `$ and $`\nu _e\nu _\mu `$ allows one to possibly distinguish between the three- and four-neutrino cases. Hence all three off-diagonal oscillation channels are useful. In the four-neutrino case it is not possible to have complete mixing and $`\delta m^2`$ parameter determinations without additional measurements in future long-baseline experiments. ## II Three-neutrino models Here we address the ability of short-baseline neutrino oscillation experiments to probe a class of three-neutrino models that can describe the results of the LSND experiment, together with the atmospheric neutrino deficit observed by the SuperK experiment. ### A Oscillation Formalism In a three-neutrino model the neutrino flavor eigenstates $`\nu _\alpha `$ are related to the mass eigenstates $`\nu _j`$ in vacuum by a unitary matrix $`U`$, $$\|\nu _\alpha =\underset{j}{}U_{\alpha j}\|\nu _j,$$ (1) with $`\alpha =e,\mu ,\tau `$ and $`j=1,2,3`$. The Maki-Nakagawa-Sakata (MNS) mixing matrix can be parameterized by $$U=\left(\begin{array}{ccc}c_{13}c_{12}& c_{13}s_{12}& s_{13}e^{i\delta }\\ c_{23}s_{12}s_{13}s_{23}c_{12}e^{i\delta }& c_{23}c_{12}s_{13}s_{23}s_{12}e^{i\delta }& c_{13}s_{23}\\ s_{23}s_{12}s_{13}c_{23}c_{12}e^{i\delta }& s_{23}c_{12}s_{13}c_{23}s_{12}e^{i\delta }& c_{13}c_{23}\end{array}\right),$$ (2) where $`c_{jk}\mathrm{cos}\theta _{jk}`$, $`s_{jk}\mathrm{sin}\theta _{jk}`$, and $`\delta `$ is the $`CP`$ non-conserving phase. Two additional diagonal phases are present in $`U`$ for Majorana neutrinos, but these do not affect oscillation probabilities. For three-neutrino models in which the results of the atmospheric and LSND experiments are explained by neutrino oscillations there are two independent mass-squared differences: $`\delta m_{\mathrm{LSND}}^2=0.3`$$`2.0`$ eV<sup>2</sup> and $`\delta m_{\mathrm{atm}}^23.5\times 10^3`$ eV<sup>2</sup>. We take $`\delta m_{32}^2=\delta m_{\mathrm{atm}}^2`$ and $`\delta m_{21}^2=\delta m_{\mathrm{LSND}}^2`$; see Fig. 1a. These mass-squared differences obey the condition $`\delta m_{\mathrm{LSND}}^2\delta m_{\mathrm{atm}}^2`$. Then the general off-diagonal vacuum oscillation probability is $$P(\nu _\alpha \nu _\beta )4\|U_{\alpha 1}\|^2\|U_{\beta 1}\|^2\mathrm{sin}^2\mathrm{\Delta }_{\mathrm{LSND}}4\mathrm{R}\mathrm{e}(U_{\alpha 2}U_{\alpha 3}^{}U_{\beta 2}^{}U_{\beta 3})\mathrm{sin}^2\mathrm{\Delta }_{\mathrm{atm}}\pm 2JS,$$ (3) where $`\mathrm{\Delta }_j`$ $``$ $`1.27\delta m_j^2(\mathrm{eV}^2)L(\mathrm{km})/E_\nu (\mathrm{GeV}),`$ (4) $`S`$ $``$ $`\left[\mathrm{sin}2\mathrm{\Delta }_{\mathrm{atm}}+\mathrm{sin}2\mathrm{\Delta }_{\mathrm{LSND}}\mathrm{sin}2(\mathrm{\Delta }_{\mathrm{LSND}}+\mathrm{\Delta }_{\mathrm{atm}})\right]`$ (5) $`=`$ $`2(\mathrm{sin}2\mathrm{\Delta }_{\mathrm{atm}}\mathrm{sin}^2\mathrm{\Delta }_{\mathrm{LSND}}+\mathrm{sin}2\mathrm{\Delta }_{\mathrm{LSND}}\mathrm{sin}^2\mathrm{\Delta }_{\mathrm{atm}}),`$ (6) and $`J`$ is the $`CP`$-violating invariant, which can be defined as $`J=Im\{U_{e2}U_{e3}^{}U_{\mu 2}^{}U_{\mu 3}\}`$. The plus (minus) sign in Eq. (3) is used when $`\alpha `$ and $`\beta `$ are in cyclic (anticyclic) order, where cyclic order is defined as $`e\mu \tau `$. For antineutrinos, the sign of the $`CP`$-violating term is reversed. For the mixing matrix in Eq. (2), $$J=s_{13}c_{13}^2s_{12}c_{12}s_{23}c_{23}\mathrm{sin}\delta .$$ (7) For recent discussions of $`CP`$ violation in neutrino oscillations, see Refs. , and . For $`L100`$ km, matter effects are very small and the vacuum formulas are a good approximation to the true oscillation probabilities. The class of scenarios that we are considering is designed to account for the large $`\nu _\mu \nu _\tau `$ mixing of atmospheric neutrinos and the small $`\nu _\mu \nu _e`$ mixing in the LSND experiment. The $`\nu _e`$ survival probability at the leading oscillation scale is given by $$P(\nu _e\nu _e)=P(\overline{\nu }_e\overline{\nu }_e)=4c_{12}^2c_{13}^2(1c_{12}^2c_{13}^2)\mathrm{sin}^2\mathrm{\Delta }_{\mathrm{LSND}}.$$ (8) Results from the BUGEY reactor experiment put an upper bound on $`P(\overline{\nu }_e\overline{\nu }_e)`$ for $`\delta m_{\mathrm{LSND}}^2>0.01`$ eV<sup>2</sup>, which provides the approximate constraint $$s_{12}^2+s_{13}^20.01.$$ (9) This leads to the conditions $$\theta _{12},\theta _{13}\theta _{23},$$ (10) with the mixing of atmospheric neutrinos, $`\theta _{23}`$, near maximal ($`\theta _{23}\pi /4`$). Hence we can take $`s_{12},s_{13}s_{23},c_{23}`$ and $`c_{12}c_{13}1`$. The off-diagonal oscillation probabilities are (to leading order in the small mixing angle parameters) $`P(\nu _e\nu _\mu )`$ $``$ $`4\|s_{12}c_{23}+s_{13}s_{23}e^{i\delta }\|^2\mathrm{sin}^2\mathrm{\Delta }_{\mathrm{LSND}}2s_{12}s_{13}c_\delta \mathrm{sin}2\theta _{23}\mathrm{sin}^2\mathrm{\Delta }_{\mathrm{atm}}`$ (11) $`+s_{12}s_{13}s_\delta \mathrm{sin}2\theta _{23}S,`$ (12) $`P(\nu _e\nu _\tau )`$ $``$ $`4\|s_{12}s_{23}s_{13}c_{23}e^{i\delta }\|^2\mathrm{sin}^2\mathrm{\Delta }_{\mathrm{LSND}}+2s_{12}s_{13}c_\delta \mathrm{sin}2\theta _{23}\mathrm{sin}^2\mathrm{\Delta }_{\mathrm{atm}}`$ (13) $`s_{12}s_{13}s_\delta \mathrm{sin}2\theta _{23}S,`$ (14) $`P(\nu _\mu \nu _\tau )`$ $``$ $`4\|s_{12}c_{23}+s_{13}s_{23}e^{i\delta }\|^2\|s_{12}s_{23}s_{13}c_{23}e^{i\delta }\|^2\mathrm{sin}^2\mathrm{\Delta }_{\mathrm{LSND}}`$ (15) $`+\mathrm{sin}^22\theta _{23}\mathrm{sin}^2\mathrm{\Delta }_{\mathrm{atm}}+s_{12}s_{13}s_\delta \mathrm{sin}2\theta _{23}S,`$ (16) where $`c_\delta =\mathrm{cos}\delta `$ and $`s_\delta =\mathrm{sin}\delta `$. The expressions for the antineutrino channels are obtained by changing the sign of the $`S`$-term in each case. A representative scenario, 1B1 in Ref. , that we will use as an example has the parameters $`\delta m_{32}^2`$ $`=`$ $`3.5\times 10^3\mathrm{eV}^2,\mathrm{sin}^22\theta _{23}=1.0`$ (17) $`\delta m_{21}^2`$ $`=`$ $`0.3\mathrm{eV}^2,\mathrm{sin}^22\theta _{12}=\mathrm{sin}^22\theta _{13}=0.015,`$ (18) with $`\delta `$ a varied parameter. More generally, $`\delta m_{\mathrm{LSND}}^2`$ in the range 0.3–2.0 eV<sup>2</sup> is allowed with $`\nu _\mu \nu _e`$ oscillation amplitude given by $$4\|s_{12}c_{23}+s_{13}s_{23}e^{i\delta }\|^2\left(\frac{0.06\mathrm{eV}^2}{\delta m_{\mathrm{LSND}}^2}\right)^2.$$ (19) In a short-baseline experiment, $`L/E`$ should optimally be chosen such that $`\mathrm{\Delta }_{\mathrm{LSND}}1`$, in which case $`\mathrm{\Delta }_{\mathrm{atm}}\mathrm{\Delta }_{\mathrm{LSND}}/1001`$; e.g., for $`\delta m_{\mathrm{LSND}}^2=0.3`$ eV<sup>2</sup> and $`E_\nu =14`$ GeV we might choose $`L=45`$ km, giving: $$\mathrm{\Delta }_{\mathrm{atm}}=0.014\left(\frac{\delta m_{\mathrm{atm}}^2}{3.5\times 10^3\mathrm{eV}^2}\right)\left(\frac{L}{45\mathrm{km}}\right)\left(\frac{14\mathrm{GeV}}{E_\nu }\right).$$ (20) Note that in the forward direction $`E_\nu =14`$ GeV is the average $`\nu _\mu `$ energy for stored unpolarized $`\mu ^{}`$ with $`E_\mu =20`$ GeV. Thus in the probability equations above, $`s_{12}`$, $`s_{13}`$, and $`\mathrm{\Delta }_{\mathrm{atm}}`$ are all small parameters, at the few percent level or less. Then to leading order in $`\mathrm{\Delta }_{\mathrm{atm}}`$ $$S4\mathrm{\Delta }_{\mathrm{atm}}\mathrm{sin}^2\mathrm{\Delta }_{\mathrm{LSND}}.$$ (21) Therefore, in these scenarios, the dominant contribution to $`P(\nu _e\nu _\mu )`$ and $`P(\nu _e\nu _\tau )`$ comes from the leading oscillation (the $`\mathrm{sin}^2\mathrm{\Delta }_{\mathrm{LSND}}`$ term), but the dominant contribution to $`P(\nu _\mu \nu _\tau )`$ comes from the subleading oscillation (the $`\mathrm{sin}^2\mathrm{\Delta }_{\mathrm{atm}}`$ term); these results are summarized in Table I. In each case, the dominant term is $`CP`$-conserving and proportional to the product of two small parameters. The $`CP`$-violating contribution in each case is $`2s_{12}s_{13}s_\delta \mathrm{sin}2\theta _{23}\mathrm{sin}2\mathrm{\Delta }_{\mathrm{atm}}\mathrm{sin}^2\mathrm{\Delta }_{\mathrm{LSND}}`$, which is a product of three small parameters (assuming $`\delta `$ is not small), and hence is smaller than the $`CP`$-conserving contribution (see Table I). In all, there are six parameters to be determined in the expressions in Eqs. (12)–(16): the three mixing angles, the phase $`\delta `$, and two independent mass-squared differences. However, although there are six off-diagonal measurements possible with $`\mu `$ and $`\tau `$ detection \[the three in Eqs. (12)–(16) plus the corresponding antineutrino channels\], the leading and next-to-leading terms in the expressions for these oscillation probabilities are only sensitive to five independent quantities: $`\|s_{12}c_{23}+s_{13}s_{23}e^{i\delta }\|`$, $`\|s_{12}s_{23}s_{13}c_{23}e^{i\delta }\|`$, $`s_{12}s_{13}s_\delta `$, $`\mathrm{sin}^2\mathrm{\Delta }_{\mathrm{LSND}}`$, and $`\mathrm{sin}2\theta _{23}\mathrm{\Delta }_{\mathrm{atm}}`$. For the parameter ranges we are considering, $`\theta _{12},\theta _{13}>\mathrm{\Delta }_{\mathrm{atm}}`$; then the $`\mathrm{sin}^2\mathrm{\Delta }_{\mathrm{LSND}}`$ term in Eq. (16) can be comparable to one of the other terms in $`P(\nu _\mu \nu _\tau )`$, but this term still depends on a subset of these same five independent quantities. The K2K , MINOS , ICANOE and OPERA long-baseline experiments will measure the parameters in the leading term of the $`\nu _\mu \nu _\tau `$ probability, $`\theta _{23}`$ and $`\delta m_{\mathrm{atm}}^2`$, and MiniBooNE will measure the parameters in the leading term of the $`\nu _e\nu _\mu `$ probability, $`\|s_{12}c_{23}+s_{13}s_{23}e^{i\delta }\|`$ and $`\delta m_{\mathrm{LSND}}^2`$. Therefore, only two independent quantities will remain to be measured in short baseline experiments: (i) the amplitude of the leading oscillation in the $`\nu _e\nu _\tau `$ probability, and (ii) the subleading CPV term, which has the same magnitude for each off-diagonal channel and can be determined by a comparison of neutrino to antineutrino rates. Hence a combination of short-baseline measurements with the results of the other accelerator-based experiments would allow all of the parameters in this three-neutrino scenario to be determined. Measurements of the other short-baseline off-diagonal oscillation probabilities may then be used to check the consistency of the result and/or improve the accuracy of the parameter determinations. One can also measure the $`\nu _\mu \nu _\mu `$ survival probability, which to leading order in small quantities can be written $$P(\nu _\mu \nu _\mu )14A(1A)\mathrm{sin}^2\mathrm{\Delta }_{\mathrm{LSND}}\mathrm{sin}^22\theta _{23}\mathrm{sin}^2\mathrm{\Delta }_{\mathrm{atm}},$$ (22) where $$As_{23}^2s_{13}^2+c_{23}^2s_{12}^2+2s_{23}c_{23}s_{12}s_{13}c_\delta .$$ (23) In short-baseline experiments, both oscillatory terms in Eq. (22) are second order in small quantities. A measurement of $`P(\nu _\mu \nu _\mu )`$ could be used in conjunction with other short-baseline measurements to make a complete determination of the oscillation parameters without the use of other data. However, since $`A`$ and $`\mathrm{\Delta }_{\mathrm{atm}}`$ are small, the deviations of $`P(\nu _\mu \nu _\mu )`$ from unity are also small, and the normalization of the beam flux would need to be known to high precision for this to be a useful measurement. An examination of Eqs. (12)–(16) and (21) shows that only the relative sign of $`\delta `$ and $`\delta m_{32}^2`$ may be determined in short-baseline measurements for the distances and oscillation parameter values that we are considering. Long-baseline experiments at a neutrino factory should be able to determine the sign of $`\delta m_{32}^2`$ since there is a significant dependence of the matter effect on the sign of $`\delta m_{32}^2`$ for $`L2000`$ km . ### B Results Figure 2 shows the ratio of the rate of tau production for stored $`\mu ^+`$ to that for stored $`\mu ^{}`$ for the $`\nu _e\nu _\tau `$ and $`\nu _\mu \nu _\tau `$ channels, versus the $`CP`$ phase $`\delta `$ for several values of baseline length $`L`$, with oscillation parameters given by Eq. (18). Statistical errors are also shown, assuming $`10^{20}`$ kt-decays (corresponding, for example, to three years of running with $`10^{20}`$ useful muon decays per year and a 1 kt detector having 33% tau detection efficiency). The event rate calculations are done according to the method outlined in Ref. . For the choice of parameters in Eq. (18), $`s_{12}`$ and $`s_{13}`$ are larger than $`\mathrm{\Delta }_{\mathrm{atm}}`$ and the $`\nu _\mu \nu _\tau `$ channel shows the largest relative $`CP`$-violating effect (see the last column of Table I). The last column in Table I shows that the relative size of the $`CP`$-violating effect in the $`\nu _\mu \nu _\tau `$ channel decreases with increasing $`\mathrm{\Delta }_{\mathrm{atm}}`$ (i.e., with increasing $`L/E`$). Also, for very small $`L/E`$, when $`\mathrm{\Delta }_{\mathrm{LSND}}1`$, $`S0`$ to leading order in $`\mathrm{\Delta }_{\mathrm{atm}}`$ and $`\mathrm{\Delta }_{\mathrm{LSND}}`$, and the $`CP`$ violation becomes negligible. Therefore, for any given set of oscillation parameters, there will be an optimum $`L/E`$ that maximizes the $`CP`$ violation effects in the $`\nu _\mu \nu _\tau `$ channel. Figure 3a shows the ratio of $`\overline{\nu }_\mu \overline{\nu }_\tau `$ event rates (from $`\mu ^+`$ decays) to $`\nu _\mu \nu _\tau `$ event rates (from $`\mu ^{}`$ decays), $`R_{\mu \tau }`$, versus baseline for 20 GeV muons, for oscillation parameters given by Eq. (18) and three values of $`\delta `$ ($`90^{}`$, $`0^{}`$, and $`90^{}`$). Approximate statistical errors are shown for $`10^{20}`$ kt-decays. The figure clearly shows that there is one distance that maximizes the size of the $`CP`$-violation effect, which in this case (20 GeV muons) is about $`L=45`$ km. This optimal distance decreases slowly with increasing $`\delta m_{32}^2`$. For $`\delta m_{32}^2`$ in the range $`2.5`$$`4.5\times 10^3`$ eV<sup>2</sup>, we find that the optimal $`L`$ is in the range $`40`$$`50`$ km for $`\delta m_{\mathrm{LSND}}^2=0.3`$ eV<sup>2</sup>. The optimal $`L`$ scales inversely with $`\delta m_{\mathrm{LSND}}^2`$, and for stored muon energies well above the tau threshold may be approximated by $$L_{\mathrm{opt}}45\mathrm{km}\left(\frac{0.3\mathrm{eV}^2}{\delta m_{\mathrm{LSND}}^2}\right)\left(\frac{E_\mu }{20\mathrm{GeV}}\right);$$ (24) e.g., for $`\delta m_{\mathrm{LSND}}^2=2`$ eV<sup>2</sup> the best sensitivity to $`CP`$ violation is obtained with $`L6`$ km. The distance from Fermilab to Argonne is about 30 km, which would be optimal for $`E_\mu =20`$ GeV and $`\delta m_{\mathrm{LSND}}^2=0.45`$ eV<sup>2</sup>. Similar results for an optimal distance for $`CP`$ violation in $`\nu _\mu \nu _\tau `$ in the context of four-neutrino models have been reported in Ref. . Given Eq. (19) (the LSND constraint on the $`\nu _\mu \nu _e`$ oscillation amplitude), $`CP`$ violation is maximized when $`\theta _{12}=\theta _{13}`$ since $`J`$ is proportional to the product of $`s_{12}`$ and $`s_{13}`$. Figure 3b shows $`R_{\mu \tau }`$ for unequal $`\theta _{12}`$ and $`\theta _{13}`$: $`\mathrm{sin}^22\theta _{12}=0.0336`$ and $`\mathrm{sin}^22\theta _{13}=0.0038`$, which for $`\delta =0`$ gives the same LSND result as $`\mathrm{sin}^22\theta _{12}=\mathrm{sin}^22\theta _{12}=0.015`$. The $`CP`$–violation effects for $`\theta _{12}\theta _{13}`$ are not as dramatic as with $`\theta _{12}=\theta _{13}`$, but still may be observable with $`10^{20}`$ kt-decays. An examination of Table I shows that the relative size of the $`CP`$ violation in the $`\nu _e\nu _\mu `$ or $`\nu _e\nu _\tau `$ channels increases with $`L/E`$. However, once $`\mathrm{\Delta }_{\mathrm{atm}}`$ is of order unity or larger the $`\mathrm{sin}2\mathrm{\Delta }_{\mathrm{atm}}`$ term in $`S`$ averages to zero, washing out the $`CP`$ violation; this does not happen until $`L`$ is much larger than $`100`$ km. Since the flux falls off like $`1/L^2`$, the statistical uncertainty increases roughly like $`L`$, and as long as $`\mathrm{\Delta }_{\mathrm{atm}}1`$ a wide range of distances have comparable sensitivity to $`CP`$ violation in the $`\nu _e\nu _\mu `$ or $`\nu _e\nu _\tau `$ channels. Figure 4a shows the ratio $`R_{e\mu }`$ of $`\nu _e\nu _\mu `$ event rates (from $`\mu ^+`$ decays) to $`\overline{\nu }_e\overline{\nu }_\mu `$ event rates (from $`\mu ^{}`$ decays) versus baseline for oscillation parameters given by Eq. (18), with $`E_\mu =20`$ GeV. The statistical errors correspond to $`2\times 10^{21}`$ kt-decays (which could be obtained, for example, by three years of running with $`10^{20}`$ useful muon decays per year, and a 10 kt detector having a 67% muon efficiency). A 4 GeV minimum energy cut has been made on the detected muon. Although $`R_{e\mu }`$ is not as sensitive to $`CPV`$ effects as $`R_{\mu \tau }`$, the increased statistics (due to a larger overall rate resulting from the use of a larger detector for muons) make $`\nu _e\nu _\mu `$ another attractive channel for $`CP`$ violation. Figure 4b shows similar results for $`R_{e\tau }`$, the ratio of $`\nu _e\nu _\tau `$ event rates (from $`\mu ^+`$ decays) to $`\overline{\nu }_e\overline{\nu }_\tau `$ event rates (from $`\mu ^{}`$ decays). The statistical errors correspond to $`10^{20}`$ kt-decays. It is evident from the figure that the $`\nu _e\nu _\tau `$ channel is not as useful for detecting $`CP`$ violation (primarily because of the reduced event rate in tau detection), although this channel is the most sensitive to the $`\nu _e\nu _\tau `$ oscillation amplitude. The combination of parameters $`\|s_{12}s_{23}s_{13}c_{23}e^{i\delta }\|`$ in the $`\nu _e\nu _\tau `$ amplitude is also present in the $`\mathrm{sin}^2\mathrm{\Delta }_{\mathrm{LSND}}`$ term of $`P(\nu _\mu \nu _\tau )`$ \[see Eq. (16)\], along with an additional factor involving small mixing angles; hence, it would be more difficult to measure $`\|s_{12}s_{23}s_{13}c_{23}e^{i\delta }\|`$ in the $`\nu _\mu \nu _\tau `$ channel. Fig. 5 shows, for various $`CP`$-violating cases with $`E_\mu =20`$ GeV, the statistical significance (number of standard deviations) that the ratios $`R_{\mu \tau }`$, $`R_{e\mu }`$, and $`R_{e\tau }`$ deviate from their expected values for the $`CP`$-conserving case. While both the $`\nu _e\nu _\mu `$ and $`\nu _\mu \nu _\tau `$ channels provide good sensitivity near the optimal $`L`$, $`\nu _e\nu _\mu `$ is more sensitive for a wider range of $`L`$, especially for values of $`\delta `$ that do not give maximal $`CP`$ violation. An additional potential advantage of the $`\nu _e\nu _\mu `$ channel is that, in principle, lower energy neutrino factories can be used since there is no need to be above the tau-lepton production threshold. However, as the energy of muons in a neutrino factory decreases, the sensitivity to CP violation also decreases (see Fig. 6); this is especially true if a lower bound is imposed on the energy of the detected muon. In principle, measurements can be made with varying $`L/E`$, either by using experiments at more than one baseline, or by using a measure of the neutrino energy (for example, the total observed event energy) at a fixed baseline. However, since there are only five independent quantities in the leading- and subleading-order probabilities in Eqs. (12)–(16), such measurements still cannot completely determine the three-neutrino parameter set; this can only be done at short baselines by a measurement of the subsubleading terms in the off-diagonal probabilities, or of the subleading terms in the diagonal $`\nu _\mu \nu _\mu `$ probability \[see Eq (22)\], both of which would be very challenging experimentally. Similar conclusions apply for additional measurements involving electron detection. To summarize, when used in conjunction with long-baseline measurements from K2K, MINOS, ICANOE, and OPERA, and with results from MiniBooNE, all six of the three-neutrino oscillation parameters can in principle be determined with short-baseline measurements at a neutrino factory. The short-baseline measurements would determine two parameters not determined by other accelerator-based experiments, and would provide a consistency check by independently measuring three other parameters also measured by other accelerator-based experiments. The $`\nu _e\nu _\tau `$ channel is most sensitive to the quantity $`\|s_{12}s_{23}s_{13}c_{23}e^{i\delta }\|`$. The $`\nu _e\nu _\mu `$ channel provides good sensitivity to $`CP`$ violation over a wide range of $`L`$ when a large muon detector is used, and the $`\nu _\mu \nu _\tau `$ channel is most useful for detecting $`CP`$ violation near the optimal $`L`$ for the parameter choice illustrated. ## III Four-neutrino models ### A Oscillation formalism Four-neutrino models are required to completely describe the solar, atmospheric, and LSND data in terms of oscillations, since a third independent mass-squared difference is necessary. A fourth neutrino must be sterile, i.e., have negligible interactions, since only three neutrinos are measured in $`Z\nu \overline{\nu }`$ decays . Following Ref. , we label the fourth mass eigenvalue $`m_0`$. Given the pattern of masses $`m_1`$, $`m_2`$, and $`m_3`$ from the three neutrino case in Sec. II, there is a preferred choice for the scale of $`m_0`$ that can fit all of the data, including constraints from accelerator experiments, namely, $`m_0`$ must be nearly degenerate with $`m_1`$, so that there are two pairs of nearly degenerate states separated by a mass gap of about 1 eV ; see Fig. 1b. Then $`\delta m_{10}^2`$ governs the oscillation of solar neutrinos, $`\delta m_{32}^2`$ governs the oscillation of atmospheric neutrinos, and $`\delta m_{21}^2`$, $`\delta m_{31}^2`$, $`\delta m_{20}^2`$, and $`\delta m_{30}^2`$ all contribute to the LSND oscillations. Six mixing angles and three (six) phases are needed to parameterize the mixing of four Dirac (Majorana) neutrinos; only three of these phases are measurable in neutrino oscillations. Thus three additional mixing angles, which we label $`\theta _{01}`$, $`\theta _{02}`$, and $`\theta _{03}`$, and two additional phases are required in extending the three-neutrino phenomenology to the four neutrino case. The simplest situation, which occurs in most explicit four-neutrino models, is that large mixing occurs only between the nearly degenerate pairs; then the four-neutrino mixing matrix can be parametrized as $$U=\left(\begin{array}{cccc}c_{01}& s_{01}^{}& s_{02}^{}& s_{03}^{}\\ & & & \\ s_{01}& c_{01}& s_{12}^{}& s_{13}^{}\\ & & & \\ c_{01}(s_{23}^{}s_{03}+c_{23}s_{02})& s_{01}^{}(s_{23}^{}s_{03}+c_{23}s_{02})& c_{23}& s_{23}^{}\\ +s_{01}(s_{23}^{}s_{13}+c_{23}s_{12})& c_{01}(s_{23}^{}s_{13}+c_{23}s_{12})& & \\ & & & \\ c_{01}(s_{23}s_{02}c_{23}s_{03})& s_{01}^{}(s_{23}s_{02}c_{23}s_{03})& s_{23}& c_{23}\\ s_{01}(s_{23}s_{12}c_{23}s_{13})& +c_{01}(s_{23}s_{12}c_{23}s_{13})& & \end{array}\right),$$ (25) where $`s_{jk}`$ is here defined as $`\mathrm{sin}\theta _{jk}e^{i\delta _{jk}}`$, and the $`\delta _{jk}`$ are the six possible phases for Majorana neutrinos. We set $`\delta _{12}=\delta _{23}=\delta _{02}=0`$ without loss of generality, since only three phases are measurable in neutrino oscillations. In this four-neutrino scenario, the parameters $`\delta m_{32}^2`$, $`\delta m_{21}^2`$, $`\theta _{23}`$, $`\theta _{12}`$ and $`\theta _{13}`$ have the same roles as in the three-neutrino scenario in the previous section; the phase $`\delta _{13}`$ can be identified with $`\delta `$ in the three-neutrino case. The angle $`\theta _{01}`$ describes the mixing of the fourth flavor of neutrino with the state it is nearly degenerate with, $`\nu _e`$; together $`\delta m_{10}^2`$ and $`\theta _{01}`$ can take on the approximate values appropriate to any of the solar neutrino oscillation solutions (small angle MSW, large angle MSW, LOW, or vacuum). The remaining mixing angles $`\theta _{02}`$ and $`\theta _{03}`$ describe mixing of the fourth neutrino with the two states $`\nu _\mu `$, $`\nu _\tau `$ in the other nearly-degenerate pair. Although models with pure oscillations to sterile neutrinos are disfavored for the solar and atmospheric data, models with mixed oscillations to sterile and active neutrinos , i.e., non-negligible $`\theta _{02}`$ and $`\theta _{03}`$, are presumably also possible. Finally, $`\delta _{01}`$ and $`\delta _{03}`$ are extra phases that may be observable in neutrino oscillations with four neutrinos. In the limit that terms involving the solar mass-squared difference can be ignored, the expressions for the oscillation probabilities in Eqs. (12) and (14) remain the same. However, the $`\nu _\mu \nu _\tau `$ probability becomes $`P(\nu _\mu \nu _\tau )`$ $``$ $`4\|(s_{12}c_{23}+s_{13}s_{23}e^{i\delta _{13}})(s_{12}s_{23}s_{13}c_{23}e^{i\delta _{13}})`$ (26) $`+(s_{02}c_{23}+s_{03}s_{23}e^{i\delta _{03}})(s_{02}s_{23}s_{03}c_{23}e^{i\delta _{03}})\|^2\mathrm{sin}^2\mathrm{\Delta }_{\mathrm{LSND}}`$ (27) $`+\mathrm{sin}^22\theta _{23}\mathrm{sin}^2\mathrm{\Delta }_{\mathrm{atm}}+(s_{12}s_{13}\mathrm{sin}\delta _{13}+s_{02}s_{03}\mathrm{sin}\delta _{03})\mathrm{sin}2\theta _{23}S.`$ (28) For $`P(\overline{\nu }_\mu \overline{\nu }_\tau )`$, the sign of the $`S`$-term is reversed. ### B Discussion When combined with measurements from K2K, MINOS, ICANOE, OPERA, and MiniBooNE, the three parameters $`\theta _{12}`$, $`\theta _{13}`$ and $`\delta _{13}`$ can in principle be determined by short-baseline measurements of the $`\nu _e\nu _\tau `$ amplitude and the $`CPV`$ term in the $`\nu _e\nu _\mu `$ channel. This is similar to the three-neutrino case in Sec. II. However, measurements of $`\nu _\mu \nu _\tau `$ and $`\overline{\nu }_\mu \overline{\nu }_\tau `$ will only give partial information on $`\theta _{02}`$, $`\theta _{03}`$, and $`\delta _{03}`$. The $`\nu _\mu \nu _\mu `$ survival probability can also be measured, which in the limit that the solar mass-squared difference can be ignored is $$P(\nu _\mu \nu _\mu )14A(1A)\mathrm{sin}^2\mathrm{\Delta }_{\mathrm{LSND}}\mathrm{sin}^22\theta _{23}\mathrm{sin}^2\mathrm{\Delta }_{\mathrm{atm}},$$ (29) where $$A4\left[s_{23}^2(\|s_{03}\|^2+\|s_{13}\|^2)+c_{23}^2(s_{02}^2+s_{12}^2)+2s_{23}c_{23}\mathrm{Re}(s_{02}s_{03}+s_{12}s_{13})\right].$$ (30) Both oscillatory terms in Eq. (29) are second order in small quantities in short-baseline experiments. In principle $`P(\nu _\mu \nu _\mu )`$ could be used as an additional measurement; however, unless $`\theta _{02}`$ and $`\theta _{03}`$ are larger than $`\theta _{12}`$ and $`\theta _{13}`$, the deviations of $`P(\nu _\mu \nu _\mu )`$ from unity are very small, and the normalization of the beam flux would have to be known to high precision for this to be useful. The parameters $`s_{12}`$, $`s_{13}`$ and $`s_\delta `$ determined by short baseline measurements of the $`\nu _e\nu _\tau `$ and $`\nu _e\nu _\mu `$ channels give predictions for the $`\nu _\mu \nu _\tau `$ channel at short baselines that can differ for the three- and four-neutrino cases \[Eqs. (16) and (28), respectively\]. If the three-neutrino predictions for $`\nu _\mu \nu _\tau `$ are found to substantially disagree with the experimental measurements, then the disagreement would provide evidence for the existence of a fourth neutrino. The absence of such a disagreement would indicate that either there are only three neutrinos or that $`\theta _{02}`$ and $`\theta _{03}`$ are significantly smaller than $`\theta _{12}`$ and $`\theta _{13}`$. In general, the sensitivity of the measurements of $`\nu _\mu \nu _\tau `$ and $`\overline{\nu }_\mu \overline{\nu }_\tau `$ to $`\theta _{02}`$, $`\theta _{03}`$, and $`\delta _{03}`$ in the four-neutrino case are similar to the sensitivities to $`\theta _{12}`$, $`\theta _{13}`$, and $`\delta `$ in the three-neutrino case. For example, given the parameters in Eq. (18) and $`\delta _{13}=0`$ (i.e., no $`CP`$ violation in the $`\nu _e\nu _\mu `$ and $`\nu _e\nu _\tau `$ channels), if $`s_{02}=s_{12}`$ and $`s_{03}=s_{13}`$, then the four-neutrino predictions would be given by Figs. 3a and 3b, where $`\delta _{03}`$ takes on the values of $`\delta `$ in the figures; the corresponding three-neutrino predictions would be given by the $`\delta =0`$ curves. Larger values of $`s_{02}`$ and $`s_{03}`$ could give a much larger $`CP`$-violating effect, provided that $`\delta _{03}`$ was not small. If both $`\delta _{03}`$ and $`\delta _{02}`$ were nonzero, their $`CP`$-violating effects could add together either constructively or destructively. Hence, larger $`CP`$-violating effects are possible in the $`\nu _\mu \nu _\tau `$ channel in the four-neutrino case than with three neutrinos, or $`CP`$ violation may be present in $`\nu _\mu \nu _\tau `$ when it is not present in the $`\nu _e\nu _\mu `$ or $`\nu _e\nu _\tau `$ channels (or vice versa), unlike the three-neutrino case, as illustrated in Fig. 7. Even if there is no $`CP`$ violation, effects of the angles $`\theta _{02}`$ and $`\theta _{03}`$ could be seen in the $`\mathrm{sin}^2\mathrm{\Delta }_{\mathrm{LSND}}`$ term in Eq. (28), as illustrated in Fig. 8. Here, sensitivities to $`\theta _{02}`$ and $`\theta _{03}`$ tend to increase with decreasing $`L`$, due to the increased flux at shorter distances. Proof of the existence of a fourth neutrino does not exclude the possibility that there may be more than four neutrinos. Strictly speaking, an inconsistency between the measurement of $`\nu _\mu \nu _\tau `$ oscillations and the three-neutrino predictions would imply only that there are four or more neutrinos. In a model with four or more neutrinos, there are many more mixing angles and phases in the neutrino mixing matrix, and it would not be possible to determine them all from these oscillation measurements. However, when combined with results at other baselines, most of the four-neutrino parameter set could be determined . ## IV Summary In scenarios designed to describe the LSND oscillation results, our results show that short-baseline neutrino factory measurements can in principle determine all of the three-neutrino parameters provided the results are used together with future results from other accelerator-based experiments. The $`\nu _e\nu _\tau `$ channel is most sensitive to one combination of parameters in the $`CP`$-conserving terms. The $`\nu _e\nu _\mu `$ channel provides good sensitivity to $`CP`$ violation over a wide range of $`L`$ (20–100 km for $`E_\mu =20`$ GeV), assuming that a large muon detector is used, e.g. 10 kt. The $`\nu _\mu \nu _\tau `$ channel is also sensitive to $`CP`$ violation for a more restricted range of $`L`$, and may be used to explore whether more than three neutrinos exist. For four or more neutrinos, the $`CP`$-violating effect in $`\nu _\mu \nu _\tau `$ may be either enhanced or reduced by the additional mixing parameters. If MiniBooNE confirms the LSND oscillation results, then it will be important to measure the rates in all three appearance modes, as well as to search for $`CP`$ violation in $`\nu _e\nu _\mu `$ and $`\nu _\mu \nu _\tau `$, to obtain indirect evidence for the existence of sterile neutrinos. Neutral current measurements would complement the charge current studies of this paper in the search for sterile neutrinos. ###### Acknowledgements. VB thanks the Aspen Center for Physics for hospitality during the course of this work. This research was supported in part by the U.S. Department of Energy under Grant Nos. DE-FG02-94ER40817, DE-FG02-95ER40896 and DE-AC02-76CH03000, and in part by the University of Wisconsin Research Committee with funds granted by the Wisconsin Alumni Research Foundation.
warning/0007/astro-ph0007164.html	ar5iv	text	# Abundance gradients and their evolution in the Milky Way disk ## 1 Introduction Abundance gradients constitute one of the most important observational constraints for models of the evolution of the Milky Way disk. The existence of such gradients is now well established, through radio and optical observations of HII regions, stars and planetary nebulae (see Henry and Worthey 1999 for a review of the abundance profiles in the Milky Way as well as in external galaxies). An average gradient of dlog(X/H)/dR$``$0.06 dex kpc<sup>-1</sup> is observed in the Milky Way for O, S, Ne and Ar, while the gradients for C and N are, perhaps, slightly larger (Smartt 2000). Observations of open clusters support the above picture (Friel 1995, 1999; Carraro et al. 1998), albeit with large uncertainties. The existence of those gradients offers the opportunity to test theories of disk evolution and stellar nucleosynthesis. Indeed, the magnitude of the observed gradients in the Milky Way disk is rather large: metal abundances in the inner disk are larger than those in the outer disk by a factor of $``$10. This suggests that the role of the Galactic bar in inducing large scale radial mixing and therefore flatenning the abundance profile (e.g. Friedli and Benz 1995) has been rather limited; alternatively, the Galactic bar is too young ($`<`$1 Gyr) to have brought any important modifications to the gaseous and abundance profiles of the disk. In any case, the magnitude of the observed gradients puts important constraints on disk models invoking radial inflows (see e.g. Portinari and Chiosi 2000 and references therein). Even in the case of simpler models (with no radial inflows), the abundance gradients constrain the various parameters, like the local timescales of star formation and infall (e.g. Matteucci and Francois 1989, Prantzos and Aubert 1995, Mollà et al. 1997) or any variation of the stellar Initial Mass Function properties with metallicity (e.g. Chiappini et al. 2000). On the other hand, profiles of abundance ratios across the disk constrain the nature of various elements and isotopes (i.e. “secondaries” vs. “primaries”, “odd-Z” vs. “even-Z”) and their nucleosynthesis sites (e.g. Prantzos et al. 1996). A large number of chemical evolution models has been developped in the 90ies, aiming to explore one or more of the above issues, either with radial inflows (e.g. Götz and Köppen 1992, Chamcham and Tayler 1994, Tsujimoto et al. 1995, Firmani et al. 1996, Thon and Meusinger 1998, Portinari and Chiosi 2000) or without radial flows (Ferrini et al. 1994, Prantzos and Aubert 1995, Giovagnoli and Tosi 1995, Carigi 1996, Chiapinni et al. 1997, Boissier and Prantzos 1999, Chang et al. 1999). Some promising chemodynamical evolutionary models for the Milky Way have also been developed by a few groups (Steinmetz and Müller 1994, Samland et al. 1997, Berczik 1999). For the simplest of those models (namely those with no radial flows), some convergence has been recently reached between the various groups, at least concerning the basic ingredients (e.g. Tosi 2000): i) necessity of substantial infall and of radially varying timescales for the infall and the star formation, and ii) no need for varying IMF or strong galactic winds. Despite this agreement, some important differences still exist between the various models. The most important concerns the history of the abundance profiles: were they steeper or flatter in the past? The former is suggested by models of e.g. Prantzos and Aubert (1995), Mollà et al. (1997), Allen et al. (1998), Boissier and Prantzos (1999), while the latter is supported by models of Tosi (1988) and Chiappini et al. (1997). The situation is not settled observationally either. Estimated ages of open clusters and planetary nebulae of various types span a large fraction of the age of the Galaxy. Observations of the abundances of those objects across the Milky Way disk could, in principle, provide some information on the past history of the abundance gradients. In practice, however, observational uncertainties are too large at present to allow conclusions on that matter (e.g. Carraro et al. 1998, Maciel 1999). Several observational studies on the abundance profiles of the Milky Way disk were performed in the late 90ies. They concerned a variety of elements (He, C, N, O, Ne, Mg, Al, Si, S, Ar and Fe), observed in HII regions (Simpson et al. 1995, Rudolph et al. 1997, Afflerbach et al. 1997, Deharveng et al. 2000), B-stars (Smartt and Rolleston 1997, Gummersbach et al. 1998) and planetary nebulae (Maciel and Köppen 1994, Maciel and Quireza 1999). Although they leave the question of the history of abundance profiles unsettled, these studies allow for some important tests of stellar nucleosynthesis theories, since they provide absolute abundances, and abundance ratios along the Milky Way disk. In a previous work (Boissier and Prantzos 1999, herefater BP99) we developed a simple model for the chemical and spectrophotometric evolution of the Milky Way disk. With few assumptions (basically concerning the radial variation of the infall and star formation timescales), the model reproduces all the major obsevational constraints for the Galaxy: total amounts of gas and stars, supernovae rates, radial profiles of gas, stars, star formation rate and oxygen abundance, luminosities and scalelengths in various wavelength bands. The success of the model does not guarantee, of course, its correctness or its uniqueness. However, it offers a sound basis for a more detailed exploration of issues related to the chemical evolution of the Milky Way. Two such issues are explored here: i) how successful are currently available stellar yields in reproducing the wealth of recently observed abundance profiles in the Galaxy? are Wolf-Rayet (massive, mass losing) stars or intermediate mass stars (IMS) important in shaping the abundance profiles, and for which elements? ii) is the evolution of abundance gradients predicted by the model compatible with observations of available tracers? are there any new, potentially testable, predictions of the model in that respect? The plan of the paper is as follows: In Sect. 2, we present briefly the basic ingredients and the underlying assumptions of the BP99 model; in particulrar, we present in some detail the metallicity dependent yields of Woosley and Weaver (1995) and Maeder (1992), on which much of our work is based. Model results for abundance gradients of He, C, N, O, Ne, Mg, Al, Si, S, Ar and Fe and their evolution are given in Sect. 3. Sect. 4 constitutes the main body of the paper. In sect. 4.1 we present a compilation of observational data from HII regions, B-stars and planetary nebulae in the Milky Way disk. Detailed comparison between model predictions and current observed disk profiles is made in Sect. 4.2. In Sect. 4.3, we further compare the behaviour of model profiles of abundance ratios to observations and draw some conclusions on the adopted yields. Then, based on the abundance data for objects of different ages we discuss the evolution of abundance gradients for elements O, Ne, S and Ar (Sect. 4.4). The results are summarized in Sect. 5. ## 2 Model The adopted model for the chemical evolution of the Milky Way disk is described in detail in BP99. In Sect. 2.1 we briefly recall the main features of the model. In Sect. 2.2 we present in some more detail the only novel ingredient with respect to BP99, namely the metallicity dependent yields of Woosley and Weaver (1995, hereafter WW95) for intermediate mass elements and of Maeder (1992, hereafter M92) for He, C, N, O (in BP99 only yields for stars of solar metallicity are used). ### 2.1 Description of the model The galactic disk is considered as an ensemble of concentric, independently evolving rings, progressively built up by infall of primordial composition. The assumption of infall is traditionally based upon the need to explain the locally observed metallicity distribution of long-lived stars, which cannot be explained by the simple “closed-box” model (leading to the well-known “G-dwarf problem”). However, the recent work of Blitz et al. (1999) gives observational support to this idea, showing that the Milky Way and M31 are currently accreting substantial amounts of gas ($``$ 1 M/yr) in the form of high velocity clouds of low metallicity. The infall rate is assumed to be exponentially decreasing in time, i.e. $$f(t,R)=A(R)e^{t/\tau (R)}$$ (1) with a characteristic timescale $`\tau (R_0`$) = 7 Gyr in the solar neighborhood ($`R_0`$ = 8kpc), in order to reproduce the local G-dwarf metallicity distribution. $`\tau (R)`$ is assumed to increase outwards, from $`\tau `$(R = 2 kpc) = 1 Gyr to $`\tau `$(R = 17 kpc) = 12 Gyr. This radial dependence of the timescale of the infall rate $`f(R)`$ is simulating the inside-out formation of galactic disks and, combined with the adopted SFR $`\mathrm{\Psi }(R)`$ (Eq. 3) allows to reproduce the observed current profiles of gas, oxygen abundance and SFR in the Milky Way disk (see BP99 and Sects. 2.3 and 4 below). The coefficient $`A(R)`$ is obtained by the requirement that at time T = 13.5 Gyr the current mass profile of the disk $`\mathrm{\Sigma }(R)`$ is obtained, i.e. $$_0^Tf(t,R)=\mathrm{\Sigma }(R)$$ (2) with $`\mathrm{\Sigma }(R)e^{R/R_G}`$ and a scalelength $`R_G`$ = 2.6 kpc for the Milky Way disk. The chemical evolution of each zone is followed by solving the appropriate set of integro-differential equations, without the Instantaneous Recycling Approximation. The adopted stellar Initial Mass Function (IMF) is a multi-slope power-law between 0.1 M and 100 M from the work of Kroupa et al. (1993), leading to a Return Fraction R = 0.32. The star formation rate (SFR) is locally given by a Schmidt-type law, i.e. proportional to some power of the gas surface density $`\mathrm{\Sigma }_g`$: $`\mathrm{\Psi }\mathrm{\Sigma }_g^{1.5}`$, according to the observations of Kennicutt (1998). It varies with galactocentric radius $`R`$, as: $$\mathrm{\Psi }(t,R)=\alpha \mathrm{\Sigma }_g(t,R)^{1.5}V(R)R^1$$ (3) where $`V(R)`$ is the circular velocity at radius $`R`$. This radial dependence of the SFR is suggested by the theory of star formation induced by density waves in spiral galaxies (e.g. Wyse and Silk 1989). Since $`V(R)constant`$ in the largest part of the disk, this is equivalent to $`\mathrm{\Psi }(R)\mathrm{\Sigma }_g(R)^{1.5}R^1`$. The efficiency $`\alpha `$ of the SFR in Eq. 3 is fixed by the requirement that the local gas fraction $`\sigma _g(R_0)`$ 0.2, is reproduced at T = 13.5 Gyr. We assume that the “rings” of the disk are evolving independently from one another. This (over)simplification ignores in general the possibility of radial inflows in gaseous disks, resulting e.g. by viscosity or by infalling gas with specific angular momentum different from the one of the underlying disk; in both cases, the resulting redistribution of angular momentum leads to radial mass flows. The magnitude of the effect is difficult to evaluate, because of our poor understanding of viscosity and our ignorance of the kinematics of the infalling gas. Models with radial inflows have been explored in the past (Mayor and Vigroux 1981; Lacey and Fall 1985; Clarke 1989; Chamcham and Tayler 1994). It turns out that for some combinations of the parameters of infall, radial inflow and SFR, acceptable solutions are obtained, i.e. the current radial profiles of various quantities are successfully reproduced (see, e.g. Portinari and Chiosi 2000 for a recent overview of the problem). However, at the present stage of our knowledge introduction of radial inflows in the models would imply more free parameters than observables. For simplicity reasons we stick to the model of “independently evolving rings” for the Milky Way disk. ### 2.2 Yields of massive stars An important ingredient in our study of abundance gradients are the stellar yields of various elements. Most of the intermediate mass elements studied here are produced by massive stars, with the exception of some CNO isotopes that are also produced by intermediate mass stars. We consider no yields from intermediate mass stars in this work; in the line of Goswami and Prantzos (2000), concerning the evolution of the halo+local disk, our explicit purpose is to check to what extent massive stars can account for observations of intermediate mass elements and for which elements the contribution of intermediate mass stars is mandatory. We use the metallicity dependant yields of WW95, which are given for stars of mass M = 12, 13, 15, 18, 20, 22, 25, 30, 35 and 40 M and metallicities Z/Z= 0, 10<sup>-4</sup>, 10<sup>-2</sup>, 10<sup>-1</sup> and 1. In Fig. 1 we present the WW95 yields, folded with the Kroupa et al. (1993) IMF. They are presented as overproduction factors, i.e. the yields (ejected mass of a given element) are divided by the mass of that element initially present in the part of the star that is finally ejected: $$<F>=\frac{_{M1}^{M2}Y_i(M)\mathrm{\Phi }(M)𝑑M}{_{M1}^{M2}X_{,i}(MM_R)\mathrm{\Phi }(M)𝑑M}$$ (4) where: $`\mathrm{\Phi }(M)`$ is the IMF, $`M1`$ and $`M2`$ the lower and upper mass limits of the stellar models (12 M and 40 M, respectively). $`Y_i(M)`$ are the individual stellar yields and $`M_R`$ the mass of the stellar remnant. Adopting $`X_{,i}`$ in Eq. (4) creates a slight inconsistency with the definition of the overpoduction factor given above, but it allows to visualize the effects of metallicity in the yields of secondary and odd-Z elements. As can be seen from Fig. 1: i) most of the intermediate mass elements are nicely co-produced (within a factor of 2) by solar metallicity stars; ii) the “odd-even effect”, favoring the production of odd-nuclei at high metallicities, is clearly present; iii) the yields of Ne and Mg show, curiously, some dependence on metallicity (not as large as the one of the odd-elements Na and Al, but still enough to lead to some interesting abundance patterns, as we shall see in Sect. 3); iv) He, C, N, Sc, V and Ti are underproduced relative to Oxygen. He, C and N clearly require another source (intermediate mass stars and/or Wolf-Rayet stars, see Prantzos et al. 1994 and Sect. 4.2), while the situation is less clear for the elements, Sc, V and Ti (see Goswami and Prantzos 2000). The calculations of WW95 did not consider any mass loss during stellar evolution. Thus, they probably underestimated the yields of several elements that are expelled by the intense winds of massive stars, i.e. He, N and C. The effect of stellar winds is stronger when the stellar metallicity is larger. Maeder (1992) found that stars with $`M>`$30 M and Z$`>`$0.1 Z eject considerably larger amounts of He, N and C during their lifetime than their lower metallicity counterparts; for that reason, less matter is left in the He-core to be processed into Oxygen. At lower metallicities, the effect of stellar winds is negligible and stars of all masses evolve almost at constant mass. Elements heavier than Oxygen are produced in the subsequent, very rapid, stages of stellar evolution (after core He exhaustion) and their yields are not directly affected by the intensity of the mass loss. However, the structure of the stellar core may be affected by the loss of mass and this may also affect the final yields of heavy elements (e.g. Woosley, Langer and Weaver 1993). In Fig. 2 we present the yields of M92 for He, N, C and O as a function of stellar mass; they are given for two metallicities (Z/Z= 0.05 and 1, respectively) and are compared to the corresponding yields of WW95. The aforementionned effect of metallicity-dependent stellar winds on the yields of stars with M$`>`$30 M is clearly seen. In Sect. 3.4 we shall explore the effect of those yields on the abundance gradients in the disk. To account for the additional source of Fe-peak elements, required to explain the observed decline of O/Fe abundance ratio in the disk (e.g. Goswami and Prantzos 2000), we utilise the recent yields of SNIa from the exploding Chandrashekhar-mass CO white dwarf models W7 and W70 of Iwamoto et al. (1999). These are updated versions of the original W7 model of Thielemann et al. (1986), calculated for metallicities Z = Z(W7) and Z = 0 (W70), respectively. In this model, the deflagration is starting in the centre of an accreting white dwarf, burns $``$ half of the stellar material in Nuclear Statistical Equilibrium and produces $``$ 0.7 $`M_{}`$ of <sup>56</sup>Fe ( in the form of <sup>56</sup>Ni). These SNIa models lead to an oveproduction of Ni, but the evolution of this element will no be considered here. It should be emphasised that the evolution of the SNIa rate is not well known, and hardly constrained by observations. For the purpose of this work, we shall adopt the formalism of Matteucci and Greggio (1986) for the rate of SNIa, adjusting it as to have them appearing locally after the first Gyr, i.e. at a time when \[Fe/H\]$``$ $``$1 in the solar neighborhood. However, use of that same formalism along the disk will lead to different O/Fe abundance ratios at T = 13.5 Gyr as we shall see in Sect. 3, i.e. the final Fe gradient will be different from the one of oxygen (see also Prantzos and Aubert 1995). ### 2.3 Results for the Milky Way disk As described in detail in BP99, the simple model presented in Sect. 2.1 can readily account for the evolution of the solar neighborhood, reproducing quite successfully the main observational constraints (age-metallicity relationship, metallicity distribution of long-lived F-stars, current local surface densities of stars, gas, star formation and supernova rates). Also, in Goswami and Prantzos (2000) it is shown that the use of the WW95 yields for massive stars and of the Iwamoto et al. (1999) yields for SNIa leads to a successful agreement between the gaseous composition of the model at an age of 9 Gyr and the observed solar one. A few exceptions concern the elements He, C and N (which are underproduced) and Ni (which is overproduced, because of the adopted SNIa yields). The model predictions for the disk are equally successful, at least to a first order. Indeed, the adopted combination of SFR (Eq. 3) and infall rate (Eq. 1) leads to final profiles of gas and SFR that are in fair agreeement with the observed ones, as can be seen in Fig. 3. The stellar profile is also in agreement with observations, but it is essentially determined by the boundary condition of Eq. 2. However, the adopted inside-out formation scheme of the disk leads naturally to different scalelengths in the B-band (reflecting mostly the SFR profile in the past $``$ 1 Gyr) and the K-band (reflecting the total stellar population, cumulated over T = 13.5 Gyr). As shown in BP99, the corresponding scale-lengths ($``$ 4 kpc in the B-band and $``$ 2.6 kpc in the K-band, respectively) are in fair agreement with observations. Moreover, the model also reproduces reasonably well the total current SFR and supernova rates as well as the total luminosities in various wavelenght bands. This is a rather encouraging success, since the number of the constraints is much larger than the number of the parameters. The success of this simple model encourages us to use it for a thorough study of the various abundance gradients in the Milky Way disk. ## 3 Model results for abundance gradients We calculated the evolution of the abundances of all elements between H and Zn with the WW95 metallicity dependent yields in all the zones of our model disk. In Fig. 4 we present the results concerning all elements with measured abundance gradients in the Milky Way. We notice that, since we did not include yields from IMS in our calculation, our results for He, C and N represent rather lower limits. We shall discuss the comparison to observations in the next section, where we shall also explore the role of massive, mass losing stars to the abundance profiles of those elements. Here we focus on the model results, which may be summarised as follows: i) Final values (at T = 13.5 Gyr) of the abundances at R<sub>0</sub> = 8 kpc are $``$ solar for O, Al, Si, S, Ar and Fe; they are slightly lower than solar for Ne and Mg; and they are considerably below solar for C and N. It may appear surprising that Ne and Mg do not reach their solar values, but this is a consequence of the adopted metallicity dependent yields of WW95 (Fig. 1): the average (i.e. over the disk’s age) overproduction factors of Ne and Mg are lower than the one of e.g. O. It is not clear whether such a dependence is physical or just an artifact of the WW95 models. As for C and N, it is clear that massive stars with no mass loss cannot be the main sources of those elements (see also Sect. 4). ii) The most prominent feature of the model is the prediction that abundance gradients flatten with time. This is a generic feature of all models forming the galactic disk “inside-out”. Indeed, in that case, there is a rapid increase of the metal abundance at early times in the inner disk, leading to a steep abundance gradient. As time goes on, star formation “migrates” to the outer disk, producing metals there and flattening the abundance gradient. iii) The final abundance profile (T = 13.5 Gyr) is, in general, flatter in the inner disk. As already described in Prantzos and Aubert (1995) this is due to the fact that in those regions the large populations of low-mass, long-lived stars that are formed early on in galactic history reject a lot of metal-poor gas at the end of their evolution, which dilutes the metal abundances; this effect is absent in the outer regions, where there are not very old stellar populations. Notice that the flattening seems to be more important in the case of N, Ne, Mg and Al. These elements show some metallicity dependence in their yields (at least according to WW95) which is difficult to understand in the case of Ne and Mg, but expected in the case of the secondary N and of the odd-Z Al. Since WW95 give yields only up to stellar metallicities of Z = Z, for higher metallicities we use their Z = Zyields. This means that in the inner disk, where metallicities higher than solar are reached, we underestimate the production of any element with metallicity dependent yield. If appropriate yields were used, the Al and N profiles in the inner disk would be steeper, not flatter, than the one of O. The magnitude of the current abundance gradients in the Milky Way disk is one of the most important constraints in the models of the evolution of our Galaxy. Most of the proposed models reproduce it fairly well (e.g. Tosi 2000 and references therein), at least when no radial inflows are included. However, equally important is the question of the evolution of those gradients and, in particular, whether they flatten or steepen with time. We shall confront our models to the data and to other theoretical works in the next section. Here, we present the evolution of our model gradients in the 4$``$14 kpc region for a few selected elements (Fig. 5). All gradients were systematically larger in the past. The gradient of Fe is slightly larger than the one of O, because our adopted prescription for the SNIa rate produces a smaller O/Fe ratio in the inner disk than in the outer one. The gradient of secondary N is always steeper than the one of O, but since we do not include N production from intermediate mass stars or WR stars in this calculation, this result serves merely for illustration purposes. ## 4 Comparison to observations In this section we compare our results to the observed abundance profiles of various elements across the galactic disk. Tracers of abundance profiles include emission-line objects (HII regions and planetary nebulae), as well as stars and stellar associations (B-stars and open clusters, respectively). In most cases, those tracers are relatively young objects (HII regions, B-stars and planetary nebulae of type I), younger than $``$ 1 Gyr; they provide then information about the current abundance profile of the Milky Way disk. In other cases (planetary nebulae of type II or III, open clusters), the objects involved are several Gyr old and provide information about the past status of the disk (albeit with much larger uncertainties than in the former case). In Sect. 4.1 we present briefly the available observational data from various sources. In Sect. 4.2 we compare the data for young objects to our results at T = 13.5 Gyr; we show, in particular, how the results for C (and to a much lesser extent, N and O) may be affected by the M92 metallicity-dependent yields of massive, mass losing stars. In Sect. 4.3 we discuss the profiles of the corresponding abundance ratios of elements to Oxygen. In Sect. 4.4 we discuss the history of the abundance profiles, comparing our results to observations of old objects (of rather uncertain ages). ### 4.1 Observational data One of the first comprehensive optical surveys of HII region abundances was performed by Shaver et al. (1983). They have found strong abundance gradients for N/H, O/H and Ar/H in the range of galactocentric distances R<sub>G</sub> = 5 $``$ 12 kpc. Those results were later confirmed by the surveys of Fich & Silkey(1991) and Vílchez & Esteban (1996), extending the data up to distances of R$`{}_{G}{}^{}`$ 17 kpc. Those authors suggested that the N/H and O/H gradients show a tendency for flatenning between 11 and 18 kpc. Taking advantage of the small extinction in the infrared, Simpson et al. (1995) and Afflerbach et al. (1997) studied objects towards the center of the Galaxy. Simpson et al. (1995) observed 12 HII regions in R<sub>G</sub> = 0 $``$ 10 kpc and found a somewhat better fit with a step function than with a smoothly decreasing one. On the contrary, after adding five outer Galaxy HII regions (R<sub>G</sub> = 13 $``$ 17 kpc), Rudolph et al. (1997) concluded that the single slope profiles for N/H and S/H are more likely than step functions. The existence of a gradient in the oxygen abundance profile obtained from HII regions is not in doubt at present, but its magnitude has been recently challenged by the work of Deharveng et al. (2000), who report a value about half as large as previous measurements. Observations of B-type stars in young clusters and associations in the late 90ies confirmed the existence of an abundance gradient. Smartt & Rolleston (1997) found an oxygen abundance gradient similar to the one obtained in HII regions (except for the work of Deharveng et al. 2000), in the galactocentric distance range R<sub>G</sub> = 6 $``$ 18 kpc. This was confirmed by Gummersbach et al. (1998) who detected significant abundance gradients for C, O, Mg, Al and Si in galactocentric distances R<sub>G</sub> = 5 $``$ 14 kpc. Several works have been devoted to the study of the abundance patterns of planetary nebulae (PN, see e.g. Maciel and Quiroza 1999 and references therein). According to a classification scheme originally suggested in Peimbert (1978) and revised in Pasquali and Perinotto (1993), planetary nebulae are subdivided into types I, II and III. PNI have He/H $`>`$ 0.125 and log(N/O) $`>`$ $``$0.3, are located in the galactic thin disk and their progenitors are probably associated with stars in the mass range 2.5 $``$ 8 M. PNII belong also to the thin disk, are not particularly enriched in nitrogen or He and are thought to evolve from stars in the 1.2 $``$ 2.5 M range. Finally, PNIII are associated with the galactic thick disk and thought to originate from low mass progenitors, in the 1 $``$ 1.2 M range. The estimated progenitor masses imply ages of $`<`$1 Gyr for PNI, 1 $``$ 8 Gyr for PNII and $`>`$8 Gyr for PNIII, respectively. These differences should allow, in principle, to use PN as tracers of the chemical evolution of the galactic disk. However, the mass and age differences between the various PN types are not quite well defined, as we shall see in Sect. 4.3. Moreover, all PN are expected to be auto-enriched in products of the CN cycle (i.e. N-rich and C-poor), while PNI (especially those originating from relatively massive progenitors, in the 5 $``$ 8 M range) are probably auto-enriched in products of the NO cycle (N-rich and O-poor). For those reasons, we shall not consider at all in the following the abundances of C and N from PN of all types and the abundances of O from PNI. For the purpose of this work, we shall use PNI (along with HII regions and B-stars) as tracers of the young disk in Sects. 4.2 and 4.3. PNII and PNIII will be discussed in Sect. 4.4. Table 1 summarizes the currently available observational data on the abundance profiles across the Milky Way disk. Two points should be noticed: 1) Contrary to other authors, we did not consider data on Fe from open clusters, since the situation is rather uncertain at present: for instance, in the solar neighborhood, open cluster data show no evolution of the Fe abundance with age (e.g. Friel 1999), contrary to the familiar age-metallicity relation suggested by data for F-stars (Edvardsson et al. 1993). In view of that dicrepancy, we chose not to consider Fe in this work (although we do calculate its abundance profile, as shown in Fig. 4). 2) Data in Table 1 are adopted directly from the references listed and no attempt has been made to homogenize them by recalculating the abundances in a consistent way. In general, differences in techniques and atomic data employed produce an abundance scatter which is smaller than observational uncertainties in the line strengths. Therefore, we believe that direct comparison between our models and those inhomogeneous data can still provide statistically meaningful results. ### 4.2 Current disk abundance profiles The observational data of Table 1 concerning “young” objects (HII regions, B-stars, PNI) are plotted in Fig. 6, along with our model results, obtained with the WW95 yields, at T = 13.5 Gyr. In Fig. 7 we also plot the same data for He, C, N and O and we show the corresponding results obtained with both WW95 and M92 yields; as discussed in Sect. 2.2, the metallicity dependenence of massive star winds has an effect on the yields of those elements. Since the prescriptions for the radial dependence for infall and SFR of the BP99 model adopted here were such as to reproduce the observed oxygen abundance profile, we start the description of our results with this element. Oxygen: As can be seen from Table 1, data from both HII regions and B-stars suggest an abundance gradient dlog(O/H)/dR $``$ $``$0.07 dex kpc<sup>-1</sup> (see the recent review by Smarrt 2000 on B-stars). Maciel and Quiroza (1999) include PNII (despite the fact that these objects are, in principle, older than 1 Gyr) and find an average gradient of dlog(O/H)/dR $``$ $``$0.065 dex kpc<sup>-1</sup>for all tracers of the “young” population. However, Deharveng et al. (2000), after a consistent analysis of their own data on HII regions, as well as of those of previous works, conclude that the gradient should be $``$ 40% smaller than generally thought, i.e. $``$0.04 dex kpc<sup>-1</sup>. The data of Deharveng et al. (2000) concern the galactocentric distance R<sub>G</sub> = 5 $``$ 15 kpc, i.e. the same as the one of the work of Shaver et al. (1983) on HII regions or studies on B-stars; the difference in the resulting abundance gradient cannot then be attributed in the studied galactocentric distance range. As stressed by Deharveng et al. (2000), their derived O/H abundances depend strongly on their two-temperature HII region model (one temperature for the high excitation O<sup>++</sup> zone and another temperature for the low excitation O<sup>+</sup> zone); an alternative HII region model (with temperature fluctuations) would lead to larger O/H abundances. As explained in Sect. 2, our model is based on the radial variation of the SFR (Eq. 3) and infall rate (Eq. 1). A different prescription would lead to different results for the abundance gradients (as in e.g. Prantzos and Aubert 1995, who used a radially independent infall time-scale and found a smaller gradient of dlog(O/H)/dR $``$ $``$0.03 dex kpc<sup>-1</sup>). Also, by adopting the same prescriptions, but assuming radial inflows, one would obtain smaller values for the abundance gradient (e.g. Tsujimoto et al. 1995, Portinari and Chiosi 2000). Obviously, the magnitude of the abundance gradient is crucial in fixing the parameters entering the current phenomenological chemical evolution models. In particular, a strong gradient does not support the idea of important radial inflows induced by a galactic bar (such as those found e.g. in the calculations of Friedli and Benz 1995); the Milky Way bar must then have formed relatively recently and/or played a negligible role in driving gaseous flows in the disk. The only published successful chemodynamical model for the Milky Way (Samland et al. 1997) predicts an oxygen abundance gradient that flattens considerably in the outer disk. Despite some claims for such a flatenning (Vilchez and Esteban 1996), no convincing observational evidence exists up to now (see the discussion in Deharveng et al. 2000). Our model predicts no flatenning of the gradient in the outer disk. It does predict a small flatenning in the inner disk. As explained in Prantzos and Aubert (1995), this is due to the fact that all metal abundances are diluted in the inner disk by the late ejection of metal-free material by the numerous low-mass, long-lived stars of the first stellar generations; the ejection rate of that material at late times is larger than the metal production rate, since most of the gas in the inner disk has been consumed at that time. This is not the case in the outer disk, which is formed late. As can be seen in Fig. 7, the M92 yields lead to results for oxygen that are almost indistinguishable from those of WW95. Only in the inner (more metal-rich) zones a slight difference is obtained in the oxygen abundance profile. As already discussed in Prantzos et al. (1994) the metallicity dependent stellar winds have a negligible impact on the oxygen yields. Finally, another aspect of the observed oxygen abundances (Fig. 6) has been discussed in several places: all “young” objects in the local disk have lower oxygen abundances (in the range O/H $``$ 3.2$``$5 10<sup>-4</sup>) than the Sun ((O/H) = 6.8$`{}_{0.9}{}^{}{}_{}{}^{+1.0}`$ 10<sup>-4</sup>, Grevesse & Sauval 1998). For instance, Gies and Lambert (1992) derive an oxygen abundace O/H$``$ 4.8 10<sup>-4</sup> for local B stars, while Esteban et al.(1998) give O/H = 5.2$`{}_{0.7}{}^{}{}_{}{}^{+0.9}`$ 10<sup>-4</sup> for Orion nebula. This is difficult to understand in the framework of conventional models of chemical evolution, where metallicity increases monotonically with time. The idea that the Sun was born in the metal-rich inner disk and subsequently migrated outwards (Wielen et al. 1996), has been recently rejected by Binney and Sellwood (2000), on the grounds of dynamical arguments. Thus, at present, there is no satisfactory explanation for the “super-metallicity” of the Sun relative to young objects in the solar neighborhood. Carbon: There are very few studies of carbon abundances, either in HII regions or B-stars. Based on only two nearby HII regions (Orion and M17), Peimbert et al.(1992) first derived an abundance gradient of dlog(C/H)/dR $``$ $``$0.08$`\pm `$0.02 dex kpc<sup>-1</sup>in the solar vicinity. Adding another HII region (M8) to that sample, Esteban et al. (1999) obtain $``$0.133$`\pm `$0.022 dex kpc<sup>-1</sup>. However, both those studies concern a very limited range of galactocentric distances (6 $``$ 9 kpc), and the derived abundance gradients cannot be considered as representative of the disk as a whole. On the other hand, the study of Gummersbach et al. (1998) suggests a rather small C gradient for B-stars in the range R<sub>G</sub> = 6 $``$ 12 kpc: dlog(C/H)/dR $``$ $``$0.045 dex kpc<sup>-1</sup>; this is half as steep as the oxygen abundance gradient derived for those same B-stars in that study. If the Gummersbach et al. (1998) value is correct, it should then be difficult to understand the corresponding disk profile of the C/O ratio, which should decrease in the inner disk. Indeed, observations of low mass stars in the solar neighborhood show that C/O increases with metallicity (e.g. Gustafsson et al. 1999 and references therein), and we would expect to see the same effect in the inner, metal-rich disk. Hibbins et al. (1998) performed a differential analysis of C and N abundances of B-stars in the outer disk (10 $``$ 17 kpc). They found that C correlates with N, but not with O. However, their absolute values of C/H for nearby stars are systematically lower by a factor 2 $``$ 3 than those of other studies; they suggest then that their data (and data coming from different samples, in general) should be used with caution in studies of galactic chemical evolution. In summary, the carbon abundance profile in the Milky Way disk is poorly determined. The few available data sets cover limited galactocentric distance ranges and are not derived in a consistent way. If all the data are plotted, as in Fig. 6 and 7, a carbon abundance gradient larger than the one of oxygen appears; but this may well be an artifact. We note that in his recent review, Smartt (2000) suggests a gradient of dlog(C/H)/dR $``$ $``$0.07 dex kpc<sup>-1</sup>, based on unpublished work of Rolleston et al. (2000) on B-stars. Our model results (Fig. 6) show a carbon abundance gradient of dlog(C/H)/dR $``$ $``$0.06 dex kpc<sup>-1</sup>, i.e. comparable to the one of oxygen, when the WW95 yields are used. However, the absolute value of the local carbon abundance at the Sun’s birth is underproduced in that case (by a factor of $``$2); another C source is then required. Intermediate mass stars (IMS) are well known net producers of carbon, but in recent years, several works (Prantzos et al. 1994, Carigi 1994, Gustaffson et al. 1999, Henry et al. 2000) suggested that the M92 metallicity dependent yields of massive stars fit better the observed C/O abundance patterns in extragalactic HII regions and low-mass stars in the solar neighborhood. This is confirmed in Fig. 7, where it can be seen that the use of the M92 yields allows indeed to reproduce the absolute abundance of carbon in the solar neighborhood, with no need for a contribution by IMS. In that case, the resulting C abundance gradient is much steeper: $``$0.086 dex kpc<sup>-1</sup>. Clearly, the abundance profiles of carbon vs. oxygen in the Milky Way disk are crucial in any attempt to evaluate the role of massive star winds in the production of carbon. Nitrogen: There is a large number of works concerning N abundances in HII regions; the N abundance profile seems to be steeper, in general, than the corresponding O profile (see Table 1). Observations of B-stars (Gummersbach et al. 1998) support this conclusion, but the error bars are much larger in that case. Smarrt (2000) suggests a gradient of dlog(N/H)/dR $``$ $``$0.08$`\pm `$0.002, i.e. compatible with the one of oxygen. Our models also produce a steeper gradient of N (with respect to that of O), with the yields of both WW95 (Fig. 6) and M92 (Fig. 7). As can be seen in Fig. 7, the differences resulting from the use of those two sets of yields are rather small (slightly larger than in the case of oxygen, but considerably smaller than in the case of carbon). Two important points should be noticed: i) The absolute value of the N abundance obtained with both sets of yields is too low compared with observations in the solar neighborhood and in the disk. Contrary to the case of carbon, the M92 yields are not sufficient to account for the current abundance of N in the Milky Way. Unless the N yields are seriously underestimated in both WW95 and M92 (which is improbable, since they consist of the sum of the initial C+N+O), another N source is required. Intermediate mass stars are the obvious candidate, but their yields are notoriously difficult to estimate (in view of the many uncertainties concerning mass loss rates, “hot-bottom burning”, etc.). Our calculations clearly show what is the magnitude of the expected contribution of IMS, in order to complement the (presumably better understood) yields of massive stars: on average, IMS have to produce 3-4 times more N than massive stars. ii) We stress again (see also Sect. 3) that the flatenning of the N abundance profile obtained in our calculations for the inner disk is due to the fact that, since yields of WW95 and M92 are available only up to Z = Z, for higher metallicities we use the yields at Z = Z; thus, in our models N is produced in the inner disk (where high metallicities are developed) with the same yield always, i.e. as a primary element and its abundance follows the one of oxygen. Clearly, this is an artifact of the calculation, due to the lack of appropriate input data; if N were treated correctly, as a secondary, its abundance profile would not flatten in the inner disk \[Notice: we could simply scale the N yields of WW95 with metallicity, but it would be more difficult to do the same thing with the M92 yields, since the effect of the stellar wind on the yield is not simply proportional to metallicity\]. Magnesium, Aluminium, Silicon: The abundance profiles of those elements have been studied by Gummesrbach et al. (1998), who observed B-stars in galactocentric distances 5 $``$ 14 kpc. In view of the large uncertainties (see Table 1), all those gradients can be considered as compatible with the oxygen gradient of $``$0.07 dex kpc<sup>-1</sup>, at the 1 $`\sigma `$ level; this is particularly true for Mg. On the other hand, taken at face value, the Al profile seems to be considerably steeper than the one of Si ($``$0.045 dex kpc<sup>-1</sup> vs. $``$0.107 dex kpc<sup>-1</sup>, respectively). This is rather surprising, taking into account that Al is an odd-Z element and its yield depends slightly on metallicity (the “odd-even” effect, see Fig. 1). Because of this effect, our models produce a steeper gradient for Al than for Si. The latter is comparable to the one of oxygen, as expected. We note that in his recent review Smarrt (2000) suggests gradients of similar magnitude for Al and Si ($``$0.05$`\pm `$0.02 dex kpc<sup>-1</sup> and $``$0.06$`\pm `$0.01 dex kpc<sup>-1</sup>, respectively), based on B-star data. In that case the problem is alleviated, although not completely solved. We also note that low-mass stars in the solar neighborhood do not exhibit the theoretically expected behaviour of the Al/O abundance pattern (e.g. Goswami and Prantzos 2000 and references therein). It may well be then that the odd-even effect has been overestimated in the WW95 yields of Al. In summary, the observed abundance profiles of Al and Si run opposite to theoretical expectations, but in view of their large uncertainties, it is difficult to draw firm conclusions on stellar nucleosynthesis. Neon, Sulphur, Argon: These elements have been observed through their emission lines in both HII regions and PN of all types. In general, HII region abundances suggest a radial profile similar to the one of oxygen for all three elements, i.e. a gradient of $``$0.07 dex kpc<sup>-1</sup> (the flat S profile of Shaver et al. 1983 is an exception to this general agreement). Data from PNI and PNII support gradients of this magnitude also for S and Ar, but suggest a smaller gradient for Ne ($``$0.036 dex kpc<sup>-1</sup>, according to Maciel and Quiroza 1999). Our model leads to S and Ar abundance gradients similar to the one of O ($``$0.05 dex kpc<sup>-1</sup> in both cases). In the case of Ne, the (unexcpected) small metallicity dependence of the WW95 yields (Fig. 1) leads to a steeper profile: dlog(Ne/H)/dR $``$ $``$0.08 dex kpc<sup>-1</sup>, exactly as in the case of Mg. In summary, the observed S and Ar abundance profiles are in agreement with theoretical expectations, while Ne remains problematic, both theoretically and observationally. Helium: For the sake of completeness, we present here the He profiles of our models. As can be seen in Fig. 7, there is no difference between the profiles obtained with the WW95 and M92 yields. In both cases a flat He/H profile is obtained. However, the lack of IMS yields in our model does not allow to draw conclusions about He. The observational situation concerning the He abundance profile is not clear. The flat profile of the He<sup>+</sup>/H<sup>+</sup> ratio found by Shaver et al. (1983) was confirmed by the recent work of Deharveng et al. (2000). However, the true He/H gradient depends also on the unknown amount of the corresponding neutral species. The data on Fig. 6 and 7 are from the work of Gummersbach et al. (1998) on B-stars. The error bars are quite large ($`\pm `$0.3 dex for log(He/H)) and do not allow to draw any conclusion on the existence of an abundance gradient for He. Fig. 8 summarizes the discussion of this section, concerning observed and calculated abundance profiles in the Milky Way disk. If the values of Smartt (2000) are adopted for the C and Si gradients from B-stars, then there is satisfactory agreement between the model and observations for all elements but Al (taking into account error bars); He and N require another source to account for their absolute abundances, even if the gradient has the correct value. But the main “message” of Fig. 8 is that homogeneous data sets are required for a meaningful comparison between theory and observations. ### 4.3 Abundance ratios along the Galactic disk Abundance ratios between metals are, in principle, more reliable tracers of the chemical evolution than absolute abundances, since they do not depend on the star formation efficiency and they allow to identify if e.g. an element has a secondary origin, whether it is produced in long-lived sources (low-mass stars or SNIa) or whether it is affected by the “odd-even” effect, etc. These properties have been widely used as diagnostic tools of the history of the halo + solar neighborhood system (se e.g. Pagel 1997). In practice, the situation is complicated due to various theoretical and observational uncertainties (e.g. Goswami and Prantzos 2000 and references therein), except for a few trivial cases. This is also the case for the Milky Way disk, as can be seen in Fig. 9: of the nine element abundance ratios to oxygen displayed in the figure, not a single one shows any clear trend with galactocentric distance. Although in most cases this is expected on theoretical grounds (e.g. for Ne, Mg, Si, S, Ar, which are primary elements), for others (N and, to a smaller extent, He and Al) this is rather surprising. Even worse is the large scatter obtained for all elements and at all galactocentric distances. There is a striking difference with F stars in the solar neighborhood, which display very small dispersion in their abundance ratios (Edvardsson et al. 1993). The inhomogeneity of the data sets displayed in Fig. 9 contributes certainly to that scatter. However, even in cases where only one tracer is involved (B-stars for Mg, Al and Si), the dispersion is rather large: it is about twice as large as the typical uncertainty of individual abundance ratios, estimated here to be $`\pm `$0.2 dex (see also Henry and Worthey 1999). A meaningful comparison between theory and observations is difficult in such conditions. However, Fig. 9 allows to draw some conclusions: i) The theoretical abundance ratios of Si/O, S/O and Ar/O are $``$ constant across the disk, as expected for primaries; their absolute value is always $``$ solar and compatible with available observations. The WW95 yields reproduce well the current average abundance ratios of those elements. ii) The theoretically expected “odd-even” effect for Al/O is not manifested in the available observational data, i.e. no trend of Al/O with galactocentric distance (or metallicity) is observed. The WW95 yields reproduce correctly the current local Al/O ratio ($``$ solar), and lead to a small dependence on metallicity (or radius), because of the “odd-even” effect. iii) The WW95 yields underproduce the observed Ne/O and Mg/O ratios (as anticipated from Fig. 1); this is also the case for the solar neighborhood, as discussed in Goswami and Prantzos (2000). The discrepancy is not large, taken into account the various uncertainties of stellar nucleosynthesis calculations (see Prantzos 2000 for a review) as well as in observational data. Our Ne/O data come mainly from PN I, where this ratio is usually determined under the assumption that Ne/O = Ne<sup>++</sup>/O<sup>++</sup>. The Ne<sup>++</sup>/O<sup>++</sup> values probably are upper limits to the Ne/O value for central stars hotter than 50,000 degrees, due to the presence of charge exchange reaction O<sup>++</sup> \+ H<sup>0</sup> $``$ O<sup>+</sup> \+ H<sup>+</sup> that allows some O<sup>+</sup>to coexist with Ne<sup>++</sup>; this is particularly the case for PNs of type I with low density (Peimbert et al. 1995). Therefore the best values from type I PNs are given by the lower envelope of the data presented in Fig. 9 (M. Peimbert, private communication). On the other hand, the WW95 yields of Ne and Mg manifest an unexplained dependence on metallicity. Put together, these discrepancies point to some problems in the WW95 yields of Mg and Al. One possibility is that the extent of the C-shell (where both elements are mainly synthesized) is underestimated by the Ledoux criterion for convection employed in the WW95 calculations. iv) The WW95 yields underproduce the local C/O ratio, while the M92 yields of massive stars reproduce it fairly well (with no need for a C contribution by IMS). However, if WR stars are the main producers of C, one expects this ratio to decrease with galactocentric distance (as indicated by the thick curve for C/O in Fig. 9), and this trend is not seen in currently available data. On the other hand, if IMS stars produce the bulk of C as primary, no variation of the C/O ratio with galactocentric distance is expected. It should be noted that gradients of C/O (in fact, variations of C/O with O/H) have been observed in extragalactic HII regions (e.g. Garnett et al. 1999) and are also expected in the case of the Milky Way disk. Thus, accurate determination of the C/O ratio across the Milky Way disk is crucial in determining the roles of IMS and Wolf-Rayet stars in carbon production. \[Note: the flattening of the theoretical C/O profile in the inner disk is due to our use of M92 yields for Z = Zeven when the metallicity is larger than solar; in principle, the C/O ratio should continue increasing in the inner disk. C and O yields for metallicities higher than solar are required for a proper calculation\]. v) Both WW95 and M92 yields underproduce the N/O ratio in the disk and lead to a N/O profile declining with radius. IMS stars are expected to be the main N producers in the disk. The dispersion of presently available data on N/O in the disk does not allow to conclude whether they produce N as primary (i.e. through hot-bottom burning), or as a secondary. As with the case of C, observations of N/O in extragalactic HII regions show that N behaves as secondary at high metallicities (Henry et al. 2000). This does not seem to be the case in the inner regions of the Milky Way disk, at least with currently available data. vi) The observed He/O ratio in B-stars is higher than the corresponding solar value at all galactocentric distances. This is obviously related to the fact that the O abundances of those young objects are (somewhat surprisingly) systematically lower than solar. Since no satisfactory explanation for that exists up to now (at least in the framework of conventional chemical evolution models), we do not expect our results to match the average observed value of He/O. On the other hand, the B-star data show a flat He/O profile (as a result of the similar He/H and O/H gradients shown on Fig. 6 and 7). In our calculations, the amount of He produced by massive stars alone (either with the WW95 or the M92 yields) is negligible relative to the primordial one of the infalling gas. Thus, the final He/H profile is flat (Fig. 6 and 7) and the final He/O profile (Fig. 9) is decreasing in the inner disk; obviously, the contribution of intermediate mass stars is mandatory to account for He observations. ### 4.4 Evolution of abundance gradients As already mentioned in Sect. 4.1, abundances in planetary nebulae of various types (I, II and III) allow, in principle, to follow the history of the abundance profile of the Milky Way disk. However, the progenitor masses and lifetimes of PNs are not well known. For instance, Allen et al. (1998) assume that both PNII and PNIII have progenitors of similar mass range (M$`<`$2 M) and classify them in terms of their kinematical properties: PNII have peculiar velocities $`v_P<`$ 60 km/s and PNIII have $`v_P>`$ 60 km/s. To account for these properties as well as for the relative frequency of the various PN types, Allen et al. (1998) propose a “dynamical” model for the Milky Way disk, simulating orbital diffusion. They find good agreement with observed gradients of PN (classified according to that scheme), provided that not all IMS go through the PN stage. In that same paper, they also explore the “conventional” scheme, classifying PN in terms of progenitor mass: 1.3$`<`$M/M$`<`$8.4 and $`\tau <`$ 3 Gyr for PNI, 0.9$`<`$M/M$`<`$1.3 and 3$`<`$ $`\tau <`$ 9 Gyr for PNII and 0.8$`<`$M/M$`<`$0.9 and $`\tau >`$ 9 Gyr for PNIII. They find that their model fails to reproduce observations of the various PN types in that case. In our model we have no information on the kinematical properties of the gas or the stars. We adopt then the classification scheme of Pasquali and Perinoto (1993) for PN, already presented in Sect. 4.1. We are aware that this scheme is not necessarily the most appropriate one (uncertainties in ages may be larger than stated and kinematics may play a non negligible role), but we adopt it for the sake of comparison to our model results. This comparison is presented in Fig. 10. It can be seen that: i) The model describes relatively well the O, Ne, S and Ar abundance gradients of young objects (left hand panels, column A), as already discussed in the previous sections. However, the observed scatter at all galactocentric distances is much larger than the range of values given by the model for the last 1 Gyr (the shaded region in Fig. 10). It should be noted that the observed scatter ($``$ 0.6 dex) is larger than the corresponding one for young stars in the solar neighborhood (e.g. Garnett and Kobulnicky 2000). Probably, systematic effects and uncertainties in distance estimates contribute largely to the observed scatter of “young” tracers in Fig. 10. ii) Model predictions are compatible with observations for O and Ar in PNII, as can be seen in the middle panels of Fig. 10 (column B). In fact, there is very little difference between the data for PNII and those for PNI. Predictions for S are slightly above the data points, while predictions for Ne are clearly below the corresponding data. The reason for the latter discrepancy is obviously the unexplained underproduction of Ne in the low metallicity yields of WW95, already discussed in Sects. 2.2 and 4.3. In the other 3 cases (O, S and Ar) there is fairly good agreement between theory and observations concerning all properties of the observed abundance profiles (absolute values, gradient, scatter). iii) Observations show that, on average, abundances in PNIII are lower and present a larger dispersion than in PNII. Both features are relatively well reproduced by our models, as can be seen on the right hand panels in Fig. 10 (column C), with the exception of Ne (for the reasons already mentioned in the previous paragraph). Our models, as some other models of this kind (Mollá et al. 1997, Portinari and Chiosi 1999) suggest that abundance gradients are steeper for older objects. Other works (e.g. Chiappini et al. 1999) reach an opposite conclusion, although they are based on the same assumption, namely an “inside-out” formation of the disk; clearly, the adopted time-scales for star formation and infall play an important role in this discrepancy. Although our model results are compatible with available data on PN, the large scatter in those data does not allow to conclude on the temporal variation of the gradients (see Maciel and Quireza 1999). However, our model makes a new, and perhaps testable, prediction: if the observed abundance scatter at a given galactocentric distance is to be attributed, at least partially, to intrinsic age differences between the tracers (e.g. PN) then we expect that scatter to be smaller in the inner disk than in the outer one (as indicated by the extent of the shaded areas in Fig. 10). ## 5 Summary Based on the successful chemical evolution model for the Milky Way disk developped by BP99, we calculated the corresponding abundance profile of elements up to the Fe-peak. For that purpose we used the metallicity dependent yields of WW95 for massive stars evolving at constant mass, along with those of Iwamoto et al. (1999) for SNIa. We also explored the yields of M92 for massive mass losing stars, concerning the elements He, C, N and O. We deliberately ignored yields from intermediate mass stars (IMS), our purpose being to check to what extent such stars are indeed required to explain observed abundance patterns. However, we did take into account ejecta from IMS, assuming that their net yield is zero for all elements; as explained in Goswami and Prantzos (2000), this procedure is crucial to ensure that the absolute metal abundances, expressed as X/H, are correctly evaluated. We compiled a large sample of observational data, concerning abundance profiles of He, C, N, O, Ne, Mg, Al, Si, S and Ar in the Milky Way disk (Table 1). Most of the tracers are young objects (B-stars, HII regions, PNI), while PNII and PNIII abundances trace earlier stages of the Galaxy evolution (but with considerable uncertainties in the corresponding ages). Contrary to other authors, we did not consider data on Fe from open clusters, since the situation is rather uncertain at present (see Sect. 4.1). We simply note that our prescription for the rate of SNIa (major Fe producers), leads to a Fe abundance profile slightly steeper than the one of O in the disk (Fig. 4 and 5). The main results of the model may be summarized as follows: (i) We obtain abundance gradients for all elements from He to Zn. For primary elements, like O, the theoretical value of the abundance gradient at the present epoch (T = 13.5 Gyr) is: dlog(X/H)/dR $``$0.06 dex kpc<sup>-1</sup>. This value is compatible with most observational data concerning young objects in the disk (see Fig. 8 and Table 1). We note, however, the “puzzling” conclusion of Deharveng et al. (2000), pointing to an oxygen gradient smaller by about 40% than the commonly accepted value; if their conclusion is confirmed, our model parameters (the ratio of star formation to infall timescale as a function of galactocentric radius) should be appropriately modified. For secondary or “odd-Z” elements (like N and Al, respectively) we obtain slightly larger gradients; this is also the case for C when the M92 yields are used. Such values are marginally compatible with available observations. We find that the M92 yields can account for the totality of carbon production across the disk with no need for a contribution by IMS, in agreement with other recent studies (Prantzos et al. 1994, Carigi 1994, Gustafsson et al. 1999, Henry et al. 2000a, 2000b). On the contrary, the majority of N is produced by another source, most probably IMS. (ii) Current observations show no trend of the abundance ratio X/O with galactocentric distance (or metallicity), for any element X. This is unexpected in the case of N and C, which do show such a trend in extragalactic HII regions where their abundance ratio to oxygen increases with metallicity (e.g. Henry and Worthey 1999). Our model shows a flat profile of S/O, Si/O and Ar/O (in agreement with observations) and a slowly decreasing ratio of Ne/O, Mg/O and Al/O with galactocentric distance (due to a “problematic” or overestimated metallicity dependence of the WW95 yields for those elements). They also show that C/O should decrease with galactocentric distance in the case of M92 yields (but this should not be the case if C is mainly produced as a primary by IMS). Clearly, precise observations of the C/O ratio across the disk are required in order to decide about its main production site. Similar conclusions hold for N: the M92 yields cannot account for its abundance profile, but is not clear whether primary or secondary production in IMS is the dominant mechanism. Observations of extragalactic HII regions suggest that the latter mechanism dominates in high metallicity regions (Henry et al. 2000b); if this is true, then an important N/O gradient should be observed in the Milky Way disk. (iii) The evolution of abundance gradients provides a strong constraint for Galactic chemical evolution models. Despite a considerable amount of observational and theoretical work (e.g. Henry and Worthey 1999; Tosi 2000 and references therein), the question has not yet been settled. Observations of planetary nebulae (PN) of different ages could, in principle, help in that respect (Pasquali and Perinoto 1993, Maciel and Quireza 1999), but the ages and distances attributed to PN are not sufficiently well known at present. Our model predicts a steady flattening of the gradients with time, due to the adopted “inside-out” formation scheme for the disk. Such an evolution is also found in other works (e.g. Mollá et al. 1997, Portinari and Chiosi 1999) using similar assumptions. Our model makes also a testable prediction: the abundance scatter must be smaller in the inner disk than in the outer regions, if the observed dispersion between PN abundances at a given galactocentric distance is due to their intrinsic age differences. In summary, we have shown that current massive star yields, combined with the BP99 chemical evolution model, reproduce fairly well most of the observed abundance gradients in the Milky Way disk. The assumptions of the model affect mostly points (i) and (iii) above, while the adopted yields affect point (ii). Some problems remain with C and N, and (to a smaller extent) Al. However, the major question of such studies, namely the evolution of abundance gradients, is far from being settled yet; a much more precise determination of ages, distances and abundances of planetary nebulae will be required for that. ###### Acknowledgements. J.L. Hou acknowledges the warm hospitality of the IAP(Paris, France). We are grateful to the referee, Dr. M. Peimbert, for his suggestions that improved the paper. This work was made possible thanks to the support of China Scholarship Council(CSC) and National Natural Sciences Foundation of China.
warning/0007/nucl-ex0007013.html	ar5iv	text	# One-neutron removal reactions on neutron-rich psd-shell nuclei ## Abstract A systematic study of high energy, one-neutron removal reactions on 23 neutron-rich, psd–shell nuclei ($`Z=59,A=1225`$) has been carried out. The longitudinal momentum distributions of the core fragments and corresponding single-neutron removal cross sections are reported for reactions on a carbon target. Extended Glauber model calculations, weighted by the spectroscopic factors obtained from shell model calculations, are compared to the experimental results. Conclusions are drawn regarding the use of such reactions as a spectroscopic tool and spin-parity assignments are proposed for <sup>15</sup>B, <sup>17</sup>C, <sup>19-21</sup>N, <sup>21,23</sup>O, <sup>23-25</sup>F. The nature of the weakly bound systems <sup>14</sup>B and <sup>15,17</sup>C is discussed. PACS: 25.60.-t, 25.60.Gc, 27.20.+n, 27.30.+t KEYWORDS: one-neutron removal, momentum distributions, $`\sigma _{1n}`$, Glauber model. , , , , , , , , , , , , , , , , , thanks: Corresponding author: orr@caelav.in2p3.frthanks: Present address: FSU, Tallahassee, USA.thanks: Present address: LPC, Caen, France.thanks: Present address: IPN, Orsay, France.thanks: Present address: INFN, Catania, Italy. Fragment momentum distributions have long been recognised as signatures of the large spatial extent of the valence nucleons in halo nuclei . Recently measurements of one-nucleon removal<sup></sup><sup></sup>*The term “knockout”, which has been employed to refer to such reactions , is not adopted here as it has long been used for $`(p,2p)`$ and $`(e,e^{}p)`$ reactions, the description of which is very different from that of absorption and diffraction in one-nucleon removal. reactions on light targets have been proposed as a spectroscopic tool for high-energy radioactive beams . This approach has arisen from the development of reaction calculations in which the strong absorption limit and core excited states are accounted for . More specifically, the integrated cross sections are related to spectroscopic factors using an extended version of the spectator core model , whilst the momentum distributions are derived in the opaque limit of the Serber model . To date, this approach has been applied to a few near dripline and halo nuclei . In this Letter the results of an investigation of high-energy one-neutron removal reactions over a broad range of light, neutron-rich psd-shell nuclei are reported. The goals of the work were twofold. Firstly, to explore the evolution in structure, and the manner in which it is manifested in the core fragment observables, from near stability to dripline and halo systems. Secondly, for many of the near stable nuclei the ground state structure is well established and, consequently, it has been possible to test the validity of one-neutron removal reactions as a spectroscopic tool. In the following, measurements of the core fragment longitudinal momentum distributions and integrated cross sections resulting from reactions on a C target are presented. Comparison is made for both observables to the results of extended Glauber type calculations incorporating second order noneikonal corrections to the JLM parameterisation of the optical potential . In the case of those systems with unknown, or poorly defined ground state structures, probable spin-parity assignments have been made. The secondary beams were produced via the fragmentation on a 490 mg/cm<sup>2</sup> thick C target of an intense ($``$1$`\mu `$Ae) 70 MeV/nucleon <sup>40</sup>Ar<sup>17+</sup> beam provided by the GANIL coupled cyclotron facility. The reaction products were collected and selected according to magnetic rigidity using the SISSI device coupled with the alpha-shaped beam analysis spectrometer. A mean rigidity of 2.880 Tm was selected to allow for the transmission of nuclei from <sup>12</sup>B to <sup>25</sup>F with energies in the range of 43 – 71 MeV/nucleon (Table 1). The energy spread in the secondary beams, as defined by the spectrometer acceptances, was $`\mathrm{\Delta }`$E/E=2%. The measurements of the momentum distributions and one-neutron removal cross sections were performed using the SPEG spectrometer . Owing to the large energy spread in the secondary beam, SPEG was operated in a dispersion matched energy-loss mode for which a resolution in the momentum measurements of $`\delta `$p/p = 3.5$`\times `$10<sup>-3</sup> was obtained. Importantly the large angular acceptances of the spectrometer (4 in the vertical and horizontal planes) provided for complete collection of the core fragments, obviating any ambiguities in the integrated cross sections and longitudinal momentum distributions that would arise from limited transverse momentum acceptances . Furthermore, the broad momentum acceptance of the spectrometer ($`\mathrm{\Delta }`$p/p=7%) allowed the momentum distributions for one-neutron removal on all the nuclei of interest to be obtained in a single setting (B$`\rho _{SPEG}`$=2.551 Tm). A secondary C reaction target of thickness 170 mg/cm<sup>2</sup> was employed for the measurements described here (the results obtained with a Ta target will be reported elsewhere ). Ion identification at the focal plane of SPEG was achieved using the energy loss derived from a gas ionisation chamber and the time-of-flight between a thick plastic stopping detector and the cyclotron radio-frequency. Additional identification information was provided by the residual energy measurement furnished by the plastic detector and the time-of-flight with respect to a thin-foil microchannel plate detector located at the exit of the beam analysis spectrometer. Two large area drift chambers straddling the focal plane of SPEG were employed to determine the angles of entry of each ion and, consequently, allowed the focal plane position spectra to be reconstructed. The momentum of each particle was then derived from the reconstructed focal plane position. Calibration in momentum was achieved by removing the reaction target and steping the mixed secondary beam of known rigidity along the focal plane. This procedure also facilitated a determination of the efficiency across the focal plane for the collection of the reaction products. The intensities of the various components of the secondary beam were measured in runs taken with the secondary reaction target removed and the spectrometer set to the same rigidity as the beamline. These were calibrated in terms of the primary beam current, which was recorded continuously throughout the experiment using a non-interceptive beam monitor. Checks were also provided by the counting rates in the microchannel at the exit of the beam analysis spectrometer and a second located just upstream of the secondary reaction target. Typical secondary beam intensities ranged from $``$600 <sup>15</sup>C/s to $``$1 <sup>25</sup>F/s. The longitudinal momentum distributions for the core fragments arising from one-neutron removal are displayed in figure 1 and the extracted widths (FWHM in the projectile frame) are summarised in Table 1. The widths were derived from Gaussian fits to the central regions of each distribution. The effects arising from the target (straggling etc), efficiency along the focal plane and instrumental resolution have been taken into account in deriving the final values. The corresponding one-neutron removal cross sections are listed in Table 1 and displayed in figure 2. The uncertainties quoted include the contributions from both the statistical uncertainty and that arising from the determination of the secondary beam intensity ($``$7%). A number of features are immediately apparent on inspection of figures 1 and 2. Firstly, the crossing of the N=8 shell and N=14 sub-shell closures are associated with a marked reduction in the widths of the core momentum distributions (viz, <sup>14,15</sup>B, <sup>15</sup>C, <sup>23</sup>O and <sup>24,25</sup>F). Secondly, with respect to the neighbouring isotopes, <sup>14</sup>B and <sup>15</sup>C exhibit enhanced one-neutron removal cross sections. The former effects arise from the large $`\nu `$2s<sub>1/2</sub> admixtures expected in the ground states of the Z=4-6, N=9 isotones (see below), which may also persist for N=10, as suggested by recent studies of <sup>14</sup>Be . A narrowing of the momentum distributions may also be expected for N=15 and 16 as in a simple shell model picture the valence neutrons occupy the $`\nu `$2s<sub>1/2</sub> orbital. In general terms, the enhanced cross sections may be attributed to a combination of weak binding (<sup>14</sup>B: S<sub>n</sub>= 0.97 MeV; <sup>15</sup>C: S<sub>n</sub>= 1.22 MeV) and the large $`\nu `$2s<sub>1/2</sub> admixtures in the ground state wavefunctions, which may be related to extended valence neutron density distributions, as discused below. As noted in Table 1, the present measurements may be compared to those made for <sup>14</sup>B and <sup>15,17,18</sup>C . While agreement is found for the momentum distributions, the integrated cross sections are systematically some 3-5 times higher than those reported at similar energies using the A1200 fragment separator . Analysis of the transverse momentum distributions obtained in the present experiment demonstrate that the rather limited acceptances of the A1200 are the origin of this discrepancy . In the case of <sup>14</sup>B, the present results and those of ref. , also obtained using a high acceptance spectrometer, are in good accord. In order to make a more quantitative analysis of the measurements and examine the utility of such reactions as a spectroscopic tool, extended Glauber type calculations have been carried out. The calculations, the principal features of which follow refs. A similar spectator core description and treatment of core excited states was developed earlier by Sagawa et al. in a study of inclusive momentum distributions following neutron removal on <sup>11</sup>Be ., include absorption (or stripping) and diffractive (or elastic) one-nucleon breakup. An important feature is that the S-matrices describing these processes have been derived from the microscopic interaction of Jeukenne, Lejeune and Mahaux (JLM) within an eikonal approximation employing noneikonal corrections . As discussed by Bonaccorso and Carstoiu and Tostevin , such microscopic potentials are much better adapted to the intermediate energy range than optical limit or global parameterisations . In addition to the cross sections, the core longitudinal momentum distributions have been computed within this framework, as opposed to the black disk approximation of ref. . A detailed description of the calculations, together with the results obtained for the transverse momentum distributions and with a Ta target, will be presented elsewhere . In terms of structure, overlaps were calculated between the ground state wavefunctions of the projectiles ($`J^\pi `$) and the core states ($`I_c^\pi `$) coupled to a valence neutron ($`nlj`$). The single-particle wavefunctions were defined within a Woods-Saxon potential with fixed geometry ($`r_0`$=1.15 fm, $`a_0`$=0.5 fm for Z=5 and 6; $`r_0`$=1.2 fm, $`a_0`$=0.6 fm for Z=7-9) with the depth adjusted to reproduce the effective binding energy ($`S_n^{eff}`$) which was fixed as the sum of the single-neutron separation energy and the excitation energy of the core state. The cross section to populate a given core final state is then, $$\sigma (I_c^\pi )=\underset{nlj}{}C^2S(I_c^\pi ,nlj)\sigma _{sp}(nlj,S_n^{eff})$$ (1) where $`C^2S`$ is the spectroscopic factor for the removed neutron with respect to the core state and $`\sigma _{sp}`$ is the cross section for removal of the neutron by absorption ($`\sigma _{abs}`$), diffraction ($`\sigma _{diff}`$) and Coulomb dissociation (only $``$7 mb in the most favourable cases – <sup>14</sup>B and <sup>15</sup>C ). The total inclusive one-neutron removal cross section ($`\sigma _{1n}^{Glauber}`$) is then the sum over the cross sections to all core states. Similarly, the inclusive core momentum distribution is the sum of all core state momentum distributions, weighted by the corresponding cross sections. Within the framework of the spectator core description used here, excitation of the core in the reaction and final-state interactions are neglected. The spectroscopic factors employed here have been calculated with the shell model code OXBASH using the WBP interaction within the 1p-2s1d configuration space. Where known, the experimentally established spin-parity assignments and core excitation energies have been used. In all other cases the shell model predictions have been assumed. The resulting cross sections and momentum distributions are displayed in Table 1 and figures 1 and 2. The breakdown of the calculated cross sections over the core states for each nucleus is detailed in ref. ; as examples, and to aid in the following discussion, the results are listed for <sup>14</sup>B and <sup>15,17</sup>C in Table 2. As the momentum distributions reflect the orbital angular momentum of the removed neutron, the calculated distributions have been normalised to the peak number of counts to facilitate the comparison (figure 1). For all the nuclei observed, including those with known structure, very good agreement is found between the calculated and measured distributions and cross sections, with the exception of <sup>22</sup>F, where the cross section is underestimated. Consequently, spin-parity assignments, derived from the shell model predictions, have been proposed for <sup>15</sup>B, <sup>17</sup>C, <sup>19-21</sup>N, <sup>21,23</sup>O, <sup>23-25</sup>F (Table 1). In the case of <sup>24</sup>F, a 3<sup>+</sup> or 1<sup>+</sup> assignment appears possible based on the present data . The decay study of Reed et al. suggests, however, that the former is the most likely , in line with the shell model predictions. Of particular interest amongst the nuclei investigated here are <sup>14</sup>B and <sup>15,17</sup>C, which, based on the relatively weak binding of the valence neutrons and measurements of the core momentum distributions and one-neutron removal cross sections, have been suggested to be one-neutron halo systems . As may be seen in figures 1 and 2 and Table 2, the momentum distributions and cross sections for <sup>14</sup>B and <sup>15</sup>C are well reproduced by the present calculations employing the spectroscopic factors derived from the shell model, in which the ground state wavefunctions are predominately a 2s<sub>1/2</sub> valence neutron coupled to the core (<sup>13</sup>B and <sup>14</sup>C) in the ground state, as suggested by decay studies and single neutron-transfer experiments . In the case of <sup>17</sup>C, a spin-parity assignement of 3/2<sup>+</sup> is favoured, whereby the ground state configuration is predominately a 1d<sub>5/2</sub> valence neutron coupled to the <sup>16</sup>C core 2$`{}_{1}{}^{}{}_{}{}^{+}`$ state. This confirms the suggestion of Bazin et al. and the calculations of Ren et al. , and is supported by the recent observation of the 1.76 MeV gamma-rays de-exciting the 2$`{}_{1}{}^{}{}_{}{}^{+}`$ state in <sup>16</sup>C following one-neutron removal on <sup>17</sup>C . Such a structure, with a high $`S_n^{eff}`$ (2.49 MeV) and a valence neutron angular momentum of l=2, excludes the possibility of any halo structure developing as evidenced by measurements of the total reaction cross section . Moderate enhancements, however, have been observed in the total reaction cross section measurements for <sup>14</sup>B . Together with the ground state structure deduced from the present experiment and refs. , it seems probable that a spatially extended valence neutron density distribution does occur; although the one-neutron binding energy of nearly 1 MeV will supress the development of a distribution as large as that found in the more weakly bound one-neutron halo nuclei <sup>11</sup>Be and <sup>19</sup>C. Detailed measurements of the total reaction cross section over a range of energies would thus be of particular interest in mapping out the density distribution of <sup>14</sup>B. In the case of <sup>15</sup>C the situation is unclear, with measurements of the total reaction cross section exhibiting no effect and small enhancements . Despite the predominately $`\mathrm{l}`$=0 character of the valence neutron, the higher neutron binding energy of <sup>15</sup>C (S<sub>n</sub>=1.22 MeV) should restrict further the spatial extent of the neutron density distribution. Interestingly, very recent measurements of the charge-changing cross sections for the C isotopes exhibit an increase for <sup>15</sup>C . In summary, a systematic investigation of one-neutron removal reactions has been carried out on a series of neutron-rich psd-shell nuclei. The longitudinal momentum distributions and corresponding single-neutron removal cross sections for the core fragments were measured using a high acceptance spectrometer. Extended Glauber model calculations, coupled with spectroscopic factors derived from shell model calculations employing the WBP interaction, reproduce well the momentum distributions and cross sections. On this basis spin-parity assignments have been proposed for <sup>15</sup>B, <sup>17</sup>C, <sup>19-21</sup>N, <sup>21,23</sup>O, <sup>23-25</sup>F. Given the ground state configurations deduced here and measurements of the total reaction cross section, it is suggested that <sup>14</sup>B presents a moderately extended valence neutron density distribution. This does not appear to be the case for <sup>17</sup>C, whilst <sup>15</sup>C exhibits contradictory behaviour. In more general terms it is concluded that high energy one-nucleon removal reactions represent a powerful spectroscopic tool far from stability. Moreover it has been demonstrated that coupled with a high acceptance, broad range spectrograph, such reactions offer a means to survey structural evolution over a wide range of isospin in a single experiment. The current development of large area, highly segmented, multi-element Ge-arrays should further enhance the sensitivity of such studies. Acknowledgements The support provided by the staffs of lpc and ganil is gratefully acknowledged. Discussions with J.A. Tostevin are also acknowledged, as is the guidance provided by B.A. Brown into the intricacies of the shell model and the assistance provided by G. Martínez in preparing the experiment. This work was funded under the auspices of the in2p3-cnrs (France) and epsrc (United Kingdom). Additional support from the Human Capital and Mobility Programme of the European Community (contract n CHGE-CT94-0056) and the GDR Noyaux Exotiques (cnrs) is also acknowledged.
warning/0007/math-ph0007021.html	ar5iv	text	# On the application of one M.G.Krein’s result to the spectral analysis of Sturm-Liouville operators. (April, 2000 ) ## Abstract Discovered by M.G.Krein analogy between polynomials orthogonal on the unit circle and generalized eigenfunctions of certain differential systems is used to obtain some new results in spectral analysis of Sturm-Liouville operators. Section A. In this section we remind some results obtained by M.G.Krein in his famous article . In this paper author develops the ”theory of polinomials, orthogonal on the positive half-line”. And this polinoms are constructed from exponents rather than from the powers of independent variable. It’s well known that there are many ways to construct the system of orthogonal polinomials on the unit circle. One of them is to start from the moments matrix. This way was chosen by M.G.Krein to obtain his results for positive half-line. Let’s assume that $`H(t)=\overline{H(t)}`$ \- function summable on each segment $`(r,r)`$. Proposition. If for any continuous $`\phi (t)`$ the following inequaliy holds $$\underset{0}{\overset{r}{}}\left\|\phi (s)\right\|^2𝑑s+\underset{0}{\overset{r}{}}\underset{0}{\overset{r}{}}H(ts)\phi (t)\overline{\phi (s)}𝑑t𝑑s0$$ (1) for each $`r>0`$, so, and in this case only, there exists the non-decreasing function $`\sigma (\lambda )`$ $`(\lambda R,\sigma (0)=0,\sigma (\lambda 0)=\sigma (\lambda )`$ $`)`$, such that $`{\displaystyle \underset{\mathrm{}}{\overset{\mathrm{}}{}}}{\displaystyle \frac{d\sigma (\lambda )}{1+\lambda ^2}}`$ $`<`$ $`\mathrm{}`$ $`{\displaystyle \underset{0}{\overset{t}{}}}(ts)H(s)𝑑s`$ $`=`$ $`{\displaystyle \underset{\mathrm{}}{\overset{\mathrm{}}{}}}(1+{\displaystyle \frac{i\lambda t}{1+\lambda ^2}}e^{i\lambda t}){\displaystyle \frac{d\sigma (\lambda )}{\lambda ^2}}+(i\gamma {\displaystyle \frac{sign(t)}{2}})t`$ where $`\gamma `$ is real constant. If in addition we presume that the equality in (1) is possible for $`\phi =0`$ only then the Hermit kernel $`H(ts),(0t,sr)`$ has Hermit resolvent $`\mathrm{\Gamma }_r(s,t)=\overline{\mathrm{\Gamma }_r(s,t)}`$ that satisfies the relation $$\mathrm{\Gamma }_r(t,s)+\underset{0}{\overset{r}{}}H(tu)\mathrm{\Gamma }_r(u,s)𝑑u=H(ts)(0s,tr)$$ The continuous analogues of polinomials orthogonal on the unit circle are defined by the formula $`P(r,\lambda )`$ $`=`$ $`e^{i\lambda r}(1{\displaystyle \underset{0}{\overset{r}{}}}\mathrm{\Gamma }_r(s,0)e^{i\lambda s}𝑑s),`$ $`P_{}(r,\lambda )`$ $`=`$ $`1{\displaystyle \underset{0}{\overset{r}{}}}\mathrm{\Gamma }_r(0,s)e^{i\lambda s}𝑑s,r0.`$ Using the well known properties of resolvents we obtain the following system $$\begin{array}{ccc}\frac{dP(r,\lambda )}{dr}& =& i\lambda P(r,\lambda )\overline{A(r)}P_{}(r,\lambda ),\\ \frac{dP_{}(r,\lambda )}{dr}& =& A(r)P(r,\lambda ),\end{array}$$ (2) where $`A(r)=\mathrm{\Gamma }_r(0,r).`$ Proposition. For each finite $`f(x)L^2(R^+)`$ we have the following equality $$f_2^2=\underset{\mathrm{}}{\overset{\mathrm{}}{}}\left\|F_P(\lambda )\right\|^2𝑑\sigma (\lambda ),\mathrm{where}F_P(\lambda )=\underset{0}{\overset{\mathrm{}}{}}f(r)P(r,\lambda )𝑑r.$$ Consequently we have the isometric mapping $`U_P`$ from $`L^2(R^+)`$ into $`L^2(\sigma ,R).`$ Theorem. The mapping $`U_P`$ is unitary if and only if the following integral diverges (equals to $`\mathrm{}`$ ) $$\underset{\mathrm{}}{\overset{\mathrm{}}{}}\frac{\mathrm{ln}\sigma ^{^{}}(\lambda )}{1+\lambda ^2}𝑑\lambda .$$ (3) Theorem. The following statements are equivalent (1) The integral (3) is finite. (2) At least for some $`\lambda `$, $`\mathrm{}\lambda >0`$ the integral $$\underset{0}{\overset{\mathrm{}}{}}\left\|P(r,\lambda )\right\|^2𝑑r$$ (4) converges. (3) At least for some $`\lambda `$ ($`\mathrm{}\lambda >0`$ ) the function $`P_{}(r,\lambda )`$ is bounded. (4) On any compact set in the open upper half-plane integral (4) converges uniformly. That is equivalent to the existence of uniform limit $`\mathrm{\Pi }(\lambda )=lim_r\mathrm{}P_{}(r,\lambda ).`$ It’s easy to verify that in cases $`A(r)L^1(R^+),A(r)L^2(R^+)`$ the conditions (1)-(4) are satisfied. What is more, in the first case measure $`\sigma `$ is continuously differentiable with certain estimates for its derivative. Consider $`E(r,\lambda )=e^{i\lambda r}P(2r,\lambda )=\mathrm{\Phi }(r,\lambda )+i\mathrm{\Psi }(r,\lambda ).`$ Let $`E(r,\lambda )=\overline{E(r,\lambda )}=\mathrm{\Phi }(r,\lambda )i\mathrm{\Psi }(r,\lambda ).`$ From (2) we infer that $`{\displaystyle \frac{d\mathrm{\Phi }}{dr}}`$ $`=`$ $`\lambda \mathrm{\Psi }a(r)\mathrm{\Phi }+b(r)\mathrm{\Psi },\mathrm{\Phi }(0,\lambda )=1;`$ $`{\displaystyle \frac{d\mathrm{\Psi }}{dr}}`$ $`=`$ $`\lambda \mathrm{\Phi }+b(r)\mathrm{\Phi }+a(r)\mathrm{\Psi },\mathrm{\Psi }(0,\lambda )=0.`$ where $`a(r)=2\mathrm{}A(2r),b(r)=2\mathrm{}A(2r).`$ Proposition. The mapping $`U_E:f(r)F_E(\lambda )=\underset{\mathrm{}}{\overset{\mathrm{}}{}}f(r)E(r,\lambda )𝑑r,`$ defined on the finite functions $`f(r)L^2(R),`$ generates the unitary operator from $`L^2(R)`$ onto $`L^2(\sigma ,R)`$. The trivial case $`a=b=0`$ yields $`E(r,\lambda )=e^{i\lambda r},\mathrm{\Psi }(r,\lambda )=\mathrm{sin}(r\lambda ),\mathrm{\Phi }(r,\lambda )=\mathrm{cos}(r\lambda ).`$ In case when $`H(t)`$ is real, the function $`\sigma (\lambda )`$ is odd. Consequently $`b(r)=0.`$ Assuming that $`H(t)`$ is absolutely continuous, we have that $`\mathrm{\Phi }`$ and $`\mathrm{\Psi }`$ are solutions of the equations $$\begin{array}{ccc}\mathrm{\Psi }^{^{\prime \prime }}q\mathrm{\Psi }+\lambda ^2\mathrm{\Psi }=0,& \mathrm{\Psi }(0)=0,& \mathrm{\Psi }^{^{}}(0)=\lambda ;\\ \mathrm{\Phi }^{^{\prime \prime }}q_1\mathrm{\Phi }+\lambda ^2\mathrm{\Phi }=0,& \mathrm{\Phi }(0)=1,& \mathrm{\Phi }^{^{}}(0)+a(0)\mathrm{\Phi }(0)=0,\end{array}$$ (5) where $`q_1(x)=a^2(x)a^{^{}}(x)`$ and $`q(x)=a^2(x)+a^{^{}}(x).`$ Section B. Consider the Sturm-Liouville operator on the half-line with Dirichlet boundary condition at zero $$l(u)=u^{^{\prime \prime }}+qu,u(0)=0.$$ (6) Let’s assume that real-valued $`q(x)`$ admits the following representation $`q(x)=a^2(x)+a^{^{}}(x)`$ where $`a(x)`$ is absolutely continuous function on the half-line. That means that $`a(x)`$ is the solution of the Ricatti equation. Consider also the corresponding differential Dirac-type system (see or ) $$\{\begin{array}{cc}\mathrm{\Phi }^{^{}}(x,\lambda )=& \lambda \mathrm{\Psi }(x,\lambda )a(x)\mathrm{\Phi }(x,\lambda )\\ \mathrm{\Psi }^{^{}}(x,\lambda )=& \lambda \mathrm{\Phi }(x,\lambda )+a(x)\mathrm{\Psi }(x,\lambda )\end{array},$$ (7) where $`\mathrm{\Phi }(0,\lambda )=1,\mathrm{\Psi }(0,\lambda )=0.`$ From the result stated in Section A it follows that the spectral measure $`\rho (\lambda )`$ of problem (6) is connected with the spectral measure $`\widehat{\sigma }(\lambda )`$ of system (7) <sup>1</sup><sup>1</sup>1See the definition of spectral measure for the differential system in . Here the function $`\widehat{\sigma }(\lambda )`$ is connected with $`\sigma (\lambda )`$ from the section A by the relation: $`\widehat{\sigma }(\lambda )=2\sigma (\lambda ).`$ by the following relation $$\rho (t)=2\underset{0}{\overset{\sqrt{t}}{}}\alpha ^2𝑑\widehat{\sigma }(\alpha ).$$ (8) From this and results stated in Section A we can infer one very simple but significant corollary Corollary. If $`q(x)`$ is real-valued function such that $$\underset{xR}{sup}\underset{x}{\overset{x+1}{}}\left\|q(s)\right\|^2𝑑s<\mathrm{},$$ (9) the improper integral $`W(x)=\underset{x}{\overset{\mathrm{}}{}}q(s)𝑑s`$ exists and satisfies the condition $`W(x)L^2(R^+),`$ then the absolutely continuous part of spectrum of operator $`H_h`$, generated by differential expression $`l(u)=u^{^{\prime \prime }}+qu`$ and boundary condition$`u(0)=hu^{^{}}(0),(hR\mathrm{})`$ fills the whole positive half-line. Proof. Let’s consider (7) with $`a(x)=W(x).`$ The spectral measure of system (7) with chosen $`a(x)`$ has the needed property. Consequently the Sturm-Liouville operator with potential $`q^{}(x)=a^{^{}}+a^2`$ and Dirichlet boundary condition also has the a.c. component which fills the whole positive half-line. But initial potential $`q(x)`$ differs from $`q^{}(x)`$ by $`L^1(R^+)`$ term only. Consequently, Kuroda’s theorem guarantees that operator $`H_h`$ has the needed property for $`h=0`$. But it means that this statement is true for any $`h`$ since the essential support of a.c. component doesn’t depend on $`h.`$ That follows, for example, from the subordinate solutions theory.$`\mathrm{}`$ In the next theorems of this section we will show how Krein’s results will help to establish asymptotics for generalized eigenfunctions and analyze the spectrum of some Sturm-Liouville operators. Theorem 1. If $`q(x)`$ is real-valued function such that $$\underset{xR}{sup}\underset{x}{\overset{x+1}{}}\mathrm{min}\{0,q(s)\}𝑑s>\mathrm{},$$ (10) the improper integral $`W(x)=\underset{x}{\overset{\mathrm{}}{}}q(s)𝑑s`$ exists and satisfies the condition $`\left\|W(x)\right\|\frac{\gamma }{x+1},(0<x,0<\gamma <1/4),`$ then operator $`H`$ generated by (6) is non-negative and $$\left\|\underset{0}{\overset{\mathrm{}}{}}\frac{\mathrm{ln}\rho ^{^{}}(\lambda )}{\sqrt{\lambda }(1+\lambda )}𝑑\lambda \right\|<\mathrm{}.$$ (11) What is more for a.e. positive spectral parameter the generalized eigenfunctions of differential expression $`l(u)=u^{^{\prime \prime }}+qu`$ has the following asymptotic $`u(x,\lambda ,\alpha )=C(\lambda ,\alpha )\mathrm{sin}(x\lambda +\phi (\lambda ,\alpha ))+\overline{o}(1),u(0,\lambda ,\alpha )=\mathrm{cos}(\alpha ),u^{^{}}(0,\lambda ,\alpha )=\mathrm{sin}(\alpha ).`$ The spectrum of operator $`H`$, generated by (6), is purely absolutely continuous on the positive half-line. Proof. We will separate the proof on two parts. 1. Reducing to the system. Let’s consider Ricatti equation $`q(x)=a^2(x)+a^{^{}}(x)`$ . We will find solution which is absolutely continuous, tends to zero at the infinity and belongs to the class $`L^2(R^+)`$. Integrating we will get the following nonlinear integral equation $`\underset{x}{\overset{\mathrm{}}{}}a^2(s)𝑑sa(x)=W(x).`$ Our goal is to study operator $`B`$: $`Bf(x)=\underset{x}{\overset{\mathrm{}}{}}f^2(s)𝑑sW(x)`$ that acts in complete metric space $`\mathrm{\Omega }`$ of measurable functions which admit the estimate $`\left\|g(x)\right\|\frac{\ae }{x+1},(\ae =\frac{1\sqrt{14\gamma }}{2}).`$ The metric is introduced by the formula $`\rho (g_1,g_2)=esssup_{x0}\{(x+1)\|g_1(x)g_2(x)\|\}.`$ To use the contraction operators principle one should verify that the following conditions hold 1. $`B`$ is acting from $`\mathrm{\Omega }`$ to $`\mathrm{\Omega }`$, 2. $`B`$ is contraction operator. It is not difficult to show that the both conditions are satisfied. Really $$\left\|Bf\right\|\frac{\gamma }{x+1}+\underset{x}{\overset{\mathrm{}}{}}\frac{\ae ^2}{(s+1)^2}𝑑s=\frac{\ae ^2+\gamma }{x+1}=\frac{\ae }{x+1},$$ $$\left\|Bg_1Bg_2\right\|\underset{x}{\overset{\mathrm{}}{}}\left\|g_1g_2\right\|\left\|g_1+g_2\right\|𝑑s\rho (g_1,g_2)\underset{x}{\overset{\mathrm{}}{}}\frac{2\ae }{(s+1)^2}𝑑s=\frac{2\ae }{x+1}\rho (g_1,g_2),$$ that means $`\rho (Bg_1,Bg_2)2\ae \rho (g_1,g_2)`$ which implies the contraction property since $`2\ae <1`$. Thus we have the single fixed point $`a(x)\mathrm{\Omega }`$, so that $`Ba=a.`$ Certainly this function satisfies the Ricatti equation as well. So it suffices to use Proposition and (8) to obtain (11). 2. Absence of singular component. Consider the system (7). From Section A we know that function $$P(x,\lambda )=\mathrm{exp}(i\lambda \frac{x}{2})\left(\mathrm{\Phi }(x/2,\lambda )+i\mathrm{\Psi }(x/2,\lambda )\right)$$ (12) satisfies the following system $$\{\begin{array}{ccc}\frac{dP}{dx}& =& i\lambda PAP_{},\\ & & \\ \frac{dP_{}}{dx}& =& AP.\end{array}$$ Where $`P(0,\lambda )=P_{}(0,\lambda )=1,`$ and $`A(x)=\frac{1}{2}a(\frac{x}{2}).`$ Let’s introduce the following function $`Q(x)=e^{i\lambda x}P(x).`$ So we will have $$\{\begin{array}{ccc}\frac{dQ}{dx}=& Ae^{i\lambda x}P_{},& \\ & & \\ \frac{dP_{}}{dx}=& Ae^{i\lambda x}Q.& \end{array}$$ $`P_{}(0,\lambda )=Q(0,\lambda )=1.`$ It’s easy to see that $`Q=\overline{P_{}}`$. Consequently $`Q(x)=1\underset{0}{\overset{x}{}}A(s)e^{i\lambda s}\overline{Q(s)}𝑑s;`$ $`\left\|Q(x)\right\|1+\underset{0}{\overset{x}{}}\left\|A(s)\right\|\left\|Q(s)\right\|𝑑s.`$ Gronuol Lemma yields the estimate $`\left\|Q(x)\right\|\mathrm{exp}\left(\underset{0}{\overset{x}{}}\left\|A(s)\right\|𝑑s\right)\left(\frac{x+2}{2}\right)^\ae .`$ From (5) it follows that $`u(x,\lambda )=\frac{\mathrm{\Psi }(x,\lambda )}{\lambda }`$ $`(\lambda 0)`$ satisfies the conditions $`u^{^{\prime \prime }}+qu=\lambda ^2u,`$ $`u(0,\lambda )=0,`$ $`u^{^{}}(0,\lambda )=1`$. From (12) we have $`\left\|u(x,\lambda )\right\|\frac{(x+1)^\ae }{2^\ae \lambda },`$ $`(\ae <1/2`$ $`).`$ The similar estimate can be proved for linear independent solution $`v(x,\lambda )`$ such that $`v(0,\lambda )=1,v^{^{}}(0,\lambda )=0.`$ Indeed it suffices to consider second equation of (5) letting $`q_1=q`$, solve equation $`q=b^2b^{^{}}`$ and repeat the same arguments to obtain the desired inequality for linear independent solution $`w(x,\lambda )`$ which satisfies the condition $`w^{^{}}(0)+b(0)w(0)=0`$. Since $`v(x,\lambda )`$ is linear combination of $`u(x,\lambda )`$ and $`w(x,\lambda )`$ we have the needed estimate. In the same way derivatives $`u^{^{}},v^{^{}}`$ can be estimated. It follows from the inequality $`\|Q^{^{}}\|=\|A\|\|Q\|\frac{\ae }{2^\ae (2+x)^{1\ae }}`$ and (12). Consequently from the constancy of Wronskian $`W(u,v)`$ and equivalence of $`\underset{x}{\overset{x+1}{}}u^2(s,\lambda )𝑑s`$ and $`\underset{x}{\overset{x+1}{}}u_{}^{^{}}{}_{}{}^{2}(s,\lambda )𝑑s`$ <sup>2</sup><sup>2</sup>2It’s the only place in the whole proof where we use the condition (10). we get the inequalities $`C_2(\lambda )x^{12\ae }`$ $``$ $`{\displaystyle \underset{0}{\overset{x}{}}}u^2(s,\lambda )𝑑sC_1(\lambda )x^{2\ae +1},`$ $`C_2(\lambda )x^{12\ae }`$ $``$ $`{\displaystyle \underset{0}{\overset{x}{}}}v^2(s,\lambda )𝑑sC_1(\lambda )x^{2\ae +1},`$ where the constants $`C_1,C_2`$ are positive. So we can find $`\zeta >0`$ so that $`\left(\underset{0}{\overset{x}{}}u_1^2(s,\lambda )𝑑s\right)\left(\underset{0}{\overset{x}{}}u_2^2(s,\lambda )𝑑s\right)^\zeta \mathrm{}`$ for any two linearly independent solutions $`u_1,u_2`$ of equation from (6). The refined subordinacy theory , yields that there is $`\eta (\lambda )>0`$ so that $`\left(D_\eta \rho \right)(\lambda ^2)=\overline{lim}_{\epsilon 0}`$ $`\frac{\rho (\lambda ^2\epsilon ,\lambda ^2+\epsilon )}{\left(2\epsilon \right)^\eta }=0.`$ Consequently $`\rho `$ gives zero weight to every $`\mathrm{\Omega }R^+`$ with $`dim\mathrm{\Omega }=0.`$ On the other hand we will prove that $`u_1,u_2`$ might be unbounded at the infinity only on the $`\lambda `$ set with the zero Hausdorff dimension. Consider $`Q(x,\lambda ):Q^{^{}}=Ae^{i\lambda x}\overline{Q(x)},Q(0)=1.`$ So $`Q(x)`$ $`=`$ $`1{\displaystyle \underset{0}{\overset{x}{}}}A(s)e^{i\lambda s}\overline{Q(s)}𝑑s=1+{\displaystyle \underset{0}{\overset{x}{}}}\left({\displaystyle \underset{s}{\overset{\mathrm{}}{}}}A(\tau )e^{i\lambda \tau }𝑑\tau \right)^{^{}}\overline{Q(s)}𝑑s=`$ $$1+\left(\underset{s}{\overset{\mathrm{}}{}}A(\tau )e^{i\lambda \tau }𝑑\tau \right)\overline{Q(s)}\|_{s=0}^{s=x}\underset{0}{\overset{x}{}}\left(\underset{s}{\overset{\mathrm{}}{}}A(\tau )e^{i\lambda \tau }𝑑\tau \right)\overline{Q^{^{}}(s)}𝑑s,$$ (13) where $`\lambda `$ is such that integral $`\underset{0}{\overset{\mathrm{}}{}}A(\tau )e^{i\lambda \tau }𝑑\tau `$ converges. The last term in (13) can be rewritten as $`\underset{0}{\overset{x}{}}\left(\underset{s}{\overset{\mathrm{}}{}}A(\tau )e^{i\lambda \tau }𝑑\tau \right)A(s)e^{i\lambda s}Q(s)𝑑s.`$ Finally we get $$Q(x)=J(x)\underset{0}{\overset{x}{}}\left(\underset{s}{\overset{\mathrm{}}{}}A(\tau )e^{i\lambda \tau }𝑑\tau \right)A(s)e^{i\lambda s}Q(s)𝑑s,$$ where $`J(x)=C(\lambda )+o(1)\overline{Q(x)}`$. From argument used in (Theorem 1.3, 1.4 ) it follows that the set $`\mathrm{\Xi }`$ of $`\lambda `$ for which $`\left(\underset{s}{\overset{\mathrm{}}{}}A(\tau )e^{i\lambda \tau }𝑑\tau \right)A(s)e^{i\lambda x}L^1(R^+)`$ has zero Hausdorff dimension. But one can easily verify that this fact leads to the boundedness of $`Q(x)`$ at the infinity for $`\lambda \mathrm{\Xi }.`$ From formula (12) it follows that generalized eigenfunctions $`u(x,\lambda )`$ which correspond to the Dirichlet boundary condition at zero are bounded at the infinity for $`\lambda \mathrm{\Xi }.`$ At the same way we can prove that the linear independent solution $`v(x,\lambda )`$ is bounded at the infinity if $`\lambda \mathrm{{\rm Y}}`$ $`(dim\mathrm{{\rm Y}}=0).`$ Consequently results obtained by Stolz guarantee that the support of singular measure has the zero Hausdorff dimension. Meanwhile as we have shown above spectral measure gives the zero weight to any set of zero Hausdorff measure. So the singular spectrum is absent.$`\mathrm{}`$ Remark 1. We could have got rid of the condition (10) and prove the absence of positive eigenvalues by making use of Hardy inequality for the equation $`Q(x)=\underset{x}{\overset{\mathrm{}}{}}A(s)e^{i\lambda s}\overline{Q(s)}𝑑s`$ which follows from $`Q^{^{}}=Ae^{i\lambda x}\overline{Q}`$ and $`Q(\mathrm{})=0`$. Really, it would mean the absence of nonzero eigenvalues for differential system (7). By (8) we would have the absence of positive eigenvalues for (6). Remark 2. Conditions of Theorem 1 are often fulfilled for potentials that oscillate at the infinity (see , and bibliography there ). We used method different from common ones such as modified Prüfer transform, $`I+Q`$ asymptotic integration and so on. Remark 3. If we consider potentials which satisfy the estimate $`\left\|q_\epsilon (x)\right\|\frac{\epsilon }{x+1}`$, then the point spectrum may occur on $`[0,\frac{4\epsilon ^2}{\pi ^2}]`$ (see , ) for any $`\epsilon >0`$. The situation for potentials considered in Theorem 1 is different. Really, the von Neumann -Wigner example $`q_{vNW}=8\frac{\mathrm{sin}2x}{x}+O(\frac{1}{x^2})`$ shows that the condition $$\left\|\underset{x}{\overset{\mathrm{}}{}}q_\gamma (s)𝑑s\right\|\frac{\gamma }{x+1}$$ (14) doesn’t guarantee the absence of the positive eigenvalues for $`\gamma `$ large enough. Meanwhile, for small $`\gamma (0<\gamma <1/4)`$ the singular spectrum disappears on the whole positive half-line. <sup>3</sup><sup>3</sup>3It’s interesing to find out is the constant $`1/4`$ optimal or not. Remark 4. From the proof of the Theorem 1 we can infer that the asymptotics for generalized eigenfunctions may not be true on the set of zero Hausdorff dimension only. And this statement holds even without the constraint (10). In the following simple statement is proved If $`q(x)`$ is continuous function which admits the representation $`q=\frac{W}{x},WL^1(1,\mathrm{})`$ then equation $`\phi ^{^{\prime \prime }}+q\phi =k^2\phi `$ doesn’t have nontrivial square integrable solutions for $`k0.`$ What is more, if $`\left\|W(x)\right\|C\|x\|^{1\epsilon }`$ then for any $`k0`$ there exists the solution $`\phi (x,k)`$ of the same equation which has the following asymptotic $`\left\|\phi (x,k)e^{ixk}\right\|C(k)\left\|x\right\|^\epsilon ,x1,`$where $`\|C(k)\|<C(k_0)`$ for $`\left\|k\right\|k_0>0.`$ We improved this result in power scale in some extent using method from . Theorem 2. If $`q(x)`$ is real-valued function such that the improper integral $`W(x)=\underset{x}{\overset{\mathrm{}}{}}q(s)𝑑s`$ exists and satisfies the condition $`W(x)L^p(R^+)L^2(R^+),(1<p<2),`$ then for a.e. positive spectral parameter the generalized eigenfunctions has the following asymptotic $`u(x,\lambda ,\alpha )=C(\lambda ,\alpha )\mathrm{sin}(x\lambda +\phi (\lambda ,\alpha ))+\overline{o}(1),whereu(0,\lambda ,\alpha )=\mathrm{cos}(\alpha ),u^{^{}}(0,\lambda ,\alpha )=\mathrm{sin}(\alpha ).`$ Proof. Let’s consider the system (7) with $`a(x)=W(x).`$Introducing $`P(x,\lambda )`$ by (12) and $`Q(x,\lambda )=e^{i\lambda x}P(x)`$ we have the following equation for $`Q(x,\lambda ):`$ $$Q^{^{}}=Ae^{i\lambda x}\overline{Q},$$ (15) where $`Q(0,\lambda )=1,A(x)=\frac{1}{2}a(\frac{x}{2}).`$ So we can see that $`Q(x,\lambda )=1\underset{0}{\overset{x}{}}A(s)\overline{P(s,\lambda )}𝑑s=1\overline{\underset{0}{\overset{x}{}}A(s)P(s,\lambda )𝑑s}.`$Since $`P(x,\lambda )`$ play the role of orthogonal polynomials, we can expect that the analogue of Menchoff theorem for orthonormal systems will guarantee the convergence of the last integral almost everywhere. But for our purpose it’s more convenient not to derive the generalization of Menchoff results for $`P(x,\lambda )`$ system but to give the reference to the very recent results of M.Christ, A.Kiselev Really, for the solution of equation (15) we have the following formal series: $`Q(x,\lambda )`$ $`=`$ $`Q_{\mathrm{}}(\lambda )(1+{\displaystyle \underset{x}{\overset{\mathrm{}}{}}}A(s_1)e^{i\lambda s_1}ds_1+\mathrm{}`$ $`+{\displaystyle \underset{x}{\overset{\mathrm{}}{}}}A(s_1)e^{i\lambda s_1}{\displaystyle \underset{s_1}{\overset{\mathrm{}}{}}}A(s_2)e^{i\lambda s_2}\mathrm{}{\displaystyle \underset{s_{j1}}{\overset{\mathrm{}}{}}}A(s_j)e^{(1)^j\lambda s_j}ds_j\mathrm{}ds_1+\mathrm{}).`$ The convergence of this kind of series for a.e. $`\lambda `$ w.r.t. Lebesgue measure was established in for $`A(s)L^p(R^+),1p<2.`$ By the methods of this paper we can show that $`Q(x,\lambda )`$ satisfies the equation (15) for a.e. $`\lambda .`$ Thus we have that a.e. $`Q(x,\lambda )`$ is bounded at the infinity. Let’s consider now the Sturm-Liouville operator on the half-line with potential $`q^{}(x)=a^{^{}}+a^2`$ with Dirichlet boundary condition at zero. From (5) and (12) we can infer that the generalized eigenfunctions of this equation which satisfy the Dirichlet condition at zero are bounded at the infinity for a.e. positive spectral parameter. It’s obvious that $`q^{}(x)`$ can be represented in the following form $`q^{}(x)=T^{^{}}+T^2`$, where $`T=W`$. Repeating the same argument for Dirac-type system (7) with $`a(x)=T(x)`$ we see, that generalized eigenfunctions, which satisfy the conditions $`\mathrm{\Phi }(0)=1,`$ $`\mathrm{\Phi }^{^{}}(0,\lambda )+T(0)\mathrm{\Phi }(0,\lambda )=0`$ are bounded at the infinity as well. At the same way we can show that the whole transfer matrix of Sturm-Liouville equation with potential $`q^{}(x)`$ is bounded at the infinity for a.e. positive spectral parameter. But $`q^{}(x)=W^2W^{^{}}=W^2+q,`$ so it differs from the initial potential $`q`$ by the $`W^2L^1(R^+)`$ term only. It’s easy to show, that if the transfer matrix of S.L. operator is bounded for some $`\lambda ,`$ so the transfer matrix of operator obtained by the $`L^1(R^+)`$ perturbation is bounded also. Consequently the transfer matrix of the initial S.L. operator is bounded for a.e. positive spectral parameter. $`\mathrm{}`$ Remark 1. The fulfillment of the conditions of Theorem 2 doesn’t mean that the singular component of spectrum is absent. Moreover the von Neumann- Wigner potential illustrates that at least the positive eigenvalue may occur. Now we will show that conditions of Theorem 2 are satisfied under some assumption imposed on the $`\mathrm{cos}`$ -transform of potential. From then on we suppose that $`q(x)`$ is such that (A) its $`\mathrm{cos}`$ -transform $`\widehat{q(\omega )}=lim_N\mathrm{}\underset{0}{\overset{N}{}}q(x)\mathrm{cos}(\omega x)𝑑x`$exists in $`L_{loc}^1(R^+)`$ sence. (B) $`\frac{2}{\pi }lim_N\mathrm{}\underset{0}{\overset{N}{}}\widehat{q(\omega )}\mathrm{cos}(\omega x)𝑑\omega =q(x)`$ in $`L_{loc}^1(R^+).`$ Theorem 3. If $`q(x)`$ is such that (A) and (B) are satisfied and $`\widehat{q(\omega )}=\widehat{q(0)}+\widehat{\phi (\omega )}`$ where $`\frac{\widehat{\phi (\left\|\omega \right\|)}}{\omega }L^{\epsilon ,2}(R)`$ for some positive $`\epsilon `$, then the asymptotics from the Theorem 2 is true. Proof. Let’s consider even infinitely smooth function $`\widehat{\chi (\omega )}`$ such that $`\widehat{\chi (\omega )}=\{\begin{array}{cc}1,& \|\omega \|1/2;\\ 0,& \|\omega \|1..\end{array}`$ Then $`\widehat{q(\omega )}=\widehat{q(0)}\widehat{\chi (\omega )}+\widehat{\psi (\omega )}`$ where $`\widehat{\psi (\omega )}=\widehat{q(0)}\widehat{[1\chi (\omega )]}+\widehat{\phi (\omega ).}`$ It’s easy to see that $`\frac{\widehat{\psi (\left\|\omega \right\|)}}{\omega }L^{\epsilon ,2}(R).`$ And it suffices to prove that S.L. operator with potential $`\psi (x)=\frac{2}{\pi }\underset{0}{\overset{\mathrm{}}{}}\widehat{\psi (\omega )}\mathrm{cos}(\omega x)𝑑\omega `$ has the transfer matrix bounded at the infinity for a.e. positive spectral parameter. Really, the initial potential $`q(x)=\psi (x)+\chi (x)\widehat{q(0)}`$ where $`\chi (x)=\frac{2}{\pi }\underset{0}{\overset{\mathrm{}}{}}\widehat{\chi (\omega )}\mathrm{cos}(\omega x)𝑑\omega =\frac{1}{\pi }\underset{\mathrm{}}{\overset{\mathrm{}}{}}\widehat{\chi (\omega )}\mathrm{cos}(\omega x)𝑑\omega L^1(R^+).`$ So the transfer matrix of initial operator would be bounded at the infinity for a.e. positive spectral parameter. Consider $`\widehat{L(\omega )}=\frac{\widehat{\psi (\left\|\omega \right\|)}}{\omega }.`$ $`W(x)`$ $`=`$ $`\underset{N\mathrm{}}{lim}{\displaystyle \underset{x}{\overset{N}{}}}\psi (s)𝑑s={\displaystyle \frac{2}{\pi }}\underset{N\mathrm{}}{lim}{\displaystyle \underset{x}{\overset{N}{}}}\left\{\underset{T\mathrm{}}{lim}{\displaystyle \underset{0}{\overset{T}{}}}\widehat{\psi (\omega )}\mathrm{cos}(\omega s)𝑑\omega \right\}𝑑s=`$ $`=`$ $`{\displaystyle \frac{2}{\pi }}\underset{N\mathrm{}}{lim}\underset{T\mathrm{}}{lim}{\displaystyle \underset{0}{\overset{T}{}}}\widehat{\psi (\omega )}\left\{{\displaystyle \underset{x}{\overset{N}{}}}\mathrm{cos}(\omega s)𝑑s\right\}𝑑\omega ={\displaystyle \frac{1}{\pi }}\underset{N\mathrm{}}{lim}\underset{T\mathrm{}}{lim}{\displaystyle \underset{T}{\overset{T}{}}}\widehat{L(\omega )}(\mathrm{sin}(N\omega )\mathrm{sin}(x\omega ))𝑑\omega =`$ $`=`$ $`\sqrt{{\displaystyle \frac{2}{\pi }}}\mathrm{}\widehat{L(x)},`$ where $`L(x)`$ is from $`L^p(R^+)`$ (for some $`p<2`$ ) as the Fourier transform of $`L^{\epsilon ,2}(R)`$ function and so the arguments from the Theorem 2 are applicable. $`\mathrm{}`$ We see that roughly speaking these conditions are fulfilled if $`\widehat{q(\omega )}`$ is relatively smooth near the zero and admits some bounds at the infinity. Remark 1. Local condition in zero of Theorem 3 is satisfied if $`\widehat{q(\omega )}`$ is from $`W^{1+\epsilon ,2}(0,\delta )`$ for some positive $`\delta >0.`$ Remark 2. In paper it was considered the dependence of absence of singular component on certain interval on the local smoothness of Fourier transform. Remark 3. From the corollary and method used in theorem 3 it follows that if $`q(x)`$ is such that (A) and (B) are satisfied, $`\widehat{q(\omega )}=\widehat{q(0)}+\widehat{\phi (\omega )}`$ where $`\left\|\widehat{\phi (\left\|\omega \right\|)}\right\|C\|\omega \|^{1/2+\epsilon }`$ in the vicinity of zero for some positive $`\epsilon `$ and $`\frac{\widehat{\phi (\left\|\omega \right\|)}}{\|\omega \|+1}L^2(R),`$ then the essential support of spectral measure of operator is $`R^+.`$ Section C. In this section we will discuss the dependence of spectral measure $`\sigma (\lambda )`$ on the coefficient $`A(x)`$ of system (2). In fact function $`A(x)`$ plays the role of sequence $`a_n`$ for polynomials orthogonal on the unit circle (see ). We will see that for the Dirac-type systems the situation is not so simple. The basic reason is the possible oscillation of $`A(x)`$. The following Lemma is true Lemma. If measurable bounded function $`A(x)`$ is such that $$\underset{x}{\overset{\mathrm{}}{}}e^sA(s)𝑑s=\overline{o}(e^x),A(x)e^x\underset{x}{\overset{\mathrm{}}{}}A(s)e^s𝑑sL_1(R^+)$$ then conditions (1)-(4) from section A are satisfied. Proof. Consider system (2) with $`\lambda =i`$. If $`P=e^xQ`$ then we have $$\begin{array}{ccc}Q^{^{}}& =& Ae^xP_{}\\ P_{}^{^{}}& =& Ae^xQ\end{array}$$ (17) Consequently $`P_{}(x,i)`$ $`=`$ $`1{\displaystyle \underset{0}{\overset{x}{}}}A(s)e^s𝑑s+{\displaystyle \underset{0}{\overset{x}{}}}A(s)e^sP_{}(s,i){\displaystyle \underset{s}{\overset{\mathrm{}}{}}}A(\xi )e^\xi 𝑑\xi 𝑑s{\displaystyle \underset{x}{\overset{\mathrm{}}{}}}A(\xi )e^\xi 𝑑\xi {\displaystyle \underset{0}{\overset{x}{}}}A(s)e^sP_{}(s,i)𝑑s`$ And now it suffices to use the standard argument. Let $`M_n=\mathrm{max}_{x[0,n]}\|P_{}(x,i)\|=\|P_{}(x_n,i)\|`$ So $`M_n1+C+M_n\underset{0}{\overset{\mathrm{}}{}}\|A(s)\|e^s\left\|\underset{s}{\overset{\mathrm{}}{}}A(\xi )e^\xi 𝑑\xi \right\|𝑑s+\overline{o}(1)M_n`$ If the whole integral in the last formula is less then $`1`$ then $`M_n`$ is bounded. Otherwise we should start to solve the equations (17) not from zero but from some other point $`x_0`$ for which this condition is satisfied. Example. $`A(x)=(x^2+1)^\alpha \mathrm{sin}(x^\beta )`$ where $`\alpha ,\beta >0.`$ One can easily verify that conditions of Lemma are satisfied if $`2\alpha +\beta /2>1`$. Meanwhile $`A(x)L_2(R^+)`$ if and only if $`\alpha >1/4`$. Nevertheless for nonpositive $`A(x)`$ with bounded derivative the condition $`A(x)L^2(R^+)`$ is necessary for (1)-(4) from Section A to be true. Proposition. If one of the conditions (1)-(4) is true, $`A(x)0`$ and $`A^{^{}}(x)`$ is bounded then $`A(x)`$ is from $`L^2(R^+).`$ Proof. Really, since $`A(x)0`$ both $`P`$ and $`Q`$ are not less then 1. Consequently if one of (1)-(4) holds then $`P_{}(x,i)`$ is bounded and as it follows from (17) $`\underset{0}{\overset{x}{}}\|A(s)\|e^s\underset{s}{\overset{x}{}}\left\|A(\xi )\right\|e^\xi 𝑑\xi 𝑑s`$ is bounded as well. But we have the inequality $$\underset{0}{\overset{x}{}}\|A(s)\|e^s\underset{s}{\overset{x}{}}\left\|A(\xi )\right\|e^\xi 𝑑\xi 𝑑se^1\underset{0}{\overset{x1}{}}\|A(s)\|\underset{s}{\overset{s+1}{}}\left\|A(\xi )\right\|𝑑\xi 𝑑sC\underset{0}{\overset{x1}{}}\left\|A(s)\right\|^2𝑑s.$$ which concludes the proof of proposition. The latter inequality follows from the boundedness of $`A^{^{}}(x).`$ $`\mathrm{}`$ Function $`A(x)`$ from the example above with $`\alpha =1/4,1<\beta 3/2`$ satisfies the conditions of Lemma (consequently (1)-(4) holds ), has bounded derivative but is not from $`L^2(R^+).`$ The explanation is that this function is not nonpositive. We would like to conclude the paper with two open problems the first of which is much more difficult then the second one. Open problems. 1. Prove that Theorem 2 holds for $`WL^2(R^+)`$. 2. Prove that the presence of a.c. component on the half-line pertains to those potentials which Fourier transform is from $`L^2`$ near the zero. Specifically, if $`q`$ admits the Fourier transform $`\widehat{q}`$ such that $`\widehat{q}L_{loc}^2(R)`$ and $`\frac{\widehat{q}}{\|\omega \|+1}`$$`L^2(R),`$ then the a.c. part of the spectrum fills the whole positive half-line. This conjecture seems reasonable at least with some additional constraints since we can represent $`\widehat{q}=\widehat{q_1}+\widehat{q_2}`$. Where the $`\widehat{q_1}`$ is localized near the zero and is from $`L^2`$ so the methods of paper works. The other function $`\widehat{q_2}`$ is such that Theorem 3 can be applied. Acknowledgment. Author is grateful to C.Remling for attention to this work. e-mail address: saden@cs.msu.su
warning/0007/math0007093.html	ar5iv	text	# 1 Introduction ## 1 Introduction A well-known conjecture in the theory of Vassiliev invariants is that these invariants are dense in the space of all numerical knot invariants. This was posed as a problem in as follows: given any numerical knot invariant $`\varphi :L𝐐,`$ does there exist a sequence of Vassiliev invariants $`\{v_i:L𝐐,i=2,3,4,\mathrm{}\}`$ such that $$\underset{i\mathrm{}}{lim}v_i(L)=\varphi (L)$$ In this note, we find approximations by Vassiliev invariants for the coefficients of the Jones polynomial and all specializations of the HOMFLY and Kauffman polynomials. Consequently, we obtain approximations of some other link invariants. This note is organized as follows: In Section 2, we show that every Jones coefficient is the limit of a sequence of Vassiliev invariants. In Section 3, given any $`d`$, for any Jones polynomial of degree bounded by $`d`$, we find an explicit finite formula for its coefficients in terms of Vassiliev invariants. In Section 4, we find an explicit infinite approximation for any Jones coefficient. In Section 5, we extend the results to any specialization of the HOMFLY and Kauffman polynomials: we find the finite formula for polynomials of bounded degree and the infinite formula for all polynomials. In Section 6, we find approximations by Vassiliev invariants for some link invariants arising from the homology of branched covers of links. In Section 7, we discuss some conjectures related to approximations by Vassiliev invariants. ## 2 Approximating Jones coefficients by Vassiliev invariants: Existence theorem Let $`J_L(t)`$ denote the Jones polynomial of a knot $`L`$. Suppose $`J_L(t)=a_mt^m+\mathrm{}+a_0+\mathrm{}+a_nt^n`$, where $`a_m`$ and $`a_n`$ are nonzero. We call the degree of the Laurent polynomial $`d=\mathrm{max}(m,n)`$. In , it was shown that if we let $`t=e^x`$, $$J_L(e^x)=\underset{i=0}{\overset{\mathrm{}}{}}\left(\frac{1}{i!}\underset{k=m}{\overset{n}{}}k^ia_k(L)\right)x^i=\underset{i=0}{\overset{\mathrm{}}{}}v_i(L)x^i$$ (0.1) then $`v_i`$ is a Vassiliev invariant of order $`i`$. Henceforth, we will refer to $`\{v_i\}`$ as the Vassiliev invariants obtained from the coefficients of the expansion above. This can be reformulated in terms of the following infinite matrix: $$\left(\begin{array}{ccccccccc}\mathrm{}& 1& 1& 1& 1& 1& 1& 1& \mathrm{}\\ \mathrm{}& 3& 2& 1& 0& 1& 2& 3& \mathrm{}\\ \mathrm{}& (3)^2& (2)^2& (1)^2& 0& 1^2& 2^2& 3^2& \mathrm{}\\ & \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\end{array}\right)\left(\begin{array}{c}\mathrm{}\\ a_1\\ a_0\\ a_1\\ \mathrm{}\end{array}\right)=\left(\begin{array}{c}v_0\\ v_1\\ 2!v_2\\ 3!v_3\\ \mathrm{}\end{array}\right)$$ (0.2) Recall that a Vandermonde matrix has the following form (see, e.g., ): $$V(x_1,\mathrm{},x_r)=\left(\begin{array}{ccc}1& \mathrm{}& 1\\ x_1& \mathrm{}& x_r\\ x_1^2& \mathrm{}& x_r^2\\ \mathrm{}& & \mathrm{}\\ x_1^{r1}& \mathrm{}& x_r^{r1}\end{array}\right)$$ $$detV(x_1,\mathrm{},x_r)=\underset{i<j}{}(x_jx_i)$$ Thus, $`V(x_1,\mathrm{},x_r)`$ is invertible if $`x_ix_j`$ for all $`ij`$. The matrix in $`(\text{0.2})`$ is a Vandermonde matrix for every finite square block which contains the first row. From the resulting system of linear equations, we obtain the following existence theorem: ###### Theorem 0.1 Given any knot $`L`$, let $`J_L(t)`$ be the Jones polynomial, and $`a_i`$ be its $`i^{\mathrm{th}}`$ coefficient. Then for each $`i`$, $`a_i(L)`$ is the limit of a sequence of Vassiliev invariants. Proof. For any coefficient $`a_i`$, we will define a sequence of Vassiliev invariants $`\alpha _{1,i},\alpha _{2,i},\mathrm{}`$ and show that $`lim_n\mathrm{}\alpha _{n,i}(L)=a_i(L)`$. We now let $`t=e^x`$ and consider the expansion (0.1). If $`nd`$, we obtain the following system of linear equations: $$\{\begin{array}{ccccc}\hfill a_n& +\mathrm{}+a_1+a_0+a_1+\mathrm{}+& a_n\hfill & =\hfill & v_0\hfill \\ \hfill (n)a_n& +\mathrm{}+(1)a_1+a_1+\mathrm{}+& na_n\hfill & =\hfill & v_1\hfill \\ \hfill \mathrm{}& & & & \mathrm{}\hfill \\ \hfill (n)^ka_n& +\mathrm{}+(1)^ka_1+a_1+\mathrm{}+& n^ka_n\hfill & =\hfill & v_k\hfill \end{array}$$ (0.3) with $`2n+1`$ variables $`a_n,\mathrm{},a_n`$, and $`k+1`$ equations. When $`k=2n`$, the system of linear equations has a unique solution since the coefficient matrix is an invertible finite block of the Vandermonde matrix from (0.2). Denote the solution by the vector $`(\alpha _{n,n},\mathrm{},\alpha _{n,1},\alpha _{n,0},\alpha _{n,1},\mathrm{},\alpha _{n,n})`$. We claim $`lim_n\mathrm{}\alpha _{n,i}(L)=a_i(L)`$. The claim follows immediately from the following lemma: ###### Lemma 0.1 For any knot $`L`$, $`\alpha _{n,i}(L)=a_i(L)`$ for all $`nd`$. When $`nd`$, $`J_{L,n}=J_L`$. We let $`\alpha _n(L)=(\alpha _{n,n}(L),\mathrm{},\alpha _{n,0}(L),\mathrm{},\alpha _{n,n}(L))`$ be given by the coefficients of $`J_L(t)`$. Therefore, $`\alpha _n(L)`$ is the unique solution satisfying the above system of linear equations. If we now fix $`i`$, for all $`nd`$, it follows that $`\alpha _{n,i}(L)=a_i(L)`$. This completes the proof of the theorem. ###### Corollary 0.1 For any knot $`L`$ and any fixed complex number $`z`$, $`J_L(z)`$ is a limit of Vassiliev invariants. Proof. For each $`n`$, let $$\alpha _n(L)=(\mathrm{},0,\alpha _{n,n}(L),\mathrm{},\alpha _{n,1}(L),\alpha _{n,0}(L),\alpha _{n,1}(L),\mathrm{},\alpha _{n,n}(L),0,\mathrm{})$$ be the sequence of infinite vectors as defined in Theorem 0.1. Let $$g_n(L)=\alpha _{n,n}(L)z^n+\mathrm{}+\alpha _{n,0}(L)+\mathrm{}+\alpha _{n,n}(L)z^n.$$ Then $`g_n`$ is a Vassiliev invariant since it is a linear combination of such invariants. By Lemma 0.1, when $`nd`$, we have $`\alpha _{n,i}(L)=a_i(L)`$ for all $`i`$. Thus $`g_n(L)=J_L(z)`$. ## 3 Approximating Jones coefficients by Vassiliev invariants: Bounded degree case For any given $`d`$ (in particular, for any given knot), we can obtain explicit solutions to the linear system (0.3) and obtain a formula for all $`\alpha _{n,i}`$. We will use (0.2) and compute the inverse of the $`(2d+1)\times (2d+1)`$ Vandermonde matrix which is symmetric about the column with zeros. For this, we need the following generating function: ###### Definition 0.1 $$f_{d,n}(v)=\underset{\genfrac{}{}{0pt}{}{jn}{j=d}}{\overset{d}{}}\frac{vj}{nj}=\frac{(1)^{n+d}}{(d+n)!(dn)!}\underset{\genfrac{}{}{0pt}{}{jn}{j=d}}{\overset{d}{}}(vj)$$ The proof of the following proposition is immediate from the definition: ###### Proposition 0.1 For any $`m𝐙`$ such that $`dmd`$, $$f_{d,n}(m)=\{\begin{array}{cc}1\hfill & \text{if }m=n\hfill \\ 0\hfill & \text{if }mn\hfill \end{array}$$ ###### Theorem 0.2 For any Jones polynomial of a knot of degree $`d`$, $$a_n=\underset{i=0}{\overset{2d}{}}f_{d,n}^{(i)}(0)v_i,\mathrm{where}f_{d,n}(v)=\underset{\genfrac{}{}{0pt}{}{jn}{j=d}}{\overset{d}{}}\frac{vj}{nj}$$ In other words, $`\alpha _{d,n}=\underset{i=0}{\overset{2d}{}}f_{d,n}^{(i)}(0)v_i`$. Proof. Let $`c_{n,j}=\frac{1}{j!}f_{d,n}^{(j)}(0)`$, the $`j^{th}`$ coefficient of the polynomial $`f_{d,n}(v)`$. $$\left(\begin{array}{ccc}c_{d,0}& \mathrm{}& c_{d,2d}\\ \mathrm{}& & \mathrm{}\\ c_{d,0}& \mathrm{}& c_{d,2d}\end{array}\right)\left(\begin{array}{ccccc}1& \mathrm{}& 1& \mathrm{}& 1\\ d& \mathrm{}& 0& \mathrm{}& d\\ \mathrm{}& & \mathrm{}& & \mathrm{}\\ (d)^{2d}& \mathrm{}& 0& \mathrm{}& d^{2d}\end{array}\right)=$$ $$\left(\begin{array}{ccc}f_{d,d}(d)& \mathrm{}& f_{d,d}(d)\\ \mathrm{}& & \mathrm{}\\ f_{d,d}(d)& \mathrm{}& f_{d,d}(d)\end{array}\right)=I_{(2d+1)\times (2d+1)}$$ Therefore, $$\left(\begin{array}{c}a_d\\ \mathrm{}\\ a_0\\ \mathrm{}\\ a_d\end{array}\right)=\left(\begin{array}{ccc}c_{d,0}& \mathrm{}& c_{d,2d}\\ \mathrm{}& & \mathrm{}\\ c_{d,0}& \mathrm{}& c_{d,2d}\end{array}\right)\left(\begin{array}{c}v_0\\ v_1\\ 2!v_2\\ \mathrm{}\\ (2d)!v_{2d}\end{array}\right)$$ With some extra notation, we can state the theorem more succinctly. Let $`𝒱_i`$ be the vector space of Vassiliev invariants spanned by $`v_i`$ from $`(\text{0.1})`$. Consider the underlying vector space of the polynomial algebra $`𝐐[v]`$ with basis $`\{v^0,v^1,v^2,\mathrm{}\}`$. Let $`E`$ be the vector space isomorphism $`E:𝐐[v]𝒱_i`$, where $`E(v^i)=i!v_i`$. We therefore obtain: $$a_n=E(f_{d,n}(v))$$ ###### Example 0.1 As in Section 2, let $`\alpha _{d,n}`$ denote the $`n^{th}`$ coefficient of the Jones polynomial of degree $`d`$. $$\begin{array}{ccc}\hfill \alpha _{2,0}& =& \frac{1}{2!2!}(45v_2+v_4)=\frac{1}{2!2!}E\left((v2)(v1)(v+1)(v+2)\right)\hfill \\ & & \\ \hfill \alpha _{3,0}& =& \frac{1}{3!3!}(3649v_2+14v_4v_6)=\hfill \\ & & \\ & =& \frac{1}{3!3!}E\left((v3)(v2)(v1)(v+1)(v+2)(v+3)\right)\hfill \\ & & \\ \hfill \alpha _{2,1}& =& \frac{1}{1!3!}(4v_1+4v_2v_3v_4)=\frac{1}{1!3!}E\left((v2)(v)(v+1)(v+2)\right)\hfill \\ & & \\ \hfill \alpha _{3,1}& =& \frac{1}{2!4!}(36v_1+36v_213v_313v_4+v_5+v_6)=\hfill \\ & & \\ & =& \frac{1}{2!4!}E\left((v3)(v2)(v)(v+1)(v+2)(v+3)\right)\hfill \end{array}$$ ###### Remark 0.1 From Theorem 0.2, we obtain a formula for any $`v_i(L)`$ in terms of $`v_0(L),\mathrm{},v_{2d}(L)`$. Namely, $$v_i(L)=\frac{1}{i!}\underset{k=m}{\overset{n}{}}k^ia_k(L)=\underset{j=0}{\overset{2d}{}}\left(\frac{1}{i!}\underset{k=m}{\overset{n}{}}k^if_{d,k}^{(j)}(0)\right)v_j(L)$$ Similarly, in it was shown, without an explicit formula, that for any knot $`L`$, there exists $`N,`$ such that all $`v_i(L)`$ are determined by $`v_0(L),\mathrm{},v_N(L)`$. ## 4 Approximating Jones coefficients by Vassiliev invariants: Infinite case In this section, we formally let $`d\mathrm{}`$ to find the correct formula for the coefficients of an arbitrary Jones polynomial of a knot, and then prove that the resulting series of Vassiliev invariants converges. We also extend the vector space isomorphism $`E:𝐐[[v]]𝒱_i`$, where $`E(v^i)=i!v_i`$. We first consider $`a_0`$: $$f_{d,0}(v)=\underset{\genfrac{}{}{0pt}{}{j0}{j=d}}{\overset{d}{}}\frac{vj}{j}=\underset{\genfrac{}{}{0pt}{}{j0}{j=d}}{\overset{d}{}}\left(1\frac{v}{j}\right)=\underset{j=1}{\overset{d}{}}\left(1\frac{v}{j}\right)\left(1+\frac{v}{j}\right)=\underset{j=1}{\overset{d}{}}\left(1\frac{v^2}{j^2}\right)$$ We recall the Weierstrass product factorization for entire functions: $$\underset{j=1}{\overset{\mathrm{}}{}}\left(1\frac{z^2}{j^2}\right)=\frac{\mathrm{sin}\pi z}{\pi z}$$ We define $`f_{\mathrm{},0}(v)=lim_d\mathrm{}f_{d,0}(v)`$, so we obtain $$f_{\mathrm{},0}(v)=\frac{\mathrm{sin}\pi v}{\pi v}=1\frac{(\pi v)^2}{3!}+\frac{(\pi v)^4}{5!}\frac{(\pi v)^6}{7!}+\mathrm{}$$ Formally (we prove convergence below), we obtain the following beautiful formula: $$a_0=E(f_{\mathrm{},0}(v))=E(\frac{\mathrm{sin}\pi v}{\pi v})=v_0\frac{\pi ^2}{3}v_2+\frac{\pi ^4}{5}v_4\frac{\pi ^6}{7}v_6+\mathrm{}$$ (0.4) We now consider $`a_n`$: $$f_{d,n}(v)=\underset{\genfrac{}{}{0pt}{}{jn}{j=d}}{\overset{d}{}}\frac{vj}{nj}=\underset{j=d}{\overset{n1}{}}\frac{vj}{nj}\underset{j=n+1}{\overset{d}{}}\frac{vj}{nj}$$ Let $`k=nj`$ in the first product, and $`k=jn`$ in the second product, so we obtain $$f_{d,n}(v)=\underset{k=1}{\overset{d+n}{}}\frac{v+kn}{k}\underset{k=1}{\overset{dn}{}}\frac{vkn}{k}=\underset{k=1}{\overset{dn}{}}(1\frac{(vn)^2}{k^2})\underset{k=dn+1}{\overset{d+n}{}}(1+\frac{vn}{k})$$ The second product is finite: let $`l=kd`$, then we obtain $`_{l=n+1}^n(1+\frac{vn}{l+d})`$. For any $`n`$, as $`d\mathrm{}`$ we can easily see that this product converges to $`1`$. The first product converges for all $`n`$ by the same argument as above. This suggests the following theorem: ###### Theorem 0.3 $$a_n=E(f_{\mathrm{},n}(v))=\underset{i=0}{\overset{\mathrm{}}{}}f_{\mathrm{},n}^{(i)}(0)v_i,\mathrm{where}f_{\mathrm{},n}(v)=\{\begin{array}{cc}\frac{\mathrm{sin}\pi (vn)}{\pi (vn)}\hfill & \text{if }vn\hfill \\ 1\hfill & \text{if }v=n\hfill \end{array}$$ Proof. For any given knot $`L`$, the Jones polynomial has finite degree $`d`$. $$\underset{i=0}{\overset{\mathrm{}}{}}f_{\mathrm{},n}^{(i)}(0)v_i(L)=\underset{i=0}{\overset{\mathrm{}}{}}f_{\mathrm{},n}^{(i)}(0)\left(\frac{1}{i!}\underset{k=d}{\overset{d}{}}k^ia_k(L)\right)$$ $$=\underset{k=d}{\overset{d}{}}a_k(L)\left(\underset{i=0}{\overset{\mathrm{}}{}}\frac{1}{i!}f_{\mathrm{},n}^{(i)}(0)k^i\right)$$ $$=\underset{k=d}{\overset{d}{}}a_k(L)f_{\mathrm{},n}(k)=a_n(L)$$ ###### Example 0.2 $$a_1=v_1+2v_2+(3!\pi ^2)v_3+(4!4\pi ^2)v_4+(5!+\pi ^420\pi ^2)v_5+\mathrm{}$$ ###### Corollary 0.2 Any Jones coefficient of a knot can be approximated by Vassiliev invariants $`\stackrel{~}{v}_i`$ of order $`i`$: $$a_n=\underset{i\mathrm{}}{lim}\stackrel{~}{v}_i,\mathrm{where}\stackrel{~}{v}_i=\underset{j=0}{\overset{i}{}}f_{\mathrm{},n}^{(j)}(0)v_j$$ ###### Remark 0.2 Jones coefficients can be shown not to be Vassiliev invariants by considering twist sequences . Let $`T_m`$ be the $`(2,2m+1)`$-torus knot. By Theorem 2.2.1 of , the restriction of any Vassiliev invariant to the sequence $`\{T_m\}`$ is a polynomial in $`m`$. However, $`J_{T_m}(t)=t^m(t^{2m+1}\mathrm{}+t^3t^21)`$. Thus for $`n0`$, $`a_n(T_m)=0`$ for $`m<\frac{n1}{3}`$ or $`m>n`$. If $`a_n(T_m)`$ were a polynomial in $`m`$, it would be zero on $`T_m`$, which is clearly false. Similarly, we can take mirror images of $`T_m`$ to show that for all $`n<0`$, $`a_n`$ is not a Vassiliev invariant. (See also .) Trapp also showed that if a sequence of Vassiliev invariants converges uniformly for all knots, then the limit is also of finite type. Because the Jones coefficients are not of finite type, the pointwise limits above cannot be uniformly convergent for all knots. Indeed, the proof of Theorem 0.3 requires us to first choose a particular knot. ###### Remark 0.3 The function $`f_{\mathrm{},n}(v)`$ is not unique, because the infinite Vandermonde matrix can have infinitely many left inverses. Since every Jones polynomial has finite degree $`d`$, but Vassiliev invariants may be nonzero for arbitrary orders, we can view the infinite matrix as a linear operator $`_{i=\mathrm{}}^{\mathrm{}}𝐐_{i=0}^{\mathrm{}}𝐐`$, so it can have infinitely many left inverses, but no right inverse. The proof of Theorem 0.3 only requires that $`f_{\mathrm{},n}(v)`$ has a Taylor expansion about zero, and that $`f_{\mathrm{},n}(m)=\delta _{m,n}`$, the Kronecker pairing for all $`m,n𝐙`$. If we also insist that $`f_{\mathrm{},n}(v)=lim_d\mathrm{}f_{d,n}(v)`$, then the generating function depends on how we select invertible finite blocks to exhaust the infinite Vandermonde matrix. ###### Remark 0.4 Given an infinite sequence $`\{v_n(L)\}`$, we can also approximate the degree $`d`$ of the Jones polynomial by functions of finite type invariants: $$d=\underset{n\mathrm{}}{lim}\sqrt[n]{\underset{k=d}{\overset{d}{}}k^na_k(L)}=\underset{n\mathrm{}}{lim}\sqrt[n]{n!\|v_n\|}$$ ###### Remark 0.5 All of the results above carry over with slight modifications to links. The Jones polynomial of a link may be a Laurent polynomial times $`\sqrt{t}`$, so instead of (0.2), the matrix and resulting formulas appear as a special case of Theorem 0.4. ## 5 Approximations of coefficients of specializations of HOMFLY and Kauffman polynomials The Jones polynomial is a specialization of both the HOMFLY and Kauffman two-variable polynomials, $`H_L(a,z)`$ and $`F_L(a,z)𝐙[a^{\pm 1},z^{\pm 1}]`$. We consider an infinite sequence of one-variable specializations of the HOMFLY polynomial. The same result and proof applies to the Kauffman polynomial as well. Let $`N𝐙\backslash \{0\}`$. $$a=t^{N/2},z=t^{1/2}t^{1/2}H_L^N(t)Z[t^{\pm \frac{1}{2}}]$$ (0.5) The Jones polynomial is obtained at $`N=2`$. Now, suppose $`H_L^N(t)=b_m^Nt^{m/2}+\mathrm{}+b_0^N+\mathrm{}+b_n^Nt^{n/2}.`$ Let $`d=\mathrm{max}(m,n)`$. In , it was shown that if we let $`t=e^x`$, $$H_L^N(e^x)=\underset{i=0}{\overset{\mathrm{}}{}}\left(\frac{1}{i!}\underset{k=m}{\overset{n}{}}\left(\frac{k}{2}\right)^ib_k^N(L)\right)x^i=\underset{i=0}{\overset{\mathrm{}}{}}v_i^N(L)x^i$$ (0.6) then $`v_i^N`$ is a Vassiliev invariant of order $`i`$. ###### Theorem 0.4 For any link $`L`$ and any $`N𝐙\backslash \{0\}`$, let its $`N^{th}`$ HOMFLY polynomial be $`H_L^N(t)=_{n=d}^db_n^Nt^{n/2}`$. Let $`E^N(v^i)=i!v_i^N`$. Then, $$b_n^N=E^N(f_{d,n}(v))=\underset{i=0}{\overset{2d}{}}f_{d,n}^{(i)}(0)v_i^N,\mathrm{where}f_{d,n}(v)=\underset{\genfrac{}{}{0pt}{}{jn}{j=d}}{\overset{d}{}}\frac{2vj}{nj}$$ $$b_n^N=E^N(f_{\mathrm{},n}(v))=\underset{i=0}{\overset{\mathrm{}}{}}f_{\mathrm{},n}^{(i)}(0)v_i^N,\mathrm{where}f_{\mathrm{},n}(v)=\{\begin{array}{cc}\frac{\mathrm{sin}\pi (2vn)}{\pi (2vn)}\hfill & \text{if }2vn\hfill \\ 1\hfill & \text{if }2v=n\hfill \end{array}$$ Proof. The proof is just a modification of the proof for the Jones polynomial. Instead of (0.2), we have $$\left(\begin{array}{ccccccc}\mathrm{}& 1& 1& 1& 1& 1& \mathrm{}\\ \mathrm{}& 1& 1/2& 0& 1/2& 1& \mathrm{}\\ \mathrm{}& (1)^2& (1/2)^2& 0& (1/2)^2& 1^2& \mathrm{}\\ & \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\end{array}\right)\left(\begin{array}{c}\mathrm{}\\ b_1^N\\ b_0^N\\ b_1^N\\ \mathrm{}\end{array}\right)=\left(\begin{array}{c}v_0^N\\ v_1^N\\ 2!v_2^N\\ 3!v_3^N\\ \mathrm{}\end{array}\right)$$ To find the inverse of both the infinite matrix and any invertible finite block, we just need the following: ###### Proposition 0.2 For any $`m𝐙`$ such that $`dmd`$, $$f_{d,n}(\frac{m}{2})=\{\begin{array}{cc}1\hfill & \text{if }m=n\hfill \\ 0\hfill & \text{if }mn\hfill \end{array}$$ and similarly for $`f_{\mathrm{},n}(\frac{m}{2})`$ for any $`m𝐙`$. ###### Corollary 0.3 For any link $`L`$, any fixed complex number $`z`$, and any $`N𝐙\backslash \{0\}`$, $`H_L^N(z)`$ is a limit of Vassiliev invariants. Proof. The proof is the same as the proof of Corollary 0.1. ###### Remark 0.6 Following , let $`F_L(a,z)`$ denote the Dubrovnik version of the Kauffman polynomial. To obtain one-variable specializations, for $`N𝐙\backslash \{0\}`$, set $`a=t^N,z=tt^1`$ and let $`F_L^N(t)`$ denote the $`N^{th}`$ Kauffman polynomial. Theorem 0.4 and Corollary 0.3 carry over to $`F_L^N(t)`$. ## 6 Approximations of other link invariants Let $`\mathrm{\Sigma }_n(L)`$ denote the n-fold branched cover of a link $`L`$. Let $`Q_L(x)`$ be the specialization at $`a=1`$ of the standard Kauffman polynomial. This link polynomial satisfies the skein relation $`Q_{L_+}+Q_L_{}=x(Q_{L_0}+Q_L_{\mathrm{}})`$ and $`Q(\text{unknot})=1`$ . A lot of information about $`H_1(\mathrm{\Sigma }_n(L))`$ can be obtained by evaluating link polynomials of $`L`$ at special values. In fact, it seems reasonable to conjecture that $`H_1(\mathrm{\Sigma }_2(L),𝐙)`$ is determined by $`Q_L`$ . We summarize the results below. All the evaluations can be found in , except for $`Q_L(2\mathrm{cos}\frac{2\pi }{5})`$ which is given in and . For the Jones polynomial, the most interesting values to evaluate are $`t=e^{\frac{2\pi i}{r}}`$, where $`r`$ is a positive integer. In the table, $`\mathrm{}`$ is the number of components of $`L`$, Arf$`(L)`$ is the Arf invariant of $`L`$, and $`w𝐙_4`$ is the Witt class of the Seifert form mod 3 of $`L`$. Because the Jones polynomial is a specialization of the HOMFLY polynomial, the table also gives evaluations of the HOMFLY polynomial. Another interesting value not listed above is $`H_L(i,i)=(i\sqrt{2})^{dim(H_1(\mathrm{\Sigma }_3,𝐙_2))}`$. Let us also recall that $`\|H_1(\mathrm{\Sigma }_n,𝐙)\|=\|\mathrm{\Delta }_L(r_i)\|`$, where $`r_i`$’s are the $`n`$th roots of unity. For the $`Q`$-polynomial, the interesting values are at $`x=2\mathrm{cos}\frac{2\pi }{r}=q+q^1`$, where $`q=e^{\frac{2\pi i}{r}}`$. Here $`t_5`$ is $`0`$ or $`1`$, and can be written in terms of the Seifert form mod $`5`$ . ###### Theorem 0.5 Let $`\varphi `$ be any of the following knot invariants: $`detL,\|H_1(\mathrm{\Sigma }_2,𝐙_p)\|`$ where $`p=3,5`$, $`\|H_1(\mathrm{\Sigma }_3,𝐙_2)\|`$, $`\|H_1(\mathrm{\Sigma }_n,𝐙)\|`$. Then (a) $`\varphi `$ is not a Vassiliev invariant. (b) $`\varphi `$ is a limit of functions of Vassiliev invariants. Proof. (a) As remarked in , a knot invariant $`\varphi `$ is not a Vassiliev invariant if there is a knot $`K`$ with $`\varphi (K\mathrm{\#}K^{})\varphi (\text{unknot})`$ for all knots $`K^{}`$. Now let $`\varphi `$ be, for example, the order of $`H_1(\mathrm{\Sigma }_2,𝐙_p)`$. Let $`K`$ be a knot with $`\varphi (K)=p`$ (e.g., any 2-bridge knot for which $`\mathrm{\Sigma }_2(K)`$ is the lens space $`L_{p,q}`$). Then $`\varphi (K\mathrm{\#}K^{})=\varphi (K)\varphi (K^{})\varphi (\text{unknot})`$. Similarly, $`\|H_1(\mathrm{\Sigma }_3,𝐙_2)\|`$ and $`\|H_1(\mathrm{\Sigma }_n,𝐙)\|`$ are not Vassiliev invariants. (b) The following equations come from the tables and comments above. Together with Corollary 0.1 and Corollary 0.3, they imply that $`\varphi `$ is a limit of functions of Vassiliev invariants. $$\begin{array}{ccc}\|H_1(\mathrm{\Sigma }_2,𝐙)\|\hfill & =& J_L(1)\hfill \\ \|H_1(\mathrm{\Sigma }_2,𝐙_3)\|\hfill & =& \|J_L(e^{\frac{\pi i}{3}})\|^2\hfill \\ \|H_1(\mathrm{\Sigma }_2,𝐙_5)\|\hfill & =& \|Q_L(2\mathrm{cos}(2\pi /5))\|^2\hfill \\ \|H_1(\mathrm{\Sigma }_3,𝐙_2)\|\hfill & =& \|H_L(i,i)\|^2\hfill \\ \|H_1(\mathrm{\Sigma }_n,𝐙)\|\hfill & =& \|\mathrm{\Delta }_L(r_i)\|,\text{ where }r_i\text{’s are the }n\text{th roots of unity.}\hfill \end{array}$$ Note from (0.5), we obtain that $`H_L(i,i)`$ is $`H_L^N(t)`$, where $`N=9,t=e^{i\pi /3}`$. As in Section 5, let $`F_L^N(t)`$ denote the $`N^{th}`$ Dubrovnik Kauffman polynomial. By a change of variables, $`Q_L(x)=(1)^{\mathrm{}}F_L(i,ix)`$ . Thus, by Remark 0.6 we obtain that $`Q_L(2\mathrm{cos}(2\pi /5))`$ is $`(1)^{\mathrm{}}F_L^N(t)`$, where $`N=5,t=ie^{2i\pi /5}`$. Since the coefficients of the Alexander-Conway polynomial are Vassiliev invariants , by an argument similar to Corollary 0.1 we obtain that $`\mathrm{\Delta }_L(r_i)`$ is a limit of Vassiliev invariants. Thus, $`\|\mathrm{\Delta }_L(r_i)\|`$ is a limit of the absolute value function of Vassiliev invariants. ###### Remark 0.7 For all $`\varphi `$ except for $`\varphi =\|H_1(\mathrm{\Sigma }_n,𝐙)\|`$, we can show that $`\varphi `$ is actually a limit of Vassiliev invariants. For example, in the case of $`\|H_1(\mathrm{\Sigma }_2,𝐙_3)\|`$, it follows from Corollary 0.1 and the equation: $`\|H_1(\mathrm{\Sigma }_2,𝐙_3)\|=J_L(e^{\frac{\pi i}{3}})J_L(e^{\frac{\pi i}{3}})`$, since the two factors on the right are complex conjugates. For other functions, e.g., $`\varphi =dim(H_1(\mathrm{\Sigma }_2,𝐙_3))`$, the argument above can only show that $`\varphi `$ is a limit of functions of Vassiliev invariants. Note that in the case of $`\|H_1(\mathrm{\Sigma }_n,𝐙)\|`$, the $`n`$-th roots of unity all appear in conjugate pairs, thus the product $`\mathrm{\Delta }_L(r_i)`$ can be arranged in conjugate pairs. It follows that the absolute value sign is not needed, except for $`\mathrm{\Delta }(1)\mathrm{\Delta }(1)`$. ###### Remark 0.8 Another knot invariant which is not a Vassiliev invariant, but is a limit of Vassiliev invariants is $`tri(L)`$, the number of 3-colorings of $`L`$. This follows from a result of Przytycki : $`tri(L)=3\|J_L(e^{\frac{\pi i}{3}})\|^2`$. ## 7 Conclusion Here, we make some final remarks on approximations by Vassiliev invariants. Our work is motivated by the following two equivalent conjectures (see ): Conjecture 7.1. Vassiliev invariants separate knots. That is, for any two knots $`K_1`$ and $`K_2`$, there is a Vassiliev invariant $`v`$ with $`v(K_1)v(K_2)`$. Conjecture 7.2. Every knot invariant is a limit of Vassiliev invariants. That is, for any knot invariant $`\varphi `$, there is a sequence of Vassiliev invariants $`v_n`$ with $`\varphi (L)=lim_n\mathrm{}v_n(L)`$ for all knots $`L`$. We have verified that a number of knot invariants (e.g., coefficients of link polynomials) are indeed limits of Vassiliev invariants. Some other knot invariants (e.g., the degree of the Jones polynomial) are proved to be limits of functions of Vassiliev invariants. In light of this, we propose: Conjecture 7.3. Every knot invariant is a limit of functions of Vassiliev invariants. That is, for any knot invariant $`\varphi `$, there is a sequence of functions $`f_n`$ and Vassiliev invariants $`v_n`$ with $`\varphi (L)=lim_n\mathrm{}f_n(v_n)(L)`$ for all knots $`L`$. In general, if $`f`$ is an analytic function, $`v`$ is a Vassiliev invariant, then it is not hard to show that $`f(v)`$ is a limit of Vassiliev invariants. Consequently, a limit of analytic functions of Vassiliev invariants is in fact a limit of Vassiliev invariants. However, this is not clear if the functions are not analytic functions. This is the case, for example, for the degree of the Jones polynomial, where the functions $`f_n`$ are the $`n`$th root function. One good aspect of Conjecture 7.3 is that it is easier to verify than Conjecture 7.2 for a given knot invariant, but is still strong enough to imply Conjecture 7.1. Therefore by , $$\mathrm{Conjecture}\mathrm{\hspace{0.17em}7.2}\mathrm{Conjecture}\mathrm{\hspace{0.17em}7.3}\mathrm{Conjecture}\mathrm{\hspace{0.17em}7.1}\mathrm{Conjecture}\mathrm{\hspace{0.17em}7.2}.$$ Acknowledgements We would like to thank Dror Bar-Natan, Xiao-Song Lin, Ted Stanford, and Ed Swartz for helpful discussions. The first author was partially supported by NSF grant DMS-98-03518. The second author was partially supported by NSF grant DMS-97-29992 while visiting the Institute for Advanced Study. References
warning/0007/math0007114.html	ar5iv	text	# Three alternating sign matrix identities in search of bijective proofs ## 1 Introduction Alternating sign matrices (ASMs) are square matrices of 0s, 1s, and $`1`$s with row and column-sums equal to 1 and with the restriction that the non-zero entries alternate signs across each row and down each column. An example is $$\left(\begin{array}{ccccc}0& 1& 0& 0& 0\\ 1& 1& 0& 1& 0\\ 0& 1& 0& 1& 1\\ 0& 0& 0& 1& 0\\ 0& 0& 1& 0& 0\end{array}\right)$$ These are rich combinatorial objects with connections to many problems in algebraic combinatorics (see , , ). They also have many different representations. The representation that was used in Kuperberg’s proof of the counting function for alternating sign matrices and Zeilberger’s proof of the refined alternating sign matrix conjecture is the six-vertex model of statistical mechanics. These are directed graphs in which each vertex has in-degree two and out-degree two, and boundary conditions that the vertical arrows along the top and bottom are directed out, horizontal arrows along the left and right are directed in, as in the following directed graph. $$\begin{array}{ccccccccccc}& & & & & & & & & & \\ & & & & & & & & & & \\ & & & & & & & & & & \\ & & & & & & & & & & \\ & & & & & & & & & & \\ & & & & & & & & & & \\ & & & & & & & & & & \\ & & & & & & & & & & \\ & & & & & & & & & & \\ & & & & & & & & & & \\ & & & & & & & & & & \end{array}$$ To actually make this a directed graph on 25 vertices, we can identify the $`i`$th up-arrow along the top row with the $`i`$th right arrow along the left edge, and similarly identify bottom and right arrows. This is called a six-vertex model because there are six possible configurations at each vertex. We shall describe a vertex as horizontal if both in-edges are horizontal, vertical if both in-edges are vertical, and otherwise southwest, northwest, northeast, or southeast, according to the direction of the sum of the four vectors represented by the four adjacent edges. It should be noted that the sum of all vertical vectors is zero, as is the sum of all horizontal vectors. It follows that there will always be an equal number of southwest and northeast vertices, and an equal number of southeast and northwest vertices. Our example of a six-vertex model corresponds to our example of an alternating sign matrix. Each 1 in the ASM corresponds to a horizontal vertex, each $`1`$ to a vertical vertex, and the 0s to the other vertices. This is a bijection because once the positions of the horizontal and vertical vertices are known, all other vertices are uniquely determined. The six-vertex model is not the only insightful representation, but it is very suggestive, especially because there is also a natural connection between ASMs and complete directed graphs or tournaments. It would be very useful to have a direct bijective connection between ASMs and tournaments. In explaining the bijection that we seek, we shall also present two other related identities that cry out for bijective proofs. ## 2 The $`\lambda `$-determinant The first two identities that I wish to present arise from the $`\lambda `$-determinant of Robbins and Rumsey . This is based on the Desnanot–Jacobi adjoint matrix theorem , that was used by Dodgson to create his algorithm for evaluating determinants. Given a square matrix $`M`$, we let $`M_j^i`$ denote $`M`$ with row $`i`$ and column $`j`$ deleted. We then have that $$detM=\frac{detM_1^1detM_n^ndetM_1^ndetM_n^1}{detM_{1,n}^{1,n}}.$$ (1) If we define the determinant of an empty matrix ($`0\times 0`$) to be 1 and the determinant of the $`1\times 1`$ matrix $`(a)`$ to be $`a`$, then equation (1) can be used as a recursive definition of the determinant. A natural one-parameter generalization of the determinant arises if we use the same initial conditions and replace the minus sign in the numerator of the recursive step by $`+\lambda `$: $$\underset{\lambda }{det}(M)=\frac{\underset{\lambda }{det}\left(M_1^1\right)\underset{\lambda }{det}\left(M_n^n\right)+\lambda \underset{\lambda }{det}\left(M_1^n\right)\underset{\lambda }{det}\left(M_n^1\right)}{\underset{\lambda }{det}\left(M_{1,n}^{1,n}\right)}.$$ (2) The following generalization of the Vandermonde determinant evaluation follows by induction. ###### Proposition 1 $$\underset{\lambda }{det}\left(x_i^{nj}\right)=\underset{1i<jn}{}(x_i+\lambda x_j).$$ (3) If we expand a few $`\lambda `$-determinants, an interesting pattern emerges: $`\underset{\lambda }{det}\left(\begin{array}{ccc}a& b& c\\ d& e& f\\ g& h& i\end{array}\right)`$ $`=`$ $`aei+\lambda (bdi+afh)+\lambda ^2(bfg+cdh)+\lambda ^3ceg+\lambda (1+\lambda )bde^1fh,`$ $`\underset{\lambda }{det}\left(\begin{array}{cccc}a& b& c& d\\ e& f& g& h\\ i& j& k& l\\ m& n& o& p\end{array}\right)`$ $`=`$ $`\mathrm{}+\lambda ^3(1+\lambda )bef^1hkn+\lambda ^3(1+\lambda )^2cfg^1hij^1kn+\mathrm{}.`$ The monomials in roman letters that correspond to permutation matrices are each multiplied by $`\lambda `$ raised to the inversion number of the permutation. The other monomials in roman letters that appear, such as $`cfg^1hij^1kn`$, correspond to alternating sign matrices, in this case $$\left(\begin{array}{cccc}0& 0& 1& 0\\ 0& 1& 1& 1\\ 1& 1& 1& 0\\ 0& 1& 0& 0\end{array}\right).$$ Each of these monomials is multiplied by a power of $`\lambda `$ and a power of $`1+\lambda `$. Let $`𝒜_n`$ be the set of $`n\times n`$ ASMs. Given $`A=(a_{ij})𝒜_n`$, we define its inversion number, $`(A)`$, to be $$(A)=\underset{i<k,j>l}{}a_{ij}a_{kl}.$$ We define $`N(A)`$ to be the number of $`1`$s in $`A`$. The following characterization of the $`\lambda `$-determinant was published by Robbins and Rumsey in 1986 . ###### Proposition 2 $$\underset{\lambda }{det}(m_{ij})=\underset{A𝒜_n}{}\lambda ^{(A)N(A)}(1+\lambda )^{N(A)}\underset{i,j=1}{\overset{n}{}}m_{ij}^{a_{ij}}.$$ (6) Zeilberger has given a bijective proof of equation (1). It would be desirable to have a direct proof of Proposition 2 by finding a similar proof of equation (2) when the $`\lambda `$-determinant is defined by the right side of Proposition 2. ###### Problem 1 Find a direct, bijective proof of the following identity. Within each summation, the range of indices for the alternating sign matrices $`B`$ and $`C`$ is specified by the product term. $`{\displaystyle \underset{(B,C)𝒜_n\times 𝒜_{n2}}{}}\lambda ^{(B)+(C)N(B)N(C)}(1+\lambda )^{N(B)+N(C)}{\displaystyle \underset{i,j=1}{\overset{n}{}}}m_{ij}^{b_{ij}}{\displaystyle \underset{i,j=2}{\overset{n1}{}}}m_{ij}^{c_{ij}}`$ $`=`$ $`{\displaystyle \underset{(B,C)𝒜_{n1}\times 𝒜_{n1}}{}}\lambda ^{(B)+(C)N(B)N(C)}(1+\lambda )^{N(B)+N(C)}`$ $`\times {\displaystyle \underset{\genfrac{}{}{0pt}{}{1i<n}{1j<n}}{}}m_{ij}^{b_{ij}}{\displaystyle \underset{\genfrac{}{}{0pt}{}{1<in}{1<jn}}{}}m_{ij}^{c_{ij}}`$ $`+\lambda {\displaystyle \underset{(B,C)𝒜_{n1}\times 𝒜_{n1}}{}}\lambda ^{(B)+(C)N(B)N(C)}(1+\lambda )^{N(B)+N(C)}`$ $`\times {\displaystyle \underset{\genfrac{}{}{0pt}{}{1i<n}{1j<n}}{}}m_{ij}^{b_{ij}}{\displaystyle \underset{\genfrac{}{}{0pt}{}{1<in}{1<jn}}{}}m_{ij}^{c_{ij}}.`$ ## 3 Directed Graphs If we combine Propositions 1 and 2, we get that $$\underset{1i<jn}{}(x_i+\lambda x_j)=\underset{A𝒜_n}{}\lambda ^{(A)N(A)}(1+\lambda )^{N(A)}\underset{i,j=1}{\overset{n}{}}x_i^{(nj)a_{ij}}$$ (7) It is worth noting that analogs of this identity for other root systems have been found by Okada . The left side of equation (7) can be interpreted as a sum over the set of tournaments on $`n`$ vertices, $`calT_n`$. Each binomial $`x_i+\lambda x_j`$ corresponds to the edge between vertices $`i`$ and $`j`$. If the edge is directed from $`i`$ to $`j`$, we choose $`x_i`$. If it is directed from $`j`$ to $`i`$, we choose $`\lambda x_j`$. Each tournament corresponds to a monomial in which the power of $`x_i`$ is $`\omega (i)`$, the out-degree of vertex $`i`$, and the power of $`\lambda `$ is $`U(T)`$, the number of upsets in the tournament: $`j>i`$ and $`ji`$: $$\underset{1i<jn}{}(x_i+\lambda x_j)=\underset{T𝒯_n}{}\lambda ^{U(T)}\underset{i=1}{\overset{n}{}}x_i^{\omega (i)}.$$ We shall use the six-vertex model to interpret the right side of equation (7). We begin with the following observations which are explained below. ###### Proposition 3 Let $`A`$ be an $`n\times n`$ ASM. In the corresponding six-vertex model * the number of horizontal vertices is $`n+N(A)`$, * the number of vertical vertices is $`N(A)`$, * the number of southwest or northeast vertices is $`(A)N(A)`$, * the number of southeast or northwest vertices is $`\left(\genfrac{}{}{0pt}{}{n}{2}\right)(A)`$. The number of vertical vertices is immediate from the bijection, and there most be one more 1 than $`1`$ in each row. A southwest vertex corresponds to a 0 of the ASM for which there is a 1 above it in its column (due north) with no other non-zero entries in between, and a 1 to its left in its row (due west) with no other non-zero entries in between. $$\begin{array}{ccccc}& & & & 1\\ & & & & 0\\ & & & & \mathrm{}\\ & & & & 0\\ 1& 0& \mathrm{}& 0& 0(\mathrm{SW})\mathrm{or}1\end{array}$$ The inversion number is the number of such pairs of 1’s: pairs of 1’s for which there are only 0s in the positions that are both due east of the lower 1 and strictly south and west of the upper 1, and there are only 0s in the positions that are both due south of the upper 1 and strictly north and east of the lower 1. The entry in the unique position due east of the lower 1 and due south of the upper 1 must be either a 0, corresponding to a southwest vertex, or a $`1`$. The remaining observations follow from the equality of the number of southwest and northeast vertices, the equality of the number of southeast and northwest vertices, and the fact that there are $`n^2`$ vertices in all. If we let $`SW(A)`$, $`SE(A)`$, and $`V(A)`$ denote, respectively, the number of southwest, southeast, and vertical vertices in $`A`$ and $`SW_i(A)`$, $`SE_i(A)`$, and $`V_i(A)`$ the number of southwest, southeast, or vertical vertices, respectively, in column $`i`$ of $`A`$, then the right side of equation (7) can be written as $$\underset{A𝒜_n}{}\lambda ^{SW(A)}(1+\lambda )^{V(A)}\underset{i=1}{\overset{n}{}}x_i^{SW_i(A)+SE_i(A)+V_i(A)}$$ Equation (7) is equivalent to $$\underset{T𝒯_n}{}\lambda ^{U(T)}\underset{i=1}{\overset{n}{}}x_i^{\omega (i)}=\underset{A𝒜_n}{}\lambda ^{SW(A)}(1+\lambda )^{V(A)}\underset{i=1}{\overset{n}{}}x_i^{SW_i(A)+SE_i(A)+V_i(A)}$$ (8) This suggests a natural bijection between tournaments on $`n`$ vertices and six-vertex models on $`n^2`$ vertices in which we have chosen a direction (left or right) at each vertical vertex. Each vertex in the six-vertex model that has an in-edge from the north will define an out-edge of the tournament. Call this vertex of the six-vertex model an initiating vertex. If an initiating vertex is southwest, there is an out-edge to the left, and the corresponding edge in the tournament will contribute to the upset number. If the initiating vertex is southeast, there is an out-edge to the right, and the corresponding edge in the tournament will not contribute to the upset number. If the initiating vertex is vertical, we have a choice of taking either the left or right out-edge. The left choice contributes one to the upset number of the tournament; the right choice contributes nothing. ###### Problem 2 Find a bijective proof of equation (8). ## 4 The Izergin-Korepin Determinant Evaluation Kuperberg’s proof of the alternating sign matrix conjecture and Zeilberger’s proof of the refined conjecture rest on the following determinant evaluation of Izergin , described in Korepin, Bogoliubov, and Izergin’s Quantum Inverse Scattering Method . ###### Proposition 4 Given $`A𝒜_n`$, let $`(i,j)`$ be the vertex in row $`i`$, column $`j`$ of the corresponding six-vertex model, and let $`H`$, $`V,`$ $`SE,`$ $`SW,`$ $`NE,`$ $`NW`$ be, respectively, the sets of horizontal, vertical, southeast, southwest, northeast, and northwest vertices. For indeterminants $`a`$, $`x_1,\mathrm{},x_n`$, and $`y_1,\mathrm{},y_n`$, we have that $`det\left({\displaystyle \frac{1}{(x_i+y_j)(ax_i+y_j)}}\right){\displaystyle \frac{_{i,j=1}^n(x_i+y_j)(ax_i+y_j)}{_{1i<jn}(x_ix_j)(y_iy_j)}}`$ (9) $`=`$ $`{\displaystyle \underset{A𝒜_n}{}}(1)^{N(A)}(1a)^{2N(A)}a^{\left(\genfrac{}{}{0pt}{}{n}{2}\right)(A)}`$ $`\times {\displaystyle \underset{(i,j)V}{}}x_iy_j{\displaystyle \underset{(i,j)NESW}{}}(ax_i+y_j){\displaystyle \underset{(i,j)NWSE}{}}(x_i+y_j).`$ As Lascoux has pointed out , the right way to understand this identity is as an extension of Cauchy’s $$det\left(\frac{1}{(x_i+y_j)}\right)\underset{i,j=1}{\overset{n}{}}(x_i+y_j)\underset{1i<jn}{}(x_ix_j)^1(y_iy_j)^1=1.$$ (10) This is true by inspection. The determinant times the product over $`i,j`$ is an alternating polynomial in the $`x_i`$ and in the $`y_j`$. Since any alternating polynomial is divisible by the Vandermonde product, the left side of this equality is a symmetric polynomial in the $`x_i`$, and it is a symmetric polynomial in the $`y_j`$. The degree in $`x_1`$ of this polynomial is zero, and the constant can be checked by induction. Applying this same reasoning to the left side of equation (9), we see that it is a symmetric polynomial in the $`x_i`$ and in the $`y_j`$. Its degree in $`x_1`$ is $`n1`$. On the right, we also have a polynomial in $`x_1`$ of degree $`n1`$. We need only check that these two sides agree for $`n`$ values of $`x_1`$. By induction, they agree at $`x_1=y_1/a`$. If we can show that the right side is symmetric in the $`y_j`$, then the identity is proven. Symmetry follows from Baxter’s triangle-to-triangle relation which was used by Izergin to prove that $`{\displaystyle \underset{A𝒜_n}{}}{\displaystyle \underset{(i,j)H}{}}x_i(1a){\displaystyle \underset{(i,j)V}{}}(y_j)(1a)`$ $`\times {\displaystyle \underset{(i,j)NESW}{}}(ax_iy_j){\displaystyle \underset{(i,j)NWSE}{}}(x_iy_j)a^{1/2}`$ is symmetric in the $`x_i`$, and it is symmetric in the $`y_j`$. Among the corollaries of Proposition 4, we can set $`a=1`$ to get Borchardt’s permanent-determinant identity: $$det\left(\frac{1}{(x_i+y_j)^2}\right)\frac{_{i,j=1}^n(x_i+y_j)^2}{_{1i<jn}(x_ix_j)(y_iy_j)}=\mathrm{perm}\left(\frac{1}{x_i+y_j}\right)\underset{i,j=1}{\overset{n}{}}(x_i+y_j),$$ (11) where $$\mathrm{perm}(a_{ij}):=\underset{\sigma 𝒮_n}{}\underset{i=1}{\overset{n}{}}a_{i,\sigma (i)}.$$ We can set $`a=\omega :=e^{2\pi i/3}`$, $`x_j=\omega q^j,`$ and $`y_j=q^{1j}`$, evaluate the determinant, and then take the limit as $`q1`$ to get the number of ASMs of a given size: $$\underset{j=0}{\overset{n1}{}}\frac{(3j+1)!}{(n+j)!}=\left\|𝒜_n\right\|.$$ (12) If we set $`a=1`$, then the matrix for which we take the determinant is $`(1/(x_i^2y_j^2))`$, which can be evaluated using Cauchy’s formula, equation (10). The left side of equation (9) simplifies to $$(1)^{n(n1)/2}\underset{1i<jn}{}(x_i+x_j)(y_i+y_j).$$ From equation (8), each of these Vandermonde-type products can be written as a sum over alternating sign matrices. We let Ein$`{}_{i}{}^{}(A)`$ be the number of vertices in row $`i`$ with an in-edge from the left, Nin$`{}_{j}{}^{}(A)`$ be the number of vertices in column $`j`$ with an in-edge from below. Replacing $`y_j`$ by $`y_j`$ and multiplying each side by $`x_1\mathrm{}x_n`$, the case $`a=1`$ is equivalent to the identity $`{\displaystyle \underset{(B,C)𝒜_n\times 𝒜_n}{}}2^{N(B)+N(C)}{\displaystyle \underset{i=1}{\overset{n}{}}}x_i^{\mathrm{Ein}_i(B)}{\displaystyle \underset{j=1}{\overset{n}{}}}y_j^{\mathrm{Nin}_j(C)}`$ (13) $`=`$ $`{\displaystyle \underset{A𝒜_n}{}}(1)^{(A)N(A)}4^{N(A)}{\displaystyle \underset{(i,j)H}{}}x_i{\displaystyle \underset{(i,j)V}{}}y_j{\displaystyle \underset{(i,j)NE}{}}(x_i+y_j)`$ $`\times {\displaystyle \underset{(i,j)SW}{}}(x_iy_j){\displaystyle \underset{(i,j)NW}{}}(x_i+y_j){\displaystyle \underset{(i,j)SE}{}}(x_iy_j)`$ ###### Problem 3 Find a bijective proof of equation (13).
warning/0007/hep-ex0007017.html	ar5iv	text	# 1 Introduction ## 1 Introduction QCD predicts differences between quark and gluon jets. These are due to the different relative probabilities for a gluon and a quark to radiate an additional gluon, given by the SU(3) group constants $`\mathrm{C}_\mathrm{A}=3`$ and $`\mathrm{C}_\mathrm{F}=4/3.`$ The various measurements of these differences from $`\mathrm{e}^+\mathrm{e}^{}`$ collider experiments are found to agree with the theoretical predictions: gluon jets are observed to have higher mean particle multiplicity, a softer fragmentation function and a wider angular spread than light quark jets . The QCD prediction of an enhancement of particle multiplicity in gluon jets with respect to quark jets is independent of the particle species, except for some small corrections. This prediction can be tested by measuring the rates of identified particles in gluon and quark jets. These measurements are also necessary for a better understanding of fragmentation processes and hadronisation models. The models presently used are mainly the string and cluster models implemented in the Jetset and Herwig Monte Carlo generators. Both models are based on the parton shower approach, Jetset using a leading log perturbative QCD calculation, and Herwig a next-to-leading log calculation. Some measurements of identified particle rates in quark and gluon jets have already been made by the LEP experiments , and the ratios of rates in quark and gluon jets have been determined. The ratios of rates of $`\mathrm{K}^+,\mathrm{K}^0,`$ and protons were found to be consistent with the ratio of rates of the average charged particle multiplicity and also consistent with the ratio of rates determined from Monte Carlo simulations. The ratio of rates for $`\mathrm{\Lambda }`$ was measured by OPAL and found to be larger than the Monte Carlo expectation. The ratio of the $`\eta `$ meson production rate in 3-jet events to the $`\eta `$ production rate in 2-jet events was observed by the L3 experiment to be larger than the Monte Carlo expectation. It was suggested that this difference could be caused by an enhanced $`\eta `$ meson production in gluon jets. A confirmation of this result would suggest that, in addition to the QCD predicted enhancement in gluon jets with respect to quark jets at equal jet energies, other $`\eta `$ meson sources might exist in gluon jets. Such sources could be production of glueballs and their decay to isoscalar mesons as has been foreseen in some theoretical models . However, the ALEPH experiment observed no evidence for an enhancement of $`\eta `$ mesons in gluon jets in excess of the Monte Carlo expectation. Experimentally, results from the comparison of charged particle or identified particle rates in quark and gluon jets are highly dependent on event topologies (i.e. the jet localisation with respect to other jets in the event) and even more dependent on the jet energies. Furthermore, comparison between the experimental results and the QCD predictions is complicated by the use of jet-finding algorithms. Analytic calculations do not employ jet finders to assign particles to jets. To cope with these energy and topological dependences, a transverse momentum-like scale $`\mathrm{Q}_{\mathrm{jet}}`$ (see Equation 2) has been proposed and was used for a comparison of the mean charged particle multiplicity in quark and gluon jets . In this paper, the rates of $`\pi ^0,\eta `$ and $`\mathrm{K}^0`$ in quark and gluon jets are measured for the first time as a function of the scale $`\mathrm{Q}_{\mathrm{jet}}`$. A phenomenological formula for the charged particle rate $`\mathrm{N}_\mathrm{q}(\mathrm{Q}_{\mathrm{jet}})`$ in quark jets, written as a second order polynomial in $`\mathrm{ln}(\mathrm{Q}_{\mathrm{jet}}),`$ and a phenomenological formula for the charged particle rate $`\mathrm{N}_\mathrm{g}(\mathrm{Q}_{\mathrm{jet}})`$ in gluon jets, written as a linear function of $`\mathrm{N}_\mathrm{q}(\mathrm{Q}_{\mathrm{jet}}),`$ were used to describe simultaneously the observed charged particle multiplicities in quark and in gluon jets. The functions $`\mathrm{N}_\mathrm{q}(\mathrm{Q}_{\mathrm{jet}})`$ and $`\mathrm{N}_\mathrm{g}(\mathrm{Q}_{\mathrm{jet}})`$, with all parameters set to the values obtained from a fit to the observed charged particle distributions, were used as a model to compare to the measured rates of identified particles in quark and gluon jets. For each comparison, only one overall normalisation factor, common to $`\mathrm{N}_\mathrm{q}(\mathrm{Q}_{\mathrm{jet}})`$ and $`\mathrm{N}_\mathrm{g}(\mathrm{Q}_{\mathrm{jet}}),`$ is allowed to vary. This method provides a model independent way of testing that the enhancement of particle production in gluon jets with respect to quark jets is independent of the particle species. ## 2 Data selection ### 2.1 Selection of hadronic $`𝐙^\mathrm{𝟎}`$ decays The present analysis was based on the full hadronic event sample collected at centre-of-mass energies at and near the $`\mathrm{Z}^0`$ peak by the OPAL detector from 1991 to 1995. This corresponded to about 4 million hadronic $`\mathrm{Z}^0`$ decays. A full description of the OPAL detector can be found in . Standard OPAL selection criteria were applied for track and electromagnetic cluster selection . Tracks were required to have: at least 20 measured points in the jet chamber, a measured momentum greater than $`0.10\mathrm{GeV}`$, an impact parameter $`\mathrm{d}_0`$ in the $`\mathrm{r}\varphi `$ plane smaller than 2 cm, a z position within 25 cm of the interaction point and a measured angle with respect to the beam axis of greater than $`20^{}.`$ Electromagnetic clusters were required to have an energy greater than $`0.1\mathrm{GeV}`$ if they were in the barrel part of the detector (i.e. $`\mathrm{cos}\theta 0.82`$) or greater than $`0.3\mathrm{GeV}`$ if they were in the endcap part. The selected tracks and clusters not associated with tracks were fed, as four-vectors, to the jet-finding algorithms. Background from all sources was reduced to less than $`0.8\%`$ and was neglected throughout the analysis. It was reduced by requiring for each event more than 7 measured tracks, a visible energy (i.e. the sum of detected particle energies after correcting for double counting) larger than 60 $`\mathrm{GeV}`$ and an angle larger than $`25^{}`$ between the calculated thrust axis and the beam axis . ### 2.2 Event simulation Detector effects and detection efficiencies for the studied particles were evaluated using 8 million Monte Carlo hadronic $`\mathrm{Z}^0`$ decays. Events were generated using the Jetset program tuned to reproduce the global features of hadronic events as measured with the OPAL detector. About 4 million events generated by the Herwig program were also used for comparison. The generated events were processed through a full simulation of the OPAL detector and were processed using the same reconstruction and selection programs as were applied to the data. ### 2.3 Selection of 3-jet events Three jet-finding algorithms were used: Luclus , Durham and the cone jet finder. The Luclus jet finder was found to provide the best jet angular resolution, which was relevant for the present analysis. Luclus was therefore used as the reference algorithm while the two others were used for systematic comparisons. The jet algorithm was forced to resolve three jets in each hadronic event. The jet energies and momenta were then calculated by imposing energy and momentum conservation with planar massless kinematics , using the jet directions found by the jet algorithm. They are given by the cyclic relation: $$\mathrm{E}_\mathrm{i}=\frac{\sqrt{\mathrm{s}}\mathrm{sin}\theta _{\mathrm{j},\mathrm{k}}}{\mathrm{sin}\theta _{\mathrm{i},\mathrm{j}}+\mathrm{sin}\theta _{\mathrm{j},\mathrm{k}}+\mathrm{sin}\theta _{\mathrm{k},\mathrm{i}}},$$ (1) where $`\theta _{\mathrm{i},\mathrm{j}}`$ is the angle between jet i and jet j. The event was accepted as a 3-jet event if each jet contained at least 3 charged particles, had a corrected energy exceeding $`5\mathrm{GeV}`$, and pointed more than $`20^{}`$ away from the beam axis and more than $`30^{}`$ away from the direction of the other two jets. The variable $`\mathrm{Y}=(\mathrm{D}_{23}\mathrm{D}_{34})/\mathrm{E}_{\mathrm{visible}}`$ was used, where $`\mathrm{D}_{23}\mathrm{and}\mathrm{D}_{34}`$ were the Luclus jet algorithm resolution parameters ($`\mathrm{D}_{\mathrm{join}}`$ for the transition from 2 to 3 and 3 to 4 jets, respectively, and $`\mathrm{E}_{\mathrm{visible}}`$ the total visible energy. This variable measured the topological stability of the event as a 3-jet event. For larger values of Y, the events tended to be three-fold symmetric, meaning that all inter-jet angles were close to $`120^{}.`$ To select stable 3-jet events (i.e. events not close to the transition from three to four jets), only events with $`\mathrm{Y}>0.2`$ <sup>1</sup><sup>1</sup>1The corresponding Y variable for Durham was defined as $`\mathrm{Y}=\mathrm{y}_{23}\mathrm{y}_{34}`$, where y is the usual $`\mathrm{y}_{\mathrm{cut}}`$ . For the cone jet finder the inter-jet angles were considered instead. were kept for further processing. The total selected data sample contained approximately 493 000 3-jet events, which was $`12.32\%`$ of the total initial event sample. The corresponding fraction for Monte Carlo events was $`12.30\%`$ for Jetset and $`12.17\%`$ for Herwig. For the Monte Carlo events, the jet-finding algorithm was applied at the parton, hadron and detector levels (defined in Section 3.2). At each level the jet energies were corrected to satisfy the constraints of energy and momentum conservation with planar massless kinematics, after which the jets were energy ordered, the first jet being the jet with the highest energy. The matching from parton to hadron level and from hadron to detector level was done using a simple angular correspondence: a jet at hadron level was matched with only one detector jet and one parton jet, those having the minimal angular deviation with respect to the hadron jet direction. The jet energy resolution, defined as $`(\mathrm{E}_{\mathrm{jet}}^{\mathrm{parton}}\mathrm{E}_{\mathrm{jet}}^{\mathrm{detector}})/\mathrm{E}_{\mathrm{jet}}^{\mathrm{parton}},`$ was found to range from $`5\%`$ for the first jet to $`13\%`$ for the third jet. The angle between the parton jet direction and the detector jet direction was found to have an r.m.s. of $`0.07\mathrm{radians}`$ for the highest energy jet and $`0.16\mathrm{radians}`$ for the lowest energy jet. ## 3 Analysis method In this section, the jet scale $`\mathrm{Q}_{\mathrm{jet}}`$ and the jet purities are defined and the method of unfolding the average charged particle multiplicity to $`100\%`$ pure quark and gluon jets is explained. In Section 4 the measured average charged particle multiplicity in pure quark and gluon jets, as a function of the scale $`\mathrm{Q}_{\mathrm{jet}}`$, was fitted to phenomenological formulae. The purpose of the fit was to obtain an analytical shape that could be used as a reference or a model to compare in Section 6 to the corresponding shapes obtained in Section 5 for $`\pi ^0,\eta `$ and $`\mathrm{K}^0`$. Systematic effects that could have distorted the measured shape are discussed in Section 4.3. ### 3.1 Jet scale $`\mathrm{Q}_{\mathrm{jet}}`$ It has already been shown that the jet energy alone is not an adequate scale to describe the average particle multiplicity in quark and gluon jets. Coherence in QCD radiation suggests that the position of a parton with respect to other partons in the event (i.e. the event topology) should also be considered in studies of jet properties. Inter-jet coherence effects that can lead to destructive or constructive interference effects on the particle flux in the inter-jet angular region, have been predicted and observed experimentally. An example is the string effect . A transverse-momentum-like scale combining the jet energy and its angular position with respect to the other jets has been used in a phenomenological study of parton shower characteristics. This scale $`\mathrm{Q}_{\mathrm{jet}}`$ was defined as: $$\mathrm{Q}_{\mathrm{jet}}=\mathrm{E}_{\mathrm{jet}}\mathrm{sin}(\theta /2)$$ (2) where $`\theta `$ was the jet angle with respect to the closest jet. $`\mathrm{Q}_{\mathrm{jet}}`$ has already been used in an experimental study of multiplicity in $`\mathrm{e}^+\mathrm{e}^{}`$ 3-jet events . The way the scale $`\mathrm{Q}_{\mathrm{jet}}`$ incorporates the topological dependence can be seen, for example, at leading order in $`\alpha _\mathrm{s}`$ for the three parton $`\mathrm{q}\overline{\mathrm{q}}\mathrm{g}`$ vertex: colour being conserved in QCD, the gluon can be represented by a $`\mathrm{q}\overline{\mathrm{q}}`$ pair that compensates exactly the total colour charge of the initial $`\mathrm{q}\overline{\mathrm{q}}`$ pair. A mutual colour shielding occurs when the initial quark or anti-quark is close in angle to the gluon and reduces subsequent gluon radiation. This colour shielding is minimal for back-to-back partons, and the scale $`\mathrm{Q}_{\mathrm{jet}}`$ becomes equal to the parton energy for that case. The scale $`\mathrm{Q}_{\mathrm{jet}}`$ was used for the present analysis with no restriction on event topology except for the minimal inter-jet angle. The distributions of energy, $`\mathrm{E}_{\mathrm{jet}},`$ and scale, $`\mathrm{Q}_{\mathrm{jet}}`$ (normalised to the total number of analysed hadronic events), for the three energy ordered jets, are shown in Figures 1a and 1b. The Monte Carlo distributions were found to reproduce the data very well. ### 3.2 Gluon jet definition and purity estimation The jet having the smallest energy, $`\mathrm{E}_{\mathrm{jet}},`$ was found to have the highest gluon purity. The jet purities were estimated using the Monte Carlo information: the initial parton shower (parton level), the generated hadrons after the fragmentation processes (hadron level) and the reconstructed particles after the simulation of the full OPAL detector response (detector level). At the hadron level, all charged and neutral particles with lifetimes greater than $`3\times 10^{10}\mathrm{s}`$ were considered as stable particles. A jet at the detector level was considered to be a gluon jet if it matched a parton jet that did not contain either the initial quark or anti-quark from the Z decay. The jet purities were determined directly from Monte Carlo as the fraction of quark or gluon jets present in the jet sample at a fixed scale $`\mathrm{Q}_{\mathrm{jet}}`$. The purities could also be estimated from matrix element calculations. It has been shown that, for leading order QCD matrix elements, the probability for a given jet $`\{\mathrm{i}\}`$ among the jets $`\{\mathrm{i},\mathrm{j},\mathrm{k}\}`$ to be a gluon jet can be expressed as a function of the jet energies: $$\mathrm{P}_{\mathrm{i}=\mathrm{g}}\frac{\mathrm{x}_\mathrm{j}^2+\mathrm{x}_\mathrm{k}^2}{(1\mathrm{x}_\mathrm{j})(1\mathrm{x}_\mathrm{k})},$$ (3) where $`\mathrm{x}_\mathrm{i}=2\mathrm{E}_\mathrm{i}/\sqrt{\mathrm{s}}`$ and the corresponding probability for being a quark jet is $$\mathrm{P}_{\mathrm{i}=\mathrm{q}}=1\mathrm{P}_{\mathrm{i}=\mathrm{g}},$$ (4) normalised to have: $$\mathrm{P}_{1=\mathrm{q}}+\mathrm{P}_{2=\mathrm{q}}+\mathrm{P}_{3=\mathrm{q}}=2.$$ (5) The purities obtained using Monte Carlo information are shown in Figure 2 together with the purities obtained from the matrix elements. In the same figure are also shown the purities for OPAL data obtained from the matrix element formula. Very good agreement was obtained between the two methods, with less than $`2\%`$ deviation in most of the range of $`\mathrm{Q}_{\mathrm{jet}}`$ common to the second and third jet (see Figure 1 and Figure 2). Therefore, the matrix element formula was used to estimate the jet purities directly from the data. The same agreement was also observed for Herwig Monte Carlo events. ### 3.3 Unfolding to pure quark and gluon jets The measured jet samples were mixtures of quark and gluon jets, while meaningful comparison of gluon and quark jets should be performed on quantities evaluated for pure samples of gluon jets and pure samples of quark jets. This analysis used the average charged particle multiplicity for samples of jets at the same scale $`\mathrm{Q}_{\mathrm{jet}}`$ but having different gluon/quark jet purities. If $`\mathrm{N}_1(\mathrm{Q}_{\mathrm{jet}})`$ and $`\mathrm{N}_2(\mathrm{Q}_{\mathrm{jet}})`$ are the average measured charged particle multiplicities for two jet samples with different quark jet purities $`\mathrm{P}_1(\mathrm{Q}_{\mathrm{jet}})`$ and $`\mathrm{P}_2(\mathrm{Q}_{\mathrm{jet}})`$ at the same $`\mathrm{Q}_{\mathrm{jet}}`$ value, then: $$\mathrm{N}_1(\mathrm{Q}_{\mathrm{jet}})=\mathrm{P}_1(\mathrm{Q}_{\mathrm{jet}})\mathrm{N}_\mathrm{q}(\mathrm{Q}_{\mathrm{jet}})+(1\mathrm{P}_1(\mathrm{Q}_{\mathrm{jet}}))\mathrm{N}_\mathrm{g}(\mathrm{Q}_{\mathrm{jet}})$$ (6) $$\mathrm{N}_2(\mathrm{Q}_{\mathrm{jet}})=\mathrm{P}_2(\mathrm{Q}_{\mathrm{jet}})\mathrm{N}_\mathrm{q}(\mathrm{Q}_{\mathrm{jet}})+(1\mathrm{P}_2(\mathrm{Q}_{\mathrm{jet}}))\mathrm{N}_\mathrm{g}(\mathrm{Q}_{\mathrm{jet}})$$ (7) where $`\mathrm{N}_\mathrm{q}(\mathrm{Q}_{\mathrm{jet}})\mathrm{and}\mathrm{N}_\mathrm{g}(\mathrm{Q}_{\mathrm{jet}})`$ are the average particle multiplicities for $`100\%`$ pure quark and gluon jet sample. This gives: $$\mathrm{N}_\mathrm{q}(\mathrm{Q}_{\mathrm{jet}})=\frac{(1\mathrm{P}_2(\mathrm{Q}_{\mathrm{jet}}))\mathrm{N}_1(\mathrm{Q}_{\mathrm{jet}})+(1\mathrm{P}_1(\mathrm{Q}_{\mathrm{jet}}))\mathrm{N}_2(\mathrm{Q}_{\mathrm{jet}})}{\mathrm{P}_1(\mathrm{Q}_{\mathrm{jet}})\mathrm{P}_2(\mathrm{Q}_{\mathrm{jet}})}$$ (8) $$\mathrm{N}_\mathrm{g}(\mathrm{Q}_{\mathrm{jet}})=\frac{\mathrm{P}_1(\mathrm{Q}_{\mathrm{jet}})\mathrm{N}_2(\mathrm{Q}_{\mathrm{jet}})\mathrm{P}_2(\mathrm{Q}_{\mathrm{jet}})\mathrm{N}_1(\mathrm{Q}_{\mathrm{jet}})}{\mathrm{P}_1(\mathrm{Q}_{\mathrm{jet}})\mathrm{P}_2(\mathrm{Q}_{\mathrm{jet}})}.$$ (9) This unfolding was only possible in the region where the $`\mathrm{Q}_{\mathrm{jet}}`$ scales of the jet samples overlapped. The $`\mathrm{Q}_{\mathrm{jet}}`$ distributions of the energy ordered jets (see Figure 1) showed that the second and the third jets could be considered to be different samples over a common range of $`\mathrm{Q}_{\mathrm{jet}}`$ from $`6\mathrm{to}26\mathrm{GeV}`$. ## 4 Average charged particle multiplicities ### 4.1 Measurements and parametrisation The average number of charged particles, $`\mathrm{N}^{\mathrm{ch}},`$ per jet in bins of 1 $`\mathrm{GeV}`$ of the scale $`\mathrm{Q}_{\mathrm{jet}}`$ was measured for samples of the second and third jets, being respectively quark enriched and gluon enriched through jet energy ordering and having a wide common range of $`\mathrm{Q}_{\mathrm{jet}}`$. In each bin of $`\mathrm{Q}_{\mathrm{jet}}`$, the average purity was obtained from the matrix element formula. The efficiency corrections from detector to hadron level were calculated from Monte Carlo information as a function of $`\mathrm{Q}_{\mathrm{jet}}`$ for each jet sample separately. The efficiency was defined as the ratio of the average number of charged particles at the detector level for a given jet sample at a given scale $`\mathrm{Q}_{\mathrm{jet}}`$, divided by the equivalent quantity for the corresponding jet sample at the hadron level. The efficiencies were found to be approximately independent of $`\mathrm{Q}_{\mathrm{jet}}`$ for all jet samples, and the corresponding corrections to $`\mathrm{N}^{\mathrm{ch}}(\mathrm{Q}_{\mathrm{jet}})`$ were at the level of $`10\%`$. The results obtained after unfolding to $`100\%`$ pure quark and gluon jets are shown in Figure 3. The unfolding was performed for $`\mathrm{Q}_{\mathrm{jet}}>7\mathrm{GeV},`$ since the method worked for all three jet finders in this region. The average charged particle multiplicities in quark jets, $`\mathrm{N}_\mathrm{q},`$ and in gluon jets, $`\mathrm{N}_\mathrm{g},`$ were simultaneously described with phenomenological formulae. The parametrisation was given by: $$\mathrm{N}_\mathrm{q}(\mathrm{Q}_{\mathrm{jet}})=\mathrm{a}_0+\mathrm{a}_1\mathrm{ln}\mathrm{Q}_{\mathrm{jet}}+\mathrm{a}_2(\mathrm{ln}\mathrm{Q}_{\mathrm{jet}})^2$$ (10) $$\mathrm{N}_\mathrm{g}(\mathrm{Q}_{\mathrm{jet}})=\mathrm{R}_0+\mathrm{R}_1\mathrm{N}_\mathrm{q}(\mathrm{Q}_{\mathrm{jet}})$$ (11) where $`\mathrm{a}_0,\mathrm{a}_1,\mathrm{a}_2,\mathrm{R}_0\mathrm{and}\mathrm{R}_1`$ are constants. This parametrisation resembles the next-to-leading order expressions given in , but with an extra offset parameter $`\mathrm{R}_0.`$ The leading order expressions alone were found to be unable to fit the data, as was also found in and discussed in . The $`\mathrm{Q}_{\mathrm{jet}}`$ dependent average charged particle multiplicities for quark and gluon jets were fitted simultaneously to the expressions for $`\mathrm{N}_1(\mathrm{Q}_{\mathrm{jet}})\mathrm{and}\mathrm{N}_2(\mathrm{Q}_{\mathrm{jet}})`$ (Equations 6 and 7) where $`\mathrm{N}_\mathrm{q}(\mathrm{Q}_{\mathrm{jet}})\mathrm{and}\mathrm{N}_\mathrm{q}(\mathrm{Q}_{\mathrm{jet}})`$ were replaced by the expressions given in equations 10 and 11. The result of the fit yielded: $`\mathrm{a}_0=2.74\pm 0.07,\mathrm{a}_1=1.71\pm 0.05,\mathrm{a}_2=0.05\pm 0.07,\mathrm{R}_0=7.27\pm 0.52\mathrm{and}\mathrm{R}_1=2.27\pm 0.07.`$ The fitted functions are shown with the data points unfolded using equations 8 and 9 in Figure 3a, where a good description of the average charged particle multiplicities in both the gluon and the quark jets can be seen. The ratio of the average charged particle multiplicities in gluon to quark jets is shown in Figure 3b. This ratio is predicted by QCD to be the same for all particle species. A test of this prediction is the quality of the fits of the above analytical function for charged particles to the measured multiplicities of $`\pi ^0,\eta `$ and $`\mathrm{K}^0`$ in quark and gluon jets with all parameters fixed except for an overall normalisation factor. ### 4.2 Ratio of the slopes of the multiplicities From QCD perturbative calculations, the ratio of particle multiplicities in gluon and quark jets is given by: $$\mathrm{N}_\mathrm{g}(\mathrm{Q}_{\mathrm{jet}})=\mathrm{R}(\mathrm{Q}_{\mathrm{jet}})\times \mathrm{N}_\mathrm{q}(\mathrm{Q}_{\mathrm{jet}})$$ (12) where the asymptotic limit of $`\mathrm{R}(\mathrm{Q}_{\mathrm{jet}})`$ at large $`\mathrm{Q}_{\mathrm{jet}}`$ is an approximation to the QCD colour factor ratio $`\mathrm{R}=\mathrm{C}_\mathrm{A}/\mathrm{C}_\mathrm{F}=2.25.`$ At the $`\mathrm{Z}^0`$ scale, this value is lower due to sizeable higher order corrections. However, the ratio of the slopes $$\mathrm{C}(\mathrm{Q}_{\mathrm{jet}})=\frac{\mathrm{d}\mathrm{N}_\mathrm{g}/\mathrm{dQ}_{\mathrm{jet}}}{\mathrm{d}\mathrm{N}_\mathrm{q}/\mathrm{dQ}_{\mathrm{jet}}}(\mathrm{Q}_{\mathrm{jet}})$$ (13) of multiplicities in quark and gluon jets is expected to be less affected by these higher order corrections . The slope ratio C has been recently calculated using a next-to-next-to-next-to leading order (3NLO) perturbative approximation. The predicted value of $`\mathrm{C}(\mathrm{Q}_{\mathrm{jet}})`$ in $`\mathrm{Z}^0`$ decays is $`\mathrm{C}1.92.`$ From the parametrisation given in equations 10 and 11 that imposes a constant value for $`\mathrm{C}(\mathrm{Q}_{\mathrm{jet}})`$ (Equation 13), we obtained: $$\mathrm{C}(\mathrm{Q}_{\mathrm{jet}})=\mathrm{R}_1=2.27\pm 0.07(\mathrm{stat}.)\pm 0.19(\mathrm{syst}.).$$ (14) The constraint $`\mathrm{C}(\mathrm{Q}_{\mathrm{jet}})=\mathrm{constant}`$ can be released by extracting the spectrum $`\mathrm{C}(\mathrm{Q}_{\mathrm{jet}})`$ from the measured distributions $`\mathrm{N}_\mathrm{q}(\mathrm{Q}_{\mathrm{jet}})`$ and $`\mathrm{N}_\mathrm{q}(\mathrm{Q}_{\mathrm{jet}}).`$ This was done by using as an estimate of the derivative at each bin $`\mathrm{Q}_{\mathrm{jet}}`$ the slope of a fitted line to three adjacent bins centred at $`\mathrm{Q}_{\mathrm{jet}}.`$ The obtained $`\mathrm{C}(\mathrm{Q}_{\mathrm{jet}})`$ spectrum is shown in Figure 3c, and a fit to a constant yielded: $$\mathrm{C}=2.27\pm 0.09(\mathrm{stat})\pm 0.27(\mathrm{syst}).$$ The systematic uncertainty, for both methods, was mainly due to differences between the jet finders and the correlations between the unfolded values of $`\mathrm{N}_\mathrm{g}(\mathrm{Q}_{\mathrm{jet}})`$ and $`\mathrm{N}_\mathrm{q}(\mathrm{Q}_{\mathrm{jet}})`$ at each bin of $`\mathrm{Q}_{\mathrm{jet}}`$. The two values are about one standard deviation higher than the prediction $`\mathrm{C}1.92.`$ The Delphi collaboration recently presented a measurement of the ratio of slopes. Their result, $`\mathrm{C}=1.97\pm 0.10(\mathrm{stat}),`$ is about one standard deviation of the total uncertainty bellow our measurement. ### 4.3 Stability of the parametrisation of multiplicities To study systematic effects on the parametrisation of the charged particle multiplicity obtained in the previous section, several variations to the analysis procedure were applied. For each variation, new $`\mathrm{N}_\mathrm{q}(\mathrm{Q}_{\mathrm{jet}})`$ and $`\mathrm{N}_\mathrm{q}(\mathrm{Q}_{\mathrm{jet}})`$ were obtained and compared to the original distributions. At each data point, the deviation with respect to the original value is considered as a systematic error. The errors were added quadratically and were included in the error bars of the original data points shown in Figure 3. For each variation, the distributions $`\mathrm{N}_\mathrm{q}(\mathrm{Q}_{\mathrm{jet}})`$ and $`\mathrm{N}_\mathrm{q}(\mathrm{Q}_{\mathrm{jet}})`$ were also fitted to the functions of equations 10 and 11. The resulting parametrisation was compared to the original one by calculating the largest relative difference between the new and the original function values. The following systematic variations were considered: 1. The analysis was completely repeated with the Durham and cone jet finders. The parametrisation was found to agree well between the Luclus (original) and the Durham jet finders. The parametrisation obtained using the cone jet finder agreed within $`10\%`$ with the original one. 2. The analysis was repeated using a different method for estimating the efficiencies. The fully unfolded average charged particle multiplicities were obtained, as a function of $`\mathrm{Q}_{\mathrm{jet}}`$ and for each jet sample, at the hadron level and then at the detector level. The ratio of detector to hadron level was then used to correct the average multiplicities of the data after having performed the unfolding. This procedure was done with the Jetset and Herwig Monte Carlos and the corrections were found to differ by at most $`3\%.`$ The analytic formula (Equations 10 and 11) was then fitted to the average charged particle multiplicity spectra corrected separately with Jetset and Herwig. The shape was found to agree with the original one within $`2\%`$ for Jetset and $`3\%`$ for Herwig. 3. The influence of the jet purities was studied by using the purities from the Monte Carlo matching, and also by increasing the cut on D from 0.20 to 0.25. The effect on the fitted parametrisation was negligible. 4. Effects of soft particles on the measured multiplicities were studied by changing the minimum momentum required per track from 0.10 GeV to $`0.15\mathrm{GeV},`$ and the minimum number of tracks per jet was also changed from 3 to 5. The analysis was repeated and the fitted parametrisation agreed with the original one within $`3\%.`$ 5. The analysis was repeated with two different jet samples. The first jet sample was gluon enriched by requiring that two jets were tagged as b-quark jets with a neural network b-tagging method and the remaining jet was considered as a gluon jet. The second jet sample, which was quark enriched, was obtained by selecting all second jets (energy ordered) in the events having no b-tagged quark jets. The average gluon purity was $`80\%`$ for the first jet sample, and the average quark purity was $`65\%`$ for the second jet sample. The fitted parametrisation agreed with the original one within $`8\%.`$ The new average charged particle multiplicity spectra, unfolded to $`100\%`$ quark and gluon purities, agreed very well with the original spectra. Differences at each data point were considered as systematic errors; this should, in principle, also account for possible correlations between the measured $`\mathrm{N}_\mathrm{g}(\mathrm{Q}_{\mathrm{jet}})\mathrm{and}\mathrm{N}_\mathrm{q}(\mathrm{Q}_{\mathrm{jet}}).`$ ## 5 $`\pi ^0,\eta `$ and $`\mathrm{K}^0`$ meson production This section describes the reconstruction of the decay channels $`\eta 2\gamma ,\pi ^02\gamma `$ and $`\mathrm{K}_\mathrm{S}^0`$ $`\pi ^+\pi ^{}`$ for the full $`\mathrm{Z}^0`$ hadronic decay sample without any 3-jet requirement. In order to gain more confidence for the rate measurements in quark and gluon jets (see Section 6), the total inclusive rates as well as the differential rates in hadronic $`\mathrm{Z}^0`$ decays of $`\pi ^0,\eta `$ and $`\mathrm{K}^0`$ were measured and compared to previous such measurements at LEP. ### 5.1 $`\pi ^0`$ and $`\eta `$ reconstruction The reconstruction was restricted to the barrel region of the detector. A procedure using a parametrisation of the expected lateral energy distribution of electromagnetic showers was optimised to resolve as many photon candidates as possible from the overlapping energy deposits in the electromagnetic calorimeter, in the dense environment of hadronic jets. The procedure was efficient but led to a rather low purity due to reconstructing spurious photons (“fakes”). Based on Monte Carlo simulated events, the rejection of fake photons was studied using a set of five measurable variables, namely: $`\mathrm{E}_\gamma :`$ the energy of the photon candidate, $`\mathrm{E}_{\mathrm{clust}}:`$ the energy of the nearest cluster to the considered photon candidate, $`\theta _{\mathrm{clust}}:`$ the opening angle between the photon candidate and the nearest cluster, $`\theta _{\mathrm{trk}}:`$ the opening angle between the photon candidate and the closest reconstructed track, and $`\mathrm{E}_{\mathrm{trk}}:`$ the amount of energy that could be attributed to tracks in an array of 3x3 lead glass blocks around the position of the photon candidate. A large number of obvious fake photons (mostly with $`\mathrm{E}_\gamma 300`$ $`\mathrm{MeV}`$) were rejected with the following two cuts: 1. A candidate was rejected if it was found to satisfy: $$\mathrm{E}_\gamma <\mathrm{A}\mathrm{E}_{\mathrm{clust}}\mathrm{exp}\left(\left(\frac{\theta _{\mathrm{clust}}}{40\mathrm{mrad}}\right)^2\right).$$ (15) Thus, when more than one photon candidate was obtained from a single electromagnetic cluster, candidates that had small energy compared to the cluster energy were rejected if their reconstructed position was close to the cluster centre. The 40 mrad in the exponential is the average polar aperture of a lead glass block as seen from the interaction vertex. The factor A was determined empirically from the Monte Carlo sample. It was set at a value which ensured that the number of rejected photons was less than 1/10 of the number of rejected fake photons in all energy bins of $`\mathrm{E}_\gamma .`$ 2. A candidate was also rejected if it was found to satisfy: $$\theta _{\mathrm{trk}}<\mathrm{B}+\mathrm{C}\mathrm{exp}\left(\left(\frac{\mathrm{E}_\gamma }{\mathrm{E}_{\mathrm{trk}}}\right)^2\right).$$ (16) A photon candidate was likely to be fake if it had an energy, $`\mathrm{E}_\gamma `$, smaller than the electromagnetic energy which could be attributed to tracks. The factors B and C were determined from the Monte Carlo sample. They were also set at values which ensured that the number of rejected photons was less than 1/10 of the number of rejected fake photons in all energy bins of $`\mathrm{E}_\gamma .`$ The number of remaining fake photons was further reduced using a weight function $$\mathrm{W}(\mathrm{E}_\gamma ,\mathrm{E}_{\mathrm{clust}},\theta _{\mathrm{clust}},\theta _{\mathrm{trk}},\mathrm{E}_{\mathrm{trk}})$$ which was calculated for every photon candidate. The five variables were assumed uncorrelated and a likelihood ratio distribution was determined for each variable. The likelihood ratio in each bin was defined as the ratio of the number of generated photons to the total number of photon candidates. The value of W was the product of the five likelihood ratios of the bins $`\mathrm{E}_\gamma ,\mathrm{E}_{\mathrm{clust}},\theta _{\mathrm{clust}},\theta _{\mathrm{trk}}`$ and $`\mathrm{E}_{\mathrm{trk}}`$ in which the candidate was found. The discriminating power of W is shown in Figure 4, where the efficiency and purity for photons in Monte Carlo events are shown as a function of a value $`\mathrm{W}_{\mathrm{cut}}.`$ For photon candidates with $`\mathrm{W}>\mathrm{W}_{\mathrm{cut}},`$ the purity was defined as the ratio of the number of generated photons to the total number of reconstructed photons and the efficiency was defined as the fraction of generated photons which were correctly reconstructed. All possible pairs of photon candidates were then considered. Each pair was assigned a probability P for both candidates being correctly reconstructed photons, the probability being simply the product of the weights W associated to the two candidates: $$\mathrm{P}=\mathrm{W}_1\times \mathrm{W}_2,\mathrm{with}\mathrm{no}\mathrm{cut}\mathrm{on}\mathrm{W}_1\mathrm{or}\mathrm{W}_2.$$ (17) The combinatorial background consisted of a mixture of three components: wrong pairing of two correctly reconstructed photons, pairing of two fake photons and pairing of one correctly reconstructed photon with one fake photon. Choosing only photon pairs with high values of P reduces the combinatorial background to its “wrong pairing of correctly reconstructed photons” component only. It was found that requiring $`\mathrm{P}>0.1`$ removed $`60\%`$ of the total combinatorial background, with a relative loss in efficiency of $`8\%`$ and $`1\%`$ for $`\pi ^0`$ and $`\eta `$ signals, respectively. The two-photon invariant mass distribution was studied in intervals of the energy fraction $`\mathrm{x}_\mathrm{E}=\frac{\mathrm{E}_{2\gamma }}{\mathrm{E}_{\mathrm{beam}}}.`$ Because most of the true reconstructed photons came from $`\pi ^0`$ decays, an additional cut was required, in addition to the cut on the probability P, to enhance the $`\eta 2\gamma `$ signal. This cut excluded, for invariant masses $`\mathrm{M}_{2\gamma }>`$ $`300\mathrm{MeV},`$ any photon that could pair with any other photon to make an invariant mass $`\mathrm{M}_{2\gamma }<`$ $`300\mathrm{MeV}`$ with a probability $`\mathrm{P}>0.1`$. The combinatorial background could be described by a second order polynomial for all cuts on probability P, and for all $`\mathrm{x}_\mathrm{E}`$ bins. The signal was well described by a double Gaussian. The mass distributions obtained from the data over the full $`\mathrm{x}_\mathrm{E}`$ range are shown for four different cuts on the probability P in Figure 5 for $`\pi ^0,`$ and in Figure 6 for $`\eta `$. The absolute $`\pi ^0`$ and $`\eta `$ reconstruction efficiencies, for $`\mathrm{P}>0.1`$ over the entire $`\mathrm{x}_\mathrm{E}`$ range, were $`15.5\%`$ and $`7\%`$ respectively. The efficiency and the signal to background ratio were dependent on the $`\mathrm{x}_\mathrm{E}`$ interval considered. ### 5.2 $`\mathrm{K}_\mathrm{S}^0`$ reconstruction The $`\mathrm{K}_\mathrm{S}^0`$ $`\pi ^+\pi ^{}`$ reconstruction was similar to the method described in . Here it was applied to the full available $`\mathrm{Z}^0`$ hadronic sample of four million events. A track was considered to be a pion candidate if it had a transverse momentum larger than 150 MeV, had more than 20 hits in the jet chamber and had either more than 3 hits in the Z-chamber or a reconstructed end point in the jet chamber . The invariant mass of pion pairs was evaluated for pairs of oppositely charged tracks having an intersection point in the plane perpendicular to the beam axis and satisfying the following requirements: 1. the distance from the intersection point to the primary vertex had to be greater than 1 cm and less than 150 cm; 2. if the secondary vertex was reconstructed in the jet chamber, it had to be less than 5 cm from the first hit of either track; 3. if the intersection point occurred before the jet chamber, the radial distance from the track to the beam axis at the point of closest approach had to exceed 3 mm; 4. track pairs that passed the above cuts were re-fitted with the constraint that they originated from a common vertex; 5. track pairs that satisfied the photon conversion hypothesis or $`\mathrm{\Lambda }p\pi `$ hypothesis were rejected. The $`\pi ^+\pi ^{}`$ invariant mass spectrum is shown in Figure 7 for the whole measured $`\mathrm{x}_\mathrm{E}`$ region. The spectrum was studied in $`\mathrm{x}_\mathrm{E}`$ intervals. In each interval, a double Gaussian shape for the signal and a second order polynomial for the background were used to fit the $`\mathrm{M}_{\pi ^+\pi ^{}}`$ spectrum. The $`\mathrm{K}_\mathrm{S}^0`$ reconstruction efficiency was found to be $`26\%`$ for $`\mathrm{x}_\mathrm{E}<0.1`$, reducing to $`15\%`$ at higher $`\mathrm{x}_\mathrm{E}`$. ### 5.3 Inclusive $`\pi ^\mathrm{𝟎},\eta `$ and $`𝐊^0`$ rates The inclusive rate measurements were performed on the whole hadronic sample without making any 3-jet requirement. For each of the mesons $`\pi ^0,\eta `$ and $`\mathrm{K}_\mathrm{S}^0`$, the number of entries remaining after the combinatorial background subtraction was considered as the number of signal entries. This was corrected for detector and reconstruction efficiencies and for the non-measured decay channels using the Particle Data Group branching ratios. The total rates were also corrected for the inaccessible $`\mathrm{x}_\mathrm{E}`$ regions. The corrections were performed using extrapolation with the spectral shapes predicted by Jetset. They amounted to $`11\%,23\%`$ and $`0.6\%`$ of the total rates for $`\pi ^0,\eta `$ and $`\mathrm{K}^0`$ respectively. The $`\mathrm{K}^0`$ rate was the $`\mathrm{K}_\mathrm{S}^0`$ rate corrected for the non-observed $`\mathrm{K}_\mathrm{L}^0`$. The total inclusive measured rates per event were: $`\mathrm{n}_{\pi ^0}=9.871\pm 0.040(\mathrm{stat})\pm 0.39(\mathrm{syst})`$ $`\mathrm{n}_\eta =1.076\pm 0.090(\mathrm{stat})\pm 0.084(\mathrm{syst})`$ $`\mathrm{n}_{\mathrm{K}^0}=2.016\pm 0.003(\mathrm{stat})\pm 0.052(\mathrm{syst}).`$ The results are in good agreement with the values previously measured at LEP . The reconstruction method described in Section 5.1 gave an improved photon purity and higher $`\pi ^02\gamma `$ and $`\eta 2\gamma `$ reconstruction efficiencies when compared to . This led to a well controlled combinatorial background in the two-photon invariant mass spectrum. Thus, better systematic uncertainties were obtained. The statistical error on the $`\eta `$ inclusive rate was larger than the error quoted in because the $`\pi ^+\pi ^{}\pi ^0`$ channel was not included in the present analysis. Since the full available $`\mathrm{Z}^0`$ hadronic sample was analysed, the statistical error on the $`\mathrm{K}^0`$ inclusive rate was improved when compared to the value quoted in . However, the new systematic error was slightly larger, due to the inclusion of the systematic error on the Monte Carlo background and signal shapes (Section 5.4 item 1 to 3). The rate measurement was repeated in $`\mathrm{x}_\mathrm{E}`$ intervals and the obtained differential rate distributions are shown in Figure 8. The $`\pi ^0`$ measured differential rate was described well by both Jetset and Herwig Monte Carlo expectations. However, the measured $`\eta `$ and $`\mathrm{K}^0`$ spectra were harder than either Monte Carlo prediction. The discrepancy was worst for $`\eta ,`$ where the measured rate was almost double the predicted rate for $`\mathrm{x}_\mathrm{E}>0.2`$. The $`\mathrm{K}^0`$ measured rate was only about $`15\%`$ larger than the predicted rate for $`\mathrm{x}_\mathrm{E}>0.2.`$ For $`\mathrm{x}_\mathrm{E}<0.1`$ the Monte Carlo predicted rates for $`\eta `$ and $`\mathrm{K}^0`$ mesons agreed with the measured values within the error bars. These observations were in good agreement with previous OPAL results . The derived values for $`\pi ^0`$ and $`\eta `$ are not meant to supersede the former OPAL results. The new measured $`\mathrm{K}^0`$ differential rates are given in Table 1. ### 5.4 Systematic errors For the $`\pi ^0,\eta `$ and $`\mathrm{K}^0`$ the largest contribution to the systematic error came from the parametrisation of the combinatorial background and the signal shape. This contributed up to $`50\%`$ of the systematic error for some $`\mathrm{x}_\mathrm{E}`$ intervals. This error was estimated using the following procedure for each $`\mathrm{x}_\mathrm{E}`$ interval: 1. The shape of the background was measured from data by fitting a second order polynomial to the non-signal regions of the two-photon invariant mass spectrum. The procedure was repeated using different side bands. 2. The shape of the background was fixed to the shape predicted by the Monte Carlo and the background was fitted to data allowing only an overall normalisation factor to vary. 3. The shape of the signal was fixed to the shape predicted by the Monte Carlo and was fitted to data allowing only an overall normalisation factor to vary. The systematic error was taken to be the quadratic sum of all deviations from the value measured using only data. For the $`\mathrm{K}^0`$ rate, the systematic errors in each $`\mathrm{x}_\mathrm{E}`$ interval were estimated using the procedure described in to which was added the contribution from the three items above. In the case of the $`\eta ,`$ the combinatorial background was found to have a small structure at $`\mathrm{M}_{2\gamma }700\mathrm{MeV}`$ in the Monte Carlo which was not seen in the data. This structure was caused by $`\omega (780)\pi ^0\gamma 3\gamma `$ for which the rate in the OPAL-tuned Monte Carlo was twice the measured rate. To estimate the systematic error the above procedure was repeated by: 1. fitting the background shape outside the signal and the small structure regions; 2. removing from the Monte Carlo $`50\%`$ of the generated $`\omega (780)3\gamma ,`$ in which case the structure disappeared. For both the $`\pi ^0`$ and the $`\eta `$ all the following systematic variations were considered: 1. The systematic error relative to the cut on the probability P associated with each photon pair was estimated by repeating the measurements for each $`\mathrm{x}_\mathrm{E}`$ bin with different cuts on $`\mathrm{P}>0.2,\mathrm{P}>0.3\mathrm{and}\mathrm{P}>0.4`$, the original value being obtained with $`\mathrm{P}>0.1.`$ This error, which took into account signal purity and reconstruction efficiency, since they depend on P, was added quadratically to the previous error. In the worst case it contributed $`14\%`$ to the total systematic error. 2. The difference between the correction factors for detector effects obtained with Jetset and those obtained with Herwig was considered as a systematic error. It was found to contribute up to $`12\%`$ of the total systematic error. 3. The error due to the energy calibration of the electromagnetic calorimeter was estimated from Monte Carlo by shifting up and down the energy of measured electromagnetic clusters by $`2\%.`$ This had a negligible effect on the number of reconstructed photon candidates. However, the position of the peaks of the $`\pi ^0`$ and $`\eta `$ signals were shifted by about 10 MeV, and the difference in the signal extracted in each $`\mathrm{x}_\mathrm{E}`$ bin was considered as a systematic error. In the worst case, this contributed $`8\%`$ to the total systematic error. 4. Some of the mesons, mainly $`\pi ^0,`$ were produced in interactions with detector material. This effect might not have been very well modelled in the Monte Carlo. Therefore, half of the Monte Carlo prediction for these mesons was included in the uncertainty. This accounted for up to $`25\%`$ of the total systematic error, the worst being in the low $`\mathrm{x}_\mathrm{E}`$ intervals. 5. An alternative extrapolation scheme as in was used to correct for the inaccessible $`\mathrm{x}_\mathrm{E}`$ regions. This yielded slightly different corrections and these differences were included in the systematic error. ## 6 $`\pi ^0,\eta `$ and $`𝐊^0`$ production in quark and gluon jets The $`\pi ^0`$, $`\eta `$ and $`\mathrm{K}^0`$ yields in quark and gluon jets were estimated as a function of the scale $`\mathrm{Q}_{\mathrm{jet}}`$. The average number of mesons produced in the second jet and third jet samples was measured as a function of $`\mathrm{Q}_{\mathrm{jet}}`$, and then the unfolding to $`100\%`$ quark and gluon jet purities was performed in the same way as was done for charged particles. Each $`\pi ^0,\eta `$ and $`\mathrm{K}^0`$ candidate was assigned to the jet which made the smallest opening angle with respect to the total momentum direction of the meson. The $`\pi ^0,\eta `$ and $`\mathrm{K}^0`$ signals were then extracted as in the case of the inclusive rate measurement (see Section 5). To make the jets fully comparable, and because photons were reconstructed only in the barrel region of the electromagnetic calorimeter, only events that had both the second and third jet in the barrel were kept. The $`\mathrm{Q}_{\mathrm{jet}}`$ interval was divided into three bins: $`\mathrm{Q}_{\mathrm{jet}}`$ = 7 to 13, 13 to 19 and 19 to 25 GeV. The average number of mesons per jet, in each $`\mathrm{Q}_{\mathrm{jet}}`$ interval, was corrected for detector acceptance, reconstruction efficiency and for inaccessible $`\mathrm{x}_\mathrm{E}`$ regions. Efficiencies were calculated separately for each jet as a function of $`\mathrm{Q}_{\mathrm{jet}}`$. Using the average purity of each bin of $`\mathrm{Q}_{\mathrm{jet}}`$ evaluated from data using the matrix element formula, the average number of mesons per jet was unfolded to $`100\%`$ quark and gluon jet purities. ### 6.1 $`\pi ^\mathrm{𝟎}`$ production The average number of $`\pi ^0`$ produced per jet as a function of $`\mathrm{Q}_{\mathrm{jet}}`$ is shown in Figure 9a for gluon and quark jets. The ratio of multiplicities in gluon and quark jets is shown in Figure 9b. The Jetset and Herwig models were found to reproduce the data within the error bars. This is demonstrated in Figure 9c where the ratios of data to Monte Carlo are compatible with unity for both quark and gluon jets in each interval of $`\mathrm{Q}_{\mathrm{jet}}.`$ The analytical function obtained for the average charged particle multiplicity as a function of $`\mathrm{Q}_{\mathrm{jet}}`$, scaled with only one free overall normalisation factor (measured to be 0.47), were found to fit well to the rate of $`\pi ^0`$ as a function of $`\mathrm{Q}_{\mathrm{jet}}`$ in gluon and quark jets, and to the ratio between $`\pi ^0`$ production rates in gluon and quark jets. The function is shown in Figures 9a and 9b. ### 6.2 $`\eta `$ production The same analytical function, but with an overall normalisation factor of 0.047, was found to describe the $`\eta `$ rate in gluon and quark jets, as shown in Figure 10a. The ratio of $`\eta `$ multiplicities in gluon and quark jets as a function of $`\mathrm{Q}_{\mathrm{jet}}`$ is shown in Figure 10b, and was compatible with the ratio of charged particle multiplicities. The ratio was also compatible with being independent of $`\mathrm{Q}_{\mathrm{jet}}`$, with $$\frac{\mathrm{N}_\mathrm{g}^\eta }{\mathrm{N}_\mathrm{q}^\eta }=1.29\pm 0.14.$$ The measured $`\eta `$ rate was found to be slightly higher in the data than in the Monte Carlo, mainly at low $`\mathrm{Q}_{\mathrm{jet}}`$. This small disagreement was the same for both quark and gluon jets, as shown in Figure 10c, where the ratios of data to Monte Carlo for gluon and quark jets are shown. This indicates that modelling of production of $`\eta `$ is equally inadequate for both Jetset and Herwig Monte Carlos. No additional enhancement of $`\eta `$ production in gluon jets was observed. Figure 8 indicates that the measured $`\eta `$ spectrum was harder than the Monte Carlo prediction. To investigate this, the analysis was repeated for $`\mathrm{x}_\mathrm{E}0.1`$ for both $`\pi ^0`$ and $`\eta `$ mesons. Due to statistical limitations, mainly in the $`\eta `$ meson sample, a yet harder cut on $`\mathrm{x}_\mathrm{E}`$ was not appropriate. The resulting $`\pi ^0`$ production rates as a function of the scale $`\mathrm{Q}_{\mathrm{jet}}`$ for the gluon and the quark jets were in good agreement with both Jetset and Herwig expectations. In each bin of $`\mathrm{Q}_{\mathrm{jet}}`$, the ratio of data to Monte Carlo for the $`\eta `$ rate was equal for quark and gluon jets. The ratio of the $`\eta `$ production rates in data to those predicted by the Monte Carlo were found to be the same for both the quark and gluon jets. The ratio of the production rate of $`\eta `$ in gluon jets to that in quark jets was also found to be consistent with being independent of $`\mathrm{Q}_{\mathrm{jet}}`$. The measured ratio was: $$\frac{\mathrm{N}_\mathrm{g}^\eta (\mathrm{x}_\mathrm{E}>0.1)}{\mathrm{N}_\mathrm{q}^\eta (\mathrm{x}_\mathrm{E}>0.1)}=1.38\pm 0.19.$$ This result is not in contradiction with previously published results concerning an excess of high momentum $`\eta `$ mesons over Monte Carlo prediction. The excess is present equally in quark and gluon jets. ### 6.3 $`𝐊^\mathrm{𝟎}`$ production An overall scale factor of 0.093 in the analytical function of Section 4 gave the best fit to the $`\mathrm{K}^0`$ production rate in quark and gluon jets. However, as shown in Figure 11a, the data for the quark jet were systematically but not significantly, higher than the corresponding analytical shape. This is also shown in Figure 11b where the ratio of $`\mathrm{K}^0`$ production in gluon and quark jets was systematically smaller than the corresponding ratio for charged particle production, although compatible within errors. This effect could be explained by a higher probability of producing a strange meson in the fragmentation of an initial strange, charmed or b quark. The ratios data to Monte carlo of the $`\mathrm{K}^0`$ production rate in gluon and quark jets are shown in Figure 11c. The ratios are compatible with unity within errors and agree with the previous OPAL result . ### 6.4 Systematic errors The measurements of the production rates of $`\pi ^0,\eta `$ and $`\mathrm{K}^0`$ in jets were not statistically limited. The errors shown in Figures 9, 10 and 11 already include both the systematic and statistical errors added in quadrature and propagated through the unfolding formula. The systematic error ranged from $`60\%`$ to $`90\%`$ of the total quoted error, depending on the data point. It was estimated for each jet and each $`\mathrm{Q}_{\mathrm{jet}}`$ interval using the procedure described in Section 6 for $`\pi ^0`$, $`\eta `$ and $`\mathrm{K}^0`$ inclusive rate measurements. In each $`\mathrm{Q}_{\mathrm{jet}}`$ interval, the selection criteria (see Section 5) for $`\pi ^0,\eta \mathrm{and}\mathrm{K}^0`$ mesons were found to act equally on quark and gluon jets. The results were very stable against changes of the selection cuts that caused relative reconstruction efficiency variation of up to $`5\%.`$ For the $`\eta `$ case, the additional cut (Section 5.1) that excluded for invariant masses $`\mathrm{M}_{2\gamma }>`$ $`300\mathrm{MeV},`$ any photon that could pair with any other photon to make an invariant mass $`\mathrm{M}_{2\gamma }<300\mathrm{MeV}`$ with a probability $`\mathrm{P}>0.1`$, was changed: first $`\mathrm{M}_{2\gamma }<300\mathrm{MeV}`$ was moved to $`\mathrm{M}_{2\gamma }<200\mathrm{MeV}`$ and second $`\mathrm{P}>0.1`$ was replaced by $`\mathrm{P}>0.2.`$ The effect on the $`\eta `$ measured rates in quark and gluon jets was negligible($`<1\%`$). In addition, for all studied mesons: 1. The difference between the measurements with purities taken from the matrix element formula and from Monte Carlo(see Section 3) was negligible. The average jet topologies and purities were varied by changing the cut on the variable $`\mathrm{Y}`$ (see Section 2.3) such that the purities obtained with the matrix element formula still agreed with those obtained from Monte Carlo information. The analysis was repeated with $`\mathrm{Y}0.25`$ and the differences obtained in each bin of $`\mathrm{Q}_{\mathrm{jet}}`$ were considered as systematic errors and were found to contribute at most $`5\%`$ of the total quadratic sum. 2. Since the mesons were reconstructed independently of the jet-finding, and were assigned to jets by angular matching once the jets were reconstructed, very little dependence on the jet finder was expected. Indeed, this was the case for Luclus and Durham jet finders where the difference was measured to be less than $`2\%`$. However, the results obtained with the cone jet finder showed deviations of up to $`10\%`$ with respect to the two other jet finders. This was considered as a systematic error, and contributed up to $`47\%`$ of the total quadratic sum of the systematic errors. 3. The results were found to be stable in each bin of $`\mathrm{Q}_{\mathrm{jet}}`$ with respect to changes of the charged particle selection requirements. The results were also found to be stable when the number of charged particles required per jet was increased from 3 to 5. Lowering the cut on the inter-jet angle from 30 to $`20^{}`$ was correlated to the cut on the variable Y. The analysis was repeated with the minimum inter-jet angle set to $`20^{}`$ and the minimum value of Y set to 0.22. The maximum deviation obtained was less than $`2\%`$ and was added to the total systematic error. 4. Due to statistical limitations, the analysis could not be repeated on data using b-tagging to obtain gluon enriched jet samples. A systematic error, that could account for correlations between particle content in pure quark and pure gluon jets as well as for the stability of the unfolding method, was assigned using the following procedure: The analysis was repeated on four different Jetset Monte Carlo samples, each sample being as large as the full available data sample. The jet samples were selected to have different gluon and quark average purities. The purities were set based on Monte Carlo information (Section 3.2) to be: $`(40\%\mathrm{quark},60\%\mathrm{gluon}),`$ $`(25\%\mathrm{quark},75\%\mathrm{gluon}),`$ $`(60\%\mathrm{quark},40\%\mathrm{gluon}),`$ $`(75\%\mathrm{quark},25\%\mathrm{gluon})`$ for the four pairs of jet samples that were processed with the unfolding method. The largest deviation with respect to the average unfolded meson rate was considered as a systematic error. This contributed up to $`7\%`$ of the systematic error, indicating that the unfolding procedure was indeed stable against large purity variations. ## 7 Summary and conclusion Average multiplicities of $`\pi ^0`$, $`\eta `$, $`\mathrm{K}^0`$ and charged particles have been measured for gluon and quark jets as a function of a transverse momentum-like scale $`\mathrm{Q}_{\mathrm{jet}}`$. The average multiplicities were unfolded to $`100\%`$ purity by comparing two jet samples having different gluon (or quark) content for the same value of $`\mathrm{Q}_{\mathrm{jet}}`$. The $`\pi ^0,\eta `$ and $`\mathrm{K}^0`$ inclusive production rates in $`\mathrm{Z}^0`$ hadronic decays were measured and found to agree with the previous values measured at LEP. The $`\pi ^0`$ production rate was found to be well described by both Jetset and Herwig Monte Carlos for all values of $`\mathrm{x}_\mathrm{E}`$. The $`\eta `$ and $`\mathrm{K}^0`$ spectra were found to be harder than the Monte Carlo predictions, particularly in the case of $`\eta `$. The charged particle multiplicity as a function of the topological scale $`\mathrm{Q}_{\mathrm{jet}}`$ for pure quark and pure gluon jets was described by a simple phenomenological formula. The same formula, with all parameters fixed except for an overall normalisation factor, was found to provide a good fit to the $`\pi ^0,`$ $`\eta `$ and $`\mathrm{K}^0`$ production rates in gluon and quark jets. The analysis showed that for $`\pi ^0,\eta `$ and $`\mathrm{K}^0`$ mesons, there was no evidence for an enhancement in gluon jets with respect to quark jets beyond the enhancement observed for inclusive charged particles. In particular we observed no evidence for an enhancement of $`\eta `$ meson production in gluon jets, contrary to the predictions of some models for gluon jet hadronisation. We measured the ratio of the slope of the average charged particle multiplicity in gluon jets to that in quark jets. We obtained $`\mathrm{C}=2.227\pm 0.07(\mathrm{stat}.)\pm 0.19(\mathrm{syst}.)`$ for this ratio, that is about one standard deviation of the total uncertainty above the analytic next-to-next-to-next-to leading order (3NLO) prediction . ## 8 Acknowledgements We particularly wish to thank the SL Division for the efficient operation of the LEP accelerator at all energies and for their continuing close cooperation with our experimental group. We thank our colleagues from CEA, DAPNIA/SPP, CE-Saclay for their efforts over the years on the time-of-flight and trigger systems which we continue to use. In addition to the support staff at our own institutions we are pleased to acknowledge the Department of Energy, USA, National Science Foundation, USA, Particle Physics and Astronomy Research Council, UK, Natural Sciences and Engineering Research Council, Canada, Israel Science Foundation, administered by the Israel Academy of Science and Humanities, Minerva Gesellschaft, Benoziyo Center for High Energy Physics, Japanese Ministry of Education, Science and Culture (the Monbusho) and a grant under the Monbusho International Science Research Program, Japanese Society for the Promotion of Science (JSPS), German Israeli Bi-national Science Foundation (GIF), Bundesministerium für Bildung und Forschung, Germany, National Research Council of Canada, Research Corporation, USA, Hungarian Foundation for Scientific Research, OTKA T-029328, T023793 and OTKA F-023259.
warning/0007/astro-ph0007370.html	ar5iv	text	# Daemon detection experiment ## 1 Introduction. The daemon hypothesis Our Universe started from Planckian scales, and it appears only natural to suggest that the major part of its mass, i.e., the DM, is contained in primordial Planckian particles with $`M=(\pi \mathrm{}\mathrm{c}/4\mathrm{G})^{1/2}2\times 10^5`$ g and $`r_\mathrm{g}=2\mathrm{G}M/\mathrm{c}^23\times 10^{33}`$ cm . Such elementary black holes can be stable and eternal . Multidimensional ($`>`$4) theories (e.g. ) allow the existence on them of stable electric charge of up to $`Z\mathrm{e}=\mathrm{G}^{1/2}M10\mathrm{e}`$. We assume that such DArk Electric Matter Objects, i.e. daemons, carrying a charge of any sign (including possibly the zeroth one) constitute a hierarchy of populations, viz. intergalactic, Galactic halo (or crown), Galactic disk etc, with continuously increasing concentration and decreasing mean random velocity. In the disk, the DM makes up about 1/2 its total mass . If it consists of daemons, their density here is $`10^{12}\mathrm{m}^3`$. Due to their extremely low concentration, negligibly small dimensions and a giant inertia, the daemons, by themselves, constitute a non-collisional plasma and, thus, do not capture and coalesce with each other. The negative daemons’ buildup inside the Sun is capable of accounting for its energetics through catalysis of the proton fusion reactions, and for the deficiency of the electron-capture neutrinos . Because of their large mass, daemons have a giant penetrating power. That is why Markov was skeptical concerning a possibility of detecting such particles. A charged daemon falling from infinity onto the Sun acquires $`10^{23}`$ eV, while losing $`10^{19}`$$`10^{22}`$ eV in passing along the Solar diameter . Prior to becoming trapped, a daemon traverses many times the Sun along gradually contracting, strongly elongated orbits, whose perihelia lie within the Sun. If, in the course of this process, the daemon enters the Earth’s sphere of influence, its path will change slightly, and the perihelion of the orbit will leave the Sun with a high probability. This is how the daemons can populate strongly elongated Earth-crossing heliocentric orbits (SEECHOs). The flux of this population at the Earth is estimated to be $`f_{}10^3`$$`10^6\mathrm{m}^2\mathrm{s}^1`$ at $`V35`$–50 $`\mathrm{km}\mathrm{s}^1`$ . Part of this population will be gradually transferred into the near-Earth almost-circular heliocentric orbits (NEACHOs), including those crossing the Earth’s surface. If the Sun moves relative to the daemon population of the Galactic disk, these fluxes may undergo seasonal variations at the Earth. The original goal of this work was to detect the most dense (as we believed) SEECHO population. In our opinion, such an approach should be more productive than the standard searches for the strongly rarefied, Galactic halo DM population (e.g. ). One may conceive the following consequences of the negative daemon interaction with matter on the atomic and subnuclear levels: (1) When capturing (or recapturing) a nucleus, the daemon first captures the ion of the atom, so that as it is dropping to ever lower-lying levels, all the electrons of the captured ion are emitted in the Auger process. Because the binding energy of the daemon to the nucleus is measured in tens of MeV, the energy of such Auger electrons may be as high as $``$0.1–1 MeV; (2) Catalysis of the fusion of light nuclei ($`Z_\mathrm{n}<Z`$), including protons . As a result of internal conversion, the energy of fusion in the vicinity of a heavy charged particle is converted, as a rule, to the kinetic energy of the resultant nucleus; (3) Capture of heavy nuclei ($`Z_\mathrm{n}Z`$). At $`Z=10`$, it occurs in solid Fe, Zn, Sn in $``$2 $`\mu \mathrm{m}`$ at $``$50 $`\mathrm{km}\mathrm{s}^1`$, and in $``$0.1 $`\mu \mathrm{m}`$ at $``$10 $`\mathrm{km}\mathrm{s}^1`$ . This ‘poisons’ the catalytic properties of the daemon. However, straightforward estimates based on the solar energetics, if it is provided by daemon-assisted catalysis of proton fusion, suggested that the daemon should free itself of a captured heavy nucleus in $`\tau _{\mathrm{ex}}0.1`$–1 $`\mu \mathrm{s}`$, apparently as a result of the decay of the daemon-containing proton (for $`Z_\mathrm{n}>24/Z`$ the ground level of the daemon in a nucleus lies inside one of its constituent nucleons). Note that the uncertainty in the values of the parameters used in these estimates permits one to shift $`\tau _{\mathrm{ex}}`$ in either direction at least by an order of magnitude. (4) If the ‘poisoned’ daemon with $`ZZ_\mathrm{n}0`$ encounters a heavier nucleus, due to great difference in binding energies, it captures the latter while losing the previous one if at the moment it is (or becomes) noticeably lighter than the anew met nucleus (at $`ZZ_\mathrm{n}=0`$ this ‘recapture’ is an analog of the charge exchange of ions moving through a neutral gas); (5) Finally, the capture of a nucleus strongly excites its internal degrees of freedom, which should stimulate the emission of nuclear radiations (numerous nucleons and their clusters, $`\gamma `$-quanta). Nuclei with a high $`Z_\mathrm{n}`$ have a higher probability for the initiation of these processes. Thus, all the above processes stimulate the emission of radiations capable of producing scintillations, and this may be used to detect the slowly moving daemons, whose direct impact cannot generate a scintillation. Earlier, we attempted to use the catalysis of light nucleus fusion to detect the daemons. We employed an acoustic method with two Li plates . The basic idea here was that the material surrounding the trajectory should be rapidly heated by the products of the $`2{}_{}{}^{7}\mathrm{Li}{}_{}{}^{14}\mathrm{C}`$ reaction. The thermal expansion of the material should generate a sound wave, which then would be detected by piezoelectric sensors. Their output signals would trace the daemon trajectory. Our experiments revealed, however, an unexpectedly strong damping of ultrasound in Li, and the experiments were stopped. The measurements made with thick Be plates (4.5 cm), 0.12 $`\mathrm{m}^2`$ in area, coated on both sides with ZnS(Ag) scintillator, were actually a continuation of the above experiments along the same lines. It was assumed that the products of the $`2{}_{}{}^{9}\mathrm{Be}{}_{}{}^{18}\mathrm{O}`$ reaction escaping from thin near-surface regions, i.e. the points of entrance into and egress from the Be plate of a daemon moving along a SEECHO, would produce scintillations $``$1 $`\mu \mathrm{s}`$ apart. 500 h of exposure did not yield any result . The latter experiment, however, prompted us to consider other possible modes of daemon interaction with matter, particularly the possible role of their ‘poisoning’ by heavy nuclei \[Fe, Si impurities in Be, and the ZnS(Ag) nuclei\]. Shedding the nucleus poisoning the daemon, combined with the above-mentioned processes (1)–(5), suggest a variety of methods for detection of slow daemons. Because of their large number, both the Auger electrons and the radiations from heavy nuclei excited in the capture are the most efficient in this respect. ## 2 Description of the setup and assumed sequence of events triggered by a daemon These considerations served as a basis when developing an ideologically new and very simple setup of four modules. Each module contains two parallel transparent polystyrene plates 4 mm thick and $`50\times 50`$ cm in size. The distance between the plates was 7 cm. One of the polystyrene plate surface was coated by type B3-s ZnS(Ag) powder $`3.5\mathrm{mg}\mathrm{cm}^2`$ thick. Its average grain size is 12 $`\mu \mathrm{m}`$. The choice of this classical phosphor, besides its availability, simplicity in use, and a high light output, was motivated also by the fact that it consists of medium-$`Z_\mathrm{n}`$ elements. And conversely, the choice of polystyrene as a plate material was stimulated by its low $`Z_\mathrm{n}`$, in order to reduce to a minimum the possibility of heavy-nucleus poisoning of the daemon before its traversal of the ZnS(Ag) layer. Each plate was viewed from a distance of 22 cm by one FEU-167 PM tube with a dia. 100-mm photocathode. The plates were separated by a sheet of black paper. To be able to judge the essential features of the possible incoming daemon flux, the ZnS(Ag)-coated surfaces of the polystyrene plates were set facing the same side (down). As a result, the light entering the top PM tube passed also a 4-mm thick transparent polystyrene plate on the way. The polystyrene plates and the PM tubes viewing them were placed in a cubic case 51 cm on the edge made of 0.3-mm thick iron sheet with double-sided thin ($``$2 $`\mu \mathrm{m}`$) facing of tin. The upper horizontal case face was made of two sheets of black paper. The PM photocathodes were arranged flush with the horizontal case faces. All the four modules were placed side by side in one horizontal plane. The total area of the four-module detector was 1 $`\mathrm{m}^2`$. The PM tubes were powered by a voltage corresponding to their sensitivity of 10 $`\mathrm{A}\mathrm{lm}^1`$. A 4.5-k$`\mathrm{\Omega }`$ resistance served as a load. The signal from the load was supplied through a 1 mH inductance and a cable of total capacity 550 pF to an oscilloscope. The purpose of this $`LC`$ circuit was to stretch the leading edge of the pulse so as to facilitate discrimination of the long scintillations produced by heavy nonrelativistic nuclei like $`\alpha `$-particles (Heavy-Particle Scintillations — HPS; for their characterization the $`{}_{}{}^{238}\mathrm{Pu}`$ $`\alpha `$-source was used) against PM tube noise and short scintillations caused by low-mass and relativistic particles (cosmic rays etc) (the Noise-Like Scintillations — NLS). Signals from the two PM tubes of the same module were fed for comparison into two inputs of one S9-8 dual-trace digital oscilloscope. The latter was triggered by the output signal of the upper PM tube if it reached a level $`U_12.5`$ mV. The signal from the second PM tube was considered significant if it increased in 0.4–1.5 $`\mu \mathrm{s}`$ and its amplitude was $`U_20.6`$ mV. The signals from the oscilloscopes recorded during $``$100 $`\mu \mathrm{s}`$ before the trigger (i.e. with a lead) and during $`+`$100 $`\mu \mathrm{s}`$ after the trigger (i.e. with a delay) were entered into computer memory if they were seen in both traces. The system was from the outset designed to detect daemons impinging on it primarily from above. Indeed, it was assumed originally that a daemon propagating through the polystyrene plate from above would capture a heavy nucleus from the ZnS(Ag) layer coating it on the back side, and emit in the process during a certain time numerous Auger electrons and the excited nucleus radiations which produce a protracted scintillation of HPS type. If this was a daemon from the SEECHO population, it would, in passing at a velocity $`35`$$`50\mathrm{km}\mathrm{s}^1`$ in $`\tau _{\mathrm{ex}}0.1`$–1 $`\mu \mathrm{s}`$ a path $``$0.3–5 cm, cause decay of protons in the captured nucleus. Then the nucleus emits products of the proton decay (pions etc.) and, possibly, fragments of the nucleus itself. All of them also can give rise to scintillation events, which also will be detected by the top PM tube. One can in principle conceive catalytic fusion reactions of nitrogen and/or oxygen nuclei captured in air, but the range in air of the heavy products of these reactions is $`<`$0.5 cm, so that the probability of their detection after the daemon has traversed the top polystyrene plate appears low. The same applies to carbon and hydrogen, the components of the polystyrene, all the more so that in their subsequent passage through the ZnS(Ag) layer the already captured $`{}_{}{}^{12}\mathrm{C}`$ or $`{}_{}{}^{1}\mathrm{H}`$ nuclei would not have a good chance to react outside the plate, as they would be promoted to higher lying levels, and lost in the preferential capture of heavier nuclei in ZnS(Ag). On passing the 7-cm gap, the daemon traverses the second polystyrene plate with the ZnS(Ag) coating on its lower surface and enters the space bounded by a semi-cubic sheet-metal case, which fixes the PM tube at a distance of 22 cm from the scintillator. The presumed phenomena occurring here are similar to those taking place after the passage of the upper polystyrene plate, however there is more room (and time) for their manifestation. When moving from below, daemon enters the chamber viewed by the second PM tube through a 0.3 mm-thick tinned iron sheet. On exiting the Sn or Fe atom, it carries away with a high probability its captured nucleus. On the capture, some liberated energetic Auger electrons are able to escape from the sheet to excite the bottom scintillator. Afterwards, in $`\tau _{\mathrm{ex}}`$, the daemon releases the captured nucleus with emission of the products of nucleon decay. After this, the daemon captures a light nucleus in the air and emits Auger electrons, which also excite scintillations in the bottom plate. In traversing this plate in its motion upward, the daemon gets rid in it and in its ZnS(Ag) layer of the light nucleus captured in the air and recaptures a heavier ion (or nucleus) from ZnS(Ag) carrying it into the polystyrene bulk. While being in the ZnS(Ag) $``$10 $`\mu \mathrm{m}`$-size grain during $`<10^9`$ s, the daemon probably has no time to force the nucleus to emit all its radiations. So it is doubtful that some particles shed in the polystyrene plate bulk would be capable of reaching the ZnS(Ag) bottom layer to excite an intensive scintillation. When entering further the 7-cm gap between the plates or the ZnS(Ag) upper layer, the daemon (re)captures here again such nuclei with the emission of Auger electrons and nuclear radiations. These latter have to be detected by the first (upper) scintillator that triggers the oscilloscope. However keeping in mind that some of the particles are emitted after the daemon has penetrated well into the polystyrene bulk, one cannot be sure of having excited a strong scintillation. Thus, there is a variety of conceivable processes capable of creating a scintillation at the passage of a daemon. We shall not attempt to list them all here. It is also clear that far from all of the weak scintillations are recorded. The strongest are apparently the informative events initiated at a time by numerous ($``$10) Auger electrons and by nuclear particles. Therefore we can lower our electronics response level and thus we have no serious problems with numerous one-particle background event discrimination (see Fig.1 below). In any case, however, the response of our setup should be asymmetric with respect to the daemon propagation direction. ## 3 Some specifics of the experiment and its results After control tests in January-February 2000, the system was put in round-the-clock operation in March 2000, with the total exposure amounting to 700 h. Altogether, $``$6$`\times `$$`10^5`$ oscilloscope triggering events have been recorded, only $``$$`10^4`$ of which contained a signal on the second trace and were entered into the computer. About 2/3 ‘single’ triggers contain a tailing signal typical of HPS. These signals are most likely due to radioactive background decays. The remaining single triggers are of the NLS type. The double events are primarily NLS signals occurring without any time delay and coinciding in shape (delays $``$0.2 $`\mu \mathrm{s}`$). Sometimes these signals appear simultaneously even in all four modules. We assign such events to cosmic rays and neglect them. Very infrequently, once in only about twenty–thirty of all the events recorded, one of the two signals has the HPS characteristics. Interestingly, events with two and more significant signals in one channel are very rare. Therefore, while basing on the above scenarios of the events accompanying a daemon traversal of our system, one could expect numerous signals on the same oscilloscopic trace, we began with analyzing the sweeps containing one signal only. The number of events with shifted signals in both traces recorded during the month is 413. In the case of purely non-correlated stochastic generation of signals, no statistically significant clusters of events should appear in the time distribution of second-trace signals. This experiment was aimed at detecting objects moving with velocities $``$35–50 $`\mathrm{km}\mathrm{s}^1`$. We expected signals with positive shifts $`\mathrm{\Delta }t1.52.0`$ $`\mu \mathrm{s}`$ or slightly more. As always, reality introduced substantial corrections into our speculations. This relates to both the daemon population discovered by us and the details of the events (see Sec.2 above) accompanying daemon passage through our system. Fig.1 shows an $`N(\mathrm{\Delta }t)`$ distribution of second-trace signals in the time $`\mathrm{\Delta }t`$ of their shift relative to the onset of the triggering signal on the first oscilloscopic trace. This is a fairly asymmetric and non-monotonic distribution. By the $`\chi ^2`$ criterion, the C.L. of this being not a $`\mathrm{\Delta }t`$-independent distribution is not less than 99.8%. First of all, there is no noticeable clustering of events near a few $`\mu \mathrm{s}`$, as we expected to be based on the hypothesis of existence of the SEECHO population. The main feature is a maximum in the region of +30 $`\mu \mathrm{s}`$ for an average 41.3 event/bin level. It contains 62 events, which exceeds by a factor of 3.22 the statistically allowable scatter $`\sigma =\sqrt{41.3}6.43`$. Accepting only an excess of 1.22 over 2$`\sigma `$, this yields seven to eight events which can be assigned to nuclear-active objects crossing the detector. For its area of 1 $`\mathrm{m}^2`$ and an exposure time of $`2.5\times 10^6`$ s, this amounts to a total flux $`f_{}3\times 10^6\mathrm{m}^2\mathrm{s}^1`$. Our initial goal was detection of the DM objects, not their flux measurement. We possibly are not able now to reveal all the daemons traversing our system due to their partial poisoning with heavy nuclei, etc. So the real flux of daemons can reach $`f_{}5\times 10^5\mathrm{m}^2\mathrm{s}^1`$. An absence of a symmetric to $`\mathrm{\Delta }t+30`$ $`\mu \mathrm{s}`$ negative feature in $`N(\mathrm{\Delta }t)`$ distribution can be attributed to long-lasting Sn nucleus poisoning of daemons moving from below. The various cross-checks did not reveal any noticeable systematic errors in the operation of our simple measuring instrumentation. ## 4 An attempt at interpreting the results An analysis of the $`N(\mathrm{\Delta }t)`$ distribution displayed in Fig.1 permits certain conclusions both on the nature of the agent responsible for this distribution and on the character of its interaction with matter. While one cannot rule out a possibility of other interpretations, we shall try to treat the results within the daemon hypothesis. We note immediately that when compared with the 7 cm interplate distance, the position of the maximum on the time axis (+30 $`\mu \mathrm{s}`$) indicates a fairly low velocity of the maximum-forming objects. This velocity is about 2–3 $`\mathrm{km}\mathrm{s}^1`$ only. Taking into account the possible slope of the trajectories could double the velocity at most. Initially, we attempted to explain such a low velocity as due to its characterizing the population in geocentric orbits intercepting the Earth surface, which was captured from the one moving in SEECHOs . This interpretation meets, however, with a difficulty pointed out as far back as 1965 by Markov . The fact is that despite their giant penetrating ability the daemons moving with a velocity of $``$10 $`\mathrm{km}\mathrm{s}^1`$ can traverse the Earth only $``$$`10^2`$$`10^3`$ times. Their buildup in the Earth during 4.5 Byr and interaction with the material (even if only the proton decay with an energy release $``$1 BeV during $`\tau _{\mathrm{ex}}=10^5`$ s is taken into account, see below) should bring about an energy dissipation corresponding to a heat flux of $``$2$`10^5`$ $`\mathrm{erg}\mathrm{cm}^2\mathrm{s}^1`$. This figure exceeds at least by four orders of magnitude the flux emanating from the Earth’s mantle (10 $`\mathrm{erg}\mathrm{cm}^2\mathrm{s}^1`$). The assumption of the daemon velocities ranging widely in magnitude and directions comes also in conflict with the narrowness of the maximum ($`20<\mathrm{\Delta }t<40`$ $`\mu \mathrm{s}`$) in the $`N(\mathrm{\Delta }t)`$ distribution. In view of the fact that this distribution was obtained from sweeps containing only one signal, it appears that we were too optimistic by assuming that the daemon frees itself of the captured heavy nucleus and recovers the catalytic properties in as short a time as $`\tau _{\mathrm{ex}}=0.1`$–1 $`\mu \mathrm{s}`$, which is shorter than the time required to cross the 7-cm gap between the plates. As already mentioned, an analysis of the solar energetics allows considerably larger values, up to $`\tau _{\mathrm{ex}}=10^5`$$`10^4`$ s. We have thus to admit that our starting scenario of a possible sequence of events (see Sec. 2) initiated by a daemon traversal of the detector contains excess or weakly revealing processes (capture of nuclei from the air as a result, say, of the daemon shedding heavy nuclei during the time it crosses the case etc.). One could go still further and assume that the capture by the daemon of a new nucleus, which is accompanied by ejection of a large number of Auger electrons and nuclear particles, occurs in our small system (10–50 cm) only when a new nucleus is recaptured during the entry into a material with a larger atomic weight. In this case, all pieces of the puzzle fall into place, and the sequence of the events accompanying the daemon traversal of the system looks somewhat differently. To begin with, on having crossed the roof (Fe, Zn) and the floors (Mg, Al, Si, Fe, O) of our building, the daemon reduces the mass of the captured nucleus in $`10^5`$$`10^4`$ s to such an extent that, on penetrating into the top ZnS(Ag) layer, it can already capture a S or Zn nucleus (this is possibly the only time where we directly invoke the hypothesis of the decay of a daemon-containing nucleon, so that for $`Z_\mathrm{n}=26`$ $`\tau _{\mathrm{ex}}10`$–100 $`\mu \mathrm{s}`$, a figure that still can be reconciled with the estimates based on solar energetics ). The Auger electrons, nucleons, and their clusters ejected in the process excite HPSs in the scintillator. Not having enough time to reduce substantially the mass of the nucleus captured here, the ‘poisoned’ daemon enters, 7 cm thereafter, the bottom scintillator, but it does not excite it. The NLS excitation in the bottom scintillator is triggered by the energetic $``$0.1–1 MeV long-range Auger electrons, which are ejected at the recapture of still heavier Sn or Fe nuclei, when the daemon reaches the lower lid of the tinned-iron case. (If it leaves the lower part of the case through its side wall most of the recapture Auger electrons released as the daemon traverses the material move almost perpendicular to this wall, i.e., parallel to the ZnS(Ag) layer. As a result, most of the electrons ejected from the side wall do not enter the scintillator and, thus, will not excite a strong scintillation.) The above reasoning suggests also that the side walls of the upper half of the case, while ”poisoning” the daemons crossing them by Sn or Fe nuclei, leave only a solid angle of $`4\pi /6`$ ster for the daemons to pass freely into the system (through the black paper). It thus becomes clear that the distance to be taken into account is the separation of 29 cm between the top scintillator and the lower case lid, and the angular spread of trajectories of the detectable daemons is limited by a solid angle of $``$2 ster. These factors are responsible for the narrowness of the maximum at 30 $`\mu \mathrm{s}`$ and yield 10–15 $`\mathrm{km}\mathrm{s}^1`$ for the velocity. We immediately see that the latter value is in a good agreement with the velocity of the objects falling on the Earth from NEACHOs. It appears that they are possibly transferred here through perturbations by the Earth (including traversal of its material) from the population in SEECHOs, which was captured by the Sun with the help of the Earth, and which was the initial target of our search. The concentration of the objects found by us in NEACHOs is determined by the balance between their capture and subsequent ejection 1–10 Myr later, through gravitational perturbations by the Earth, into the region of other planets action (as well as out of the Solar system altogether). Because of these orbits being close to that of the Earth, the flux of the particles of this population through the Earth should exceed the flux of the NEACHO-population replenishing SEECHO population, which is exactly what is observed. An object impacting with this velocity cannot excite atoms. The fact that oscilloscopic traces typical of data presented in Fig.1 exhibit only one pulse shifted relative to the primary one, and that there is no pulse paired with the trigger (i.e. without a shift in time) suggests that the radiation emitted in interaction of the daemon with matter has a low penetrating power. Polystyrene 4 mm thick stops it completely. If these are electrons, their energy is $`<`$1 MeV. In conclusion, consider the important information that can be gained by using for the $`N(\mathrm{\Delta }t)`$ plot only the events containing HPSs in the first channel (Fig. 1). In accordance with the above analysis of the sequence of the events initiated by daemon traversal of our detector, most of these events ($``$80%) have NLSs in the second channel; this is a consequence of the recapture by daemons of nuclei in the bottom lid accompanied by the ejection of electrons, with most of the latter, after passing a distance of 22 cm in the air, impinge on the bottom ZnS(Ag) screen. While this distribution consists of 212 events only, the C.L. of its being different from $`N(\mathrm{\Delta }t)`$ = const, calculated by the $`\chi ^2`$ criterion, becomes 99.9%. All the 39 events responsible for the maximum near +30 $`\mu \mathrm{s}`$ are concentrated in this distribution. This maximum exceeds the mean level by a factor 3.86$`\sigma `$. Fig.1 displays also an HPS distribution for 10-$`\mu \mathrm{s}`$ wide bins. The absence of features in the remaining NLS distribution provides one more argument for the events of interest to us here not being the result of interference or regular instrumental malfunctions. ## 5 Discussion and main conclusions One usually searches for DM objects with $`V=200`$–300 $`\mathrm{km}\mathrm{s}^1`$, which populate the Galactic halo. We were the first to look for a much denser near-Sun population . In choosing a method capable of their detection, we counted primarily on the specific activity of the objects of this population at the nuclear and subnuclear levels rather than on the purely collisional interactions with particles of the matter, as this is done, for instance, in experimental search for much less massive (and more numerous) hypothetical WIMPs. In these and similar experiments, the discrimination used to reveal the expected signals is so tight as to practically exclude the possibility of detecting fairly infrequent but very energetic daemon signals, which become manifest in specific conditions (for instance, when entering a material with a high atomic number). Disregarding some details in the possible interpretation of the results of our experiments, which at first glance might appear very simple, it can be maintained that we have detected at a $`\mathrm{C}.\mathrm{L}.99.9`$% an indication of the existence of some highly penetrative nuclear-active cosmic radiation whose objects move with low astronomical velocities ($`V=10`$–15 $`\mathrm{km}\mathrm{s}^1`$). They appear to have an enormous penetrating power and are capable of passing near and through the Earth to populate finally NEACHOs. The main indication of the existence of this superslow radiation is the non-stochastic component in the long-time shifts of signals of two PM tubes viewing phosphor coatings on two parallel plates arranged one beneath the other. This component crowds in the 20 $`\mu \mathrm{s}`$ range. The simplicity of our detector system, however, is only apparent, in that in actual fact it is a result of a hard search that has been going on for many years . The observed intense scintillations originate from a simultaneous ejection of many ($``$10) energetic particles, and the events occurring in the detector, as we have seen, are very complex. The order in which they appear and the manifestation itself depend on a number of factors which might seem to be of minor importance. Revealing these factors and their effectively operating combinations has become possible due to our gradually refining understanding of the various aspects of the problem. If we had not used tinned iron sheets for the scintillator cases, which play a role of not passive but of the essential pieces of detector, or black paper for their top covers, or elements of asymmetry in the scintillator coating orientation, etc., the efficiency of daemon detection would have dropped strongly. On the other hand, the observed distribution with such a well pronounced, narrow shifted peak cannot be produced, say, by neutrons thermalized somewhere outside the detector, or by long-lived nuclear states excited, for instance, by cosmic rays. Obviously enough, because of the unavoidable dispersion in the neutron velocity or nuclear deexcitation times such a distribution would be fairly difficult to produce even deliberately. In our case, with a sweep $`\pm `$100 $`\mu \mathrm{s}`$ long, these effects could at best increase slightly the average background level. Judging from its manifestations and properties, the discovered radiation can be identified in a rather self-consistent manner with daemons, i.e. hypothetical primordial Planckian objects carrying an electric charge and moving in close to the Earth’s orbits. The primary source of this population is apparently the DM of the Galactic disk. The intermediate stage is a SEECHO population producing apparently a $``$35–50 $`\mathrm{km}\mathrm{s}^1`$ flux at the Earth at a level noticeably lower than $`10^4\mathrm{m}^2\mathrm{s}^1`$, which accounts for our not having yet detected it. It is a search for similar objects with a local concentration enhanced strongly by the gravitational action of the Sun and the Earth that has initiated this experiment. Several modified experiments are currently under way. ## Acknowledgements The author is greatly indebted to M.V.Beloborodyy, R.O.Kurakin, V.G.Latypov, K.A.Pelepelin for providing software and electronics maintenance.
warning/0007/math-ph0007032.html	ar5iv	text	# 1 INTRODUCTION ## 1 INTRODUCTION There has been considerable progress during the last ten years in the mathematical study of long-time existence properties of solutions of geometrically-based classical field theories. A significant portion of this work has focussed on the study of what are called wave maps in the mathematics literature and nonlinear sigma models (or, in certain special cases, chiral models) in the physics literature. These are defined as maps $`\psi `$ from a Lorentzian geometry $`(M^{m+1},\eta )`$, e.g. Minkowski space, to a Riemannian geometry $`(N^n,g)`$, e.g. a symmetric space or a compact Lie group, with $`\psi `$ being a critical point for the functional <sup>1</sup><sup>1</sup>1 Integrals over $`M^{m+1}`$ are understood to use the natural volume form compatible with the metric $`\eta `$. $`S[\psi ]={\displaystyle _{M^{m+1}}}\eta ^{\mu \nu }g_{AB}(\psi )_\mu \psi ^A_\nu \psi ^B,`$ (1) and hence satisfying the wave map equation $`\eta ^{\mu \nu }_\mu _\nu \psi ^A+\mathrm{\Gamma }^A{}_{BC}{}^{}(\psi )_\mu \psi ^B_\nu \psi ^C\eta ^{\mu \nu }=0;`$ (2) here $``$ is the (torsion-free) derivative operator determined by the metric $`\eta `$, and $`\mathrm{\Gamma }`$ are the (torsion-free) connection coefficients compatible with $`g`$. Wave maps have a well-posed Cauchy problem, and it is known that for 1+1 dimensional base geometries $`(M^{1+1},\eta )`$, every choice of smooth initial data evolves into a global smooth solution , while for 3+1 (or higher) dimensional base geometries, certain smooth initial data leads to solutions with singularities . Not yet understood is what happens in general for 2+1 dimensional base geometries. This is the “critical dimension” (see ), where global smooth solutions are expected, at least, for all smooth initial data of sufficiently small energy. While global existence results are known to hold for certain classes of rotationally-symmetric wave maps in 2+1 dimensions (without restrictions on the energy) , not much is known otherwise for critical wave maps . An interesting modification to the wave map equations can be obtained by adding torsion. This can be done in 2+1 dimensions, without adding extra dynamical fields, as follows. One fixes a pair of background fields: a closed one-form field $`v`$ on the base manifold $`M^{2+1}`$ and a non-closed two-form field $`p`$ on the target $`N^n`$. The field $`p`$ serves as a “torsion potential” in the sense that the torsion tensor on $`N^n`$ is defined as $`Q^A{}_{BC}{}^{}=3/2g^{AD}_{[D}p_{BC]}.`$ (3) A map $`\psi :M^{2+1}N^n`$ is defined to be a torsion wave map if it is a critical point for the functional $`S_{\mathrm{tor}}[\psi ]={\displaystyle _{M^{2+1}}}\left(\eta ^{\mu \nu }g_{AB}(\psi )_\mu \psi ^A_\nu \psi ^B+\lambda ϵ^{\mu \nu \sigma }v_\sigma p_{AB}(\psi )_\mu \psi ^A_\nu \psi ^B\right)`$ (4) where $`\lambda `$ is a coupling constant and $`ϵ`$ is the 2+1 volume tensor normalized with respect to $`\eta `$. The torsion wave map equation obtained from (4) is given by $`\eta ^{\mu \nu }_\mu _\nu \psi ^A+\stackrel{~}{\mathrm{\Gamma }}^A{}_{BC}{}^{}(\psi )_\mu \psi ^B_\nu \psi ^C(\eta ^{\mu \nu }+\lambda ϵ^{\mu \nu \sigma }v_\sigma )=0`$ (5) where $`\stackrel{~}{\mathrm{\Gamma }}^A{}_{BC}{}^{}=\mathrm{\Gamma }^A{}_{BC}{}^{}+Q^A_{BC}`$ (6) are the connection coefficients <sup>2</sup><sup>2</sup>2The contorsion coefficients $`\stackrel{~}{\mathrm{\Gamma }}^A_{[BC]}`$ compatible with $`g`$ as determined by the torsion are identically equal to $`Q`$. compatible with $`g`$, with torsion $`Q`$. Note that the effect of the torsion is to add the nonlinear term $`\lambda ϵ^{\mu \nu \sigma }v_\sigma Q^A{}_{BC}{}^{}(\psi )_\mu \psi ^B_\nu \psi ^C`$ (7) to the wave map equation (2). Wave maps without torsion have a conserved, symmetric stress-energy tensor arising from the functional $`S[\psi ]`$. With the addition of torsion, the corresponding symmetric stress-energy tensor obtained from the functional $`S_{\mathrm{tor}}[\psi ]`$ is no longer conserved. However, we point out that a non-symmetric stress-energy tensor can be derived by considering the variation of $`S_{\mathrm{tor}}[\psi ]`$ under infinitesimal diffeomorphisms of $`M^{2+1}`$ acting on $`\eta `$, $`ϵ`$, $`v`$, and $`\psi `$. This leads to $`T^\mu {}_{\alpha }{}^{}=\eta ^{\mu \nu }g_{AB}_\nu \psi ^A_\alpha \psi ^B1/2\delta ^\mu {}_{\alpha }{}^{}\eta _{}^{\nu \sigma }g_{AB}_\nu \psi ^A_\sigma \psi ^B`$ $`+1/2\lambda ϵ^{\mu \nu \sigma }v_\alpha p_{AB}_\nu \psi ^A_\sigma \psi ^B`$ (8) which satisfies $`_\mu T^\mu {}_{\alpha }{}^{}=1/2\lambda ϵ^{\mu \nu \sigma }_\mu \psi ^A_\nu \psi ^Bp_{AB}_\alpha v_\sigma .`$ (9) Hence $`T^\mu _\alpha `$ is conserved if $`v`$ is covariantly constant on $`(M^{2+1},\eta )`$. Furthermore, $`T^\mu _\alpha `$ reduces to the standard symmetric stress-energy tensor for wave maps without torsion when $`v`$ is set to zero. The stress-energy tensor (8) is central to investigating global existence for torsion wave maps. The critical dimension for torsion wave maps, just as for standard wave maps, is 2+1. While we do not attempt here to investigate the general class of critical torsion wave maps, we are able to prove global existence for various reductions of critical wave maps, with and without torsion, where the base geometry is Minkowski space. These reductions are defined by the invariance or equivariance of the wave map $`\psi `$ under a one-dimensional group of translations acting on $`M^{2+1}`$. More specifically, choose Cartesian coordinates $`(x,y,t)`$ for $`(M^{2+1},\eta )`$ and denote the translation group action by $`(x,y,t)(x,y+\lambda ,t)`$. Then, for any target $`N^n`$, a wave map $`\psi `$ is translation invariant if $`\psi ^A(x,y+\lambda ,t)=\psi ^A(x,y,t).`$ (10) Translation equivariant wave maps require that the target $`N^n`$ admit a translation group action. Let $`\rho ^A{}_{B}{}^{}(\lambda )`$ denote a representation of the translation group action on the base $`M^{2+1}`$ acting on the target $`N^n`$. Then a wave map $`\psi `$ is translation equivariant if $`\psi ^A(x,y+\lambda ,t)=\rho ^A{}_{B}{}^{}(\lambda )\psi ^B(x,y,t).`$ (11) Note that translation equivariance (11) reduces to translation invariance (10) when (and only when) the representation $`\rho (\lambda )`$ is chosen to be trivial, $`\rho ^A{}_{B}{}^{}(\lambda )=\delta ^A_B`$. One class of targets for which there is a natural translation group action available are Lie groups, $`G`$. For a Lie group target $`N^n=G`$, left and right multiplication on $`G`$ by a one-parameter exponential subgroup $`\mathrm{exp}(\lambda A)`$ define translation group actions, where $`A`$ is any element in the Lie algebra of $`G`$. This leads to three types of equivariance as follows. Let $`\mathrm{\Psi }`$ devote a matrix representation of the wave map $`\psi :M^{2+1}G`$ and let $`L`$ and $`R`$ be matrix representations of elements of the Lie algebra of $`G`$. Then $`\psi `$ is said to be, respectively, left-translation equivariant if $`\mathrm{\Psi }(x,y+\lambda ,t)=\mathrm{exp}(\lambda L)\mathrm{\Psi }(x,y,t),`$ (12) or right-translation equivariant if $`\mathrm{\Psi }(x,y+\lambda ,t)=\mathrm{\Psi }(x,y,t)\mathrm{exp}(\lambda R),`$ (13) or conjugate-translation equivariant if $`\mathrm{\Psi }(x,y+\lambda ,t)=\mathrm{exp}(\lambda L)\mathrm{\Psi }(x,y,t)\mathrm{exp}(\lambda R).`$ (14) Corresponding to invariant wave maps (10) and equivariant wave maps (12), (13), (14), we have the following four classes of reductions: Invariant Wave maps (Any target) $`\psi =\varphi (x,t)`$ (15) Left-Equivariant Wave maps (Lie group target) $`\mathrm{\Psi }=\mathrm{exp}(yL)\mathrm{\Phi }_L(x,t)`$ (16) Right-Equivariant Wave maps (Lie group target) $`\mathrm{\Psi }=\mathrm{\Phi }_R(x,t)\mathrm{exp}(yR)`$ (17) Conjugate-Equivariant Wave maps (Lie group target) $`\mathrm{\Psi }=\mathrm{exp}(yL)\mathrm{\Phi }_C(x,t)\mathrm{exp}(yR)`$ (18) In each case the 2+1 wave map equation for $`\psi `$ yields a 1+1 reduced equation for $`\varphi ,\mathrm{\Phi }_L,\mathrm{\Phi }_R,\mathrm{\Phi }_C`$, respectively, provided that the target geometry is suitably invariant as discussed later. We establish global existence of solutions to the Cauchy problem for the class of translation-invariant wave maps with torsion in Section 2. While the proof for these wave maps is very similar to that for $`1+1`$ wave maps with no torsion, the torsion terms do introduce some subtleties into the analysis, which we highlight. In order to prove global existence of solutions to the Cauchy problem for the three classes of translation equivariant wave maps with torsion, we find it useful to work with a frame formulation for 2 + 1 wave maps . In Section 3 we introduce the frame formulation for general targets and then proceed to relate wave map equivariance for Lie group targets to frame invariance and equivariance. In particular, our global existence theorems for equivariant wave maps have a natural formulation and proof using frames. The proof for the left equivariant, right equivariant, and conjugate equivariant wave maps with torsion is fairly similar in each case. We focus on the left equivariant case (which corresponds to invariant frames) and carry out the global existence proof in detail, in Section 4. We then briefly note in Section 5 the differences entailed in proving global existence for the other two cases. We make a few concluding remarks in Section 6. ## 2 INVARIANT WAVE MAPS WITH TORSION The translation invariance condition (10) is characterized by the wave map functions (15) being independent of $`y`$. Under this reduction the torsion wave map equation (5) becomes $`\gamma ^{\alpha \beta }_\alpha _\beta \varphi ^A+\stackrel{~}{\mathrm{\Gamma }}_{BC}^A(\varphi )_\alpha \varphi ^B_\beta \varphi ^C(\gamma ^{\alpha \beta }+\lambda v_yϵ^{\alpha \beta })=0`$ (19) where $`\gamma ^{\alpha \beta }`$ is the 1 + 1 Minkowski metric ($`\alpha `$, $`\beta `$ run over $`x`$ and $`t`$) and $`ϵ^{\alpha \beta }`$ is the 1 + 1 Levi-Civita tensor. We hereafter take $`v_y`$ to be constant, but we make no further restrictions: The target $`(N^n,g)`$ can be any Riemannian geometry, and the torsion potential $`p`$ can be any non-closed two form on $`N^n`$. Interestingly, while the torsion term $`\lambda v_yQ^A{}_{BC}{}^{}(\varphi )_\alpha \varphi ^B_\beta \varphi ^Cϵ^{\alpha \beta }`$ (20) appears in a nontrivial way in the reduced wave map equation (19), and while the stress-energy tensor (8) generally contains a torsion term, for translation invariant wave maps the torsion drops out of many of the stress-energy tensor components. We have $`T_{tt}=T_{xx}`$ $`=`$ $`{\displaystyle \frac{1}{2}}(\|_t\varphi \|^2+\|_x\varphi \|^2),`$ (21) $`T_{xt}=T_{tx}`$ $`=`$ $`_t\varphi ^A_x\varphi ^Bg_{AB},`$ (22) all of which contain no torsion, along with $`T_{yy}`$ $`=`$ $`{\displaystyle \frac{1}{2}}(\|_t\varphi \|^2\|_x\varphi \|^2)+{\displaystyle \frac{1}{2}}\lambda v_yϵ^{\alpha \beta }_\alpha \varphi ^A_\beta \varphi ^Bp_{AB},`$ (23) $`T_{xy}`$ $`=`$ $`0,`$ (24) $`T_{yx}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\lambda v_xϵ^{\alpha \beta }_\alpha \varphi ^A_\beta \varphi ^Bp_{AB},`$ (25) $`T_{ty}`$ $`=`$ $`0,`$ (26) $`T_{yt}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\lambda v_tϵ^{\alpha \beta }_\alpha \varphi ^A_\beta \varphi ^Bp_{AB}.`$ (27) Note that $`v_x`$ and $`v_t`$ do not appear in the reduced wave map equation (19); setting them to zero does not affect (19), but it does result in $`T_{yx}`$ and $`T_{yt}`$ vanishing. We now consider the Cauchy problem for translation invariant wave maps (15) with torsion. Initial data at $`t=t_0`$ consists of a pair of maps $`\widehat{\varphi }:\mathrm{\Sigma }N,\widehat{\theta }:\mathrm{\Sigma }TN`$ (28) (here $`\mathrm{\Sigma }=R^1`$ or $`S^1`$ allowing for periodic boundary conditions). A solution to the Cauchy problem is then a map $`\varphi :\mathrm{\Sigma }\times R^1M^{2+1}N`$ which satisfies (19) along with the initial conditions $`\varphi (x,t_0)=\widehat{\varphi }(x),_t\varphi (x,t_0)=\widehat{\theta }(x).`$ (29) Note that there are no constraints on the choice of initial data $`\{\widehat{\varphi },\widehat{\theta }\}`$. Global existence of initial value solutions is established by the following theorem. Theorem 1. For any smooth compact support <sup>3</sup><sup>3</sup>3 $`\widehat{\varphi }`$ is compactly supported if it is constant everywhere outside a compact region in $`\mathrm{\Sigma }`$; $`\widehat{\theta }`$ is compactly supported if it zero outside such a region. initial data, the Cauchy problem (19) and (29) has a unique smooth global solution $`\varphi (x,t)`$ for all $`tR^1`$. Proof: The PDE system (19) is manifestly hyperbolic; hence, local existence and uniqueness are immediate . To prove global existence, it is sufficient (by the usual open-closed arguments ) to show that if $`\varphi (x,t)`$ satisfies (19) on $`\mathrm{\Sigma }\times I`$, with $`I`$ a bounded open interval in $`R^1`$, then $`\varphi (x,t)`$ and all its derivatives are bounded on $`\mathrm{\Sigma }\times I`$. To show that $`\varphi `$ and its first derivatives $`_\alpha \varphi `$ are bounded, we use an argument based on stress-energy conservation (see ). From the form of the stress-energy components (21) to (27), together with the conservation equations $`_tT^t{}_{t}{}^{}+_xT_t^x=0,_tT^t{}_{x}{}^{}+_xT_x^x=0,`$ (30) we find that $`\gamma ^{\alpha \beta }_\alpha _\beta T_{tt}=0.`$ (31) It then follows from standard results (see ) for the wave equation on 1 + 1 Minkowski space that $`T_{tt}`$ is bounded on $`I`$. Thus the first derivatives of $`\varphi `$ are bounded. As a consequence of the mean value theorem and the assumed compact support of the initial data, $`\varphi `$ is then bounded as well. There are a number of ways of proceeding to argue that second and higher order derivatives of $`\varphi `$ are bounded. Here we use an argument which is adapted from Shatah based on bounding successive $`k`$th order energies $`_k(t)={\displaystyle \frac{1}{2}}{\displaystyle _\mathrm{\Sigma }}\left(\|_t_x{}_{}{}^{k}\varphi \|^2+\|_x{}_{}{}^{k+1}\varphi \|^2\right)𝑑x.`$ (32) Note that the ordinary energy $`_0(t)={\displaystyle _\mathrm{\Sigma }}T_{tt}𝑑x={\displaystyle \frac{1}{2}}{\displaystyle _\mathrm{\Sigma }}\left(\|_t\varphi \|^2+\|_x\varphi \|^2\right)𝑑x`$ (33) is bounded and independent of $`t`$, $`_0(t)=_0(t_0)`$, for smooth compact support initial data. We start by rewriting the torsion wave map equation (19) in the form $`D^\alpha V_\alpha ^A=0`$ (34) where $`V_\alpha ^A=_\alpha \varphi ^A`$, $`D^\alpha =\gamma ^{\alpha \beta }D_\beta `$, and $`D_\beta =_\beta +\mathrm{\Gamma }^A{}_{BC}{}^{}V_{\beta }^{C}+\lambda Q^A{}_{BC}{}^{}ϵ_{\beta }^{}{}_{}{}^{\alpha }V_{\alpha }^{C}`$ defines a covariant derivative operator which includes the connection with torsion. If we now apply $`D_\beta `$ to equation (34) and commute $`D_\beta `$ past the derivative operators, keeping track of the various curvature and torsion terms which arise, then we obtain a nonlinear wave equation for $`V_\alpha ^A`$: $`D^\alpha D_\alpha V_\beta {}_{}{}^{A}+P_\beta {}_{}{}^{A}(V,V,V)=0`$ (35) where $`P(V,V,V)`$ denotes an expression which is trilinear in $`V_\beta ^A`$ and involves no higher derivatives of $`\varphi ^A`$. By multiplying (35) by $`\gamma ^{\alpha \beta }g_{CA}D_tV_\alpha ^C`$, we straightforwardly derive the conservation equation $`D_t\left({\displaystyle \frac{1}{2}}\|D_tV\|^2+{\displaystyle \frac{1}{2}}\|D_xV\|^2\right)D_x\left(D_tVD_xV\right)=\stackrel{~}{P}(V,V,V)DV`$ (36) where $`\stackrel{~}{P}(V,V,V)`$ is, like $`P(V,V,V)`$, trilinear in $`V`$ with no higher derivatives of $`\varphi ^A`$. Now, integrating (36) over $`\mathrm{\Sigma }`$, we obtain $`_t_1(t)={\displaystyle _\mathrm{\Sigma }}\stackrel{~}{P}(V,V,V)DV𝑑x`$ (37) for the $`1`$st order energy defined in (32). Estimating the right hand side of (37), we find $`_t_1(t)CV_{L^6}^3D_\alpha V_{L^2}`$ $`CV_{L^6}^3\sqrt{_1(t)}`$ (38) and hence $`_t\sqrt{_1}CV_{L^6}^3.`$ (39) It follows from Sobolev inequalities that $`V_{L^6}CV_{L^2}^{2/3}DV_{L^2}^{1/3},`$ (40) so we have $`_t\sqrt{_1}CV_{L^2}^2DV_{L^2}C_0\sqrt{_1}.`$ (41) Since $`_0`$ is bounded, it follows from (41) that $`\sqrt{_1(t)}Ce^{kt}`$ (42) which bounds $`_1(t)`$, and therefore bounds the $`L^2`$ norm of $`DV`$. Hence $`_\alpha ^2\varphi _{L^2}`$ is bounded. To bound $`_2(t)`$, we start from the wave equation (35) for $`V`$ and repeat the previous argument. Setting $`W_{\beta \gamma }{}_{}{}^{A}:=D_\beta V_\gamma ^A`$, we derive $`D^\alpha D_\alpha W_{\beta \gamma }{}_{}{}^{A}+R_{\beta \gamma }{}_{}{}^{A}(V,V,W)=0`$ (43) where $`R(V,V,W)`$ is bilinear in $`V`$, linear in $`W`$, and involves no other derivatives of $`\varphi `$. From (43) we obtain the conservation equation $`D_t\left({\displaystyle \frac{1}{2}}\|D_tW\|^2+{\displaystyle \frac{1}{2}}\|D_xW\|^2\right)D_x\left(D_tWD_xW\right)=\stackrel{~}{R}(V,V,W)DW`$ (44) where $`\stackrel{~}{R}(V,V,W)`$ has the same properties as $`R(V,V,W)`$. Integrating over $`\mathrm{\Sigma }`$ and estimating, we obtain $`_t_2(t)CV_{L^8}^2W_{L^4}\sqrt{_2(t)}.`$ (45) Since $`V`$ and $`W=DV`$ are $`L^2`$ bounded, by applying Sobolev inequalities to (45) we obtain $`_t\sqrt{_2}C\sqrt{_2}`$ (46) Hence we have that $`_2(t)`$ is bounded and therefore so is the $`L^2`$ norm of $`DW`$. Thus, since $`DW=DDV`$, it follows that $`_\alpha ^3\varphi _{L^2}`$ is bounded. The argument proceeds to all successively higher orders and we thereby determine that all derivatives of $`\varphi `$ are $`L^2`$ bounded. It follows from Sobolev embedding that all derivatives of $`\varphi `$ are pointwise bounded, which completes the proof of Theorem 1. ## 3 EQUIVARIANT WAVE MAPS WITH TORSION <br>AND THE FRAME FORMULATION We begin by setting up a frame formulation for wave maps with and without torsion. (See also ). We first choose a frame basis $`\{e_a^A\}(a=1,\mathrm{},n)`$ for the target geometry $`(N^n,g)`$ and let $`e_A^a(\psi )`$ denote the frame associated to $`\psi `$. We now define the “frame fields” $`K_\mu ^a:=e_A^a(\psi )_\mu \psi ^A`$ (47) where $`\{e_A^a\}`$ are the components of the dual basis to $`\{e_a^A\}`$. These frame fields $`K_\mu ^a`$ may be viewed either as the pull-back of the dual frame $`\{e_A^a\}`$ from $`N^n`$ to $`M^{2+1}`$ along the map $`\psi :M^{2+1}N^n`$, or as the frame components of the wave map gradient on the tangent space of the target geometry $`(N^n,g)`$. In any case, it follows from (47) that $`K`$ satisfies the identity $`_{[\nu }K_{\mu ]}^a=1/2C_{bc}{}_{}{}^{a}(\psi )K_\nu ^bK_\mu ^c`$ (48) where $`C_{bc}^a`$ are the frame commutator coefficients defined by $`[e_b,e_c]=C_{bc}{}_{}{}^{a}e_{a}^{}.`$ (49) Moreover, one verifies that if $`\psi `$ satisfies the wave map equation (2), then $`K`$ satisfies $`^\mu K_\mu ^a=C^a{}_{bc}{}^{}(\psi )K_\nu ^bK_\mu ^c\eta ^{\nu \mu }`$ (50) where $`C^a{}_{bc}{}^{}:=g^{ad}C_{db}{}_{}{}^{e}g_{ec}^{}`$, with $`g^{ad}:=e_A^ae_B^bg^{AB}`$ and $`g_{ab}:=e_a^Ae_b^Bg_{AB}`$; or if $`\psi `$ satisfies the torsion wave map equation (5), then $`K`$ satisfies $`^\mu K_\mu ^a=C^a{}_{bc}{}^{}(\psi )K_\nu ^bK_\mu ^c\eta ^{\nu \mu }\lambda ϵ^{\sigma \nu \mu }v_\sigma Q^a{}_{bc}{}^{}(\psi )K_\nu ^bK_\mu ^c`$ (51) where $`Q^a{}_{bc}{}^{}:=e_A^aQ^A{}_{BC}{}^{}e_{b}^{B}e_c^C`$. Up to this point in setting up the frame formulation, we have made no restrictions on the choice of the target or on the nature of the wave maps. We now focus on equivariant wave maps (12) to (14) and their corresponding frame formulations, so we assume the target geometry to be a Lie group $`G`$. While $`K`$ can be defined for any frame basis on $`G`$, the frame field equations are simplest if we require that $`\{e_a^A\}`$ be a left-invariant basis for $`G`$. It then follows that the commutator coefficients $`C_{bc}^a`$ are independent of $`\psi `$ and are constant. If we make the further restrictions that the metric $`g`$ be a left-invariant tensor on $`G`$, $`g_{AB}=e^a{}_{A}{}^{}e_{}^{b}{}_{B}{}^{}g_{ab}^{},`$ (52) and the torsion potential $`p`$ be a left invariant two-form on $`G`$, $`p_{AB}=e^a{}_{A}{}^{}e_{}^{b}{}_{B}{}^{}p_{ab}^{},`$ (53) so that the components $`g_{ab}`$ and $`p_{ab}`$ are constant, then the coefficients $`C^a{}_{bc}{}^{}(\psi )`$ are independent of $`\psi `$ and constant while so are the frame components $`Q^a{}_{bc}{}^{}(\psi )`$ as well; in particular, we have $`Q^a{}_{bc}{}^{}=3/2g^{ad}p_{e[d}C_{bc]}^e`$ (54) with $`C_{bc}{}_{}{}^{a}=2e_A^ae_{[b\|}^B_Be_{\|c]}^A.`$ (55) Remark 1: Every nonabelian Lie group admits both a left-invariant metric $`g`$ and a left-invariant two-form $`p`$. However, for semi-simple Lie groups $`G`$, if $`G`$ has dimension three then all left-invariant two-forms $`p`$ are necessarily closed, and consequently $`Q=0`$ so there is no torsion. This is not the case if $`G`$ has larger dimension. In particular, a non-closed left-invariant two-form $`p`$ and hence non-zero torsion $`Q`$ is admitted by all nonabelian semi-simple Lie groups $`G`$ other than the three-dimensional ones (namely $`SU(2)`$ and its real forms $`SO(3)`$, $`SO(1,2)`$, $`SO(2,1)`$). See Proposition A in the appendix. We now find that, assuming the restrictions just noted, we can write equations (48), (50) and (51) strictly in terms of the frame fields $`K`$, with no explicit $`\psi `$ dependence: $`_{[\nu }K^a{}_{\mu ]}{}^{}=1/2C_{bc}{}_{}{}^{a}K_{\nu }^{}{}_{}{}^{b}K_{\mu }^{c},`$ (56) $`^\mu K_\mu ^a=C^a{}_{bc}{}^{}K_{}^{b}{}_{\nu }{}^{}K_{\mu }^{c}\eta ^{\nu \mu },`$ (57) $`^\mu K_\mu {}_{}{}^{a}=C^a{}_{bc}{}^{}K_{}^{b}{}_{\nu }{}^{}K_{\mu }^{c}\eta ^{\nu \mu }\lambda ϵ^{\sigma \nu \mu }v_\sigma Q^a{}_{bc}{}^{}K_{}^{b}{}_{\nu }{}^{}K_{\mu }^{c}.`$ (58) The field equations (56) and(57) together are a self-contained PDE system for $`K`$ which is equivalent to the wave map equation (2); the field equations (56) and(58) likewise are a self-contained PDE system for $`K`$ which is equivalent to the wave map equation with torsion (5). Note that the system with torsion reduces to the system without torsion when $`\lambda =0`$. Proposition 1. Let $`(M^{2+1},\eta )`$ be a Lorentzian geometry, and let $`N^n=G`$ be a Lie group target. 1. Suppose that $`\psi ^A`$ is a solution of the torsion wave map equation (5). Then $`K_\mu ^a`$ defined by (47) satisfies the field equations (56) and (58). 2. Suppose that $`K_\mu ^a`$ is a solution of the field equations (56) and (58). If $`M^{2+1}`$ is simply connected, then there exists a torsion wave map $`\psi ^A`$, satisfying equation (5), which is related to $`K_\mu ^a`$ by (47). Proof: To prove part (1), we first note that for $`K^a_\mu `$ given by (47), the field equation (56) is an identity. We then verify that, through the torsion wave map equation (5), the substitution of (47) for $`K_\mu ^a`$ satisfies the field equation (58). For the converse, to prove part (2), we note the field equation (56) shows that $`K_\mu ^a`$ can be viewed as a Lie-algebra valued connection one-form on the trivial bundle $`M^{2+1}\times G`$, with zero curvature. Since the bundle is trivial and the manifold $`M^{2+1}`$ is assumed to be simply connected, there exists a global parallel section. Correspondingly, there exists a smooth map $`U:M^{2+1}G`$ (called a “gauge transformation”) in terms of which we have $`K_\mu ^a=(U^1_\mu U)^Ae_A^a(I)`$ (59) where $`I`$ is the identity element of $`G`$. Now let $`\psi ^A`$ denote $`U`$ written in terms of local coordinates on $`G`$. It follows that $`U^1`$ pulls back $`e_A^a(I)`$ to $`e_A^a(\psi )`$, at the Lie group element specified by $`U`$. Hence we have $`(U^1_\mu U)^Ae_A^a(I)=e_A^a(\psi )_\mu \psi ^A.`$ (60) Combining (60) with (59), we obtain equation (47). Then by substituting (47) into the field equation (58), we verify that $`\psi `$ satisfies (5). We note that, independent of their usefulness for the study of wave maps, these field theories in terms of $`K`$ viewed as a Lie-algebra valued one-form field on $`M^{2+1}`$ have some interest as a nonlinear generalization of Maxwell’s equations. Indeed, for the abelian case $`C_{ab}{}_{}{}^{c}=0`$, the field equations (56) and (57) are exactly Maxwell’s equations in 2+1 dimensions, while the field equations (56) and (58) are a modification of Maxwell’s equations by adding torsion. This relationship is explored elsewhere . Since we will use frame fields to study translation equivariant wave maps, we now characterize frame fields which correspond to the three classes of equivariant wave maps (12), (13), (14). We begin with the following definitions of invariant and equivariant frame fields under a translation group action. Invariant Frame Field: $`K(x,y+\lambda ,t)=K(x,y,t)`$ (61) Equivariant Frame Field: $`K(x,y+\lambda ,t)=\mathrm{exp}(\lambda A)K(x,y,t)\mathrm{exp}(\lambda A)`$ (62) Here $`A`$ is an element of the Lie algebra of the target Lie group $`G`$, and $`(x,y,t)`$ are standard coordinates for the Minkowski space base geometry $`(M^{2+1},\eta )`$. Geometrically, the translation equivariant group action (62) on $`K`$ arises via the pull-back of the dual frame components $`\{e_A^a\}`$ under right multiplication in $`G`$ by the one-parameter exponential subgroup generated from the Lie algebra element $`R`$. When $`R=0`$ this group action reduces to the translation invariant group action (61) on $`K`$. (Alternatively, note that the translation invariant group action arises directly by left multiplication in $`G`$ since the dual frame is left-invariant.) Based on Proposition 1, the correspondence between invariant/equivariant frame fields and wave maps is summarized by the following two results. Proposition 2. 1. If $`\psi `$ is left equivariant (12), then the corresponding frame field $`K`$ is invariant (61). 2. If $`\psi `$ is right equivariant (13), then the corresponding frame field $`K`$ is equivariant (62), with the components $`K_y^a`$ constant. 3. If $`\psi `$ is conjugate equivariant (14), then the corresponding frame field $`K`$ is equivariant (62). The proof of these correspondences amounts to a direct calculation using a matrix representation for $`\psi `$ and $`K`$. There are straightforward converse correspondences as well. Proposition 3. 1. If K is invariant (61), then the corresponding wave map $`\psi `$ is left equivariant (12). 2. If K is equivariant (62), then the corresponding wave map $`\psi `$ is conjugate equivariant (14). 3. If K is equivariant (62) with the components $`K_y^a`$ constant, then the corresponding wave map $`\psi `$ is right equivariant (13). Proof: Let $`U`$ denote a the matrix representation of the wave map $`\psi `$ corresponding to $`K`$. We first prove part (1). It follows from the definition of frame field invariance, together with relation (59), that $`U`$ satisfies $`_y(U^1_\mu U)=0.`$ (63) Integrating the $`y`$ component of this equation, and then multiplying both sides by $`U`$, we obtain the linear matrix ordinary differential equation $`_yU(x,y,t)=U(x,y,t)f(x,t)`$ (64) where $`f`$ is an arbitrary Lie-algebra matrix valued function (independent of $`y`$). The general solution to (64) is $`U(x,y,t)=\mathrm{exp}(yA(x,t))V(x,t)`$ (65) where $`V`$ is an arbitrary Lie-algebra nonsingular matrix valued function (independent of $`y`$), and $`A:=VfV^1`$. We now impose the $`x,t`$ components of equation (63). Calculating $`U^1_tU`$ with $`U`$ from (65), we find $`U^1_tU=V^1_tV+V^1\mathrm{exp}(yA)(_t\mathrm{exp}(yA))V,`$ (66) so $`0`$ $`=`$ $`_y(U^1_tU)=V^1\mathrm{exp}(yA)_tA\mathrm{exp}(yA)V,`$ (67) which implies that $`_tA=0.`$ (68) Similarly, working with $`U^1_xU`$ and imposing $`0=_y(U^1_xU)`$ we determine that $`_xA=0.`$ (69) Thus $`A`$ must be a constant Lie-algebra valued matrix, which we denote $`L`$; then (65) becomes $`U(x,y,t)=\mathrm{exp}(yL)V(x,t).`$ (70) Condition (12) immediately follows, so the wave map corresponding to $`K`$ is left equivariant. We now prove part (2). From the definition of frame equivariance, there is a $`y`$-independent Lie-algebra matrix valued field $`f_\mu (x,t)`$ and a constant Lie-algebra matrix $`A`$ such that $`K_\mu (x,y,t)=\mathrm{exp}(yA)f_\mu (x,t)\mathrm{exp}(yA).`$ (71) Hence, from relation (59), $`U(x,y,t)`$ must satisfy $`U^1_\mu U=\mathrm{exp}(yA)f_\mu \mathrm{exp}(yA).`$ (72) The $`y`$-component of this equation yields $`(_yU)\mathrm{exp}(yA)=U\mathrm{exp}(yA)f_y,`$ (73) and after some manipulation we obtain the linear matrix ODE $`_yW=W(f_yA)`$ (74) where $`W(x,y,t):=U(x,y,t)\mathrm{exp}(yA)`$. The general solution to (74) is $`W(x,y,t)=\mathrm{exp}(yB(x,t))V(x,t)`$ (75) where $`V`$ is an arbitrary Lie-algebra matrix valued function, and $`B`$ is defined as $`B:=V(f_yA)V^1.`$ (76) Working with the other components of equation (72) we derive $`_tW=Wf_t`$ (77) and $`_xW=Wf_x.`$ (78) Then rearranging (77) and using (75), we obtain $`f_t(x,t)`$ $`=`$ $`W^1_tW=V^1_tV+V^1\mathrm{exp}(yB)_t(\mathrm{exp}(yB))V.`$ (79) Since both $`f_t`$ and $`V`$ are independent of $`y`$, if we take $`_y`$ of both sides of equation (79) we have $`0`$ $`=`$ $`_y(V^1\mathrm{exp}(yB)_t(\mathrm{exp}(yB))V)`$ (80) $`=`$ $`V^1\mathrm{exp}(yB)_tB\mathrm{exp}(yB)V`$ which implies that $`_tB=0`$. Similarly, using (78), we find that $`_xB=0`$. Hence $`B`$ is a constant Lie-algebra matrix, which we denote $`L`$. Thus, after combining (75) with the definition $`W=U\mathrm{exp}(yA)`$, we see that $`U(x,y,t)=\mathrm{exp}(yL)V(x,t)\mathrm{exp}(yA)`$ (81) so $`U(x,y,t)`$ is conjugate equivariant (14) with $`R=A`$. Finally, we prove part (3). From (81) we have $`K_y=U^1_yU=\mathrm{exp}(yA)V^1LV\mathrm{exp}(yA)+A`$ (82) which is assumed to be constant. By differentiating with respect to $`y`$, we obtain $`[V^1LV,A]=0`$, and hence (81) becomes $`K_y=V^1LV+A`$. Thus, it follows that $`B:=V^1LV`$ defines a constant Lie-algebra matrix which commutes with $`A`$. We then have $`U(x,y,t)`$ $`=`$ $`\mathrm{exp}(yL)V(x,t)\mathrm{exp}(yA)`$ (83) $`=`$ $`V(x,t)\mathrm{exp}(yB)\mathrm{exp}(yA)=V(x,t)\mathrm{exp}(yR)`$ where $`R:=A+B`$. Hence $`U(x,y,t)`$ is right equivariant (13). As a consequence of Propositions 2 and 3, we can prove global existence of solutions to the Cauchy problem for the three classes of translation equivariant wave maps (with or without torsion) by using invariant or equivariant frame fields. We do this first for the invariant frame fields in the next section. Our analysis makes essential use of the wave map stress-energy tensor (8). Through the relation (47) for $`K`$ in terms of $`\psi `$, we obtain $`T^\mu {}_{\alpha }{}^{}=\eta ^{\mu \nu }K_\nu ^aK_\alpha ^bg_{ab}1/2\delta ^\mu {}_{\alpha }{}^{}\eta _{}^{\nu \sigma }K_\nu ^aK_\sigma ^bg_{ab}+1/2\lambda ϵ^{\mu \nu \sigma }v_\alpha p_{ab}K_\nu ^aK_\sigma ^b.`$ (84) One verifies that, for solutions $`K`$ of (56) and (58) in which $`(M^{2+1},\eta )`$ is Minkowski space, this non-symmetric stress-energy tensor satisfies the conservation equation $`_\mu T^\mu {}_{\alpha }{}^{}=1/2\lambda ϵ^{\mu \nu \sigma }p_{ab}K_\mu ^aK_\nu ^b_\alpha v_\sigma .`$ (85) Hereafter we specialize to the situation where $`v`$ is constant on $`M^{2+1}`$. This makes the analysis of the field equations considerably simpler. In particular, the stress-energy is strictly conserved, $`_\mu T^\mu {}_{\alpha }{}^{}=0`$. ## 4 GLOBAL EXISTENCE FOR INVARIANT FRAME FIELD EQUATIONS WITH TORSION By definition (61) of translation invariance for frame fields, the component functions $`K_\mu ^a`$ are independent of $`y`$. Then, adopting the convenient notation $`E^a:=K_x^a,H^a:=K_y^a,B^a:=K_t^a,`$ (86) we find that the translation-invariant frame field equations take the form $`_xH^a=`$ $`C_{bc}{}_{}{}^{a}E_{}^{b}H^c`$ (87) $`_tE^a=`$ $`_xBC_{bc}{}_{}{}^{a}B_{}^{b}E^c`$ (88) $`_tH^a=`$ $`C_{bc}{}_{}{}^{a}B_{}^{b}H^c`$ (89) $`_tB^a=`$ $`_xE^aC^a{}_{bc}{}^{}(B^bB^cE^bE^cH^bH^c)`$ (90) $`\lambda Q^a{}_{bc}{}^{}(v_yB^bE^cv_xB^bH^c+v_tE^bH^c)`$ for the functions $`\{E^a(x,t),H^a(x,t),B^a(x,t)\}`$. Note that, in this system of field equations, (87) is a constraint equation while (88) to (90) are evolution equations. Initial data at $`t=t_0`$ for the Cauchy problem is specified by choosing (on $`\mathrm{\Sigma }=R^1`$ or $`S^1`$ allowing for periodic boundary conditions) Lie-algebra valued functions $`\{\widehat{E}^a(x),\widehat{H}^a(x),\widehat{B}^a(x)\}`$ which satisfy the constraint $`_x\widehat{H}^a=C_{bc}{}_{}{}^{a}\widehat{E}_{}^{b}\widehat{H}^c.`$ (91) A solution to the Cauchy problem is then a set of fields $`\{E^a(x,t)`$, $`H^a(x,t)`$, $`B^a(x,t)\}`$ satisfying (88) to (90) and the initial conditions $`E^a(x,t_0)=\widehat{E}^a(x),H^a(x,t_0)=\widehat{H}^a(x),B^a(x,t_0)=\widehat{B}^a(x).`$ (92) To show that the Cauchy problem is well-posed, we note that well-posedness is known for the wave map equation without torsion , which is equivalent to the system (87) to (90) up to the addition of the torsion terms involving $`\lambda `$. These terms do not involve any derivatives of the fields and hence do not effect the well-posedness. Alternatively, we note that, up to such terms, the system is equivalent to the Maxwell equations in 2+1 dimensions, which constitute a well-posed system. It follows that the system (87) to (90) is well-posed and, moreover, is first-order hyperbolic. In this section we prove global existence of smooth solutions to the Cauchy problem for the 1+1 field equations (87) to (90). The proof relies on the use of the stress-energy tensor (84) along with light cone estimates. To proceed we write out the components of the stress-energy tensor (8) in terms of $`E^a,H^a,B^a`$. Using the coordinates $`(x,y,t)`$ for $`M^{2+1}`$ we have $`T_{xx}={\displaystyle \frac{1}{2}}\left(E_x^2E_y^2+B^2\right)+\lambda v_xH^aB^bp_{ab}`$ (93) $`T_{yy}={\displaystyle \frac{1}{2}}\left(E_x^2+E_y^2+B^2\right)+\lambda v_yB^aE^bp_{ab}`$ (94) $`T_{tx}=EB+\lambda v_xH^aE^bp_{ab}`$ (95) $`T_{xt}=EB+\lambda v_tH^aB^bp_{ab}`$ (96) $`T_{ty}=HB+\lambda v_yH^aE^bp_{ab}`$ (97) $`T_{yt}=HB+\lambda v_tB^aE^bp_{ab}`$ (98) $`T_{xy}=EH+\lambda v_yH^aB^bp_{ab}`$ (99) $`T_{yx}=EH+\lambda v_xB^aE^bp_{ab}`$ (100) where $`E^2:=E^aE^bg_{ab}`$ and $`EB=E^aB^bg_{ab}`$, etc.. For derivation of light cone estimates, it is useful to work with null components of the stress-energy tensor. We introduce null coordinates which mix $`t`$ and $`x`$ (but not $`y`$): $`\begin{array}{ccc}\mathrm{}=t+x& & x=\frac{1}{2}(\mathrm{}+n)\\ & & \\ n=t+x& & t=\frac{1}{2}(\mathrm{}n)\end{array}.`$ (104) Then we find (for the components we will need): $`T_{\mathrm{}\mathrm{}}`$ $`=`$ $`K_{\mathrm{}}^2+\lambda v_{\mathrm{}}H^aK_{\mathrm{}}^bp_{ab}`$ (105) $`=`$ $`{\displaystyle \frac{1}{4}}(B+E)^2+{\displaystyle \frac{\lambda }{4}}(v_t+v_x)H^a(B^b+E^b)p_{ab}`$ $`T_{nn}`$ $`=`$ $`K_n^2\lambda v_nH^aK_n^bp_{ab}`$ (106) $`=`$ $`{\displaystyle \frac{1}{4}}(B+E)^2{\displaystyle \frac{\lambda }{4}}(v_t+v_x)H^a(B^b+E^b)p_{ab}`$ $`T_\mathrm{}n`$ $`=`$ $`{\displaystyle \frac{1}{2}}K_y^2\lambda v_nK_{\mathrm{}}^aH^bp_{ab}`$ (107) $`=`$ $`{\displaystyle \frac{1}{2}}E_y^2+{\displaystyle \frac{\lambda }{4}}(v_t+v_x)H^a(B^b+E^b)p_{ab}`$ $`T_n\mathrm{}`$ $`=`$ $`{\displaystyle \frac{1}{2}}K_y^2+\lambda v_{\mathrm{}}K_n^aH^bp_{ab}`$ (108) $`=`$ $`{\displaystyle \frac{1}{2}}E_y^2{\displaystyle \frac{\lambda }{4}}(v_t+v_x)H^a(B^b+E^b)p_{ab}.`$ For these components the stress-energy conservation equation (85) has the null component form $`_nT_{\mathrm{}\mathrm{}}+_{\mathrm{}}T_n\mathrm{}=0,`$ (109) $`_{\mathrm{}}T_{nn}+_nT_\mathrm{}n=0.`$ (110) These equations are essential for the derivation of the light cone estimates we will need. Also important for our analysis is the energy function $`(t)`$ $`=`$ $`{\displaystyle _\mathrm{\Sigma }}T_{tt}𝑑x`$ (111) $`=`$ $`{\displaystyle _\mathrm{\Sigma }}({\displaystyle \frac{1}{2}}(E{}_{}{}^{2}+H{}_{}{}^{2}+B^2)+\lambda v_tH^aE^bp_{ab})dx`$ We note that for certain values of the coupling constant $`\lambda `$, the energy $`(t)`$ can be negative, and it therefore does not in general control the $`L^2`$ norm of $`E^a`$, $`H^a`$, or $`B^a`$. However, for sufficiently small $`\lambda `$, there is a constant $`k>0`$ such that $`{\displaystyle \frac{1}{k}}(E_x^2+E_y^2)E_x^2+E_y^2+2\lambda v_tH^aE^bp_{ab}k(E_x^2+E_y^2)`$ (112) and hence the energy is positive, so that $`(t)`$ does consequently control $`E_{L^2}`$,$`H_{L^2}`$ , and $`B_{L^2}`$. We assume henceforth that $`\lambda `$ is sufficiently small for this to be the case. <sup>4</sup><sup>4</sup>4It is sufficient that $`\lambda `$ satisfy $`\|\lambda \|1/\sqrt{\|v_t\|\|p\|}`$ where $`\|p\|^2=\|p_{ab}p_{cd}g^{ac}g^{bd}\|`$. We now state our main results. Let $`\mathrm{\Sigma }`$ denote $`R^1`$ or $`S^1`$, and introduce coordinates $`(x,t)`$ for $`\mathrm{\Sigma }\times R^1M^{2+1}`$. Fix constants $`v_t,v_x,v_y`$. Let $`G`$ be a Lie group with $`C_{bc}^a`$ denoting the Lie-algebra commutator structure tensor. Fix on the Lie algebra of $`G`$ a positive definite metric tensor $`g_{ab}`$ (it need not necessarily be compatible with the commutator) and a skew-tensor $`p_{ab}`$. Let $`Q^a_{bc}`$ be the tensor defined by (54). Theorem 2. Let $`\lambda `$ be a small constant.<sup>4</sup> For any smooth compact support initial data (92) satisfying (91), the Cauchy problem (87) to (90) has a unique smooth global solution $`\{E^a(x,t),H^a(x,t),B^a(x,t)\}`$ for all $`tR^1`$ Combining this result with Propositions 2 and 3 from Section 3, we have a corresponding result for wave maps. Theorem 3. The Cauchy problem for left-translation equivariant Lie group wave maps (12), with or without torsion, has a unique smooth global solution for all smooth compact support initial data. Proof of Theorem 2: Local existence and uniqueness of smooth solutions of the PDE system (87) to (90) follows from standard results (see, for example,) for first-order hyperbolic systems in 1+1 dimensions. In order to prove global existence, it is sufficient by the usual “open-closed” arguments to establish the following: For $`\{E^a(x,t)`$, $`H^a(x,t)`$, $`B^a(x,t)\}`$ satisfying equations (87) to (90) for $`tI`$, with $`I`$ a bounded open interval in $`R^1`$, each component of these fields is bounded for $`tI`$, as are all orders of their derivatives. We prove this boundedness result as follows: ### Step 1: Conserved Energy It follows from the stress-energy conservation equation $`_tT_{tt}_xT_{xt}=0`$ that the energy $`(t)`$ satisfies $`\frac{d}{dt}(t)=(t)`$ where $`(t):={\displaystyle _\mathrm{\Sigma }}_xT_{xt}dx=T_{xt}\|_\mathrm{\Sigma }`$ (113) is the flux. If we are working on $`\mathrm{\Sigma }=S^1`$, then $`\mathrm{\Sigma }`$ is empty, so $`(t)=0`$. If instead $`\mathrm{\Sigma }=R^1`$, then we note that as a consequence of hyperbolicity of the system (87) to (90), the fields $`\{E^a,H^a,B^a\}`$ have compact support on $`\mathrm{\Sigma }`$ for all $`tI`$, and hence $`(t)=0`$. Thus, the energy is conserved, $`(t)=(t_0)`$, for all $`tI`$. As we noted earlier, the energy controls the $`L^2`$ norm of the fields $`\{E^a`$, $`H^a`$, $`B^a\}`$, so long as $`\lambda `$ is sufficiently small (as assumed in the theorem). Hence we have $`E^a_{L^2(\mathrm{\Sigma })}<k,H^a_{L^2(\mathrm{\Sigma })}<k,B^a_{L^2(\mathrm{\Sigma })}<k`$ (114) for some constant $`k`$ (depending on $`(t_0)`$), for all $`tI`$. ### Step 2: Bounded $`H^a`$ In the system (87) to (90), the field $`H^a`$ enters in a different way from $`E^a`$ and $`B^a`$, since the evolution equation (89) for $`H^a`$ involves no spatial derivative terms, and the constraint equation (87) has $`_xH^a`$ appearing, but no spatial derivatives of $`E^a`$ or $`B^a`$. Consequently, we treat $`H^a`$ differently from the other two fields: we first show that $`H^a`$ is bounded, and then use this in showing that $`E^a`$ and $`B^a`$ are bounded. To start, we integrate the absolute values of both sides of the constraint equation (87) over $`\mathrm{\Sigma }`$, obtaining $`{\displaystyle _\mathrm{\Sigma }}\|_xH^a\|𝑑x={\displaystyle _\mathrm{\Sigma }}\|C_{bc}{}_{}{}^{a}E_{}^{b}H^c\|𝑑x.`$ (115) Since $`C_{bc}^a`$ is constant, there exists a constant $`k_1`$ such that $`\|C_{bc}{}_{}{}^{a}E_{}^{b}H^c\|\|C_{bc}{}_{}{}^{a}\|\|E^b\|\|H^c\|k_1(E^2+H^2)`$ (116) by standard algebraic inequalities. It follows from (115) and (116) together with the bounds (114) that $`{\displaystyle _\mathrm{\Sigma }}\|_xH^a\|𝑑xk_2`$ (117) for a constant $`k_2`$. Combining (117) with the mean value theorem, we obtain controls on the spatial variation of $`E_y(x,t)`$ for any fixed time $`t`$. In particular, for any $`x_1,x_2\mathrm{\Sigma }`$ with fixed $`t`$, we have $`\|H^a(x_2,t)H^a(x_1,t)\|`$ $`=`$ $`\|{\displaystyle _{x_1}^{x_2}}_xH^a(x,t)dx\|`$ (118) $``$ $`{\displaystyle _{x_1}^{x_2}}\|_xH^a(x,t)\|𝑑xk_2.`$ If we are working on $`\mathrm{\Sigma }=R^1`$, we can choose $`x_1`$ outside the support of $`H^a(x,t)`$ for all $`tI`$, and therefore it follows from (118) that $`\|H^a(x,t)\|k_2`$ for all $`(x,t)\mathrm{\Sigma }\times I`$. Hence, $`H^a(x,t)`$ is bounded on $`\mathrm{\Sigma }\times I`$. If instead we are working on $`\mathrm{\Sigma }=S^1`$, we need to do more to bound $`H^a(x,t)`$. Consider $`_{S^1}H^a(x,t)𝑑x`$, which is the spatial average of $`H^a`$ on $`S^1`$. From the fundamental theorem of calculus, and from the evolution equation (89), we obtain (for $`tI`$) $`{\displaystyle _{S^1}}H^a(x,t)𝑑x`$ $`=`$ $`{\displaystyle _{t_0}^t}{\displaystyle \frac{d}{ds}}{\displaystyle _{S^1}}H^a(x,s)𝑑x𝑑s+{\displaystyle _{S^1}}H^a(x,t_0)𝑑x`$ $`=`$ $`{\displaystyle _{t_0}^t}{\displaystyle _{S^1}}C_{bc}{}_{}{}^{a}B_{}^{b}(x,s)H^c(x,s)𝑑x𝑑s+{\displaystyle _{S^1}}H^a(x,t_0)𝑑x.`$ Next, using standard quadratic algebraic inequalities, we note that $`_{S^1}C_{bc}{}_{}{}^{a}B_{}^{b}H^c𝑑x`$ is bounded in terms of the energy, $`\|{\displaystyle _{S^1}}C_{bc}^aB^bH^c𝑑x\|k_3(t)=k_3(t_0)`$ (120) for some constant $`k_3`$. Hence, $`_{t_0}^t_{S^1}C_{bc}{}_{}{}^{a}B_{}^{b}(x,s)H^c(x,s)𝑑x𝑑s`$ is bounded above and below, $`\|{\displaystyle _{t_0}^t}{\displaystyle _{S^1}}C_{bc}{}_{}{}^{a}B_{}^{b}(x,s)H^c(x,s)𝑑x𝑑s\|(tt_0)k_3(t_0)k_4`$ (121) for some constant $`k_4`$, for all $`tI`$. Then since $`_{S^1}H^a(x,t_0)𝑑x`$ involves initial data only, it also is bounded above and below. Therefore, from (LABEL:LoneSEy) we have that $`\|{\displaystyle _{S^1}}H^a(x,t)𝑑x\|k_5`$ (122) and so the average of $`H^a`$ over $`S^1`$ is bounded above and below, for all $`tI`$. Combining this result with the spatial variance control (118), we conclude that $`H^a(x,t)`$ is bounded (above and below) on $`\mathrm{\Sigma }\times I`$. ### Step 3: Bounded $`E^a`$ and $`B^a`$ While standard $`1+1`$ light cone arguments do not directly apply to the system (87) to (90), a modified argument can be used with the pointwise bounds on $`H^a`$ achieved in Step 2. Using the null form of the stress-energy conservation laws (109)-(110), along with the expressions (105)-(108) for the stress-energy components, we have $`_n(B+E)^2`$ $`=`$ $`2_{\mathrm{}}H^2\lambda (v_t+v_x)p_{ab}_n\left(H^a(B^b+E^b)\right)`$ (123) $`+\lambda (v_t+v_x)p_{ab}_{\mathrm{}}\left(H^a(B^b+E^b)\right),`$ $`_{\mathrm{}}(B+E)^2`$ $`=`$ $`2_nH^2\lambda (v_t+v_x)p_{ab}_n\left(H^a(B^b+E^b)\right)`$ (124) $`+\lambda (v_t+v_x)p_{ab}_{\mathrm{}}\left(H^a(B^b+E^b)\right).`$ We use the field equations (87) to (90) to remove all of the derivatives which appear on the right-side of these equations. Thus $`_n(B+E)^2`$ $`=`$ $`2H^bH^c(B^a+E^a)C_{abc}`$ (125) $`2\lambda (v_t+v_x)H^a(B^b+E^b)(B^c+E^c)Q_{bca}`$ and $`_{\mathrm{}}(B+E)^2`$ $`=`$ $`2H^bH^c(B^a+E^a)C_{abc}`$ (126) $`2\lambda (v_t+v_x)H^a(B^b+E^b)(B^c+E^c)Q_{bca}.`$ It is convenient here to let $`\alpha ^a:=B^a+E^a`$ and $`\beta ^a:=B^a+E^a`$, and so we have $`_n\alpha ^2=2C_{abc}H^bH^c\alpha ^a2\lambda (v_t+v_x)Q_{bca}H^a\beta ^b\alpha ^c`$ (127) and $`_{\mathrm{}}\beta ^2=2C_{abc}H^bH^c\beta ^a2\lambda (v_t+v_x)Q_{bca}H^a\beta ^b\alpha ^c.`$ (128) Since $`C_{abc}`$, $`Q_{bca}`$, $`\lambda `$, $`v_t`$ and $`v_x`$ are constant, and since $`H^a`$ is bounded on $`\mathrm{\Sigma }\times I`$, we immediately have the following estimates for the right-sides of (127) and (128): $`_n\alpha ^2k_6\sqrt{\alpha ^2}+k_7\sqrt{\alpha ^2}\sqrt{\beta ^2}`$ (129) and $`_{\mathrm{}}\beta ^2k_8\sqrt{\beta ^2}+k_9\sqrt{\alpha ^2}\sqrt{\beta ^2}`$ (130) with some constants $`k_6`$, $`k_7`$, $`k_8`$, and $`k_9`$. We now apply a light cone argument to the differential inequalities (129) and (130). First, choose an arbitrary point $`(\widehat{x},\widehat{t})`$ in $`\mathrm{\Sigma }\times I`$ to the future of the initial surface $`\mathrm{\Sigma }`$, so $`\widehat{t}>t_0`$, and integrate $`\alpha ^2`$ back along the light ray parallel to $`_n`$ via (129) and also integrate $`\beta ^2`$ back along the light ray parallel to $`_{\mathrm{}}`$ via (130). This yields $`\alpha ^2(\widehat{x},\widehat{t})\alpha ^2(\widehat{x}+\widehat{t}t_0,t_0)+k_6{\displaystyle _{t_0}^{\widehat{t}}}\sqrt{\alpha ^2(\widehat{x}+\widehat{t}s,s)}𝑑s`$ $`+k_7{\displaystyle _{t_0}^{\widehat{t}}}\sqrt{\alpha ^2(\widehat{x}+\widehat{t}s,s)}\sqrt{\beta ^2(\widehat{x}+\widehat{t}s,s)}𝑑s`$ (131) and $`\beta ^2(\widehat{x},\widehat{t})\beta ^2(\widehat{x}+t_0\widehat{t},t_0)+k_8{\displaystyle _{t_0}^{\widehat{t}}}\sqrt{\beta ^2(\widehat{x}\widehat{t}+s,s)}𝑑s`$ $`+k_9{\displaystyle _{t_0}^{\widehat{t}}}\sqrt{\alpha ^2(\widehat{x}\widehat{t}+s,s)}\sqrt{\beta ^2(\widehat{x}\widehat{t}+s,s)}𝑑s.`$ (132) Next, take the supremum of these expressions over $`\mathrm{\Sigma }`$. Letting $`\widehat{\alpha }^2(t):=sup_{x\mathrm{\Sigma }}\alpha ^2(x,t)`$ and $`\widehat{\beta }^2(t):=sup_{x\mathrm{\Sigma }}\beta ^2(x,t)`$, we obtain from (131) $`\widehat{\alpha }^2(\widehat{t})`$ $``$ $`\widehat{\alpha }^2(t_0)+k_6\underset{x\mathrm{\Sigma }}{sup}{\displaystyle _{t_0}^{\widehat{t}}}\sqrt{\alpha ^2(x,s)}𝑑s`$ (133) $`+k_7\underset{x\mathrm{\Sigma }}{sup}{\displaystyle _{t_0}^{\widehat{t}}}\sqrt{\alpha ^2(x,s)}\sqrt{\beta ^2(x,s)}𝑑s`$ $``$ $`\widehat{\alpha }^2(t_0)+k_6{\displaystyle _{t_0}^{\widehat{t}}}\sqrt{\widehat{\alpha }^2(s)}𝑑s+k_7{\displaystyle _{t_0}^{\widehat{t}}}\sqrt{\widehat{\alpha }^2(s)}\sqrt{\widehat{\beta }^2(s)}𝑑s`$ $``$ $`\widehat{\alpha }^2(t_0)+k_{10}(\widehat{t}t_0)^{1/2}\left({\displaystyle _{t_0}^{\widehat{t}}}\widehat{\alpha }^2(s)𝑑s\right)^{1/2}`$ $`+k_{11}\left({\displaystyle _{t_0}^{\widehat{t}}}\widehat{\alpha }^2(s)𝑑s\right)^{1/2}\left({\displaystyle _{t_0}^{\widehat{t}}}\widehat{\beta }^2(s)𝑑s\right)^{1/2}`$ where the last step is a consequence of the Holder inequality. If we define $`a(t):={\displaystyle _{t_0}^t}\widehat{\alpha }^2(s)𝑑s`$ (134) and $`b(t):={\displaystyle _{t_0}^t}\widehat{\beta }^2(s)𝑑s`$ (135) then (133) can be written as (with $`\widehat{t}`$ replaced by $`t`$) $`{\displaystyle \frac{d}{dt}}a(t)a(t_0)+k_{10}(tt_0)^{1/2}a^{1/2}(t)+k_{11}a^{1/2}(t)b^{1/2}(t).`$ (136) Similarly, from (132), we derive $`{\displaystyle \frac{d}{dt}}b(t)b(t_0)+k_{12}(tt_0)^{1/2}b^{1/2}(t)+k_{13}a^{1/2}(t)b^{1/2}(t).`$ (137) We want to show $`a(t)`$ and $`b(t)`$ are bounded functions of $`t`$ by applying a Gronwall type argument to the coupled inequalities (136),(137). It is useful first to divide by $`a^{1/2}(t)`$ in (136) and by $`b^{1/2}(t)`$ in (137), yielding $`{\displaystyle \frac{d}{dt}}a^{1/2}(t)a(t_0)a^{1/2}(t)+k_{10}(tt_0)^{1/2}+k_{11}b^{1/2}(t),`$ (138) $`{\displaystyle \frac{d}{dt}}b^{1/2}(t)b(t_0)b^{1/2}(t)+k_{12}(tt_0)^{1/2}+k_{13}a^{1/2}(t).`$ (139) We estimate the term $`a(t_0)a^{1/2}(t)`$ by using the fact that $`a(t)`$ is a monotonic increasing function of $`t`$, due to positivity of $`\widehat{\alpha }^2`$ in (134). Thus, $`a(t_0)a^{1/2}(t)`$ is bounded by $`a^{1/2}(t_0)`$. In addition, we note the term $`k_{10}(tt_0)^{1/2}`$ is bounded since $`tI`$ is bounded. We thereby obtain $`{\displaystyle \frac{d}{dt}}a^{1/2}(t)k_{14}+k_{11}b^{1/2}(t).`$ (140) Similarly, we obtain $`{\displaystyle \frac{d}{dt}}b^{1/2}(t)k_{15}+k_{13}a^{1/2}(t).`$ (141) Adding (140) and (141), and defining $`c(t):=a^{1/2}(t)+b^{1/2}(t)`$, we derive $`{\displaystyle \frac{d}{dt}}c(t)k_{18}+k_{17}c(t).`$ (142) Gronwall’s inequality immediately applies to (142), and so we determine that $`c(t)`$ is bounded for all $`tI`$. Then $`a^{1/2}(t)`$ and $`b^{1/2}(t)`$, which are positive, are bounded. Returning to the inequalities (136)-(137), it follows that $`\frac{d}{dt}a(t)`$ and $`\frac{d}{dt}b(t)`$ are each bounded. Hence, from the definitions of $`a`$ and $`b`$, we obtain that $`sup_\mathrm{\Sigma }\alpha ^2`$ and $`sup_\mathrm{\Sigma }\beta ^2`$ are bounded for all $`tI`$. Since $`\alpha ^2=(B^a+E^a)^2`$ and $`\beta ^2=(B^a+E^a)^2`$, we conclude that $`B^a(x,t)`$ and $`E^a(x,t)`$ are bounded on $`\mathrm{\Sigma }\times I`$. ### Step 4: Bounded Derivatives Now that we have determined that $`E^a`$, $`H^a`$, and $`B^a`$ are bounded on $`\mathrm{\Sigma }\times I`$, we proceed to show that the first derivatives of these functions, and subsequently all higher order derivatives, are bounded on $`\mathrm{\Sigma }\times I`$. We start with $`H^a`$. From (87) and (89), it follows that since $`E^a`$, $`H^a`$, and $`B^a`$ are bounded, then $`_xH^a`$ and $`_tH^a`$ are bounded. Similarly, if the order $`n`$ derivatives of $`E^a`$, $`H^a`$, and $`B^a`$ are bounded, then it follows from (derivatives of) (87) and (89) that the order $`n+1`$ derivatives of $`H^a`$ are bounded. Hence, (by induction on $`n`$), the derivatives of $`H^a`$ to all orders are bounded. For $`E^a`$ and $`B^a`$, we use light cone arguments much like step 3, but involving a “derivative stress-energy” tensor. Specifically, let $`T_{(1)\mathrm{}\mathrm{}}={\displaystyle \frac{1}{4}}(_xB+_xE)^2+{\displaystyle \frac{\lambda }{4}}(v_t+v_n)_xH^a(_xB^b+_xE^b)p_{ab}`$ (143) $`T_{(1)nn}={\displaystyle \frac{1}{4}}(_xB+_xE)^2{\displaystyle \frac{\lambda }{4}}(v_t+v_n)_xH^a(_xB^b+_xE^b)p_{ab}`$ (144) $`T_{(1)\mathrm{}n}={\displaystyle \frac{1}{2}}(_xH)^2+{\displaystyle \frac{\lambda }{4}}(v_t+v_n)_xH^a(_xB^b+_xE^b)p_{ab}`$ (145) $`T_{(1)n\mathrm{}}={\displaystyle \frac{1}{2}}(_xH)^2{\displaystyle \frac{\lambda }{4}}(v_t+v_n)_xH^a(_xB^b+_xE^b)p_{ab}`$ (146) as defined analogously to the stress-energy components (105) to (108). Then, as a consequence of the field equations (87) to (90), we find $`_nT_{(1)\mathrm{}\mathrm{}}+_{\mathrm{}}T_{(1)n\mathrm{}}=Y_1(E,H,B;_xE,_xB)`$ (147) and $`_{\mathrm{}}T_{(1)nn}+_nT_{(1)\mathrm{}n}=Y_2(E,H,B;_xE,_xB)`$ (148) where $`Y_1`$ and $`Y_2`$ are homogeneous quadratic in $`_xE`$ and $`_xB`$, with bounded coefficients. Although (147) and (148) are not strict conservation equations, we can nevertheless proceed similarly to step 3. Through use of the field equations, we can express (147) and (148) as $`_n(_xB+_xE)^2=\stackrel{~}{Y}_1(E,H,B;_xE,_xB)`$ (149) and $`_{\mathrm{}}(_xB+_xE)^2=\stackrel{~}{Y}_2(E,H,B;_xE,_xB)`$ (150) with $`\stackrel{~}{Y}_1`$ and $`\stackrel{~}{Y}_2`$ of the same nature as $`Y_1`$ and $`Y_2`$. It then follows from standard algebraic inequalities that $`_n(_xB+_xE)^2k_{19}\|_xB\|^2+k_{20}\|_xE\|^2,`$ (151) $`_{\mathrm{}}(_xB+_xE)^2k_{21}\|_xB\|^2+k_{22}\|_xE\|^2.`$ (152) We now apply the light cone arguments of step 3 to the differential inequalities (151) and (152): Starting at an arbitrary point in $`\mathrm{\Sigma }\times I`$, we integrate (151) and (152) back to the initial surface along light rays generated by $`_n`$ and $`_{\mathrm{}}`$. Taking suprema over $`\mathrm{\Sigma }`$ and adding the resulting inequalities, we obtain $`\underset{x\mathrm{\Sigma }}{sup}[(_xB(x,t))^2+(_xE(x,t))^2]`$ $`k_{23}+k_{24}{\displaystyle _{t_0}^t}\underset{x\mathrm{\Sigma }}{sup}[(_xB(x,s))^2+(_xE(x,s))^2]ds.`$ (153) Applying the Gronwall inequality to (153) shows that $`\|_xB\|`$ and $`\|_xE\|`$ are bounded for all $`tI`$. With $`E^a`$, $`H^a`$, $`B^a`$, $`_xE^a`$ and $`_xB^a`$ bounded, it follows from the evolution equations (88) and (90) that $`_tE^a`$ and $`_tB^a`$ are bounded as well. The previous argument can be applied to all orders of derivatives of the fields. This establishes that $`E^a`$, $`H^a`$, $`B^a`$, and all of their derivatives are bounded on $`\mathrm{\Sigma }\times I`$. Global existence now follows from the usual “open-closed” arguments. ## 5 GLOBAL EXISTENCE FOR EQUIVARIANT <br>FRAME FIELD EQUATIONS WITH TORSION As discussed in Section 3, while left-equivariant wave maps (12) correspond to invariant frame fields (61), conjugate-equivariant wave maps (14) and right-equivariant wave maps (13) correspond to equivariant frame fields (62). In this section we show that global existence holds for smooth solutions to the Cauchy problem for translation equivariant frame fields. We first note that by definition of translation equivariance, $`K_\mu ^a(x,y,t)`$ can be expressed as $`K_x^a(x,y,t)=\mathrm{exp}(yR)E^a(x,t)\mathrm{exp}(yR)`$ (154) $`K_y^a(x,y,t)=\mathrm{exp}(yR)H^a(x,t)\mathrm{exp}(yR)`$ (155) $`K_t^a(x,y,t)=\mathrm{exp}(yR)B^a(x,t)\mathrm{exp}(yR)`$ (156) in terms of some Lie-algebra valued fields $`\{E^a(x,t),H^a(x,t),B^a(x,t)\}`$ which do not depend on $`y`$. Here $`R`$ is a fixed (constant) element in the Lie algebra; the left multiplication by $`\mathrm{exp}(yR)`$ combined with right multiplication by $`\mathrm{exp}(yR)`$ denotes the adjoint action of a one-parameter Lie subgroup on the Lie algebra. Substituting expressions (154) to (156) into the frame field equations with torsion (56) and (58) on Minkowski space, we obtain the following 1+1 reduced PDE system $`_xH^a=`$ $`C_{bc}{}_{}{}^{a}E_{}^{b}(H^cR^c)`$ (157) $`_tE^a=`$ $`_xB^aC_{bc}{}_{}{}^{a}B_{}^{b}E^c`$ (158) $`_tH^a=`$ $`C_{bc}{}_{}{}^{a}B_{}^{b}(H^cR^c)`$ (159) $`_tB^a=`$ $`_xE^aC_{bc}{}_{}{}^{a}H_{}^{b}R^cC^a{}_{bc}{}^{}(B^bB^cE^bE^cH^bH^c)`$ (160) $`\lambda Q^a{}_{bc}{}^{}(v_yB^bE^cv_xB^bH^c+v_tE^bH^c)`$ provided that $`C^a_{bc}`$ and $`Q^a_{bc}`$ are invariant under the adjoint action of $`\mathrm{exp}(yR)`$. We note that the only difference between these equations for translation equivariant frame fields and equations (87) to (90) for translation invariant frame fields is the presence of the commutator terms involving $`R`$. While the expressions for the field equations are changed somewhat in passing from invariant to equivariant frame fields, the expressions for the stress-energy components (93) to (100) and (105) to (108) do not change at all. (In particular, while $`K_\mu ^a`$ is not independent of $`y`$, the quadratic expressions $`K_\mu ^aK_\nu ^bg_{ab}`$ and $`K_\mu ^aK_\nu ^bp_{ab}`$ are invariant, and consequently so is $`T_{\mu \nu }`$. ) Initial data for the Cauchy problem for translation equivariant frame fields consists of Lie-algebra valued functions $`\{\widehat{E}^a(x),\widehat{H}^a(x),\widehat{B}^a(x)\}`$ on $`\mathrm{\Sigma }`$ which satisfy the constraint $`_x\widehat{H}^a=C_{bc}{}_{}{}^{a}\widehat{E}_{}^{b}(\widehat{H}^cR^c).`$ (161) A solution to the Cauchy problem is a set of fields $`\{E^a(x,t),H^a(x,t),B^a(x,t)\}`$ satisfying equations (157) to (160) and the initial conditions (92). The global existence result, and its corollary, are stated as follows. Let $`\mathrm{\Sigma }`$ denote $`R^1`$ or $`S^1`$, and introduce coordinates $`(x,t)`$ for $`\mathrm{\Sigma }\times R^1`$. Fix constants $`v_t,v_x,v_y`$. Let $`G`$ be a Lie group and let $`R^a`$ be a fixed (constant) vector in the Lie algebra of $`G`$. Assume $`G`$ admits on its Lie algebra a positive definite metric tensor $`g_{ab}`$ and a skew tensor $`p_{ab}`$ which are each invariant under the adjoint action of the Lie subgroup generated by $`R^a`$: $`g_{ae}C_{bc}{}_{}{}^{e}R_{}^{c}=g_{be}C_{ac}{}_{}{}^{e}R_{}^{c}`$ (162) $`p_{ae}C_{bc}{}_{}{}^{e}R_{}^{c}=p_{be}C_{ac}{}_{}{}^{e}R_{}^{c}`$ (163) where $`C_{bc}^a`$ denotes the Lie-algebra commutator structure tensor. Let $`Q^a_{bc}`$ be the tensor defined by (54). Theorem 4. Let $`\lambda `$ be a small constant.<sup>4</sup> For any smooth compact support initial data (92) satisfying (161), the Cauchy problem (157) to (160) has a unique smooth global solution $`\{E^a(x,t),H^a(x,t),B^a(x,t)\}`$ for all $`tR^1`$. From Propositions 2 and 3 we obtain a corresponding result for wave maps. Theorem 5. The Cauchy problems for conjugate-translation equivariant wave maps (14) and for right-translation equivariant wave maps (13), with or without torsion, have unique smooth global solutions for all smooth compact support initial data. Remark 2: Under the translation invariant form (61) for frame fields, which corresponds to left-translation equivariant (12) or translation invariant (10) wave maps, the reduction of the frame field equations and corresponding wave map equation is consistent for any Lie group target with $`(G,g,p)`$ invariant under left multiplication. However, this is not the case under the translation equivariant form (62) for frame fields, which corresponds to conjugate-translation equivariant (14) or right-translation equivariant (12) wave maps. The translation equivariance ansatz gives a consistent reduction of the frame field equations and corresponding wave map equation only if the target geometry $`(G,g,p)`$ is invariant under right multiplication by the translation group generated by the Lie algebra element $`R`$ appearing in (12) to (14) for wave maps and (62) for frame fields. We refer to this condition, given by (162) and (163), as translation invariance of the target. As shown in Proposition A in the appendix, every compact semi-simple Lie group $`G`$ admits a translation invariant geometry $`(G,g,p)`$, except that the dimension of $`G`$ must be greater than three to support a non-zero torsion $`Q`$ (see Remark 1). ### Proof of Theorem 4: The proof of Theorem 4 is very similar to that of Theorem 2. We summarize the differences (if any) in each step. ### Step 1: Conserved Energy Since the expression for the energy is unchanged and since it is conserved, there are no changes in obtaining $`L^2`$ bounds for $`E^a(x,t),H^a(x,t),B^a(x,t)`$. ### Step 2: Bounded $`H^a`$ Instead of (115), we have $`{\displaystyle _\mathrm{\Sigma }}\|_xH^a\|𝑑x{\displaystyle _\mathrm{\Sigma }}\|C_{bc}{}_{}{}^{a}E_{}^{b}H^c\|𝑑x+{\displaystyle _\mathrm{\Sigma }}\|C_{bc}{}_{}{}^{a}E_{}^{b}R^c\|𝑑x.`$ (164) The first of the two terms on the right hand side of (164) may be handled as in (116). As for the second term, we have $`{\displaystyle _\mathrm{\Sigma }}\|C_{bc}{}_{}{}^{a}E_{}^{b}R^c\|𝑑x`$ $``$ $`k_{25}{\displaystyle _\mathrm{\Sigma }}\|E^b\|𝑑x`$ (165) $`<`$ $`k_{25}\left({\displaystyle _\mathrm{\Sigma }}E^2𝑑x\right)^{1/2}`$ $``$ $`k_{26}`$ where the second inequality uses the compact support of $`E^b`$ together with the Holder inequality, and the last inequality follows from the $`L^2`$ bound on $`E^a`$. Hence we obtain $`{\displaystyle _\mathrm{\Sigma }}\|_xH^a\|𝑑xk_{27}`$ (166) analogous to (117). If $`\mathrm{\Sigma }=R^1`$, the argument leading to a pointwise bound on $`H^a(x,t)`$ for $`tI`$ proceeds exactly as in the proof of Theorem 2. If $`\mathrm{\Sigma }=S^1`$, then we need to modify the argument which begins with (LABEL:LoneSEy). We have $`{\displaystyle _{S^1}}H^a(x,t)𝑑x`$ $`=`$ $`{\displaystyle _{t_0}^t}{\displaystyle _{S^1}}C^a{}_{bc}{}^{}B_{}^{b}(x,s)\left(H^c(x,s)+R^c\right)𝑑x𝑑s`$ (167) $`+{\displaystyle _{S^1}}H^a(x,t_0)𝑑x.`$ The term $`_{t_0}^t_{S^1}C^a{}_{bc}{}^{}B_{}^{b}(x,s)R^c𝑑x𝑑s`$ can be bounded from above using the same quadratic inequality that is used in (165), and so we obtain $`\|{\displaystyle _{S^1}}H^a(x,t)𝑑x\|k_{28}`$ (168) analogous to (122). The argument for pointwise bounds on $`H^a(x,t)`$ for $`\mathrm{\Sigma }=S^1`$ can then be completed as in the proof of Theorem 2. ### Step 3: Bounded $`E^a`$ and $`B^a`$ From inequalities (129) and (130) onward, the light-cone arguments used to bound $`E^a(x,t)`$ and $`B^a(x,t)`$ in the proof of Theorem 2 work identically to bound $`E^a(x,t)`$ and $`B^a(x,t)`$ here. To arrive at (129) and (130) we use the following equations, analogous to (125) and (126), $`_n(B+E)^2`$ $`=`$ $`2H^b(H^c+R^c)(B^a+E^a)C_{abc}`$ (169) $`2\lambda (v_t+v_x)H^a(B^b+E^b)(B^c+E^c)Q_{bca}`$ $`+2\lambda (v_tv_x)C_{bc}{}_{}{}^{a}R_{}^{b}(B^c+E^c)(B^bE^b)p_{ab},`$ $`_n(B+E)^2`$ $`=`$ $`2H^b(H^c+R^c)(B^a+E^a)C_{abc}`$ (170) $`2\lambda (v_t+v_x)H^a(B^b+E^b)(B^c+E^c)Q_{bca}`$ $`+2\lambda (v_t+v_x)C_{bc}{}_{}{}^{a}R_{}^{b}(B^c+E^c)(B^b+E^b)p_{ab}.`$ Adopting the notation $`\alpha ^a:=B^a+E^a`$, $`\beta ^a:=B^a+E^a`$, these equations become $`_n\alpha ^2`$ $`=`$ $`2C_{abc}H^b(H^c+R^c)\alpha ^a`$ (171) $`2\lambda (v_t+v_x)Q_{bca}H^a\beta ^b\alpha ^c+2\lambda (v_t+v_x)R^b\alpha ^c\beta ^bC_{bc}{}_{}{}^{a}p_{ab}^{},`$ $`_{\mathrm{}}\beta ^2`$ $`=`$ $`2C_{abc}H^b(H^c+R^c)\beta ^a`$ (172) $`2\lambda (v_t+v_x)Q_{bca}H^a\beta ^a\alpha ^c+2\lambda (v_t+v_x)R^b\beta ^c\alpha ^bC_{bc}{}_{}{}^{a}p_{ab}^{}.`$ Then, noting that $`C_{abc}`$, $`Q_{abc}`$, $`\lambda `$, $`v_t`$, $`v_x`$ and $`R`$ are constant, and recalling that $`H^a`$ is bounded on $`\mathrm{\Sigma }\times I`$, we obtain $`_n\alpha ^2k_{29}\sqrt{\alpha ^2}+k_{30}\sqrt{\alpha ^2}\sqrt{\beta ^2}`$ (173) and $`_{\mathrm{}}\beta ^2k_{31}\sqrt{\beta ^2}+k_{32}\sqrt{\alpha ^2}\sqrt{\beta ^2}`$ (174) which are identical to (129) and (130). ### Step 4: Bounded Derivatives One can see in Step 2 and Step 3 that the presence of the commutator terms involving $`R`$ in the field equations (157) to (160) changes little in the arguments for boundedness, since these extra terms are easily controlled by the analogous quadratic terms appearing in the equations. The same holds true for Step 4. We can define the derivative stress-energy components just as in (143) to (146) and then obtain conservation equations similar to (147) to (148), with small modifications in the expressions $`Y_1`$ and $`Y_2`$ which appear there. These modifications are readily handled in deriving the estimate (151) and (152). The rest of the argument proceeds unchanged. Hence we obtain global existence. ## 6 CONCLUDING REMARKS The wave map global existence results we have obtained here extend previous work in two significant ways. First, our study of translation equivariant critical wave maps for Lie group targets (Theorems 3 and 5) provides a counterpart to work on rotationally equivariant critical wave maps for symmetric-space targets (see ). Second, our inclusion of torsion gives an interesting generalization of critical wave maps for arbitrary targets, which ties into current work on integrable chiral models in 2+1 dimensions in the case of Lie group targets . Furthermore, our results demonstrate the utility of the frame formulation of wave maps for Lie group targets (Proposition 1). The translation-equivariant reduction of critical wave maps studied here is motivated by this formulation and the analysis is especially straightforward in terms of frames. An important question to investigate for future work is how the frame formulation might help in understanding the unreduced critical wave map equation for general Lie group targets and symmetric-space targets. ## APPENDIX Proposition A. Let $`G`$ be a semi-simple Lie group with commutator structure tensor $`C_{bc}^a`$. 1. The Lie algebra of $`G`$ admits a translation invariant (162) positive-definite metric $`g_{ab}`$ if $`G`$ is compact. 2. The Lie algebra of $`G`$ admits a translation invariant (163) skew-tensor $`p_{ab}`$ with non-zero torsion (54) if $`G`$ is compact and has dimension greater than three. 3. If $`G`$ has dimension three then the torsion (54) is zero for every skew-tensor $`p_{ab}`$ on the Lie algebra of $`G`$. Proof of 1: If $`G`$ is compact then its Lie algebra admits an invariant positive-definite metric $`g_{ab}`$ (see, e.g. ), which satisfies $`g_{ae}C_{bc}{}_{}{}^{e}=g_{be}C_{ac}{}_{}{}^{e}.`$ (175) (In particular, the Cartan-Killing metric given by $`g_{ab}:=C_{ae}{}_{}{}^{c}C_{bc}^{}^e`$ is both invariant and positive-definite.) Hence condition (162) holds. Proof of 2 and 3: Hereafter $`g_{ab}`$ denotes the Cartan-Killing metric. We first remark that, for any $`G`$, the natural construction $`p_{ab}:=C_{ab}{}_{}{}^{d}g_{de}^{}R^e`$ (176) is easily seen to yield a translation invariant skew-tensor. But the resulting torsion tensor (54) is always zero, since $`p_{e[a}C_{bc]}{}_{}{}^{e}=C_{e[a}{}_{}{}^{d}C_{bc]}^{}{}_{}{}^{e}g_{df}^{}R^f=0`$ (177) by the Jacobi identity. In three dimensions it is easy to show that $`C_{ab}{}_{}{}^{e}g_{ec}^{}`$ must be proportional to the totally-skew Levi-Civita tensor $`ϵ_{abc}`$, while any skew-tensor $`p_{ab}`$ can be expressed in the form $`p_{ab}=ϵ_{abc}p^c`$ (178) for some vector $`p^c`$ in the Lie algebra of $`G`$. Thus, it follows that $`p_{ab}`$ must have the form (176) where $`R^e`$ is proportional to $`p^e`$, and hence from (176) and (177) we have that the torsion tensor (54) is zero. This shows that there is no torsion for any three-dimensional $`G`$ (and hence none in particular with $`p_{ab}`$ being translation invariant). Now suppose $`G`$ has dimension greater than three. In this case, $`G`$ must have rank greater than one and hence the Lie algebra of $`G`$ possesses an abelian subalgebra of dimension at least two (see, e.g. ). This allows the explicit construction of a translation invariant skew-tensor $`p_{ab}`$ as follows. Let $`p^a`$, $`q^a`$ be any two (linearly independent) commuting vectors in the Lie algebra of $`G`$, so $`p^aq^bC_{ab}{}_{}{}^{c}=0`$, and let $`p_e:=g_{ea}p^a,q_e:=g_{ea}q^a`$. Set $`R^a:=\alpha p^a+\beta q^a0`$ with constants $`\alpha ,\beta `$. Then it is straightforward to show that the skew-tensor defined by $`p_{ab}:=2p_{[a}q_{b]}`$ (179) is translation invariant as a consequence of $`p`$ and $`q`$ commuting with $`R`$. Now it remains to show that the torsion tensor given by (54) and (179) is non-zero. We have $`g_{ad}Q^d{}_{bc}{}^{}=3/2(p_eq_{[a}C_{bc]}{}_{}{}^{e}q_ep_{[a}C_{bc]}{}_{}{}^{e}).`$ (180) To show that the tensor (180) is non-zero when $`G`$ is compact, we contract (180) with the vector $`s^a=p^aq^eq_eq^ap^eq_e`$ satisfying $`s^aq_a=0`$. This yields $`s^ag_{ad}Q^d{}_{bc}{}^{}=1/2s^ap_aC_{bc}{}_{}{}^{e}q_{e}^{}.`$ (181) with $`s^ap_a=p^ap_aq^dq_d(p^aq_a)^20`$ due to positive-definiteness of $`g_{ab}`$. Moreover, since $`G`$ is semi-simple, its Lie algebra has empty center and so $`C_{bc}{}_{}{}^{e}q_{e}^{}=C_{eb}{}_{}{}^{a}q_{}^{e}g_{ca}`$ is non-zero (that is, there exists a vector $`v^b`$ so that $`C_{eb}{}_{}{}^{a}q_{}^{e}v^b0`$). Therefore, $`s^ag_{ad}Q^d_{bc}`$ is non-zero and thus so is the torsion tensor (180). ACKNOWLEDGEMENTS Various portions of this work were carried out at the Courant Institute and at the University of Washington. Partial support for this research has come from NSF grant PHY-9800732 at the University of Oregon.
warning/0007/nucl-ex0007014.html	ar5iv	text	# A consistent analysis of (e,e′p) and (d,3He) experiments ## 1 Introduction Spectroscopic factors deduced from (e,ep) reactions (see Fig. 1) are found to be substantially lower than the sum-rule limit given by the independent-particle shell model (IPSM). In contrast to this, experiments with hadronic probes such as the (d,<sup>3</sup>He) reaction, generally find spectroscopic factors that are close to values predicted by the IPSM. However, there is a strong model dependence in the extraction of these spectroscopic factors from transfer reactions (see e.g. Ref. and references therein). In this paper it will be investigated to what extent this model dependence can account for the apparent discrepancy between reported spectroscopic factors derived from (e,ep) and from (d,<sup>3</sup>He) experiments. Modern nuclear-structure calculations predict occupations for valence orbitals in the range of 60 to 90 % of the IPSM limit. The precise value and the spreading of the strength depend sensitively on the amount of short- and long-range correlations included in the calculation. Recently, it was demonstrated for the nucleus <sup>7</sup>Li that structure calculations based on a realistic nucleon-nucleon potential are indeed able to describe accurately the momentum distributions and spectroscopic factors measured with the reaction (e,ep). To put such calculations to a further test for other nuclei it is necessary to avail of accurate absolute spectroscopic factors. In this respect the existing discrepancy between spectroscopic factors deduced from the (e,ep) and (d,<sup>3</sup>He) reactions needs a detailed investigation. It is the aim of the present paper to carry out such a study and to provide a consistent set of spectroscopic factors extracted from both the (e,ep) and the (d,<sup>3</sup>He) reaction. In Section 2 the Coulomb Distorted Wave Impulse Approximation (CDWIA), which is used in the analysis of the (e,ep) experiments, and the Distorted Wave Born Approximation (DWBA) method, used for the (d,<sup>3</sup>He) reaction, are reviewed with special emphasis on the sensitivities of the spectroscopic factors to the various approximations made. In Section 3 it is investigated, which part of the bound-state wave function (BSWF) is probed by the (e,ep) and (d,<sup>3</sup>He) reactions, in order to understand the model sensitivity arising from the shape of the BSWF. In Section 4 one <sup>48</sup>Ca(e,ep) and two <sup>48</sup>Ca(d,<sup>3</sup>He) data sets are used for a detailed comparison between the (e,ep) and (d,<sup>3</sup>He) spectroscopic factors. In Section 5 a re-analysis of (d,<sup>3</sup>He) data sets for other nuclei is made, in which non-locality and finite range corrections are included and BSWF’s deduced from (e,ep) experiments are used. Conclusions are drawn in Section 6. ## 2 Description of the reactions (e,ep) and (d,<sup>3</sup>He) ### 2.1 The (e,ep) reaction In (e,ep) experiments the energy $`e_i`$ and momentum $`\stackrel{}{k}_i`$ of the initial electron, and the energies and momenta of the final electron and knocked-out proton, denoted by $`e_f`$, $`\stackrel{}{k}_f`$ and $`E_p`$, $`\stackrel{}{k}_p`$, respectively, are measured. The energy and momentum transferred by the scattered electron are denoted by : $`\omega =e_ie_f`$ and $`\stackrel{}{q}=\stackrel{}{k}_i\stackrel{}{k}_f`$. From energy and momentum conservation the missing energy $`E_m`$ and missing momentum $`\stackrel{}{p}_m`$ are determined : $`E_m=e_ie_fE_pT_{A1}=\omega E_pT_{A1}`$ $`\stackrel{}{p}_m=\stackrel{}{k}_p^{lab}(\stackrel{}{k}_i\stackrel{}{k}_f)=\stackrel{}{k}_p\stackrel{}{q},`$ (1) where $`T_{A1}`$ is the kinetic energy of the residual nucleus. The missing energy is the energy required to separate the struck proton from the target nucleus, where the final nucleus is left in the ground-state or in one of its excited states. The missing momentum is, according to the definition of Eq. (1), the proton momentum in the nucleus just before the reaction provided that there is no further interaction between the incoming electron and the initial nucleus and the outgoing electron and proton and the final nucleus. In the plane-wave impulse approximation (PWIA, see below) the (e,ep) cross section can be written as : $$\frac{\mathrm{d}^6\sigma }{\mathrm{d}E_e\mathrm{d}\mathrm{\Omega }_e\mathrm{d}E_p\mathrm{d}\mathrm{\Omega }_p}=k\sigma _{ep}S(E_m,\stackrel{}{p}_m),$$ (2) where the left-hand side represents the measured (e,ep) cross section, $`k\sigma _{ep}`$ a kinematic factor times the elementary electron-proton cross section and $`S`$($`E_m`$,$`\stackrel{}{p}_m`$) the spectral function . The spectral function is the joint probability of finding a proton with momentum $`\stackrel{}{p}_m`$ and binding energy $`E_m`$ inside the nucleus. For a transition leading to a discrete state at $`E_m`$=$`E_{tr}`$, the spectral function is written as the momentum distribution $`\rho `$($`\stackrel{}{p}_m`$) times a delta function for the energy : $$S(E_m,\stackrel{}{p}_m)=\rho (\stackrel{}{p}_m)\delta (E_mE_{tr}).$$ (3) The spectral function as given in Eq. (2) cannot be determined experimentally because the outgoing proton interacts strongly with the final nucleus (this is often called the Final State Interaction, FSI). Moreover the factorization of the six-fold differential cross section into an elementary electron-proton cross section times the spectral function does not hold any more due to the FSI and Coulomb distortions of the electron waves. In the following a theoretical basis similar to the one used in the description of the (d,<sup>3</sup>He) reaction is given to compute the (e,ep) cross section, in which the interaction between the participating particles is taken into account. The influence of various approximations made in this analysis is also investigated in order to reveal the origin of the model uncertainties on the extracted observables. Since we want to keep the formulae transparent the angular momentum and spin parts are not given in this article. The basis of the theoretical description of the reaction : $`A+e_iB+p+e_f`$ is the T-matrix formalism . This T-matrix in the prior form is defined in the following way : $$T_{if}=<\mathrm{\Psi }_f^{()}\|V_iU_i\|\psi _A\phi _i^{(+)}>,$$ (4) where $`\phi _i^{(+)}(\stackrel{}{k}_i,\stackrel{}{r})`$ is the incoming distorted electron wave with outgoing-wave boundary conditions, $`\psi _A`$ the wave function of the target nucleus, $`V_i`$ the total interaction between the incoming electron and the nucleus, from which $`U_i`$, the potential used to generate the distorted wave $`\phi _i`$, is excluded, and $`\mathrm{\Psi }_f^{()}`$ is the exact final-state wave function of the electron-proton-residual-nucleus system obeying incoming-wave boundary conditions. The distorting potential $`U_i`$ is usually taken to be the Coulomb potential arising from a uniformly charged sphere. In the distorted-wave approximation the exact final-state wave function is approximated by the product of an internal wave function for the residual nucleus $`\psi _B`$, a distorted outgoing electron-wave $`\phi _f^{()}(\stackrel{}{k}_f,\stackrel{}{r})`$, and the distorted outgoing proton-wave $`\chi _p^{()}(\stackrel{}{k}_p,\mu \stackrel{}{R})`$ (with $`\mu =(A1)/A`$). Here the displacement of the knocked-out proton from the residual nucleus $`B`$ is denoted by $`\stackrel{}{R}`$, while $`\stackrel{}{r}`$ is the displacement of the electron from the center of mass of the residual-nucleus plus proton final system (see Fig. 2a). Under these assumptions the distorted wave transition amplitude becomes : $$T_{if}^{DW}=<\phi _f^{()}(\stackrel{}{k}_f,\stackrel{}{r})\chi _p^{()}(\stackrel{}{k}_p,\mu \stackrel{}{R})\psi _B\|V_iU_i\|\psi _A\phi _i^{(+)}(\stackrel{}{k}_i,\stackrel{}{r})>.$$ (5) For light and medium-heavy nuclei the effects of the electron distortions are small ; for heavy nuclei, however, these effects become sizable. Because of the long range of the Coulomb potential it is difficult to include electron distortions in theoretical codes that compute (e,ep) momentum distributions. Up to now two approaches have been followed to deal with these distortions. In the simplest one, the Effective Momentum Approximation (EMA) , the electron momenta are replaced by effective ones. A more precise treatment of the electron distortions is employed in the work of Giusti and Pacati , in which the eikonal approach is used to expand the electron waves in powers of $`Z\alpha `$ ($`Z`$ the nuclear charge and $`\alpha `$ the fine-structure constant). For medium heavy nuclei such as the calcium isotopes this approximation is accurate enough as all higher order terms have negligible influence on the calculated cross sections . Recently, full relativistic calculations have been used to analyze momentum distributions measured with the reaction (e,ep). The deduced spectroscopic factors are different by up to 10% from those resulting from a non-relativistic analysis. However, in these analyses relativistic Hartree-Fock wave functions are used for the BSWF, which not always provide a satisfactory description of the experimental momentum distributions. Moreover, the use of relativistic optical potentials at the low proton energies ($`T_p`$ 100 MeV) employed in the presently discussed experiments may be questionable. We therefore limit the present analysis to the non-relativistic approach. The distorted proton wave, $`\chi _p^{()}(\stackrel{}{k}_p,\mu \stackrel{}{R})`$, is usually chosen to be the solution of the Schrödinger equation with the optical potential that describes elastic proton scattering off the final nucleus, $`B`$. The parameterization of the optical potential is not unique; using a different parameterization to generate the distorted proton waves results in waves identical at large distances but different in the nuclear region where the (e,ep) reaction takes place. The effect of different parameterizations on the value of the extracted spectroscopic factors has been investigated for the reaction <sup>51</sup>V(e,ep) . A model uncertainty of about 6 % was found there due to the treatment of the final-state interaction. The optical potential used for the generation of the distorted waves as well as the binding potential for the proton are local potentials. However, for fundamental reasons this potential is expected to be non-local. Perey has pointed out that the wave function of a non-local potential is systematically smaller in the nuclear interior than the wave function of the local potential that gives an equivalent description of the elastic scattering process. This non-locality correction can be taken into account effectively by multiplying the local wave function with the factor : $$F(r)=\left(1\frac{\mu _p\beta ^2}{2\mathrm{}^2}U_L(r)\right)^{1/2},$$ (6) where $`\mu _p=\frac{A1}{A}m_p`$ is the reduced proton mass, $`U_L(r)`$ the local optical potential and $`\beta `$ the range of the non-locality. The non-locality correction affects the distorted waves in the region where the potential is significantly different from zero. This is also the region where the reaction takes place so that the non-locality correction also affects the spectroscopic factors determined from knock-out reactions. The T-matrix element $`T_{if}^{DW}`$ can be written in a way that explicitly shows the nuclear-structure part : $`T_{if}^{DW}={\displaystyle 𝑑\stackrel{}{r}𝑑\stackrel{}{R}}`$ $`\phi _f^{()}(\stackrel{}{k}_f,\stackrel{}{r})\chi _p^{()}(\stackrel{}{k}_p,\mu \stackrel{}{R})\times `$ (7) $`<\psi _B\|V_iU_i\|\psi _A>\phi _i^{(+)}(\stackrel{}{k}_i,\stackrel{}{r}).`$ The matrix element $`<\psi _B\|V_iU_i\|\psi _A>`$ contains the nuclear-structure information. It involves integration over all internal coordinates $`\xi _B`$, independent of $`\stackrel{}{r}`$ and $`\stackrel{}{R}`$. The potential $`V_i`$ describes the total interaction between the electron and the nucleus, whereas in the distorting potential $`U_i`$ the part of the interaction leading to the (e,ep) channel is excluded. In this way the difference $`V_iU_i`$ is the interaction between the electron and the struck proton : $$V_iU_i=V_{ep}(\stackrel{}{r},\mu \stackrel{}{R}).$$ (8) Since $`V_{ep}`$ is independent of the internal coordinates of $`\psi _B`$, the nuclear matrix element can be factorized into $`V_{ep}`$ and the nuclear overlap integral : $$<\psi _B\|V_iU_i\|\psi _A>=V_{ep}(\stackrel{}{r},\mu \stackrel{}{R})𝑑\xi _B\psi _B^{}(\xi _B)\psi _A(\xi _B,\stackrel{}{R}).$$ (9) The integral is usually expanded into single-particle states : $$𝑑\xi _B\psi _B^{}(\xi _B)\psi _A(\xi _B,\stackrel{}{R})=\underset{n\mathrm{}jm}{}<J_BjM_Bm\|J_AM_A>\sqrt{S_{n\mathrm{}j}}\varphi _{n\mathrm{}jm}(\stackrel{}{R}),$$ (10) where $`<J_BjM_Bm\|J_AM_A>`$ is a Clebsch-Gordan coefficient, $`\sqrt{S_{n\mathrm{}j}}`$ the spectroscopic amplitude and $`\varphi _{n\mathrm{}jm}(\stackrel{}{R})`$ a normalized single particle wave function, usually referred to as the bound state wave function (BSWF). Substituting the foregoing two expressions into the transition amplitude yields : $`T_{if}^{CDWIA}=`$ $`{\displaystyle \underset{n\mathrm{}jm}{}}<J_BjM_Bm\|J_AM_A>\sqrt{S_{n\mathrm{}j}}{\displaystyle }d\stackrel{}{r}{\displaystyle }d\stackrel{}{R}\phi _f^{()}(\stackrel{}{k}_f,\stackrel{}{r})\times `$ (11) $`\chi _p^{()}(\stackrel{}{k}_p,\mu \stackrel{}{R})V_{ep}(\stackrel{}{r},\mu \stackrel{}{R})\varphi _{n\mathrm{}jm}(\stackrel{}{R})\phi _i^{(+)}(\stackrel{}{k}_i,\stackrel{}{r}),`$ which is the Coulomb Distorted Wave Impulse Approximation (CDWIA) amplitude. In the code DWEEPY that was used to calculate the momentum distributions presented in Section 4, this expression is evaluated together with the angular momentum and spin parts, which are not shown in Eq. (11). From Eq. (11) the Distorted-Wave Impulse Approximation (DWIA) amplitude is obtained by replacing the electron waves by plane waves : $`T_{if}^{DWIA}=`$ $`{\displaystyle \underset{n\mathrm{}jm}{}}<J_BjM_Bm\|J_AM_A>\sqrt{S_{n\mathrm{}j}}{\displaystyle }d\stackrel{}{r}{\displaystyle }d\stackrel{}{R}\mathrm{exp}(i\stackrel{}{k}_f\stackrel{}{r})\times `$ (12) $`\chi _p^{()}(\stackrel{}{k}_p,\mu \stackrel{}{R})V_{ep}(\stackrel{}{r},\mu \stackrel{}{R})\varphi _{n\mathrm{}jm}(\stackrel{}{R})\mathrm{exp}(i\stackrel{}{k}_i\stackrel{}{r}).`$ In order to gain some further insight into Eq. (12) and to make the connection with the PWIA expression the Coulomb potential is now used for the interaction : $`V_{ep}=\alpha /\|\mu \stackrel{}{R}\stackrel{}{r}\|`$. With this potential the integration over $`\stackrel{}{r}`$ can be performed : $`T_{if}^{DWIA}=`$ $`{\displaystyle \underset{n\mathrm{}jm}{}}<J_BjM_Bm\|J_AM_A>{\displaystyle \frac{4\pi \alpha }{\stackrel{}{q}^2}}\sqrt{S_{n\mathrm{}j}}\times `$ (13) $`{\displaystyle 𝑑\stackrel{}{R}\mathrm{exp}(i\stackrel{}{q}\mu \stackrel{}{R})\chi _p^{()}(\stackrel{}{k}_p,\mu \stackrel{}{R})\varphi _{n\mathrm{}jm}(\stackrel{}{R})},`$ where $`\stackrel{}{q}=\stackrel{}{k}_i\stackrel{}{k}_f`$. The Plane-Wave Impulse Approximation (PWIA) amplitude is obtained from the expression (13) by replacing the distorted proton waves by plane waves : $`T_{if}^{DWIA}={\displaystyle \underset{n\mathrm{}jm}{}}<J_BjM_Bm\|J_AM_A>\sqrt{S_{n\mathrm{}j}}\times `$ $`{\displaystyle 𝑑\stackrel{}{R}\mathrm{exp}(i\stackrel{}{q}\mu \stackrel{}{R})\mathrm{exp}(i\stackrel{}{k}_p\mu \stackrel{}{R})\varphi _{n\mathrm{}jm}(\stackrel{}{R})}`$ $`={\displaystyle \underset{n\mathrm{}jm}{}}<J_BjM_Bm\|J_AM_A>\sqrt{S_{n\mathrm{}j}}{\displaystyle 𝑑\stackrel{}{R}\mathrm{exp}(i\stackrel{}{p}_m\stackrel{}{R})\varphi _{n\mathrm{}jm}(\stackrel{}{R})},`$ (14) where in the last expression the proton momentum $`\stackrel{}{k}_p`$, which is the center-of-mass momentum, has been written in terms of laboratory momenta : $`\stackrel{}{k}_p=\stackrel{}{k}_p^{lab}+\stackrel{}{p}_m/(A1)`$. Expression (14) is just the Fourier transform of the BSWF. After including the angular momentum and spin parts and squaring the $`T_{if}`$-matrix element the well known expression for the (e,ep) cross section is obtained. When distortions (of the proton or electron waves) are included, the cross section cannot be factorized any more into $`\sigma _{ep}`$ and S($`E_m`$,$`\stackrel{}{p}_m`$) (see Eq. (2)). For convenience one defines the reduced cross section or distorted momentum distribution (both the experimental and the calculated one) by $$\rho ^D(\stackrel{}{p}_m)\delta (E_mE_{tr})=S^D(E_m,\stackrel{}{p}_m)=\frac{1}{k\sigma _{ep}}\frac{\mathrm{d}^6\sigma ^{CDWIA}}{\mathrm{d}E_e\mathrm{d}\mathrm{\Omega }_e\mathrm{d}E_p\mathrm{d}\mathrm{\Omega }_p}$$ (15) (compare eqs. 2 and 3). In calculating the sixfold differential cross section there is some ambiguity in the current operator to be used as the proton is off-shell . However, in the used kinematics the influence on the spectroscopic factors of the different prescriptions as given by de Forest is smaller than a few percent. ### 2.2 The (d,<sup>3</sup>He) reaction The basis of the DWBA description of transfer reactions $`A+aB+b`$, such as the (d,<sup>3</sup>He) reaction presently under study, is the transition amplitude : $$T_{\alpha \beta }=<\mathrm{\Psi }^{()}\|V_\alpha U_\alpha \|\psi _A\psi _a\chi _\alpha ^{(+)}(\stackrel{}{k}_\alpha ,\stackrel{}{r}_\alpha )>,$$ (16) where $`\alpha (\beta )`$ is the entrance (exit) channel with projectile (ejectile) $`a(b)`$ and target (final) nucleus $`A(B)`$ and $`\stackrel{}{r}_\alpha `$ the displacement of $`a`$ from $`A`$ (see Fig. 2b). The interaction $`V_\alpha `$ is the sum of two-body interactions between the nucleons of the projectile and those of the target nucleus. The wave function $`\psi _a(\psi _A)`$ is the internal wave function of the projectile (target nucleus), while $`\chi _\alpha ^{(+)}(\stackrel{}{k}_\alpha ,\stackrel{}{r}_\alpha )`$ is the solution of the Schrödinger equation for the incoming particle with the distorting potential $`U_\alpha `$, usually chosen to be an optical potential that fits the elastic scattering in channel $`\alpha `$, and $`\mathrm{\Psi }^{()}`$ is the exact wave function of the system with incoming-wave boundary conditions. In the DWBA method the following approximations are made : First the exact wave function $`\mathrm{\Psi }^{()}`$ is replaced by a product of internal wave functions of the outgoing particle $`\psi _b`$, the residual nucleus $`\psi _B`$ and a function $`\chi _\beta ^{()}(\stackrel{}{k}_\beta ,\stackrel{}{r}_\beta )`$ describing the elastic scattering of the outgoing particle off the final nucleus B. This leads to the distorted-wave transition amplitude : $`T_{\alpha \beta }^{DW}`$ $`=`$ $`<\chi _\beta ^{()}(\stackrel{}{k}_\beta ,\stackrel{}{r}_\beta )\psi _b\psi _B\|V_\alpha U_\alpha \|\psi _A\psi _a\chi _\alpha ^{(+)}(\stackrel{}{k}_\alpha ,\stackrel{}{r}_\alpha )>`$ $`=`$ $`{\displaystyle 𝑑\stackrel{}{r}_\beta 𝑑\stackrel{}{r}_\alpha \chi _\beta ^{()}(\stackrel{}{k}_\beta ,\stackrel{}{r}_\beta )<Bb\|V_\alpha U_\alpha \|Aa>\chi _\alpha ^{(+)}(\stackrel{}{k}_\alpha ,\stackrel{}{r}_\alpha )},`$ (17) where $`\stackrel{}{r}_\beta `$ is the displacement of $`b`$ from $`B`$ (see Fig. 2a). In the nuclear matrix element $`<\psi _b\psi _B\|V_\alpha U_\alpha \|\psi _A\psi _a>`$ the integration is performed over all coordinates independent of $`\stackrel{}{r}_\alpha `$ and $`\stackrel{}{r}_\beta `$. The second approximation deals with the interaction $`V_\alpha U_\alpha `$, which is replaced in the prior formalism by the interaction between the transferred nucleon and the projectile nucleons : $$V_\alpha U_\alpha \underset{i}{\overset{a}{}}V_{in}=V_{an}(\stackrel{}{r}),$$ (18) where the sum runs over all the constituents of the projectile $`a`$, $`n`$ is the nucleon to be transferred and $`\stackrel{}{r}`$ is the displacement of the pickup nucleon from the center of mass of the projectile. (In principle this interaction should be taken off-shell. In view of the many other approximations made, this point is probably not relevant and is always neglected in analyses.) The above approximations are based on the following assumptions made for the transfer reaction mechanism itself. Firstly it is assumed that the interaction that drives the reaction is weak enough so that the reaction process may be treated in first order perturbation theory. Secondly the reaction is assumed to be a one-step process; the transferred nucleon is picked up by the incoming projectile, whereas all other target nucleons do not change their state of motion. The distorted waves, $`\chi _\alpha ^{(+)}(\stackrel{}{k}_\alpha ,\stackrel{}{r}_\alpha )`$ and $`\chi _\beta ^{()}(\stackrel{}{k}_\beta ,\stackrel{}{r}_\beta )`$, are usually chosen to be the wave functions of optical potentials describing elastic scattering in the entrance and exit channels. However, as already mentioned, the same elastic scattering data can be described with different parameterizations of the optical potential, resulting in waves identical at large distances but differing in the nuclear region. Their contribution to the transition amplitude will differ accordingly and gives rise to an extra uncertainty in the spectroscopic factors deduced from transfer reactions. In contrast to the (e,ep) reaction where only one wave function enters that is generated in an optical model potential, in the (d,<sup>3</sup>He) reaction two such wave functions are entering. Moreover, the uncertainty in the optical-model wave functions of composite particles in the interior of the nucleus is appreciably larger than that for nucleons. Consequently, the uncertainties due to different possible parameterizations of the optical-model potential are larger for spectroscopic factors deduced from the (d,<sup>3</sup>He) reaction. As pointed out earlier, non-locality corrections must be applied to the wave functions obtained from the (local) optical potential. This correction affects the wave functions of the projectile and ejectile in the region where the transfer takes place. The nuclear matrix element in Eq. (17) can be expanded, along the lines given in , into the nuclear overlap integral as given in Eq. (10) and into the overlap between the projectile and ejectile : $`f(\stackrel{}{r})=<\psi _b^{}(\xi _a,\stackrel{}{r})\|V_{an}(\stackrel{}{r})\|\psi _a(\xi _a)>`$. This gives for the DWBA transition amplitude : $$T_{\alpha \beta }^{DW}\underset{n\mathrm{}jm}{}\sqrt{S_{n\mathrm{}j}}𝑑\stackrel{}{r}_\beta 𝑑\stackrel{}{r}_\alpha \chi _\beta ^{()}(\stackrel{}{k}_\beta ,\stackrel{}{r}_\beta )\varphi _{n\mathrm{}jm}(\stackrel{}{R})f(\stackrel{}{r})\chi _\alpha ^{(+)}(\stackrel{}{k}_\alpha ,\stackrel{}{r}_\alpha ).$$ (19) The evaluation of this amplitude involves a six-dimensional integral over $`\stackrel{}{r}_\alpha `$ and $`\stackrel{}{r}_\beta `$. A reduction to a more convenient three-dimensional integral is achieved in the zero-range approximation, where the effective interaction $`V_{an}(\stackrel{}{r})`$ is assumed to have a range equal to zero, so that $`f(\stackrel{}{r})`$ $`=`$ $`D_0\delta (\stackrel{}{r}),`$ $`D_0`$ $`=`$ $`{\displaystyle d\stackrel{}{r}<\psi _b^{}(\xi _a,\stackrel{}{r})\|V_{an}(\stackrel{}{r})\|\psi _a(\xi _a)>}.`$ (20) The physical meaning of this approximation is that the ejectile $`b`$ is assumed to be emitted at the same position where the absorption of the projectile $`a`$ has taken place. The effect of neglecting the finite range of the interaction is that the spectroscopic factors deduced from transfer reactions in a zero-range analysis are larger than the ones obtained from a full finite-range analysis . Full finite-range calculations are hard to perform because of the six dimensional integral in Eq. (19). However it has been shown by Buttle and Goldfarb that the effects of the finite range of the interaction can be taken into account approximately (local energy approximation, LEA) by replacing the delta function in Eq. 20 by the following radial factor : $$\mathrm{\Lambda }(r)=[1+\frac{2m_am_p}{\mathrm{}^2m_b}R_{fr}^2(E_\alpha U_\alpha (\mu r)+E_pU_p(r)E_\beta +U_\beta (r))]^1,$$ (21) in which $`R_{fr}`$ is the finite range distance. ## 3 BSWF probing functions As shown in Section 2, the transition amplitudes of both the (e,ep) and (d,<sup>3</sup>He) reaction consist of the nuclear matrix element sandwiched between the incoming and outgoing distorted ”probing” waves (see eqs. (7) and (17)). In this section it is investigated to which part of the BSWF the (e,ep) reduced cross section and the (d,<sup>3</sup>He) cross section are sensitive. Cross sections are obtained from the T-matrix elements by integrating the radial coordinate from zero to infinity. The radial sensitivity of the cross section was investigated by varying the lower radial integration bound between 0 and 10 fm and plotting these results as : $$P(r)=\frac{1}{\mathrm{\Delta }r}(\sigma _{r\mathrm{\Delta }r/2}\sigma _{r+\mathrm{\Delta }r/2}),$$ (22) where $`\sigma _x`$ reperesents the cross section calculated with a lower limit $`x`$ on the integral over the T-matrix element : $$\sigma _x\|_x^{\mathrm{}}T(r)dr\|^2.$$ (23) In this way the separate contribution of the interval $`\mathrm{\Delta }r`$ around $`r`$ to the cross section was obtained and hence the part of the BSWF to which the reaction is sensitive can be determined. For the reaction <sup>48</sup>Ca(e,ep) these calculations were performed for the transitions leading to the 1/2<sup>+</sup> ground-state and the first 3/2<sup>+</sup> excited state in <sup>47</sup>K. The BSWF shown in the upper part of Fig. 3 that was used in these calculations was generated in a Woods-Saxon well with the parameters as given in Section 4. In that section the other parameters that entered these CDWIA calculations are also given. The results of the calculation of $`P(r)`$ for the (e,ep) reaction are shown in the middle part of Fig. 3 for different values of the missing momentum. From this figure it can be seen that the (e,ep) reaction is sensitive to the whole BSWF and the largest contribution to the momentum distribution comes from those regions in $`r`$, where also $`r^2\mathrm{\Phi }_{nlj}(r)`$ is large. For the 1d<sub>3/2</sub> orbital this is the region between r=2.5 and 4.5 fm and for the 2s<sub>1/2</sub> orbital between r=0.7 and 6.7 fm. Because of the node in the 2s<sub>1/2</sub> orbital, for low missing momenta, there is a destructive contribution to the momentum distribution from the inner lobe, whereas this contribution becomes constructive for high missing momenta. For the above mentioned transitions the sensitivity to the BSWF was also determined for the (d,<sup>3</sup>He) reaction. The DWBA calculations were performed with the parameters as given in Section 4 and the same BSWF as for the (e,ep) experiment was used. The results of these calculations are presented in the lower part of Fig. 3. Here it can be seen that apart from strong interferences between the incoming and outgoing distorted waves in the interior of the nucleus, the (d,<sup>3</sup>He) reaction is most sensitive to the region between $`r`$=5 and 10 fm. Therefore, the (d,<sup>3</sup>He) reaction is not sensitive to the details of the BSWF inside the nucleus. In the region where the (d,<sup>3</sup>He) reaction is sensitive the BSWF has the global form : $`\nu \mathrm{exp}(\kappa r)`$, where $`\kappa `$ depends on the (measured) binding energy of the proton, and the normalization $`\nu `$ depends on the depth and shape of the potential that generates the BSWF. As the spectroscopic factor is the integral of the BSWF over the total radial region, one can only determine spectroscopic factors from the (d,<sup>3</sup>He) reaction by assuming some shape for the BSWF. The conclusion is that with the (e,ep) reaction the BSWF is probed in the whole radial region whereas, with the (d,<sup>3</sup>He) reaction only the exponential tail of the BSWF is probed. This tail is very sensitive to the exact shape of the used proton-binding potential. The shape of the BSWF introduces thus a large model dependence, sometimes up to 50% , in spectroscopic factors deduced from (d,<sup>3</sup>He) experiments. Given this sensitivity, it is even questionable whether ratios of spectroscopic factors for different isotopes can be determined accurately in the (d,<sup>3</sup>He) reaction , as it is not certain that the radius of the BSWF well scales with A<sup>1/3</sup>. ## 4 Analysis of the <sup>48</sup>Ca data ### 4.1 Analysis of the <sup>48</sup>Ca(e,ep) experiment The experimental (e,ep) momentum distributions were obtained with the coincidence set-up at NIKHEF . Two metal foils with a thicknesses of 7.3 and 15.0 mg/cm<sup>2</sup>, enriched to 95.2 % in <sup>48</sup>Ca were used. Reduced cross sections were obtained under parallel kinematic conditions in the range between -60 and 260 MeV/c. The electron beam energy was 440 MeV and the outgoing proton kinetic energy was 100 MeV. The experimental systematic error on the extracted distributions is 4 %. Further details can be found in Ref. . The CDWIA calculations were performed with the code DWEEPY . The proton optical-potential parameters were obtained from the work of Schwandt et al. . A non-locality correction according to the prescription of Perey (see also Eq. (6)) was applied with a range parameter $`\beta `$ of 0.85 fm. The bound state wave function was calculated in a Woods-Saxon well with a diffuseness $`a_0`$ of 0.65 fm and a Thomas spin-orbit parameter $`\lambda `$ of 25. A non-locality correction was also applied to the BSWF with a $`\beta _{nloc}`$ of 0.85 fm. The well depth $`V_0`$ and the radius parameter r<sub>0</sub> were adjusted with the separation energy as a constraint to get the best description of the measured momentum distributions. In Fig. 4 these reduced cross sections are shown for the transitions to the first three positive parity states together with the results of the CDWIA calculations, while in Table 1 the deduced spectroscopic factors and radii of the BSWF are given. ### 4.2 Re-analysis of the <sup>48</sup>Ca(d,<sup>3</sup>He) experiments For the comparison with the data from the (e,ep) experiment a re-analysis of two ($`\stackrel{}{\mathrm{d}},^3`$He) experiments was performed. The first (d,<sup>3</sup>He) experiment was performed with an incoming deuteron energy of 79.2 MeV. Further details of this experiment can be found in the original paper . In the analysis in that paper a local zero-range DWBA calculation was used for the extraction of the spectroscopic factors together with a BSWF potential well with $`r_0`$=1.25 fm, $`a_0`$=0.60 fm and $`V_0`$ adjusted to get the correct binding energy. Non-locality corrections were not applied to the BSWF. The ratio of the spectroscopic factors given in to those deduced from the present (e,ep) experiment for several discrete states is shown in Fig. 5a as a function of the excitation energy. The (d,<sup>3</sup>He) spectroscopic factors calculated this way are on the average 50% higher than those obtained from the (e,ep) experiment. In the present re-analysis, performed with the code DWUCK4 , non-locality corrections and finite range corrections via the LEA approach were included together with the BSWF obtained from the present (e,ep) experiment. The same optical model parameter sets for the deuteron and <sup>3</sup>He waves were used as in Ref. . Spectroscopic factors deduced with this re-analysis are given in Table 1. In order to estimate realistic errors on the spectroscopic factors deduced from the (d,<sup>3</sup>He) experiment the following sources of uncertainties were taken into account : i) a total experimental systematic error of 5% which includes the error on the target thickness; ii) the effect of the uncertainty (about 3%) in the rms radii obtained from the (e,ep) experiment, which yields 25% for the transition to the 1/2<sup>+</sup>, 28% for the transition to the 3/2<sup>+</sup> and 29% for the transition to the 5/2<sup>+</sup>state; iii) a 10% uncertainty due to the $`<\mathrm{d}\|^3`$He$`>`$ overlap function (the value of $`D_0`$); and iv) at least a 10% uncertainty due to different possible parameterizations of the deuteron and <sup>3</sup>He optical potentials. The latter number is taken from , where the sensitivity of the spectroscopic factors to different optical potential parameterizations was investigated for the reaction <sup>51</sup>V(d,<sup>3</sup>He)<sup>50</sup>Ti at 53 MeV. The ratio of the spectroscopic factors obtained from the present analysis of the (d,<sup>3</sup>He) data to those obtained from the (e,ep) experiment is plotted in Fig. 5b. The average ratio is one, so it is concluded that, except for two points, there is a very good agreement between the spectroscopic factors obtained from both reactions. This agreement is obtained by including non-locality and finite-range corrections in the analysis together with experimental BSWF’s obtained from the (e,ep) experiment. Finite-range corrections reduce the spectroscopic factors by about 15 % and the use of the BSWF determined from the (e,ep) analysis gives a a further reduction of 30 to 40 %. The deviation of the spectroscopic factors and the relatively small rms radius for the transition leading to the 3.42 MeV excited state in <sup>47</sup>K might be ascribed to some unresolved 1f<sub>7/2</sub> strength at 3.4 MeV but in the scarce literature on the level scheme of <sup>47</sup>K no 7/2<sup>-</sup> states have been reported so far. The reduced cross section for this transition can also be described well with a BSWF with an rms radius of 3.47 fm, which is the average value for the 1d<sub>5/2</sub> orbital obtained from all the 5/2<sup>+</sup> transitions observed in the (e,ep) experiment. This gives a 3 % lower spectroscopic factor in the (e,ep) experiment, but in the (d,<sup>3</sup>He) analysis the spectroscopic factor drops by 14 %. The deviation for the spectroscopic factor for the very weak transition at 6.87 MeV may be due to the uncertainty in the rms radius that is not well determined from the (e,ep) experiment. It is also possible that two-step processes have a different effect on the (e,ep) and (d,<sup>3</sup>He) cross sections for this weak transition. The second experiment was performed with an incoming deuteron energy of 56 MeV . Angular distributions and asymmetries were measured for the first three positive parity transitions, leading to 1/2<sup>+</sup>, 3/2<sup>+</sup>and 5/2<sup>+</sup> states in <sup>47</sup>K, see Fig. 6. The used optical-model potential parameterizations for the deuteron and <sup>3</sup>He waves were obtained from elastic deuteron and <sup>3</sup>He scattering off <sup>48</sup>Ca and are listed in Table 2. The non-locality corrections were taken into account according to the prescription of Perey (see also Eq. (6)). For the deuteron and <sup>3</sup>He wave-functions the parameters are given in the last column of Table 2. The same BSWF as obtained in the analysis of the (e,ep) experiment was employed in the calculation of the (d,<sup>3</sup>He) cross sections. Finite-range effects were included by applying the LEA correction (Eq. (21)). A finite range distance of 0.77 fm was used together with the Bassel normalization $`D_0`$=2.95 for the overlap between the deuteron and the <sup>3</sup>He ejectile. The LEA approach was compared to full finite-range calculations performed with the code DWUCK5 . In the latter calculations the D-state of the deuteron was included and a (d,<sup>3</sup>He) overlap function was used that yields $`D_0`$=2.95. The cross sections in the full finite range calculations were globally 10 to 15 % larger (and hence the spectroscopic factors smaller) than in the LEA calculation. As the used value of $`D_0`$ has also an uncertainty for convenience all DWBA calculations to be presented were performed in LEA and the presented spectroscopic factors were obtained from those. In Fig. 6 the results of DWBA calculations are shown for the transitions mentioned. Both angular distributions and asymmetries are described well with the used optical potentials and the BSWF’s obtained from the (e,ep) experiment. The spectroscopic factors extracted for these three transitions are given in Table 1. An estimate of the errors on the spectroscopic factors from this (d,<sup>3</sup>He) experiment gives typically numbers between 25 and 35 %. ## 5 (d,<sup>3</sup>He) spectroscopic factors for other nuclei A similar comparison as for <sup>48</sup>Ca has also been made for other nuclei where good (e,ep) and (d,<sup>3</sup>He) data exist. For these nuclei the (d,<sup>3</sup>He) data were re-analyzed in the same way as described above. The optical potentials were taken from the original papers, non-locality and finite-range corrections were included in the same way as for <sup>48</sup>Ca and the BSWF’s were taken from the (e,ep) work. Only pick-up from the valence shells was considered. The result of this comparison is presented in Table 3, while in Fig. 7 the spectroscopic factors expressed as a fraction of the IPSM limit are shown. The agreement between the spectroscopic factors for these transitions from the (d,<sup>3</sup>He) experiments and (e,ep) experiments is very good. The average ratio of (d,<sup>3</sup>He) over (e,ep) spectroscopic factors is 1.01 with a spread of 0.25. The error on the spectroscopic factors obtained from the (e,ep) experiments, typically 10%, was taken from their respective references and a 25% error was assigned to the spectroscopic factors from the (d,<sup>3</sup>He) experiments. The latter error is mainly due to to the large dependence of the spectroscopic factors on the rms radius of the BSWF : $`\mathrm{\Delta }S/S7\mathrm{\Delta }r_{rms}/r_{rms}`$. In view of the above conclusion that the spectroscopic factors deduced from (e,ep) and (d,<sup>3</sup>He) reactions are in agreement, and exhaust only about 60 % of the IPSM value, the question may be asked why previously applied spin-dependent sum rules for hadronic transfer reactions yielded values close to 100 %. As argued earlier by us , this sum rule, which connects stripping and pick-up strengths, is only valid if all strength for a given spin is included in the summation. This condition is clearly not fulfilled, as is known both from experimental data and from modern nuclear-structure calculations. In angular-momentum decompositions of spectral functions obtained with the reaction (e,ep), it has been shown that spectroscopic strength distributions for a given angular momentum possess long tails extending to large energies. These tails, which may contain up to 20 % of the total strength in the energy distribution, have not been included in the sum-rule analysis. Moreover, calculations of correlated nuclear matter have shown that not only the energy distributions of hole states, but also those of particle states exhibit such tails, which extend to several hundred MeV beyond the quasi-particle pole. Consequently, application of the sum rule will necessarily fail, since appreciable parts of both the hole and the particle strength are lacking in the summation. ## 6 Conclusion In this article it has been shown that spectroscopic factors obtained from (e,ep) and (d,<sup>3</sup>He) experiments are mutually consistent, provided that in the DWBA calculations for the analysis of the (d,<sup>3</sup>He) data non-locality and finite-range corrections are included together with the BSWF obtained from (e,ep) experiments. It was also shown that the (e,ep) reaction is sensitive to the whole BSWF, whereas the (d,<sup>3</sup>He) reaction is only sensitive to the exponential tail of the BSWF. This tail is very sensitive to the assumed shape of the potential well used to generate the BSWF. From the consistency of the obtained results it can be concluded that the reaction mechanism for these transitions in (e,ep) as well as in the (d,<sup>3</sup>He) reaction is understood well enough to obtain reliable nuclear structure information. From both reactions a spectroscopic strength of about 50 to 70% of the IPSM limit is found for strong valence transitions. We would like to thank Dr. N. Matsuoka for making available to us the the <sup>48</sup>Ca(d,<sup>3</sup>He) data and the deuteron and <sup>3</sup>He optical potentials, and Dr. P.D. Kunz for many useful discussions on finite-range corrections and making available the codes DWUCK4 and DWUCK5. This work is part of the research program of the National Institute for Nuclear Physics and High-Energy Physics (NIKHEF), which is made possible by the financial support from the Foundation for Fundamental Research of Matter (FOM) and the Netherlands Organization for the Advancement of Research (NWO).
warning/0007/cond-mat0007335.html	ar5iv	text	# Perturbative polydispersity: Phase equilibria of near-monodisperse systems ## I Introduction Materials whose constituent elements are much larger than atoms or simple molecules are now ubiquitous. Such substances include polymer blends, colloidal suspensions and emulsions. In the manufacture of such large components, it is impossible to produce a pure strain of truly identical particles. Instead, each colloidal latex is minutely different in size from its colleagues, and the number of monomers on each long polymer molecule inevitably varies across some distribution. These systems are said to be ‘polydisperse’, though the word is applied only loosely to the polymeric case, since not every molecule is unique. Some intriguing experimental and simulational data have been collected from polydisperse systems. For instance, unusual textures have been found in polydisperse polymer blends , and coexisting phases of hard spheres have been shown to contain unequal proportions of the various particle sizes, in both simulation and experiment . It is important, then, to understand the statistical thermodynamics of polydisperse mixtures. It is well known how to calculate the equilibrium state of a thermodynamic system . At fixed temperature, an expression must be found for the Helmholtz free energy. It is then straightforward to determine any property of the system at equilibrium: its value is that which minimises the free energy. To establish the concentrations of the system’s various constituents, the minimisation must be performed subject to constraints of global conservation. The constitution of a uniform system cannot vary, because the amount of each component is conserved, but the free energy can sometimes be lowered by partitioning into coexisting phases of differing compositions. The non-trivial part of this procedure is to establish a formula for the Helmholtz free energy. In principle, the subsequent minimisation to derive phase equilibria is simple. For a polydisperse system the story is quite different. Here, the number of components is thermodynamically large. In general, this does not greatly complicate the (already difficult) many-body problem of formulating the free energy. However, since that free energy is now a function of infinitely many conserved variables, the minimisation procedure, though formally understood , becomes intractable. In the past, many calculations have been performed to study the effects of polydispersity, using pedagogical model systems with a convenient form of Hamiltonian , species distribution , or approximate free energy . Also, approximate solutions can be found by making ad hoc guesses about the compositions of coexisting phases, and then minimising with respect to just a few variables such as the overall density of each phase . Alternatively, the minimisation can be performed numerically, by replacing the continuous distribution of concentrations with an arbitrary, finite set of variables . The above methods invariably involve some way of approximating the many-body problem inherent in the thermodynamics of the particular system. The present study differs in that we shall calculate only the effects of polydispersity. The analytical procedure was outlined recently , and is presented in more detail in the present article, and applied in new ways to solve a range of problems in the physics of polydispersity. The method is a perturbative one, whereby polydisperse phase equilibria are derived from a reference state that is monodisperse and also in multi-phase equilibrium. We make no assumptions as to the nature of the system, and require no special properties or approximations to ensure tractability of the statistical mechanics, since we shall not address the underlying many-body problem. Instead, we assume that the many-body properties of the monodisperse reference system are known (whether from experiment or theory), and derive some exact thermodynamic relations, applicable to most systems, for the changes in its properties due to the introduction of polydispersity. The aim then is to treat a polydisperse system as a perturbation to a monodisperse reference state, using the variables that characterise the polydisperse particles (for instance, size and charge) as expansion parameters. There are two major obstacles to such a theory. Firstly, the reference state is singular. To calculate thermodynamic equilibria, a knowledge of the chemical potential of each species is required. In the reference system, the population is zero for all species but one. Unfortunately, the chemical potential of a non-populated species is (as a rule) negative infinity. The second obstacle is that a narrow distribution of species is not necessarily a smooth one, and therefore may not be conducive to small-variable expansion. To proceed, then, we must isolate the badly behaved functions for exact treatment, and expand only smoothly varying quantities. The distribution of species is not expanded (contrast ), so that any such distribution may be treated, including a set of delta peaks representing a finite number of components. The canonical ensemble is used throughout, as it is of greatest practical relevance. Thus, the overall mix of species (the ‘parent distribution’) is specified a priori, and the pressure and chemical potentials are derived quantities, used to calculate the resulting distribution adopted by each coexisting phase after partitioning. Perturbative methods suggested previously have been dogged by a proliferation of variables, leading to unwieldy expressions. That problem is avoided here by relating all properties of the polydisperse system to just two or three functions that parameterise a generic free energy. An alternative approach for perturbing about a monodisperse reference is to assume that all species but one are dilute, and to use their concentrations as small expansion parameters , in the manner of a virial series. In the present study, by expanding in the width, rather than height, of the distribution of polydisperse components, we are using a reference state that correctly treats the many-particle interactions of concentrated, phase-separated systems. For narrow distributions therefore, this series is expected to converge more quickly than a virial-like expansion. An exception is at critical points, where both methods fail. The formalism and notation are set out in the next section, where a generic free energy is perturbatively expanded to low order, to derive the pressure and chemical potentials of a single polydisperse phase. These expressions are used in section III to analyse polydisperse phase equilibria. The normalised distributions, giving the relative amount of each species in each phase, are calculated in section III A. Their differences (due to fractionation) are found to obey a very simple law. For systems that are simultaneously polydisperse in two properties (e.g. size and charge), this law is shown to describe ‘convective fractionation’, whereby one property is partitioned between phases due to a driving force on the other. In section III B, the effect of polydispersity on the total number densities at coexistence is calculated. The resulting shift of a phase boundary along the density-axis of a phase diagram is shown to be proportional to the variance (standard deviation squared) of the parent distribution, in the limit of a narrow distribution. Explicit expressions are given for the cloud- and shadow-point densities, as well as more general binodals. All the derived formulae quantifying polydisperse phase equilibria are expressed in terms of a couple of basic parameters of the free energy. In case the free energy is not known for the polydisperse system in question, a method for evaluating the relevant parameters is given in section IV A for soft interaction potentials, and in section IV B for hard spheres and related systems. The method for soft potentials is extended to determine the lowest-order effect of polydispersity on correlation functions in appendix A. To demonstrate the ease of use of the formulae for phase equilibria, they are applied in section V to a Flory-Huggins model of chemically polydisperse polymers, and in section VI to fluid-fluid coexistence in colloid-polymer mixtures, where slight polydispersity is shown to widen the coexistence region. ## II Properties of a single phase Consider a polydisperse phase of overall number density $`\rho `$, whose constituent particles are distinguished by a number $`\nu `$ of properties (e.g. mass, charge, diameter, oblateness…). Each particle is characterised by an $`\nu `$-component vector $`\mathit{ϵ}`$, drawn from a narrow distribution $`p(\mathit{ϵ})`$ normalised over the whole of $`\mathit{ϵ}`$-space. Let each component of $`\mathit{ϵ}`$ represent the (dimensionless) fractional deviation of each property from some reference value. For instance, a hard sphere of radius $`r`$ in a polydisperse sample would be characterised by a scalar $`ϵ=(rr_0)/r_0`$, in terms of the reference size $`r_0`$. Since the distribution $`p(\mathit{ϵ})`$ is narrow, the components of $`\mathit{ϵ}`$ are always small numbers. To isolate the badly behaved part of the chemical potential, we shall write the phase’s free energy density (which is a functional of the distribution of species densities $`\rho p(\mathit{ϵ})`$) in two parts $$[\rho p(\mathit{ϵ})]^{\mathrm{id}}+^{\mathrm{ex}}.$$ (1) This is an exact statement, as it serves only to define $`^{\mathrm{ex}}`$ as the excess free energy density, over and above that of a polydisperse ideal gas of densities $`\rho p(\mathit{ϵ})`$. The ideal part is simply the sum, over the continuum of species, of the free energy density of an ideal gas of each species, $$^{\mathrm{id}}[\rho p(\mathit{ϵ})]=\text{d}^\nu ϵ\rho p(\mathit{ϵ})\left\{\mathrm{ln}(\rho p(\mathit{ϵ}))1\right\}.$$ (2) The chemical potentials for the continuum of species are given by a functional derivative of $``$ with respect to the density of each species, $$\mu (\mathit{ϵ})\frac{\delta }{\delta [\rho p(\mathit{ϵ})]}$$ (3) which, from Eq. (1), splits into two parts, $$\mu (\mathit{ϵ})=\mu ^{\mathrm{id}}(\mathit{ϵ})+\mu ^{\mathrm{ex}}(\mathit{ϵ}).$$ (4) By functional differentiation of Eq. (2), the ideal part of the chemical potential is $$\mu ^{\mathrm{id}}(\mathit{ϵ})=\mathrm{ln}(\rho p(\mathit{ϵ}))$$ (5) which is singular in the monodisperse limit $`p(\mathit{ϵ})\delta (\mathit{ϵ})`$, because it derives from the entropy of mixing. The other part, $`\mu ^{\mathrm{ex}}(\mathit{ϵ})`$, remains well behaved, since it describes the physics of interactions which vary little from one species to the next. The obstacles to a small-variable expansion, discussed above, have now been isolated in the badly-behaved but well-defined function $`\mu ^{\mathrm{id}}(\mathit{ϵ})`$ in Eq. (5), and the narrow, but possibly rapidly varying distribution $`p(\mathit{ϵ})`$, also appearing in Eq. (5). These quantities will be treated exactly, while the effects of interactions, embodied in the excess free energy density $`^{\mathrm{ex}}`$, are expanded in the small variable $`\mathit{ϵ}`$. ### A Expansion of the excess free energy The excess free energy density in a phase at equilibrium is a function(al) $`^{\mathrm{ex}}([p(\mathit{ϵ})],\rho )`$ of the intensive variables $`p(\mathit{ϵ})`$ and $`\rho `$. The dependence of $`^{\mathrm{ex}}`$ on temperature and other fields is suppressed in the notation, and $`^{\mathrm{ex}}`$ is assumed to be absolutely minimised with respect to any non-conserved order parameters. Any distribution $`p(\mathit{ϵ})`$ is uniquely defined by its moments $`\{\mathit{ϵ},\mathit{ϵ}\mathit{ϵ},\mathit{ϵ}\mathit{ϵ}\mathit{ϵ},\mathrm{}\}`$ which, for vectorial $`\mathit{ϵ}`$, are averages over $`p(\mathit{ϵ})`$ of outer products of $`\mathit{ϵ}`$. For instance, $`\mathit{ϵ}\mathit{ϵ}{\displaystyle \mathit{ϵ}\mathit{ϵ}p(\mathit{ϵ})\text{d}^\nu ϵ}.`$ Since $`p(\mathit{ϵ})`$ is narrow, higher moments are increasingly small. So a controlled small-variable expansion of $`^{\mathrm{ex}}`$ is $$^{\mathrm{ex}}(\rho ,\mathit{ϵ},\mathit{ϵ}\mathit{ϵ},\mathrm{})=^{\mathrm{ex}}(\rho ,0,0,\mathrm{})+\mathit{ϵ}𝑨(\rho )+\mathit{ϵ}\mathit{ϵ}:𝑩(\rho )+\mathit{ϵ}\mathit{ϵ}:𝑪(\rho )+Ø(ϵ^3)$$ (6) where $`^{\mathrm{ex}}(\rho ,0,0,\mathrm{})`$ is the free energy of a phase of monodisperse particles all characterised by $`\mathit{ϵ}=\mathrm{𝟎}`$. We shall write this as $`_\mathrm{m}^{\mathrm{ex}}(\rho )`$, where the subscript m denotes a quantity in the monodisperse reference system. The coefficients $`𝑨`$, $`𝑩`$ and $`𝑪`$ are vector and tensor functions of the overall density $`\rho `$. The notation $`Ø(ϵ^3)`$ indicates terms of third order in $`\mathit{ϵ}`$ and higher. Equation (6) is a generic expansion, not specific to any particular system. The use of a non-singular expansion, involving only integer powers of small quantities, has the status of a conjecture. An alternative motivation for the form of Eq. (6) is given in section IV. ### B Expansion of the excess chemical potential Following Eq. (3), Eq. (6) is differentiated, using the identity $`{\displaystyle \frac{\delta ^{\mathrm{ex}}}{\delta [\rho p(\mathit{ϵ})]}}{\displaystyle \frac{^{\mathrm{ex}}}{\rho }}+{\displaystyle \frac{^{\mathrm{ex}}}{\mathit{ϵ}}}{\displaystyle \frac{\delta \mathit{ϵ}}{\delta [\rho p(\mathit{ϵ})]}}+{\displaystyle \frac{^{\mathrm{ex}}}{\mathit{ϵ}\mathit{ϵ}}}{\displaystyle \frac{\delta \mathit{ϵ}\mathit{ϵ}}{\delta [\rho p(\mathit{ϵ})]}}+\mathrm{}`$ where $`\delta \mathit{ϵ}^m/\delta [\rho p(\mathit{ϵ})]=(\mathit{ϵ}^m\mathit{ϵ}^m)/\rho `$. In this notation, $`^{\mathrm{ex}}/\mathit{ϵ}`$ is a vector whose components are the derivatives of $`^{\mathrm{ex}}`$ with respect to each component of $`\mathit{ϵ}`$. Similarly, $`^{\mathrm{ex}}/\mathit{ϵ}\mathit{ϵ}`$ is a tensor. Thus we obtain a general expression, truncated at second order in $`\mathit{ϵ}`$, for the excess chemical potential $`\mu ^{\mathrm{ex}}(\mathit{ϵ})`$ of the species with property $`\mathit{ϵ}`$, in a phase characterised by $`\{\rho ,p(\mathit{ϵ})\}`$, $`\mu ^{\mathrm{ex}}(\mathit{ϵ})`$ $`=`$ $`\mu _\mathrm{m}^{\mathrm{ex}}(\rho )+\mathit{ϵ}\left(𝑨^{}(\rho )𝑨/\rho \right)+\mathit{ϵ}\mathit{ϵ}:\left(𝑩^{}(\rho )𝑩/\rho \right)+\mathit{ϵ}\mathit{ϵ}:\left(𝑪^{}(\rho )2𝑪/\rho \right)`$ (8) $`+\mathit{ϵ}\left(𝑨+2\mathit{ϵ}𝑪\right)/\rho +\mathit{ϵ}\mathit{ϵ}:𝑩/\rho +Ø(ϵ^3)`$ where $`\mu _\mathrm{m}^{\mathrm{ex}}(\rho )`$ is the excess chemical potential, $`d_\mathrm{m}^{\mathrm{ex}}/d\rho `$ of a monodisperse reference phase at density $`\rho `$. Equation (8) will be central to the calculations of phase equilibria in section III. ### C Expansion of the pressure The pressure of a polydisperse phase with density $`\rho `$ and distribution $`p(\mathit{ϵ})`$ can be found from $`P=+{\displaystyle \mu (\mathit{ϵ})\rho p(\mathit{ϵ})\text{d}^\nu ϵ}`$ which is the continuum analogue of the standard thermodynamic relation for the pressure of a mixture. Using Eqs. (1) and (4) it follows that $$P=P^{\mathrm{id}}^{\mathrm{ex}}+\rho \mu ^{\mathrm{ex}}(\mathit{ϵ})p(\mathit{ϵ})\text{d}^\nu ϵ$$ (9) where $`P^{\mathrm{id}}`$ is the ideal pressure of the polydisperse mixture. It is easily confirmed from Eqs. (2) and (5) that this is the usual ideal gas pressure $`P^{\mathrm{id}}=\rho `$. Substituting the series expressions for the excess free energy and chemical potential (Eqs. (6) and (8)) into Eq. (9), and using the pressure of the monodisperse reference system at density $`\rho `$, $`P_\mathrm{m}(\rho )=\rho _\mathrm{m}^{\mathrm{ex}}(\rho )+\rho \mu _\mathrm{m}^{\mathrm{ex}}(\rho )`$, yields the expansion for the pressure of the polydisperse phase, $$P=P_\mathrm{m}(\rho )+\mathit{ϵ}\left(\rho 𝑨^{}(\rho )𝑨\right)+\mathit{ϵ}\mathit{ϵ}:\left(\rho 𝑩^{}(\rho )𝑩\right)+\mathit{ϵ}\mathit{ϵ}:\left(\rho 𝑪^{}(\rho )𝑪\right)+Ø(ϵ^3).$$ (10) ## III Phase equilibria ### A Fractionation At coexistence between two or more phases, the distribution $`p(\mathit{ϵ})`$ of properties is typically different for each phase. We shall assume that the overall distribution of all particles in all phases is known (since the particles might have been synthesised en masse, before phase separation was induced). This will be called the parent distribution $`p_\mathrm{p}(\mathit{ϵ})`$. Without loss of generality, the origin of $`\mathit{ϵ}`$ will henceforth be set at the mean of this parent distribution, so that $`\mathit{ϵ}_\mathrm{p}0`$. The aim of this section is to calculate the distribution $`p_\alpha (\mathit{ϵ})`$ in any given phase $`\alpha `$, of $`M`$ coexisting phases. The number of particles in phase $`\alpha `$ is some fraction $`n_\alpha `$ of the total number in the system. So conservation of material is expressed as $`_\beta ^Mn_\beta =1`$. In fact, a stronger criterion than this holds. The number of particles belonging to each species is conserved. That is, $$\underset{\beta }{\overset{M}{}}n_\beta p_\beta (\mathit{ϵ})=p_\mathrm{p}(\mathit{ϵ})$$ (11) for all $`\mathit{ϵ}`$. This is a convenient point at which to introduce a special notation which will be useful later. If $`Q`$ is any property of a phase then, for that phase, let $`\widehat{\mathrm{\Delta }}[Q]`$ be its deviation from the mean of $`Q`$ over all phases, weighted by the number of particles in each. Hence, in phase $`\alpha `$, $`\widehat{\mathrm{\Delta }}_\alpha [Q]Q_\alpha {\displaystyle \underset{\beta }{\overset{M}{}}}n_\beta Q_\beta ={\displaystyle \underset{\beta }{\overset{M}{}}}n_\beta [Q]_\beta ^\alpha .`$ The notation $`[\mathrm{}]_\beta ^\alpha `$ indicates the difference between the quantity evaluated in phase $`\alpha `$ and in phase $`\beta `$. Returning to the problem of coexistence, note that, in any two of the $`M`$ coexisting phases, $`\alpha `$ and $`\beta `$ say, the chemical potential is equal, i.e$`\mu _\alpha (\mathit{ϵ})=\mu _\beta (\mathit{ϵ})`$ for all $`\mathit{ϵ}`$. Combining Eqs. (4) and (5), this implies $$\rho _\alpha p_\alpha (\mathit{ϵ})\mathrm{exp}\mu _\alpha ^{\mathrm{ex}}(\mathit{ϵ})=\rho _\beta p_\beta (\mathit{ϵ})\mathrm{exp}\mu _\beta ^{\mathrm{ex}}(\mathit{ϵ})$$ (12) which can be used to eliminate each $`p_\beta (\mathit{ϵ})`$ in Eq. (11), yielding $$p_\alpha (\mathit{ϵ})=\frac{p_\mathrm{p}(\mathit{ϵ})}{_\beta n_\beta \frac{\rho _\alpha }{\rho _\beta }\mathrm{exp}\left[\mu _\alpha ^{\mathrm{ex}}(\mathit{ϵ})\mu _\beta ^{\mathrm{ex}}(\mathit{ϵ})\right]}.$$ (13) We shall now use Eq. (8) for the excess chemical potential, truncated at first order, $`\mu ^{\mathrm{ex}}(\mathit{ϵ})=\mu _\mathrm{m}^{\mathrm{ex}}(\rho )+\mathit{ϵ}(𝑨^{}𝑨/\rho )+\mathit{ϵ}𝑨/\rho +Ø(ϵ^2)`$ and expand the exponential in Eq. (13) as $`\mathrm{exp}[\mu ^{\mathrm{ex}}(\mathit{ϵ})]_\beta ^\alpha =(1+\mathit{ϵ}[𝑨/\rho ]_\beta ^\alpha )\mathrm{exp}[\mu _\mathrm{m}^{\mathrm{ex}}+\mathit{ϵ}(𝑨^{}𝑨/\rho )]_\beta ^\alpha +Ø(ϵ^2).`$ So Eq. (13) can be written $$p_\alpha (\mathit{ϵ})=\frac{p_\mathrm{p}(\mathit{ϵ})}{𝒮_\alpha }\left(1\mathit{ϵ}\frac{_\beta [𝑨/\rho ]_\beta ^\alpha n_\beta \frac{\rho _\alpha }{\rho _\beta }\mathrm{exp}[\mu _\mathrm{m}^{\mathrm{ex}}+\mathit{ϵ}(𝑨^{}𝑨/\rho )]_\beta ^\alpha }{𝒮_\alpha }+Ø(ϵ^2)\right)$$ (14) where $`𝒮_\alpha _\beta n_\beta \frac{\rho _\alpha }{\rho _\beta }\mathrm{exp}[\mu _\mathrm{m}^{\mathrm{ex}}+\mathit{ϵ}(𝑨^{}𝑨/\rho )]_\beta ^\alpha `$. If Eq. (14) is integrated with respect to $`\mathit{ϵ}`$, then the term linear in $`\mathit{ϵ}`$ disappears, due to the condition $`\mathit{ϵ}_\mathrm{p}\mathrm{𝟎}`$. We are left with $`1=1/𝒮_\alpha +Ø(ϵ^2)`$, or $`𝒮_\alpha =1+Ø(ϵ^2)`$, which can be substituted back into Eq. (14). Finally, integrating Eq. (12) with $`\mu ^{\mathrm{ex}}(\mathit{ϵ})`$ expanded to zeroth order yields $`(\rho _\alpha /\rho _\beta )\mathrm{exp}\left[\mu _\mathrm{m}^{\mathrm{ex}}\right]_\beta ^\alpha =1+Ø(ϵ)`$, which we also substitute into Eq. (14). The result is $$p_\alpha (\mathit{ϵ})=p_\mathrm{p}(\mathit{ϵ})\left\{1\mathit{ϵ}\widehat{\mathrm{\Delta }}_\alpha [𝑨(\rho )/\rho ]+Ø(ϵ^2)\right\}.$$ (15) This is the result we wanted: an expression for the distribution of species in each phase. Not surprisingly, it is almost the same as the parent distribution (which remains unapproximated), since all species have similar properties ($`\mathit{ϵ}`$ is small for each). The lowest-order correction to $`p_\mathrm{p}(\mathit{ϵ})`$ is expressed in terms of $`𝑨`$, a (vectorial) parameter of the excess free energy (Eq. (6)), and $`\rho _\beta `$, the density of each phase. In principle, the quantities on the right hand side of Eq. (15) are known. The densities of the phases $`\rho _\alpha `$ will differ a little (by an amount $`\delta _\alpha `$, say) from the binodals of the monodisperse reference system $`\rho _\alpha ^\mathrm{m}`$, so that $`\rho _\alpha =\rho _\alpha ^\mathrm{m}+\delta _\alpha `$, and the numbers of particles $`n_\alpha `$ will have similar small changes. However, in Eq. (15) one may substitute either the monodisperse or the polydisperse values of $`\rho `$ and $`n`$, whichever are available from theory or experiment, since the small differences affect only higher-order terms. A more elegant expression is obtained if we find the difference between the normalised distributions $`p(\mathit{ϵ})`$ of two of the $`M`$ coexisting phases. From Eq. (15), $$[p(\mathit{ϵ})]_\beta ^\alpha p_\mathrm{p}(\mathit{ϵ})\mathit{ϵ}\left[𝑨/\rho \right]_\beta ^\alpha $$ (16) which, together with Eq. (11), contains the same information as Eq. (15). The arrow($``$) indicates the strict asymptotic limit as $`ϵ^2_\mathrm{p}0`$. The parameter $`𝑨/\rho `$ in Eq. (16), is the lowest-order coefficient of $`\mathit{ϵ}`$ in Eq. (8), so we can write $$[p(\mathit{ϵ})]_\beta ^\alpha p_\mathrm{p}(\mathit{ϵ})\mathit{ϵ}\left[\mathbf{}_\mathit{ϵ}\mu ^{\mathrm{ex}}(\mathit{ϵ})\|_{\mathit{ϵ}=\mathrm{𝟎}}\right]_\beta ^\alpha .$$ (17) A gradient in $`\mathit{ϵ}`$-space is denoted $`\mathbf{}_\mathit{ϵ}`$, and it is evaluated for phases $`\alpha `$, $`\beta `$ in Eq. (17) at the mean species $`\mathit{ϵ}=\mathrm{𝟎}`$. Equation (17) says that, at coexistence, the distributions separate along the steepest gradient of $`\mu _\alpha ^{\mathrm{ex}}\mu _\beta ^{\mathrm{ex}}`$ in $`\mathit{ϵ}`$-space. It is informative to take moments of Eq. (16), since these will tell us how the various properties of the particles are correlated within the phases. One might wish to examine averages of the individual components $`ϵ_i`$ of $`\mathit{ϵ}`$, such as $`ϵ_1`$ the mean value of some property such as the size of particles, or $`ϵ_1ϵ_2`$ the cross-correlation between e.g. sizes and charges of the particles. The most general moment, from which this information can be extracted, is the mean of an outer product of $`m`$ vectors $`\mathit{ϵ}`$. That is the $`m`$th rank tensor $`\mathit{ϵ}^m`$ $``$ $`{\displaystyle \underset{}{\mathit{ϵ}\mathit{ϵ}\mathrm{}\mathit{ϵ}}p(\mathit{ϵ})\text{d}^\nu ϵ}`$ (19) $`m`$ whose components are all possible correlators of order $`m`$. From Eq. (16), the difference of such a moment between two of the coexisting phases is $$\left[\mathit{ϵ}^m\right]_\beta ^\alpha =\mathit{ϵ}^{m+1}_\mathrm{p}\left[𝑨/\rho \right]_\beta ^\alpha +Ø(\mathit{ϵ}^{m+2})$$ (20) which involves the next moment of the parent distribution, and one summation (the scalar product) over the $`\nu `$ polydisperse properties. The scalar version of Eq. (20) was derived previously for substances that are polydisperse in a single property, at two-phase and multi-phase coexistence. It has been noted that for some distributions, odd moments $`\mathit{ϵ}^{2k+1}_\mathrm{p}`$ are not of order $`(2k+1)`$ in the width of the distribution but, due to cancellation in the integral (Eq. (19)), are of higher order, so that the RHS of Eq. (20) is zero to $`Ø(\mathit{ϵ}^{m+1})`$ and the unevaluated $`Ø(\mathit{ϵ}^{m+2})`$ terms are the dominant order<sup></sup><sup></sup>When $`m=2`$, Eq. (20) relates the second moment of the daughters to the skew of the parent . For highly symmetric distributions the skew is small (or vanishing), so the LHS of Eq. (20) is small, though not vanishingly so, as it is limited by the unevaluated $`Ø(\mathit{ϵ}^{m+2})`$ terms. On the other hand, for a strongly asymmetric parent, by Eq. (20), the daughters have very different second moments. Thus the skew of the parent is highly influential, contrary to the impression given in Ref. .. Specifically, this is the case for distributions which become symmetric in the narrow limit (e.g. Gaussian, Schultz). For strongly asymmetric distributions though, excepting unlikely cancellations, the $`m`$th moment is of $`m`$th order in the width. Equations (16) and (17) remain valid for all narrow distributions. As the notation of Eq. (20) is rather abstract, let us consider a specific example: a system of particles with a range of radii $`r`$ and charges $`q`$. At equilibrium the particles are partitioned into two coexisting phases denoted $`\alpha `$ and $`\beta `$, and we are interested in the difference between the average radii in these phases. The overall average radius throughout the system is called $`r_\mathrm{p}`$ and the average charge $`q_\mathrm{p}`$. Applying Eq. (20) with $`m=1`$, $`ϵ_1=r/r_\mathrm{p}1`$ and $`ϵ_2=q/q_\mathrm{p}1`$, and taking the first component gives the answer $`r_\alpha r_\beta =a\left[r^2_\mathrm{p}r_\mathrm{p}^2\right]+b\left[rq_\mathrm{p}r_\mathrm{p}q_\mathrm{p}\right]`$ with coefficients $`a=[A_1/\rho ]_\beta ^\alpha /r_\mathrm{p}`$ and $`b=[A_2/\rho ]_\beta ^\alpha /q_\mathrm{p}`$. The first term says that, in the absence of charges, the amount of size fractionation is proportional to the overall size variance (the square of the standard deviation), with a coefficient $`a`$ indicating the difference between each phase’s affinity for large particles. With charge polydispersity, the second term indicates that size fractionation can occur even if $`a=0`$ so that neither phase favours larger particles. If size and charge are cross-correlated in the parent, i.e. bigger particles tend to have bigger (or smaller) charges, then a ‘convective’ fractionation of sizes will be driven by $`b`$, the ‘driving force’ separating charges. These two contributions are additive for a system with a narrow distribution. ### B Shift of the phase boundaries So far we have established the lowest order change (with width of the distribution) of the distribution of properties in each phase, relative to the original, ‘parent’ distribution (Eq. (15)). However, a knowledge of the normalised distributions is not a complete characterisation of the coexisting phases. We would also like to know the overall densities of the phases, $`\rho _\alpha `$. For a monodisperse system, the binodal densities are independent of the number of particles in each phase. This is not so when the system is polydisperse, since the size of one phase’s share of the available particles influences the shape of its distribution of species and, hence, its thermodynamic properties. One coexistence of particular interest is the ‘cloud point’. A phase is at its cloud point when its density lies on the ‘cloud curve’ in the phase diagram. It then coexists with an infinitesimal amount of another phase, whose density is on the corresponding ‘shadow’ curve. The cloud point is important from a practical point of view, since it marks the threshold at which a separation of phases is first observable, and the resulting morphology and fractionation can have useful applications. The cloud point also inspires theoretical interest since it is relatively easy to treat in a two-phase system, where the distribution of species in the majority phase is simply equal to the parent. In the present analysis, as in the previous section, we shall calculate the conditions of general phase equilibria, not just the cloud and shadow curves. Having found the general solution, we shall return to the special case of the cloud point. To establish the phase equilibria, we shall demand that the pressure and chemical potentials are equal in all phases, using the expressions for pressure and chemical potential derived in section II. Before proceeding further, let us take moments of Eq. (15), as this will lead to some substantial simplifications. We see that the second moment, $`\mathit{ϵ}\mathit{ϵ}`$, in any given phase is, to lowest order, equal to that in the parent. However, the mean, $`\mathit{ϵ}`$, in that phase is also of second order in the parent’s width, since the mean of the parent vanishes by definition. In orders of the width of the parent, $`\sigma _\mathrm{p}\sqrt{\mathit{ϵ}\mathit{ϵ}_\mathrm{p}}`$, we have $`\mathit{ϵ}`$ $`=`$ $`\mathit{ϵ}\mathit{ϵ}_\mathrm{p}\widehat{\mathrm{\Delta }}\left[𝑨/\rho \right]+Ø(\sigma _\mathrm{p}^3)`$ (22) $`\text{and}\mathit{ϵ}\mathit{ϵ}`$ $`=`$ $`\mathit{ϵ}\mathit{ϵ}_\mathrm{p}+Ø(\sigma _\mathrm{p}^3).`$ (23) So the terms involving $`\mathit{ϵ}\mathit{ϵ}`$ in Eqs. (8) and (10) can be dropped, since they are of order $`\sigma _\mathrm{p}^4`$. Let us clarify this point. The excess chemical potentials in Eq. (8) and the pressure in Eq. (10) were calculated consistently to second order in the deviations $`\mathit{ϵ}`$ of the properties of a phase’s constituent particles from some reference. In general, then, to evaluate these thermodynamic potentials for an arbitrary phase, with respect to an arbitrary reference, terms involving $`\mathit{ϵ}\mathit{ϵ}`$ must be preserved. However, the choice $`\mathit{ϵ}_\mathrm{p}=\mathrm{𝟎}`$, to fix the arbitrary reference, allows us to drop the terms in Eqs. (8) and 10 with coefficient $`𝑪`$, for coexisting phases. This is because we have found that when several phases coexist, their mean properties differ very little from each other, so that $`\mathit{ϵ}_\alpha \mathit{ϵ}_\mathrm{p}`$ is second-order small. To calculate the densities of the coexisting phases, we first demand that the pressure $`P(\rho ,\mathit{ϵ},\mathit{ϵ}\mathit{ϵ})`$ is equal in any two phases $`\alpha `$, $`\beta `$ of the $`M`$ coexisting phases. We shall write this as $$\left[P(\rho ,\mathit{ϵ},\mathit{ϵ}\mathit{ϵ})\right]_\beta ^\alpha =0$$ (24) and use the expression for $`P(\rho ,\mathit{ϵ},\mathit{ϵ}\mathit{ϵ})`$ in Eq. (10). Equation (24) is a condition on $`\rho _\alpha `$, the binodal densities which we want to find. As mentioned above, in the limit of a narrow parent, these densities are close to the monodisperse binodals $`\rho _\alpha ^\mathrm{m}`$, differing by an amount $`\delta _\alpha `$, so that $`\rho _\alpha =\rho _\alpha ^\mathrm{m}+\delta _\alpha `$. We shall determine the small shift $`\delta _\alpha `$. In Eq. (10) we can Taylor-expand the monodisperse pressure $`P_\mathrm{m}(\rho _\alpha )=P_\mathrm{m}(\rho _\alpha ^\mathrm{m})+P_\mathrm{m}^{}(\rho _\alpha ^\mathrm{m})\delta _\alpha `$. It will transpire that $`\delta _\alpha `$ is of order $`\sigma _\mathrm{p}^2`$, so that it is sufficient to truncate this Taylor expansion at first order in $`\delta _\alpha `$. The condition of pressure balance (Eq. (24)) becomes $`[P_\mathrm{m}(\rho _\mathrm{m})+P_\mathrm{m}^{}(\rho _\mathrm{m})\delta +\mathit{ϵ}(\rho 𝑨^{}𝑨)+\mathit{ϵ}\mathit{ϵ}:(\rho 𝑩^{}𝑩)]_\beta ^\alpha =0.`$ We see that, due to the low order of expansion, the change in density $`\delta `$ affects only $`P_\mathrm{m}`$, the ‘monodisperse contribution’ to the pressure. Effectively, each of the polydisperse phases can be replaced by a monodisperse phase that is subject to an ‘external field’ ($`\mathit{ϵ}(\rho 𝑨^{}𝑨)+\mathit{ϵ}\mathit{ϵ}:(\rho 𝑩^{}𝑩)`$). In this interpretation, each (fictitious) monodisperse phase responds to the field with a density change $`\delta `$, governed by the response function $`P_\mathrm{m}^{}`$. The same interpretation can be applied to the chemical potential balance below. We can eliminate the first term because the monodisperse reference phases at densities $`\rho _\alpha ^\mathrm{m}`$ are in pressure balance, that is $`[P_\mathrm{m}(\rho _\mathrm{m})]_\beta ^\alpha =0`$. The density derivative of the pressure in the monodisperse system, $`P_\mathrm{m}^{}(\rho _\mathrm{m})`$, can be expressed in terms of the phase’s isothermal compressibility, $$\kappa \frac{1}{V}\left(\frac{V}{P}\right)_{N,T}=1/\rho \frac{dP}{d\rho }=1/\rho ^2\frac{d\mu }{d\rho }.$$ (25) Hence pressure balance dictates $$[\frac{\delta }{\kappa \rho }+\mathit{ϵ}(𝑨^{}\rho 𝑨)+\mathit{ϵ}\mathit{ϵ}:(𝑩^{}\rho 𝑩)]_\beta ^\alpha =0$$ (26) for any pair $`(\alpha ,\beta )`$ of the $`M`$ phases. Another constraint on the shifts in density $`\delta _{\alpha ,\beta }`$ can be derived from the equality of all chemical potentials between any pair of phases, which leads to Eq. (12) for the excess parts. We substitute from Eq. (8) for $`\mu ^{\mathrm{ex}}(\mathit{ϵ})`$, and further substitute $`\mu _\mathrm{m}^{\mathrm{ex}}(\rho )`$ $`=`$ $`\mu _\mathrm{m}(\rho )\mathrm{ln}\rho `$ $`=`$ $`\mathrm{ln}\rho +\mu _\mathrm{m}(\rho _\mathrm{m})+\mu _\mathrm{m}^{}(\rho _\mathrm{m})\delta `$ to yield $`\left[p(\mathit{ϵ})\mathrm{exp}\left\{\mu _\mathrm{m}(\rho _\mathrm{m})+\mu _\mathrm{m}^{}\delta +\mathit{ϵ}(𝑨^{}𝑨/\rho )+\mathit{ϵ}\mathit{ϵ}:(𝑩^{}𝑩/\rho )+\mathit{ϵ}𝑨/\rho +\mathit{ϵ}\mathit{ϵ}:𝑩/\rho \right\}\right]_\beta ^\alpha =0`$ Expanding the exponential of small quantities and integrating gives $`[e^{\mu _\mathrm{m}(\rho _\mathrm{m})}(1+\mu _\mathrm{m}^{}\delta +\mathit{ϵ}𝑨^{}+\mathit{ϵ}\mathit{ϵ}:\{𝑩^{}+𝑨𝑨/2\rho ^2\})]_\beta ^\alpha =0.`$ Now, applying Eq. (25) and the condition $`\left[\mu _\mathrm{m}(\rho _\mathrm{m})\right]_\beta ^\alpha =0`$ yields $$[\frac{\delta }{\kappa \rho ^2}+\mathit{ϵ}𝑨^{}+\mathit{ϵ}\mathit{ϵ}:(𝑩^{}+𝑨𝑨/2\rho ^2)]_\beta ^\alpha =0$$ (27) which, with Eq. (26), forms a pair of simultaneous equations in two unknowns, $`\delta _\alpha `$ and $`\delta _\beta `$. As an aside, it is interesting that, as well as fixing the $`M`$ unknowns $`\delta _\alpha `$, the $`2(M1)`$ constraints in Eqs. (26) and (27) impose extra conditions on the reference system for $`M>2`$. The extra constraints are analogous to the Gibbs phase rule: a single component, whether monodisperse or near-monodisperse, can exist as a pair of coexisting phases under a range of conditions, whereas triple or higher coexistence requires a special choice of external fields . It is intriguing that the extra conditions imposed on the monodisperse reference system constrain the values of $`𝑨(\rho )`$ and $`𝑩(\rho )`$ (defined in Eq. (6)), which, though defined in the monodisperse system, are not manifestly relevant to its phase equilibria. It is straightforward to solve the linear equations (26 and 27) to yield the value of $`\delta _\alpha `$, expressed as $$x_\alpha =\frac{[y]_\beta ^\alpha }{[\rho ^1]_\beta ^\alpha }$$ (28) in terms of $`x\delta /\kappa \rho +\mathit{ϵ}(𝑨^{}\rho 𝑨)+\mathit{ϵ}\mathit{ϵ}:(𝑩^{}\rho 𝑩)`$ and of $`y\mathit{ϵ}𝑨/\rho +\mathit{ϵ}\mathit{ϵ}:(𝑩/\rho +𝑨𝑨/2\rho ^2),`$ where expressions for $`\mathit{ϵ}`$, $`\mathit{ϵ}\mathit{ϵ}`$ can be substituted from Eqs. (III B). However, the expression is inelegant, since there are $`M`$ phases, of which $`\beta `$ is in no way special. The shift $`\delta _\alpha `$ in the $`\alpha `$-binodal should have a symmetric dependence on all other phases. Since $`\beta `$ is arbitrary, we may symmetrize Eq. (28) by averaging it over all $`\beta `$, with respect to any weight factor $`z_{\alpha \beta }`$: $`x_\alpha ={\displaystyle \frac{_\beta z_{\alpha \beta }[y]_\beta ^\alpha /[\rho ^1]_\beta ^\alpha }{_\beta z_{\alpha \beta }}}`$ For neatness, we choose $`z_{\alpha \beta }n_\beta (\rho _\alpha ^1\rho _\beta ^1)`$. The final answer for the shift $`\delta `$ of the density of oneThe subscript $`\alpha `$, denoting the phase in question, is now dropped since it applies to every quantity in Eq. (29) (except $`\mathit{ϵ}\mathit{ϵ}_\mathrm{p}`$ and operands of $`\widehat{\mathrm{\Delta }}`$). of the $`M`$ coexisting phases due to polydispersity is $$\delta =\rho \kappa \mathit{ϵ}\mathit{ϵ}_\mathrm{p}:\left((𝑨^{}\rho 𝑨)\widehat{\mathrm{\Delta }}[𝑨/\rho ]𝑩^{}\rho +𝑩+\frac{\widehat{\mathrm{\Delta }}\left[\widehat{\mathrm{\Delta }}[𝑨/\rho ]\widehat{\mathrm{\Delta }}[𝑨/\rho ]\right]2\widehat{\mathrm{\Delta }}\left[𝑩/\rho \right]}{2\widehat{\mathrm{\Delta }}[1/\rho ]}\right)$$ (29) plus terms of order $`\sigma _\mathrm{p}^3`$. As this is a lowest-order expression, the densities $`\rho `$ appearing on the R.H.S. may be taken as either the monodisperse or the polydisperse binodal values, whichever is most convenient. As yet, we do not know the values of the parameters $`𝑨`$, $`𝑩`$, $`𝑨^{}`$ and $`𝑩^{}`$ in Eq. (29). They will be discussed in Section IV. Nevertheless, the structure of Eq. (29) is informative. As stated above, $`\delta `$ is indeed of order $`\sigma _\mathrm{p}^2`$. If, as is often the case , a phase diagram is drawn for some system, showing the densities of coexisting phases against the width $`\sigma _\mathrm{p}`$ of some parent distribution, then Eq. (29) gives the leading-order Taylor expansion of the shape of the phase boundariesThis expansion of the equations of the phase boundaries is in some ways equivalent (though in a quite different formalism) to the Clausius-Clapeyron-like equation derived by Bolhuis and Kofke for the $`P\nu `$ phase diagram, where $`\nu `$ parameterises the variance of a Gaussian distribution of activities of polydisperse hard spheres. The present analysis, though restricted to narrow distributions, is more general.. By definition, they meet the ($`\sigma _\mathrm{p}=0`$)-axis at the monodisperse values $`\rho _\mathrm{m}`$, and we have shown that they do so perpendicularly, since there is no term linear in $`\sigma _\mathrm{p}`$. According to Eq. (29), the curvature of each phase boundary is proportional to the isothermal compressibility $`\kappa `$ of the phase. Hence the small-variable expansion breaks down at a critical point, where the compressibility diverges. This is not unexpected, since the position of the critical point can itself be shifted by the presence of polydispersity, for instance to a different temperature. Criticality in the monodisperse system is a coexistence between two (or more) phases of equal densities. Under the same conditions, slightly-polydisperse phases, being non-critical, will either coexist at finitely-different densities, or not demix at all. In either case, the thermodynamic state is far from the monodisperse reference, so the divergence of Eq. (29) at criticality is correct. Equation (29) exhibits a property typical of polydisperse systems. For a given parent distribution, the positions of the phase boundaries depend on the number of particles in each phase (since this affects the averages in the definition of $`\widehat{\mathrm{\Delta }}[\mathrm{}]`$). This contrasts with monodisperse systems, for which tie-lines can be drawn, connecting a set of points in the phase diagram’s forbidden region to the same coexisting end-points. It is interesting that this signature of polydispersity is evident at the lowest order in $`\sigma _\mathrm{p}`$. Hence, as soon as any shift in the phase boundaries due to polydispersity is observable in a system, that shift should vary as a tie-line is traversed. As mentioned above, a coexisting state of particular interest is that between an infinitesimal amount of one phase (the ‘shadow’ phase) and a majority phase defined to be at its ‘cloud’ point. Such a state defines the extreme boundary of the coexistence region. The movement of the cloud point<sup>§</sup><sup>§</sup>§Greek subscripts were used to denote generic phases, whereas $`c`$, $`s`$ indicate cloud and shadow phases, for which the relative numbers of constituent particles are specified., $`\delta _c`$ is of great importance, since this tells an experimenter where, in the phase diagram, to first expect phase separation to occur. Let the portion of particles in the shadow phase be the infinitesimally small number $`n_s=s`$. Then the fraction belonging to the majority, cloud phase is $`n_c=1s`$. For any property $`Q`$, $`\widehat{\mathrm{\Delta }}[Q]`$ vanishes in the majority phase as $`s0`$, since the cloud phase defines the number-weighted mean. In the cloud phase, $`\widehat{\mathrm{\Delta }}_c[Q]=Q_c(n_cQ_c+n_sQ_s)=s[Q]_s^c`$ where $`[Q]_s^c`$ denotes the difference between the values of $`Q`$ in the cloud and shadow phases, $`Q_cQ_s`$. In the shadow phase, $`\widehat{\mathrm{\Delta }}_s[Q]=(1s)[Q]_c^s.`$ Substituting these expressions into Eq. (29) gives the shift of the cloud point as $$\delta _c=\rho _c\kappa _c\mathit{ϵ}\mathit{ϵ}_\mathrm{p}:\left(𝑩_c^{}\rho _c𝑩_c+\frac{[𝑨/\rho ]_s^c[𝑨/\rho ]_s^c+2[𝑩/\rho ]_s^c}{2[1/\rho ]_s^c}\right).$$ (30) The coexisting shadow phase has a density $`\rho _s^\mathrm{m}+\delta _s`$ where $$\delta _s=\rho _s\kappa _s\mathit{ϵ}\mathit{ϵ}_\mathrm{p}:\left(𝑩_s^{}\rho _s𝑩_s+\frac{[𝑨/\rho ]_s^c[𝑨/\rho ]_s^c+2[𝑩/\rho ]_s^c}{2[1/\rho ]_s^c}+\left(𝑨_s^{}\rho _s𝑨_s\right)[𝑨/\rho ]_s^c\right).$$ (31) It is not always appropriate to express the position of a phase boundary in terms of the number density of particles, $`\rho =\rho ^\mathrm{m}+\delta `$. For instance, the phase diagram is often represented in terms of the volume fraction of particles, $`\varphi `$. This is the product of the number density and the volume of a particle. Since fractionation occurs, so that the mean volume of a particle differs from phase to phase, $`\varphi `$ and $`\rho `$ are not equivalent, and the coexistence region has a different shape in the two representations of the phase diagram. Quantities such as $`\varphi `$ are easily derived from the above results as follows. For spherical particles, for instance, using the scalar $`ϵ`$ to denote fractional deviations in particle radii: $`r_i=(1+ϵ_i)r_\mathrm{m}`$, we have $`\varphi `$ $`=`$ $`{\displaystyle \frac{4}{3}}\pi r^3\rho ={\displaystyle \frac{4}{3}}\pi r_\mathrm{m}^3(1+ϵ)^3(\rho ^\mathrm{m}+\delta )`$ (32) $`{\displaystyle \frac{\varphi }{\varphi _\mathrm{m}}}`$ $`=`$ $`1+{\displaystyle \frac{\delta }{\rho _\mathrm{m}}}+3ϵ+3ϵ^2+Ø(\sigma _\mathrm{p}^3)`$ (33) with the quantities on the RHS given in Eqs. (III B) and (29). In describing the phase diagram, one final quantity of interest is the width of the coexistence region, sometimes called the miscibility gap. In terms of $`\rho `$, this is the range of densities of the system as a whole, for which it separates into more than one phase. Say the miscibility gap is bounded by two phases, called $`\alpha `$ and $`\beta `$ in the reference system. Quenching to one edge of this gap, the cloud point of phase $`\alpha `$, will cause an infinitesimal amount of phase $`\beta `$ to form. At the other extreme, a vanishing amount of phase $`\alpha `$ coexists with majority phase $`\beta `$. So the miscibility gap is given by application of Eq. (30) alone, to establish the two cloud points of phases $`\alpha `$ and $`\beta `$. The miscibility gap $`[\rho _c]_\beta ^\alpha `$ is then given, in terms of the gap in the monodisperse reference system, by $$[\rho _c]_\beta ^\alpha =[\rho _c^\mathrm{m}]_\beta ^\alpha +\frac{\mathit{ϵ}\mathit{ϵ}_\mathrm{p}}{[1/\rho ]_\beta ^\alpha }:\left\{[\rho \kappa ]_\beta ^\alpha [𝑩/\rho ]_\beta ^\alpha +[1/\rho ]_\beta ^\alpha [\rho \kappa (𝑩^{}\rho 𝑩)]_\beta ^\alpha +[𝑨/\rho ]_\beta ^\alpha [𝑨/\rho ]_\beta ^\alpha (\rho _\alpha \kappa _\alpha +\rho _\beta \kappa _\beta )/2\right\}.$$ (34) Note that, with the gap defined to be positive, $`[1/\rho ]_\beta ^\alpha `$ is positive. In principal, the lowest-order change in the miscibility gap due to polydispersity, given in Eq. (34), may be positive or negative, depending on the system. We shall return to this topic in section V. ## IV Evaluating the parameters We have determined the state of thermodynamic equilibrium for polydisperse phases, in terms of certain properties of the monodisperse reference phases. Some of these properties, the density $`\rho `$ and compressibility $`\kappa `$, are well established for most common systems. Others, the parameters $`𝑨(\rho )`$, $`𝑨^{}(\rho )`$, $`𝑩(\rho )`$ and $`𝑩^{}(\rho )`$, require further discussion. These parameters appear in the excess free energy, Eq. (6), and we require their values in the monodisperse system (where $`\mathit{ϵ}=\mathit{ϵ}\mathit{ϵ}=0`$) at some density $`\rho `$. One might already have an expression to hand for the free energy of a particular system, in the form of Eq. (6), in which case it is trivial to identify the appropriate parameters and substitute them into the various expressions above. The results for two commonly-studied systems are given in sections V and VI. Alternatively, the required numbers may be established by experiment, or may require derivation from first principles. We shall now find some generic expressions for $`𝑨(\rho )`$, $`𝑩(\rho )`$ and, for completeness, also $`𝑪(\rho )`$, which will inform our physical interpretation of these quantities. Then, in section IV A, we shall apply those expressions, to construct a method for first-principles calculation of the parameters from thermodynamic perturbation theory. That theory is applicable whenever the Hamiltonian of the system can be differentiated with respect to the properties $`\mathit{ϵ}`$ of the particles. Though many systems meet this criterion, one system of particular theoretical and pedagogical interest, a set of hard spheres, does not. Its Hamiltonian is discontinuous, being zero for all physical configurations and rising discontinuously to infinity if any spheres overlap. However, hard spheres belong to a special class of systems (scalable systems) for which the parameter $`𝑨`$, which controls the emergence of fractionation (see Eqs. (15, 16, 20)) is calculable by an alternative method, detailed in section IV B. To find general expressions for the quantities $`𝑨`$, $`𝑩`$ and $`𝑪`$, let us Taylor expand a generic excess free energy with respect to the properties $`\mathit{ϵ}_i`$ of each constituent particle separately, where $`i`$ labels the $`N`$ particles. We expand to second order in these $`N`$ vectorial variables thus $`F^{\mathrm{ex}}=F_\mathrm{m}^{\mathrm{ex}}+{\displaystyle \underset{i=1}{\overset{N}{}}}\mathit{ϵ}_i{\displaystyle \frac{F^{\mathrm{ex}}}{\mathit{ϵ}_i}}\|_\mathrm{m}+{\displaystyle \frac{1}{2}}{\displaystyle \underset{i,j=1}{\overset{N}{}}}\mathit{ϵ}_i\mathit{ϵ}_j:{\displaystyle \frac{^2F^{\mathrm{ex}}}{\mathit{ϵ}_i\mathit{ϵ}_j}}\|_\mathrm{m}+Ø(ϵ^3)`$ where the subscript m indicates evaluation in the system where $`\mathit{ϵ}_i=\mathrm{𝟎}`$ for all $`i`$. In that system the particles are all the same, so the derivative $`F^{\mathrm{ex}}/\mathit{ϵ}_i`$ takes the same vectorial value for any particle $`i`$. Hence we may evaluate it for particle number 1, without loss of generality. Similarly, the second derivate $`^2F^{\mathrm{ex}}/\mathit{ϵ}_i\mathit{ϵ}_j`$ takes just two tensorial values, depending on whether $`i`$ and $`j`$ label the same particle (number 1, say) or different particles (1 and 2, say). Hence, dividing the above equation by the system volume $`V`$ gives the excess free energy density $$^{\mathrm{ex}}=_\mathrm{m}^{\mathrm{ex}}+\rho \mathit{ϵ}\frac{F^{\mathrm{ex}}}{\mathit{ϵ}_1}\|_\mathrm{m}+\frac{1}{2}\rho \mathit{ϵ}\mathit{ϵ}:\frac{^2F^{\mathrm{ex}}}{\mathit{ϵ}_1\mathit{ϵ}_1}\|_\mathrm{m}+\frac{1}{2}\rho \mathit{ϵ}\mathit{ϵ}:N\frac{^2F^{\mathrm{ex}}}{\mathit{ϵ}_1\mathit{ϵ}_2}\|_\mathrm{m}+Ø(ϵ^3).$$ (35) Note that, although the final term contains a factor $`N`$, it remains intensive, because $`^2F^{\mathrm{ex}}/\mathit{ϵ}_1\mathit{ϵ}_21/V`$ since particles 1 and 2 interact less as the system grows. Comparing Eq. (35) with Eq. (6), we can identify expressions for the coefficients $`𝑨`$, $`𝑩`$ and $`𝑪`$, which we also express in terms of derivatives of Eqs. (6) and (8), $`𝑨={\displaystyle \frac{^{\mathrm{ex}}}{\mathit{ϵ}}}`$ $`=`$ $`\rho {\displaystyle \frac{F^{\mathrm{ex}}}{\mathit{ϵ}_1}}=\rho {\displaystyle \frac{\text{d}\mu ^{\mathrm{ex}}(\mathit{ϵ})}{\text{d}\mathit{ϵ}}}`$ (37) $`𝑩={\displaystyle \frac{^{\mathrm{ex}}}{\mathit{ϵ}\mathit{ϵ}}}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\rho {\displaystyle \frac{^2F^{\mathrm{ex}}}{\mathit{ϵ}_1\mathit{ϵ}_1}}={\displaystyle \frac{1}{2}}\rho {\displaystyle \frac{\text{d}^2\mu ^{\mathrm{ex}}(\mathit{ϵ})}{\text{d}\mathit{ϵ}\text{d}\mathit{ϵ}}}`$ (38) $`𝑪={\displaystyle \frac{1}{2}}{\displaystyle \frac{^2^{\mathrm{ex}}}{\mathit{ϵ}\mathit{ϵ}}}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\rho N{\displaystyle \frac{^2F^{\mathrm{ex}}}{\mathit{ϵ}_1\mathit{ϵ}_2}}={\displaystyle \frac{1}{2}}\rho {\displaystyle \frac{}{\mathit{ϵ}}}{\displaystyle \frac{\text{d}\mu ^{\mathrm{ex}}(\mathit{ϵ})}{\text{d}\mathit{ϵ}}}`$ (39) with all formulae evaluated in the monodisperse reference system. To clarify the meanings of the various derivatives, let us read Eq. (IVa) from left to right. It states that the parameter $`𝑨`$, which is a property of a monodisperse phase, is the rate of change of excess free energy density as the mean property $`\mathit{ϵ}`$ of all particles in the phase is changed. The next equality asserts that $`𝑨`$ is the overall number density $`\rho `$ times the rate at which the extensive excess free energy changes when the property $`\mathit{ϵ}`$ of particle number 1 alone is changed. Finally from Eq. (IVa), this is the same as $`\rho `$ times the variation in excess chemical potential for different species. (Though different species are absent from the monodisperse system, their excess chemical potentials remain well defined.) In sections IV A and IV B two methods are developed, using the various relations in Eqs. (IV) to ascertain the values of $`𝑨`$, $`𝑩`$ and $`𝑪`$ in any given system. ### A Thermodynamic perturbation theory In this section, polydispersity is treated as a perturbation to the Hamiltonian of a system. The Hamiltonian of a polydisperse system (neglecting the trivial kinetic part) is $`H(\mathrm{\Gamma })`$, a function of the set $`\mathrm{\Gamma }`$ of particle coordinates (both positional and internal) which defines a configuration of the system. Given that a monodisperse reference system in the same configuration $`\mathrm{\Gamma }`$ has Hamiltonian $`H_\mathrm{m}(\mathrm{\Gamma })`$, we define the perturbation $`\widehat{H}(\mathrm{\Gamma })`$ by $$H(\mathrm{\Gamma })H_\mathrm{m}(\mathrm{\Gamma })+\widehat{H}(\mathrm{\Gamma }).$$ (40) By Boltzmann and Gibbs, the Helmholtz free energy of the polydisperse system is (with units $`k_BT1`$) $`F=N{\displaystyle \text{d}\mathit{ϵ}p(\mathit{ϵ})\left(\mathrm{ln}[Np(\mathit{ϵ})]1\right)}\mathrm{ln}{\displaystyle \text{d}\mathrm{\Gamma }\mathrm{exp}[H(\mathrm{\Gamma })]}`$ and that of the monodisperse system is $`F_\mathrm{m}=N[\mathrm{ln}N1]\mathrm{ln}{\displaystyle \text{d}\mathrm{\Gamma }\mathrm{exp}[H_\mathrm{m}(\mathrm{\Gamma })]}.`$ In these two equations we may set $`H`$ and $`H_\mathrm{m}`$ to zero to find expressions for the ideal part, which we subtract, leaving the excess part of the free energies, $`F^{\mathrm{ex}}=\mathrm{ln}{\displaystyle \text{d}\mathrm{\Gamma }\mathrm{exp}[H(\mathrm{\Gamma })]}+\mathrm{ln}{\displaystyle \text{d}\mathrm{\Gamma }}`$ and $`F_\mathrm{m}^{\mathrm{ex}}=\mathrm{ln}{\displaystyle \text{d}\mathrm{\Gamma }\mathrm{exp}[H_\mathrm{m}(\mathrm{\Gamma })]}+\mathrm{ln}{\displaystyle \text{d}\mathrm{\Gamma }}.`$ By applying Eq. (40), we see that the difference between these expressions is $`F^{\mathrm{ex}}F_\mathrm{m}^{\mathrm{ex}}=\mathrm{ln}\left({\displaystyle \frac{\text{d}\mathrm{\Gamma }e^{H_\mathrm{m}}e^{\widehat{H}}}{\text{d}\mathrm{\Gamma }e^{H_\mathrm{m}}}}\right)`$ where the argument of the logarithm is the thermal average of $`\mathrm{exp}(\widehat{H}(\mathrm{\Gamma }))`$ in the monodisperse system. Let us denote the thermal average of a stochastic variable $`\zeta `$ by $`\zeta `$. This is a Boltzmann-weighted average over configurations, denoted by double brackets to distinguish it from averages over the distribution of species $`p(\mathit{ϵ})`$, for which we have been using single brackets $`\mathrm{}`$. Hence, the excess free energy of a polydisperse system, in terms of a perturbation $`\widehat{H}`$ to a monodisperse system, is $$F^{\mathrm{ex}}=F_\mathrm{m}^{\mathrm{ex}}\mathrm{ln}e^{\widehat{H}}_\mathrm{m}$$ (41) where the subscript m indicates that the thermal average is performed in the monodisperse system. Equation (41) has the familiar form of thermodynamic perturbation theory but, for polydispersity, it applies only to the excess part of the free energy. The ideal parts of the mono- and polydisperse free energies differ by a non-perturbative amount. Note that we are not making any approximation of the interactions, unlike many perturbation theories which use an ideal or harmonic reference state. Hence, as with our previous analysis, Eq. (41) treats the effects arising purely due to polydispersity, in a fully interacting system. We now apply the second equality in Eq. (IVa) to the perturbative expression for the excess free energy, Eq. (41), to find $`𝑨`$ by varying the properties of a single particle in an otherwise monodisperse system. The first term in Eq. (41), $`F_\mathrm{m}^{\mathrm{ex}}`$ is a constant and therefore does not contribute. We find $`𝑨/\rho ={\displaystyle \frac{F^{\mathrm{ex}}}{\mathit{ϵ}_1}}\|_\mathrm{m}={\displaystyle \frac{}{\mathit{ϵ}_1}}\mathrm{ln}e^{\widehat{H}}_\mathrm{m}\|_\mathrm{m}`$ yielding $$𝑨/\rho =\frac{H}{\mathit{ϵ}_1}_\mathrm{m}$$ (43) where the ‘hat’ ($`\widehat{}`$) has been dropped from the Hamiltonian $`H`$ since, by definition, the perturbation $`\widehat{H}`$ is the only part that varies with $`\mathit{ϵ}_1`$. Similarly, for Eqs. (IVb) and (IVc), $`{\displaystyle \frac{2𝑩}{\rho }}={\displaystyle \frac{^2H}{\mathit{ϵ}_1\mathit{ϵ}_1}}_\mathrm{m}+{\displaystyle \frac{H}{\mathit{ϵ}_1}}_\mathrm{m}{\displaystyle \frac{H}{\mathit{ϵ}_1}}_\mathrm{m}{\displaystyle \frac{H}{\mathit{ϵ}_1}}{\displaystyle \frac{H}{\mathit{ϵ}_1}}_\mathrm{m}`$ (44) (45) $`{\displaystyle \frac{2𝑪}{\rho N}}={\displaystyle \frac{^2H}{\mathit{ϵ}_1\mathit{ϵ}_2}}_\mathrm{m}+{\displaystyle \frac{H}{\mathit{ϵ}_1}}_\mathrm{m}{\displaystyle \frac{H}{\mathit{ϵ}_2}}_\mathrm{m}{\displaystyle \frac{H}{\mathit{ϵ}_1}}{\displaystyle \frac{H}{\mathit{ϵ}_2}}_\mathrm{m}.`$ (46) (47) In principle, the derivatives of the Hamiltonian appearing in Eqs. (43), (45) and (47) are known from the microscopic physics of a given system. For instance, in a system of particles with central, symmetric, pairwise additive interactions, whose inter-particle potential is $`U(r,\mathit{ϵ}_i,\mathit{ϵ}_j)`$, the Hamiltonian is $$H=\frac{1}{2}\underset{\stackrel{i=1}{}}{\overset{N}{}}\underset{\stackrel{j=1}{i}}{\overset{N}{}}U(\|𝒓_i𝒓_j\|,\mathit{ϵ}_i,\mathit{ϵ}_j).$$ (48) In terms of the matrices of derivatives of the potential, $`𝑼_1(r)`$ $``$ $`{\displaystyle \frac{U(r,\mathit{ϵ},\mathrm{𝟎})}{\mathit{ϵ}}}\|_{\mathit{ϵ}=\mathrm{𝟎}}`$ (50) $`𝑼_{11}(r)`$ $``$ $`{\displaystyle \frac{^2U(r,\mathit{ϵ},\mathrm{𝟎})}{\mathit{ϵ}\mathit{ϵ}}}\|_{\mathit{ϵ}=\mathrm{𝟎}}`$ (51) $`𝑼_{12}(r)`$ $``$ $`{\displaystyle \frac{^2U(r,\mathit{ϵ}_1,\mathit{ϵ}_2)}{\mathit{ϵ}_1\mathit{ϵ}_2}}\|_{\mathit{ϵ}_1=\mathit{ϵ}_2=\mathrm{𝟎}}`$ (52) this yields $$𝑨=\rho ^2_0^{\mathrm{}}𝑼_1(r)g_\mathrm{m}(r)\mathrm{\hspace{0.33em}4}\pi r^2\text{d}r.$$ (54) Here, $`g_\mathrm{m}(r)=\rho (\mathrm{𝟎})\rho (𝒓)_\mathrm{m}/\rho ^2`$ is the radial distribution function in the monodisperse reference phase. To obtain Eq. (IV Aa), we have used the fact that the density of particle centres in the system at any instant is $`\rho (𝒓)=_{i=1}^N\delta ^{(3)}(𝒓𝒓_i)`$. Equations (45, 47) can similarly be evaluated for central, symmetric, pairwise additive interactions, giving $`𝑩`$ $`=`$ $`{\displaystyle \frac{\rho ^2}{2}}{\displaystyle _0^{\mathrm{}}}\left[𝑼_{11}(r)𝑼_1(r)𝑼_1(r)\right]g_\mathrm{m}(r)\mathrm{\hspace{0.17em}4}\pi r^2\text{d}r+{\displaystyle \frac{\rho ^3}{2}}{\displaystyle \text{d}^3r\text{d}^3r^{}𝑼_1(r)𝑼_1(r^{})\left[g_\mathrm{m}(r)g_\mathrm{m}(r^{})g_\mathrm{m}^{(3)}(𝒓,𝒓^{})\right]}`$ (55) $`𝑪`$ $`=`$ $`{\displaystyle \frac{\rho ^2}{2}}{\displaystyle _0^{\mathrm{}}}\left[𝑼_{12}(r)𝑼_1(r)𝑼_1(r)\right]g_\mathrm{m}(r)\mathrm{\hspace{0.17em}4}\pi r^2\text{d}r{\displaystyle \frac{3\rho ^3}{2}}{\displaystyle \text{d}^3r\text{d}^3r^{}𝑼_1(r)𝑼_1(r^{})g_\mathrm{m}^{(3)}(𝒓,𝒓^{})}`$ (57) $`+{\displaystyle \frac{\rho ^4}{2}}{\displaystyle \text{d}^3r\text{d}^3r^{}\text{d}^3r^{\prime \prime }𝑼_1(r)𝑼_1(\left\|𝒓^{}𝒓^{\prime \prime }\right\|)\left[g_\mathrm{m}(r)g_\mathrm{m}(\left\|𝒓^{}𝒓^{\prime \prime }\right\|)g_\mathrm{m}^{(4)}(𝒓,𝒓^{},𝒓^{\prime \prime })\right]}`$ where $`g_\mathrm{m}^{(3)}(𝒓,𝒓^{})\rho (\mathrm{𝟎})\rho (𝒓)\rho (𝒓^{})_\mathrm{m}/\rho ^3`$ and $`g_\mathrm{m}^{(4)}(𝒓,𝒓^{},𝒓^{\prime \prime })\rho (\mathrm{𝟎})\rho (𝒓)\rho (𝒓^{})\rho (𝒓^{\prime \prime })_\mathrm{m}/\rho ^4`$. Notice that Eqs. (IV Aa, b, c) require a knowledge of spatial correlations in the monodisperse system only. In the appendix, it is shown how general thermal averages are perturbed by polydispersity, for such systems with soft, pairwise interactions. Equations (IV A) and (IV A) allow the parameters affecting phase equilibria (to lowest order in the width of the distribution) to be calculated for any system with soft (i.e. differentiable) interactions. We shall put these expressions to use on some real systems in sections V and VI. ### B Scalable systems: The special case of hard spheres In this section we shall consider a system of prevalent theoretical interest, polydisperse hard spheres. Hard spheres interact via a potential which is zero, except for configurations where the spheres overlap, which are forbidden and hence have infinite potential energy. Hard-sphere systems attract interest because they exemplify substances with short-range repulsive interactions, while lacking a characteristic energy scale. This leads to temperature-independent behaviour which is purely entropy-driven, and engenders simplicity due to the small number of tunable parameters. Unfortunately, hard spheres are an exceptional case to which the thermodynamic perturbation theory in section IV A cannot be applied, since the internal energy is zero and the interaction potential is discontinuous. Hence the values of the parameters $`𝑨`$, $`𝑩`$ and $`𝑪`$ cannot be found by that method, though their formulations in Eqs. (IV) still hold. In this section we shall exploit a special property of the hard-sphere system to establish the exact value of the parameter $`𝑨`$ (previously approximated). In fact the method can be applied to any system of particles whose interactions are ‘scalable’ — a term which will be elucidated below. With a knowledge of $`𝑨`$, the partitioning of hard spheres between slightly-polydisperse phases is fully determined in Eq. (15), (16) or (20). Unfortunately it is not clear how such a method might be used to determine $`𝑩`$, the other parameter on which the phase boundaries depend. In section VI its value is taken from an approximate expression for the free energy of polydisperse hard spheres, of which there are many examples in the literature. To calculate $`𝑨`$, the coefficient of $`\mathit{ϵ}`$ in the free energy expansion (Eq. (6)), for polydisperse hard spheres, we use the first equality of Eq. (IVa). Let the radius $`r`$ of each sphere be measured relative to a reference length $`r_0`$ thus: $`r=(1+ϵ)r_0`$. As this is the only property that varies from sphere to sphere, the value $`ϵ`$ which characterises the particles is a scalar, so Eq. (IVa) reduces to a scalar equation. That equation states that $`A`$ is the rate at which the free energy density changes when the first moment $`ϵ`$ of size deviations is increased, while the other moments remain stationaryIn a monodisperse system, the distribution is a delta function. It is therefore not possible to vary the mean while holding the other moments constant. Nevertheless, being of higher order in the small quantity $`\mathit{ϵ}`$, the higher moments are stationary, so the partial derivative is well defined.. When this hypothetical change is made, the particles all grow by the same amount, so the system remains monodisperse. In terms of the unique radius $`r`$ of the monodisperse particles, the derivative is $`{\displaystyle \frac{^{\mathrm{ex}}}{ϵ}}=r_0{\displaystyle \frac{^{\mathrm{ex}}}{r}}`$ which will be evaluated at $`r=r_0`$. One can imagine making this change in two stages: firstly the whole system is scaled up, so that the particles, and the space between them, and the volume of the system all increase, while the concentration $`\varphi `$ (the fraction of space occupied by particles) is held constant. Then the system is compressed back to its original volume while the particles remain at their new large size, so the concentration increases. The resulting change in the extensive excess free energy ($`F^{\mathrm{ex}}=V^{\mathrm{ex}}`$) is given by $`\left({\displaystyle \frac{F^{\mathrm{ex}}}{r}}\right)_{V,N}=\left({\displaystyle \frac{F^{\mathrm{ex}}}{r}}\right)_{\varphi ,N}\left({\displaystyle \frac{F^{\mathrm{ex}}}{V}}\right)_{r,N}\left({\displaystyle \frac{V}{r}}\right)_{\varphi ,N}.`$ The last term here is due to the reversible work done against excess pressure when the system is compressed, with $`(V/r)_{\varphi ,N}=3V/r`$ being a conversion factor between length and volume. So we have $`\left({\displaystyle \frac{F^{\mathrm{ex}}}{r}}\right)_{V,N}=\left({\displaystyle \frac{F^{\mathrm{ex}}}{r}}\right)_{\varphi ,N}+3{\displaystyle \frac{P^{\mathrm{ex}}V}{r}}.`$ The first term on the right hand side is trivial to calculate because there exist no length scales (such as the range of an interaction) that remain fixed as the particles grow. Hence the growth is simply an overall scaling. This is what was meant by the term “scalable system” above. By writing the free energy in terms of the usual configurational integral over particle positions measured in units of the thermal de Broglie wavelength, it is easy to establish that this term vanishes, since the excess part does not vary with the overall scale factor. Hence, for hard spheres, $$A=3P^{\mathrm{ex}}=3(P\rho )$$ (58) where, as usual, $`k_BT`$ has been set to unity. Since, at equilibrium, $`P`$ is a constant for all phases, the coefficient in Eqs. (1620) may be written $`[A/\rho ]_\beta ^\alpha =3P[\rho ^1]_\beta ^\alpha `$ for size-polydispersity in scalable systems. Though it is not apparent that an analogous method exists for calculating $`B`$, the parameter $`C`$ is calculable from the first equality of Eq. (IVc), by taking the second derivative of excess free energy density with respect to particle size in a monodisperse scalable system. We shall not perform that calculation here. For monodisperse hard spheres, there is a transition from a crystal of volume fraction 0.545 to a fluid phase of volume fraction 0.494 which, according to the Carnahan-Starling equation of state , exists at a pressure $`6.17\pm 0.02`$, in units scaled by $`k_BT`$ and by the volume of a particle. Hence, from Eq. (58), the coefficient $`[A/\rho ]_{\mathrm{fluid}}^{\mathrm{crystal}}`$, which appears in Eq. (20), evaluates to $`3.51\pm 0.04`$. This is consistent with simulation results in which the average sizes of polydisperse hard spheres were measured in coexisting fluid and crystal phases, and the difference plotted against the variance of the overall distribution. The gradient was found to be $`3.55\pm 0.01`$, and changed by only a few percent up to polydispersities of 6% and more. ## V Example: Chemically polydisperse polymers It has been necessary to plough through a fair amount of mathematics in order to arrive at some simple and practical formulae. The convenience of this approach is that the lengthy derivations have been performed once and for all. They will not require repetition for each application. Lest the reader has lost sight of the wood for the trees, let us demonstrate how easily the formulae may be applied to a model of polymeric phase equilibria. We consider a Flory-Huggins-style description of a polymer blend. The reference system is a binary mixture of polymers (species ‘$`𝒜`$’ and ‘$``$’) of equal size $`L`$, but two different chemical types. The chemical nature of each species is parameterised by a number in the interval $`(1,1)`$, which may be interpreted as the hydrophobicity of the molecule. We make species $``$ chemically ‘neutral’ (hydrophobicity zero), while the value for species $`𝒜`$ is $`a`$. If the concentration of $`𝒜`$-polymers is $`\varphi `$ then that of $``$-polymers is $`1\varphi `$, and the mean-field free energy density is given by $$L_\mathrm{m}=\varphi \mathrm{ln}\varphi +(1\varphi )\mathrm{ln}(1\varphi )\chi \varphi _1^2$$ (59) where $`\varphi _1=a\varphi `$ is the ‘hydrophobic concentration’. The Flory-Huggins interaction parameter $`\chi `$ determines the strength of attraction between chemically similar species. This model system separates into coexisting $`𝒜`$-rich and $`𝒜`$-poor phases for $`\stackrel{~}{\chi }\chi a^2>2`$, the concentrations of which are easily found e.g. from a double-tangent construction on the free energy (Eq. (59)). The resulting phase diagram in the $`(\varphi ,\stackrel{~}{\chi })`$-plane is shown by the solid line in Fig. 1. If component $`𝒜`$ is now made chemically polydisperse, by varying slightly the constituent monomers on each polymer molecule, while component $``$ remains monodisperse and chemically neutral, acting only as a polymeric solvent, then the overall hydrophobic concentration of polymer $`𝒜`$ is the mean of a distribution: $`\varphi _1=\varphi {\displaystyle ap(a)\text{d}a}=\varphi a`$ and the free energy density becomes $$L=(1\varphi )\mathrm{ln}(1\varphi )\chi \varphi _1^2+\varphi p(a)\mathrm{ln}(\varphi p(a))\text{d}a.$$ (60) The model free energy of Eq. (60) was studied previously to demonstrate a sophisticated approximation scheme for systematically reducing the dimensionality of polydisperse phase equilibria problems . Accurate cloud (o) and shadow () points were calculated by that method, and are reproduced in Fig. 1 (where $`\stackrel{~}{\chi }\chi a_0^2`$) for a Gaussian parent with $`a_\mathrm{p}a_0=0.5`$ and polydispersity $`\sigma =8\%`$ . For comparison, our lowest-order perturbative formulae for the cloud and shadow points (Eqs. (3031)) can be evaluated with ease. Note that the last term of Eq. (60) has the form of an ideal free energy density if we identify $`\varphi `$ with $`\rho `$ in Eq. (2) (up to an irrelevant term linear in $`\rho `$). Hence, by writing $`a=a_0(1+ϵ)`$ to expand about the monodisperse reference system with $`a=a_0=0.5`$, Eq. (60) can be cast in the form of Eq. (6), with $`A`$ $`=`$ $`2\stackrel{~}{\chi }\varphi ^2`$ $`B`$ $`=`$ $`0`$ $`C`$ $`=`$ $`\stackrel{~}{\chi }\varphi ^2`$ all of which are scalar quantities since only one property (hydrophobicity) is polydisperse. The cloud and shadow points ($`\varphi _c=\varphi _c^\mathrm{m}+\delta _c`$, $`\varphi _s=\varphi _s^\mathrm{m}+\delta _s`$) are obtained by substituting these values into Eqs. (3031) together with the isothermal compressibility derived from Eq. (59), $`\kappa =(1\varphi )/(\varphi 2\stackrel{~}{\chi }\varphi ^2(1\varphi ))`$, yielding $`\varphi _c`$ $`=`$ $`\varphi _c^\mathrm{m}+2\sigma ^2\stackrel{~}{\chi }^2{\displaystyle \frac{\varphi _c^\mathrm{m}(1\varphi _c^\mathrm{m})(\varphi _c^\mathrm{m}\varphi _s^\mathrm{m})\varphi _s^\mathrm{m}}{12\stackrel{~}{\chi }\varphi _c^\mathrm{m}(1\varphi _c^\mathrm{m})}}`$ (61) $`\varphi _s`$ $`=`$ $`\varphi _s^\mathrm{m}+2\sigma ^2\stackrel{~}{\chi }^2{\displaystyle \frac{\varphi _s^\mathrm{m}(1\varphi _s^\mathrm{m})(\varphi _s^\mathrm{m}\varphi _c^\mathrm{m})(2\varphi _s^\mathrm{m}\varphi _c^\mathrm{m})}{12\stackrel{~}{\chi }\varphi _s^\mathrm{m}(1\varphi _s^\mathrm{m})}}`$ (62) where $`\varphi _c^\mathrm{m}`$, $`\varphi _s^\mathrm{m}`$ are the respective points on the monodisperse binodal (solid line in Fig. 1). Eqs. (6162) are plotted in Fig. 1 as dashed and dash-dotted lines respectively. As expected, the perturbative expansion scheme breaks down near the critical point, where the compressibility diverges and the separate reference phases disappear. This invalid regime extends over only a small range of $`\stackrel{~}{\chi }`$, and the phase boundaries rapidly converge, with increasing $`\stackrel{~}{\chi }`$, to a very accurate answer. Notice that the miscibility gap is broadened by polydispersity. Equation (34) shows that the polydisperse perturbation to the gap must be positive whenever $`B=B^{}=0`$, i.e. when there is no explicit dependence of the excess free energy on the variance of the distribution of species. We also have formulae prescribing the chemical properties of the two phases. By definition, the distribution of polymers in the cloud phase is just the same as the parent distribution that was input to the system. The shadow phase, on the other hand, has chosen its own preferred mix. The mean hydrophobicity of type-$`𝒜`$ polymers in this phase is given by Eq. (III Ba) (or equivalently Eq. (15)) as $`a_s=a_0(1+ϵ_s)=a_0(1+2\sigma ^2\stackrel{~}{\chi }(\varphi _s\varphi _c)+𝒪(\sigma ^3))`$, which is also in good agreement with exact numerical data. ## VI Example: Sterically-stabilised colloid-polymer mixtures To further illustrate the utility of the formulae derived in this article, let us make a case study of a familiar complex fluid: a colloid-polymer mixture . The phase behaviour of the monodisperse system is first briefly reviewed. A theoretical understanding of its behaviour is well-established . However, the effects of colloidal polydispersity on the system have not previously been calculated, despite the fact that the results of colloidal synthesis are inevitably slightly polydisperse. The effects of polymeric polydispersity on the system, which have been modelled previously , will not be analysed. ### A Established monodisperse behaviour In this system, small balls (typically of diameter 10 nm to 1 $`\mu `$m) of PMMA (‘Perspex’) or silica are suspended in an organic liquid. To avoid the balls (or ‘colloidal latices’) clumping due to van der Waals attraction, they are each given a short brush-like coat of ‘steric’ polymer chains, making the latices behave as hard spheres, with no long-range interactions but a very short-range strong repulsion. Hard spheres of this kind are known to have a simple phase diagram, and carry no latent heat. The only controlling parameter is the fraction of space $`\varphi `$ occupied by the spheres. A single phase transition exists, between an amorphous fluid state and an FCC crystal, which coexist at $`\varphi 0.494`$ and $`\varphi 0.545`$ respectively . The phase behaviour is enriched by the addition of free polymer chains to the mixture. Each of these chains, being very long, bunches itself into the form of a self-avoiding random walk. Such a ‘ball of string’ cannot penetrate a colloidal latex, and so interacts with the colloid particles as if it too were a hard sphere. On the other hand, when two randomly coiled polymer molecules meet, they can interpenetrate and, with a suitable choice of suspending solvent, can be made almost entirely non-interactingIn particular, the second virial coefficient of the polymeric osmotic pressure can be made to vanish, though higher virial coefficients remain ., like an ideal gas . The ideal gas pressure that they exert on the surface of an isolated colloidal hard sphere is isotropic. However, if a second hard sphere is so close to the first that the non-penetrating polymer coils cannot fit between them, then the polymeric pressure is not felt on the adjacent parts of the colloids, so they are pushed together. Asakura and Oosawa showed that this mechanism could be described in terms of an effective attraction between colloidal spheres, known as the depletion interaction. Thus, rather than considering a mixture of two components (colloid and polymer) in the solvent, one can imagine a single-component system, consisting of spheres with a hard core repulsion and a longer-range pairwise<sup>\</sup><sup>\</sup>\**A more accurate description would include corrections to the pairwise attraction, since three or four latices in close proximity exclude polymer coils from a region with non-trivial geometry. attraction. This effective attractive potential $`U(r)`$ between two latices of radius $`a`$ and centre-to-centre separation $`r`$ is simply the product of the (ideal) polymeric osmotic pressure $`\mathrm{\Pi }_p`$ and the volume of space from which the centres of polymer coils are excluded by the hard spheres . At $`r=2a`$ (contact of the two hard spheres) the potential rises discontinuously to infinity, and for $`r>2(a+r_g)`$, where $`r_g`$ is the radius of gyration of a polymer chain, there is no interaction, since the diameter of a polymer coil ($`2r_g`$) sets the range of the effective potential. The phase diagram of this system in the $`(\varphi ,\mathrm{\Pi }_p)`$-plane has been calculated for the monodisperse case , using results from the Percus-Yevick integral equation theory for hard spheres . It was found (in agreement with experiments ) to resemble the phase diagram of simple atomic substances such as argon, with $`\mathrm{\Pi }_p^1`$ playing the rôle of temperature. As well as the colloidal crystal, there are low- and high-concentration amorphous phases, dubbed ‘colloidal gas and liquid’, which meet at a critical point. We shall study the ‘fluid-fluid’ coexistence of these amorphous phases, thus avoiding problems of non-ergodicity arising in the polydisperse crystal . ### B Calculating the polydisperse phase equilibria The strength of the treatment of polydispersity in this article is that it uses, as input, properties of the monodisperse system. The phase behaviour of monodisperse colloid-polymer mixtures is well understood, so application of the new formulae for polydisperse phase equilibria will be straightforward. In addition to the monodisperse phase diagram, we require the functions $`A(\rho )`$ and $`B(\rho )`$, defined by Eq. (6), which are scalars in this case since the particle radii alone are polydisperse. Unlike the previous example in section V, we are not given an expression for the free energy of the polydisperse system. According to Eq. (IVa), $`A(\rho )`$ can be found from the monodisperse free energy, for which we use the expression given in Ref. . It is simply the derivative of the free energy with respect to the radius of monodisperse colloidal particles (scaled by the radius), at fixed number density. Hence an expression for $`A(\rho )`$ is straightforwardly found, so long as care is taken not to overlook the dependence of concentration on particle size, $`\varphi =\frac{4}{3}\pi a^3\rho `$, and to keep the polymer size $`r_g`$ fixed while taking the derivative. It is not so trivial to evaluate $`B(\rho )`$ without a prescribed polydisperse free energy. Firstly, we must specify the difference between the Hamiltonian of the system with slightly polydisperse colloid, and that of a monodisperse reference system. Let us write the Hamiltonian of the colloid-polymer mixture with polydisperse colloid as a sum of two parts: $$H=H_{\mathrm{HS}}+H_{\mathrm{dep}}^{\mathrm{eff}}$$ (63) the Hamiltonian $`H_{\mathrm{HS}}`$ of the hard-sphere colloid, plus effective interactions $`H_{\mathrm{dep}}^{\mathrm{eff}}`$ due to the depletion of polymer. Following the analysis of the monodisperse case, we assume that the resulting free energy is also separable into two parts: $$F=F_{\mathrm{HS}}+F_{\mathrm{dep}}$$ (64) where $`F_{\mathrm{HS}}`$ is approximately the free energy of polydisperse hard spheres in the absence of polymer. This is consistent with the Weeks-Chandler-Andersen (WCA) approximation which states that the spatial distribution of repulsive particles (in this case hard spheres) is not significantly altered by the introduction of additional interactions which are purely attractive (here contained in $`H_{\mathrm{dep}}^{\mathrm{eff}}`$). Accordingly, the polymeric contribution to the free energy $`F_{\mathrm{dep}}`$ will be approximated by averaging the depletion Hamiltonian $`H_{\mathrm{dep}}^{\mathrm{eff}}`$ over the pure hard sphere distribution. Note that the structure of Eq. (64) implies that the desired function $`B(\rho )`$ also splits into two terms, $$B=B_{\mathrm{HS}}+B_{\mathrm{dep}}$$ (65) The depletion interaction is soft, so the thermodynamic perturbation theory of section IV A is suitable for calculating its contribution in $`B_{\mathrm{dep}}`$. The hard sphere potential, on the other hand, is non-differentiable and therefore inappropriate for description in terms of an energetic perturbation, as noted above. Other methods will be used to deal with the hard sphere part $`B_{\mathrm{HS}}`$. Let us tackle the polymeric (depletion) contribution first. As stated above, the effective Hamiltonian (i.e. the polymeric free energy for a fixed configuration of colloidal spheres) $`H_{\mathrm{dep}}^{\mathrm{eff}}`$ is simply the product of the ideal polymeric osmotic pressure and the volume of space from which centres of polymer coils are excluded. Let us define the polymer-colloid size ratio as $`\xi r_g/a`$, and $`v_s\frac{4}{3}\pi a^3`$ to be the volume of a reference sphere with the mean radius. An isolated colloidal sphere of radius $`a(1+ϵ)`$ excludes polymer coils from a volume $`v_s(1+\xi +ϵ)^3`$. Hence, even for well-separated colloids beyond the interaction range, there is a bulk contribution to the depletion Hamiltonian, $$H_{\mathrm{bulk}}=\frac{\widehat{\mathrm{\Pi }}_p}{\xi ^3}\underset{i=1}{\overset{N}{}}(1+\xi +ϵ_i)^3$$ (66) where the polymeric osmotic pressure has been re-scaled: $`\widehat{\mathrm{\Pi }}_pv_s\xi ^3\mathrm{\Pi }_p`$, so that the $`(\varphi ,\widehat{\mathrm{\Pi }}_p)`$-plane will correspond to phase diagrams presented in Ref. for colloid-polymer mixtures. For monodisperse colloid, $`H_{\mathrm{bulk}}`$ is usually ignored, since it is a constant, independent of volume fraction $`\varphi `$. For a polydisperse system, on the other hand, it cannot be neglected, as it contains $`ϵ`$-dependence. From Eq. (45), its contribution to $`B(\rho )`$ is $$B_{\mathrm{bulk}}=3\widehat{\mathrm{\Pi }}_p\rho (1+\xi )/\xi ^3.$$ (67) When two colloidal latices, of radii $`(1+ϵ_1)a`$ and $`(1+ϵ_2)a`$, have a centre-centre separation $`(2+ϵ_1+ϵ_2)a<r<(2+2\xi +ϵ_1+ϵ_2)a`$, their individual regions of polymeric exclusion overlap, so that the total volume available to polymer increases, giving rise to the depletion attraction. It can be shown by some elementary geometry that the overlap volume of the two spheres of depletion; radii $`r_1=(1+\xi +ϵ_1)a`$ and $`r_2=(1+\xi +ϵ_2)a`$; is $`V_{\mathrm{overlap}}=\pi \left[{\displaystyle \frac{r^3}{12}}{\displaystyle \frac{(r_1^2+r_2^2)r}{2}}+{\displaystyle \frac{2(r_1^3+r_2^3)}{3}}{\displaystyle \frac{(r_1^2r_2^2)^2}{4r}}\right].`$ The resulting interaction potential $`U(r,ϵ_1,ϵ_2)=\mathrm{\Pi }_pV_{\mathrm{overlap}}`$ gives a pairwise contribution to the effective Hamiltonian, as in Eq. (48). Its derivatives, as defined in Eqs. (48), are $`U_1(r)`$ $`=`$ $`{\displaystyle \frac{3(1+\xi )}{4\xi ^3}}(2+2\xi r/a)\widehat{\mathrm{\Pi }}_p`$ (69) $`U_{11}(r)`$ $`=`$ $`{\displaystyle \frac{3}{4\xi ^3}}\left({\displaystyle \frac{r}{a}}4[1+\xi ]+2{\displaystyle \frac{a}{r}}[1+\xi ]^2\right)\widehat{\mathrm{\Pi }}_p`$ (70) which can be substituted into Eq. (IV Ab) for $`B_{\mathrm{dep}}`$. Additionally, that equation requires the pair and three-point distribution functions $`g_\mathrm{m}(r)`$, $`g_\mathrm{m}^{(3)}(𝒓,𝒓^{})`$ for the monodisperse system. At the level of the WCA approximation, these are replaced by the pure hard-sphere (HS) distribution functions, $`gg^{\mathrm{HS}}`$, for which we use the Percus-Yevick expression . As a simplifying assumption, let us also use the Kirkwood superposition approximation for the three-body correlations, $`g^{(3)}(𝒓,𝒓^{})g(r)g(r^{})g(\|𝒓𝒓^{}\|).`$ It remains only to find the hard-sphere contribution to $`B`$. Unfortunately, the methods developed here do not enable a first-principles derivation of this quantity, so we must look to other analyses of polydisperse hard spheres for an evaluation of $`B_{\mathrm{HS}}`$. In particular, we refer to the free energy expression for the polydisperse hard-sphere fluid due to Boublik, Mansoori, Carnahan, Starling and Lealand (BMCSL), which is known to reduce to the Carnahan-Starling free energy in the monodisperse limit. The application of this expression perhaps requires some clarification. The use of results from other studies of polydispersity by no means depreciates the present analysis. In cases where a polydisperse free energy is already known, one is spared the application of first-principles methods such as section IV. However, the results of section III (or equivalent methods) are still required, in order to derive (to lowest-order) the phase equilibria from that free energy. Indeed, the gap between a knowledge of the free energy and a knowledge of the phase equilibria is evidenced by the BMCSL free energy, whose regime of thermodynamic instability was only recently established . To extract the required function, the BMCSL free energy is cast in the form of Eq. (6) by expanding to order $`ϵ^2`$, at a density $`\rho =\varphi _0/v_s`$, giving $`v_s_{\mathrm{HS}}^{\mathrm{ex}}`$ $`=`$ $`{\displaystyle \frac{\varphi _0^2(43\varphi _0)}{(1\varphi _0)^2}}+ϵ{\displaystyle \frac{6\varphi _0^2(2\varphi _0)}{(1\varphi _0)^3}}+ϵ^23\varphi _0\left[{\displaystyle \frac{\varphi _0(1+3\varphi _02\varphi _0^2)}{(1\varphi _0)^3}}\mathrm{ln}(1\varphi _0)\right]`$ (72) $`+ϵ^23\varphi _0\left[{\displaystyle \frac{\varphi _0(1+2\varphi _0)(3+\varphi _0\varphi _0^2)}{(1\varphi _0)^4}}+\mathrm{ln}(1\varphi _0)\right]+Ø(ϵ^3).`$ The zeroth-order term is the excess free energy of a monodisperse hard-sphere fluid, and recovers the Carnahan-Starling equation of state. By comparison with Eq. (6), the coefficient of $`ϵ^2`$ is the required function $`v_sB_{\mathrm{HS}}`$. For definiteness, let us consider a polymer-colloid size ratio $`\xi =0.4`$. For a given polymeric osmotic pressure $`\widehat{\mathrm{\Pi }}_p`$, the coexisting concentrations $`\varphi `$ of monodisperse colloidal fluid phases are given in Ref. . These values for the concentrations of coexisting monodisperse gas ($`\varphi _\mathrm{m}^g`$) and liquid ($`\varphi _\mathrm{m}^l`$) phases are substituted into the expressions calculated above to obtain values of $`A`$ and $`B`$ for the coexisting phases. The monodisperse gas-liquid phase boundaries $`\varphi _\mathrm{m}^g`$, $`\varphi _\mathrm{m}^l`$ are shown as a solid line in Fig. 2 from the critical point at $`\widehat{\mathrm{\Pi }}_p0.41`$ to the triple point at $`\widehat{\mathrm{\Pi }}_p0.54`$. The calculated values of $`A`$ and $`B`$ and their derivatives are used in Eqs. (3031) to find the change in each phase’s cloud- and shadow-point densities due to polydispersity, which are translated into concentrations via Eq. (32). These shifts in concentration are added to the monodisperse phase boundary in Fig. 2 to give the resulting cloud and shadow points $`\varphi ^{c,s}=\varphi _\mathrm{m}^{c,s}+\delta \varphi ^{c,s}`$ for polydispersity $`\sigma =8\%`$, shown as dashed and dash-dotted curves respectively. We see from the figure that the expansion has remained well controlled (has not blown up) even for values of $`\widehat{\mathrm{\Pi }}_p`$ very close to criticality of the reference system. The theory predicts a widening of the coexistence region (the miscibility gap) so that liquid will condense from a polydisperse gas of lower concentration than in the monodisperse case. We also observe that when the gas phase exists in an infinitesimal quantity (at the shadow point), it is more concentrated than at its cloud point. The same applies to the liquid phases. Similar results are also found for different size ratios $`\xi `$ (studied for $`0.3<\xi <1`$). The above calculation of $`B(\rho )`$ for the colloid-polymer mixture involves a degree of inexactitude, since the radial distribution function is approximated by the Percus-Yevick hard-sphere distribution function, while three-point correlations are even more crudely estimated. To assess the importance of precision in the correlations, the cloud and shadow points were re-calculated using a Heaviside step function $`g(r)=\mathrm{\Theta }(r2a)`$ in place of the Percus-Yevick radial distribution. The polydispersity-induced shift in the high-density cloud point changes sign when the correlations are neglected, indicating, not surprisingly, that the correct correlations are important in determining the high-density phase boundary. Though the positions of the other phase boundaries are also affected by the change in $`g(r)`$, certain qualitative features remain unaltered, and are therefore expected to be reliable predictions of our approximate model. In particular, each shadow point remains more concentrated than the equivalent cloud point and, far from the critical point, the gases are shifted to lower concentrations by polydispersity. According to Eqs. (16) and (20), the difference between the normalised populations of the coexisting phases is proportional to $`[A/\rho ]_g^l`$ and the fractional difference in mean particle size is $`[ϵ]_g^l=[A/\rho ]_g^l\sigma _\mathrm{p}^2`$ where $`\sigma _\mathrm{p}`$ is the overall polydispersity. So the coefficient $`[A/\rho ]_g^l`$ determines the amount of partitioning or fractionation. It is plotted for the gas-liquid phase boundary (from $`\widehat{\mathrm{\Pi }}_p=0.41`$ to $`0.54`$) in Fig. 3. Note that it is positive, indicating that, on average, the liquid phase contains larger particles than the gas. The hard-sphere contribution is negative, but is outweighed by the depletion part. ## VII Summary The results derived here are exact in the limit of low polydispersity, and can therefore provide accurate information on a great many systems, whose manufacture attempts to approximate the ideal of monodisperse constituents. Additionally, the results will be of pragmatic use even in systems with a wide scatter of particle properties, in providing qualitative predictions and estimates for the effects of polydispersity. In any case, minimal effort is required to substitute values into the formulae. Equations (15) and, equivalently, (16) show that, in the narrow limit, the difference (due to fractionation) between normalised distributions in coexisting phases does not depend on the volumes of the phases, and is in fact remarkably simple, requiring only one system-dependent parameter per polydisperse property. From Eq. (20), the $`m`$th moment about the centre of the distribution differs, between coexisting phases, by an amount proportional to the $`(m+1)`$th moment of the parent<sup>††</sup><sup>††</sup>††plus terms of the order of the $`(m+2)`$th moment, which become important in a near-symmetric distribution if $`m`$ is even .. For hard spheres, the constant of proportionality has been determined (Eq. (58)), and agrees with data from simulation . The calculation is indistinguishable from the data at 2% polydispersity, and deviates by only 5% at 4% polydispersity. Additionally, at the nematic cloud point of the polydisperse Zwanzig model of hard rods , the form of Eq. (20) has been shown to hold to within 5% up to the remarkably large polydispersity of 50%. In section III A, it was shown that this simple rule can lead to ‘convective fractionation’ for multiply-polydisperse systems, whereby one property is partitioned between phases due to a driving force on another. The expressions found for the shift in the phase boundaries due to polydispersity, at a general coexistence (Eq. (29)) and at the cloud (Eq. (30)) and shadow points (Eq. (31)), show that the shift in the density of a phase is proportional to its isothermal compressibility and to the square of the polydispersity. This explains why these phase boundaries are parabolic near $`\sigma =0`$ when plotted as a function of polydispersity $`\sigma `$ in a range of models; see e.g. . The other parameters on which these shifts depend can be extracted directly from the form of the polydisperse free energy if it is known or, if not, from the Hamiltonian via the thermodynamic perturbation theory in section IV A. In the case of pairwise interactions, the relevant parameters are obtained from a knowledge of only two- and three-body correlations in the reference system, even when larger clusters are responsible for the existence of the phase transition in question. This illustrates a strength of the small-polydispersity expansion: that the underlying reference system takes care of complicated many-body interactions, so they do not need to be calculated in the analysis of the polydispersity. The application of the methods to the fluid-fluid coexistence of colloid-polymer mixtures in section VI revealed that colloidal polydispersity leads to a widening of the coexistence region, and favours larger particles living in the liquid phase. Similar methods are used in appendix A to find how polydispersity alters correlation functions. Other methods exist for calculating polydisperse phase equilibria. Whatever the theoretical formalism by which one chooses to analyse a polydisperse system, the results, if correct, will tend in the limit of low polydispersity to Eq. (16) for the degree of fractionation, and Eq. (29) for the movement of the binodals. ## VIII Acknowledgements Many thanks for informative discussions go to Peter Sollich, David Fairhurst, Michael Cates, Paul Bartlett, Peter Bolhuis, Richard Sear, Patrick Warren and Wilson Poon. The work was funded by a Royal Society of Edinburgh/SOEID personal research fellowship, and EPSRC grant number GR/M29696. ## A Thermal averages and spatial correlations In polydisperse systems, correlation functions are easier to calculate than phase boundaries, as there is no partitioning of the sample to consider. The thermal average of a stochastic quantity $`\zeta `$ in a single polydisperse phase is given by the usual perturbative expression with respect to a monodisperse reference system $$\zeta =\frac{\zeta e^{\widehat{H}}_\mathrm{m}}{e^{\widehat{H}}_\mathrm{m}}.$$ (A1) We shall require the second-order expansion of this expression, in the small quantity $`\widehat{H}`$, which must therefore be small (in units of $`k_BT`$), $$\zeta =\zeta _\mathrm{m}\left(1+\widehat{H}_\mathrm{m}\right)\left(\zeta \widehat{H}_\mathrm{m}\zeta _\mathrm{m}\widehat{H}_\mathrm{m}\right)+\frac{1}{2}\left(\zeta \widehat{H}^2_\mathrm{m}\zeta _\mathrm{m}\widehat{H}^2_\mathrm{m}\right)+Ø(\widehat{H}^2).$$ (A2) Let us restrict the discussion to a system of particles interacting via a pairwise-additive, symmetric, isotropic, central potential $`U(r,\mathit{ϵ}_1,\mathit{ϵ}_2)`$, for which the Hamiltonian is given in Eq. (48). With the Hamiltonian Taylor-expanded to second order in $`\mathit{ϵ}`$, the perturbation can be written $`\widehat{H}={\displaystyle \frac{1}{2}}{\displaystyle \underset{i=1}{\overset{N}{}}}{\displaystyle \underset{\stackrel{j=1}{i}}{\overset{N}{}}}\{(\mathit{ϵ}_i+\mathit{ϵ}_j)𝑼_1(\|𝒓_i𝒓_j\|)+{\displaystyle \frac{1}{2}}(\mathit{ϵ}_i\mathit{ϵ}_i+\mathit{ϵ}_j\mathit{ϵ}_j):𝑼_{11}(\|𝒓_i𝒓_j\|)+\mathit{ϵ}_i\mathit{ϵ}_j:𝑼_{12}(\|𝒓_i𝒓_j\|)\}+Ø(ϵ^3)`$ in terms of the derivatives of the interaction potential defined in Eqs. (48). Substituting this expression into Eq. (A2), we keep terms to second order in $`\mathit{ϵ}`$, and measure the deviations $`\mathit{ϵ}`$ with respect to the mean properties of particles in the single-phase system, so that $`\mathit{ϵ}=\mathrm{𝟎}`$. This yields $`\zeta `$ $`=`$ $`\zeta _\mathrm{m}\mathit{ϵ}\mathit{ϵ}:{\displaystyle \frac{1}{2}}{\displaystyle \underset{i,j}{}}\left\{\zeta 𝚿(\|𝒓_i𝒓_j\|)_\mathrm{m}\zeta _\mathrm{m}𝚿(\|𝒓_i𝒓_j\|)_\mathrm{m}\right\}`$ (A4) $`+\mathit{ϵ}\mathit{ϵ}:{\displaystyle \frac{1}{2}}{\displaystyle \underset{i=1}{\overset{N}{}}}{\displaystyle \underset{\stackrel{j=1}{i}}{\overset{N}{}}}{\displaystyle \underset{\stackrel{k=1}{i,j}}{\overset{N}{}}}\left\{\zeta 𝑼_1(\|𝒓_i𝒓_j\|)𝑼_1(\|𝒓_i𝒓_k\|)_\mathrm{m}\zeta _\mathrm{m}𝑼_1(\|𝒓_i𝒓_j\|)𝑼_1(\|𝒓_i𝒓_k\|)_\mathrm{m}\right\}`$ where $`𝚿(r)𝑼_{11}(r){\displaystyle \frac{1}{2}}𝑼_1(r)𝑼_1(r).`$ The thermal averages of functions of the stochastic particle positions can be re-written in Eq. (A4) in terms of the number density field in the monodisperse system, which at any instant is $`\rho (𝒓)=_{i=1}^N\delta ^{(3)}(𝒓𝒓_i)`$, where $`\delta ^{(3)}`$ is the three-dimensional Dirac delta function. Finally, the thermal average of a given stochastic variable $`\zeta `$ becomes $`\zeta `$ $`=`$ $`\zeta _\mathrm{m}{\displaystyle \frac{1}{2}}\mathit{ϵ}\mathit{ϵ}:{\displaystyle \text{d}𝒓^{}\text{d}𝒓^{\prime \prime }\left(\zeta \rho (𝒓^{})\rho (𝒓^{\prime \prime })_\mathrm{m}\zeta _\mathrm{m}\rho (𝒓^{})\rho (𝒓^{\prime \prime })_\mathrm{m}\right)𝚿(\|𝒓^{}𝒓^{\prime \prime }\|)}`$ (A6) $`+{\displaystyle \frac{1}{2}}\mathit{ϵ}\mathit{ϵ}:{\displaystyle \text{d}𝒓^{}\text{d}𝒓^{\prime \prime }\text{d}𝒓^{\prime \prime \prime }\left(\zeta \rho (𝒓^{})\rho (𝒓^{\prime \prime })\rho (𝒓^{\prime \prime \prime })_\mathrm{m}\zeta _\mathrm{m}\rho (𝒓^{})\rho (𝒓^{\prime \prime })\rho (𝒓^{\prime \prime \prime })_\mathrm{m}\right)𝑼_1(\|𝒓^{}𝒓^{\prime \prime }\|)𝑼_1(\|𝒓^{}𝒓^{\prime \prime \prime }\|)}.`$ A case of interest, for instance, is $`\zeta =\rho (\mathrm{𝟎})\rho (𝒓)`$. Then Eq. (A6) gives the perturbation, due to polydispersity, of the pair correlation function for total density as $`\rho (\mathrm{𝟎})\rho (𝒓)\rho (\mathrm{𝟎})\rho (𝒓)_\mathrm{m}={\displaystyle \frac{1}{2}}\mathit{ϵ}\mathit{ϵ}:{\displaystyle \text{d}𝒓^{}\text{d}𝒓^{\prime \prime }\left(\rho (\mathrm{𝟎})\rho (𝒓)\rho (𝒓^{})\rho (𝒓^{\prime \prime })_\mathrm{m}\rho (\mathrm{𝟎})\rho (𝒓)_\mathrm{m}\rho (𝒓^{})\rho (𝒓^{\prime \prime })_\mathrm{m}\right)𝚿(\|𝒓^{}𝒓^{\prime \prime }\|)}`$ (A7) $`+{\displaystyle \frac{1}{2}}\mathit{ϵ}\mathit{ϵ}:{\displaystyle \text{d}𝒓^{}\text{d}𝒓^{\prime \prime }\text{d}𝒓^{\prime \prime \prime }\left(\rho (\mathrm{𝟎})\rho (𝒓)\rho (𝒓^{})\rho (𝒓^{\prime \prime })\rho (𝒓^{\prime \prime \prime })_\mathrm{m}\rho (\mathrm{𝟎})\rho (𝒓)_\mathrm{m}\rho (𝒓^{})\rho (𝒓^{\prime \prime })\rho (𝒓^{\prime \prime \prime })_\mathrm{m}\right)𝑼_1(\|𝒓^{}𝒓^{\prime \prime }\|)𝑼_1(\|𝒓^{}𝒓^{\prime \prime \prime }\|)}`$ (A8) which depends, to second order in the standard deviation of species, on four- and five-point correlations in the monodisperse reference phase.
warning/0007/cond-mat0007261.html	ar5iv	text	# Application of the generalized three-dimensional Jordan-Wigner transformation to the bilayer Heisenberg antiferromagnet \[ ## Abstract We extend the definition of the Jordan-Wigner transformation to three dimensions using the generalization of ideas that were used in the two-dimensional case by one of the present authors. Under this transformation, the three-dimensional XY Hamiltonnian is transformed into a system of spinless fermions coupled to a gauge field with only two nonzero components. We calculate the flux per plaquette for the 3 elementary perpendicular plaquettes of a cubic lattice, and find that it is nonzero for only two of the plaquettes. We provide a simple interpretation for the average phase-per-plaquette being $`\pi `$ on the plaquettes where it is nonzero. Then we apply these new findings to the investigation of the Heisenberg bilayer antiferromagnet. PACS numbers: 75.10.Jm, 75.50.Ee preprint: ver. 2.2 \] In 1928, Jordan and Wigner introduced a transformation to write the spin operators $`S^{}`$ and $`S^z`$ in terms of Fermi operators in one dimension. For example, this transformation allows a mapping of the XY Hamiltonian in one dimension into a Hamiltonian of noninteracting spinless fermions. Trying to generalize this result to two dimensions proved to be difficult. But, a natural extension of the Jordan-Wigner (JW) transformation was introduced in 1993 by one of the present authors to study the interchain-coupling effect on the one-dimensional spin-1/2 Heisenberg antiferromagnets. Several other authors have given implicit ways of defining a generalized JW transformation in two dimensions. The generalization of the JW transformation has been attempted even in three dimensions by Huerta and Zanelli who introduced an implicit transformation, and by Kochmański whose transformation, however, does not preserve some of the spin-commutation relations. In this work, we follow the procedure of Ref. to extend the definition of the JW transformation in three dimensions. This paper is organized as follows. We start by reviewing the JW transformation in one and two dimensions. Then we define its three dimensional version. For the XY model, we calculate the components of the effective gauge field that results from the transformation, and the flux per plaquette as one spinless-fermion moves around it. Finally, we use this transformation to study the effect of the interlayer coupling on the Heisenberg bilayer. The JW transformation can be performed independent of any model Hamiltonian. In one dimension this is written as $`S_i^{}=c_ie^{i\varphi _i},S_i^z=c_i^{}c_i1/2.`$ (1) The phase $`\varphi _i=\pi _{j=0}^{i1}c_j^{}c_j=\pi _jw(i,j)n_j`$ where $`w(i,j)=\mathrm{\Theta }(ij)(1\delta _{i,j}),`$ (2) with $`\mathrm{\Theta }(ij)=1`$ for $`ji`$ and $`0`$ for $`j>i`$ being the (discrete) Heaviside step function. Note that the phase term in this transformation is obtained by summing over all those sites to the left of site $`i`$. Also, when we calculate spin commutation relations we find that the fermion-number operator, at the site with the lowest index and hosting one spin of the commutator, is contained only once in the resulting phase (refer to Ref. for more details). This assures that the spin commutation relations are preserved. In two dimensions, a generalized transformation based on this idea has been introduced in Ref.\[\]. This was written as: $`S_{i,j}^{}`$ $`=c_{i,j}e^{i\pi \varphi _{i,j}},S_{i,j}^z=c_{i,j}^{}c_{i,j}1/2,`$ (3) $`\varphi _{i,j}`$ $`={\displaystyle \underset{\alpha =0}{\overset{i1}{}}}{\displaystyle \underset{\beta =0}{\overset{\mathrm{}}{}}}n_{\alpha ,\beta }+{\displaystyle \underset{\beta =0}{\overset{j1}{}}}n_{i,\beta },`$ (4) where now a given site is specified by two indices, $`i`$ along the x-axis and $`j`$ along the y-axis. In this transformation, the phase at site $`(i,j)`$ is obtained by summing over all sites to the left of the vertical line passing through $`(i,j)`$ as well as all sites directly below $`(i,j)`$. This is a necessary and sufficient condition that allows the generalized JW transformation to preserve all spin commutation relations. The phase term in (4) can be written as: $`e^{i\pi _{𝐳=0}^{\mathrm{}}w(𝐱,𝐳)n(𝐳)}`$ with $`w(𝐱,𝐳)=`$ $`\mathrm{\Theta }(x_1z_1)(1\delta _{x_1,z_1})`$ (6) $`+\mathrm{\Theta }(x_2z_2)\delta _{x_1,z_1}(1\delta _{x_2,z_2});`$ $`𝐱=(x_1,x_2)`$ and $`𝐳=(z_1,z_2)`$ denoting sites on the square lattice. A necessary condition for the preservation of the spin commutation relation is that the following condition is satisfied $`e^{i\pi w(𝐱,𝐳)}=e^{i\pi w(𝐳,𝐱)}.`$ (7) A relation that is indeed fulfilled by Eq. (6). Next, consider the transformation of the XY model under (4), and calculate the flux per plaquette generated as one spinless fermion moves around an elementary plaquette. The phase change is found to be: $$(n_{i,j}n_{i+1,j})\pi .$$ (8) This result is more general than the one of Fradkin, who reported that the phase should be $`\pi `$, and also shows that Shaofeng’s claim that the flux per plaquette is zero is not true. Using the same procedure described above for one and two dimensions, we define the JW transformation in three dimensions by: $`S_{i,j,k}^{}`$ $`=c_{i,j,k}e^{i\pi \varphi _{i,j,k}},S_{i,j,k}^z=c_{i,j,k}^{}c_{i,j,k}1/2`$ (9) $`\varphi _{i,j,k}`$ $`={\displaystyle \underset{ϵ=0}{\overset{k1}{}}}{\displaystyle \underset{\alpha =0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{\beta =0}{\overset{\mathrm{}}{}}}n_{\alpha ,\beta ,ϵ}+{\displaystyle \underset{\alpha =0}{\overset{i1}{}}}{\displaystyle \underset{\beta =0}{\overset{\mathrm{}}{}}}n_{\alpha ,\beta ,k}`$ (11) $`+{\displaystyle \underset{\beta =0}{\overset{j1}{}}}n_{i,\beta ,k}`$ where now the phase is: $`\varphi (𝐱)={\displaystyle \underset{z=0}{\overset{\mathrm{}}{}}}w(𝐱,𝐳)n(𝐳)`$ (12) with $`w(𝐱,𝐳)=`$ $`\mathrm{\Theta }(x_3z_3)(1\delta _{x_3,z_3})`$ (15) $`+\mathrm{\Theta }(x_1z_1)\delta _{x_3,z_3}(1\delta _{x_1,z_1})`$ $`+\mathrm{\Theta }(x_2z_2)\delta _{x_1,z_1}\delta _{x_3,z_3}(1\delta _{x_2,z_2}),`$ $`𝐱=(x_1,x_2,x_3)`$ and $`𝐳=(z_1,z_2,z_3)`$ are now sites on the cubic lattice. It is straightforward to check again that Eq. (7) is satisfied, and that all spin commutation relations are preserved. To obtain the phase (12) we choose a plane going through the site where we want to define the transformation. Here the plane chosen is the one perpendicular to the vector $`(0,0,1)`$. Then we sum over those sites on this plane as in the two-dimensional transformation. Finally we sum over all sites that are behind this plane. As for the phase change that occurs when a spinless fermion completes a motion around a plaquette, we consider the elementary plaquettes in the $`xy`$, $`xz`$, and $`yz`$ planes. The flux per plaquette in the $`xy`$ plane turns out to be the same as in two dimensions, that is $$(n_{i,j,k}n_{i+1,j,k})\pi .$$ (16) For a spinless fermion starting at site $`(i,j,k)`$, $`n_{i,j,k}=1`$. Then for this fermion to go to site $`(i+1,j,k)`$, this later site has to be empty due to Pauli exclusion principle; thus $`n_{i+1,j,k}=0`$. In this approximation one finds that the phase is $`\pi `$. For a plaquette in the $`yz`$ plane the result is $$(n_{i,j,k+1}n_{i,j,k})\pi .$$ (17) Similarly, the phase per plaquette can be set to be approximately $`\pi `$. As for the phase change around an elementary plaquette in the $`xz`$ plane, it is found to be identically zero. The results for these fluxes per plaquette suggest that as a spinless fermion moves around a given plaquette, it will couple to a gauge field if the plaquette belongs to plane $`xy`$ or $`yz`$, but no coupling to a gauge field occurs in the remaining plane $`xz`$. This thus suggests that the effective magnetic field to which the spinless fermions couple has only two components. To show this, we consider the XY model written for coupled layers: $`H_{XY}`$ $`={\displaystyle \frac{J}{2}}{\displaystyle \underset{i,j,k}{}}S_{i,j,k}^{}S_{i+1,j,k}^++{\displaystyle \frac{J}{2}}{\displaystyle \underset{i,j,k}{}}S_{i,j,k}^{}S_{i,j+1,k}^+`$ (19) $`+{\displaystyle \frac{J_{}}{2}}{\displaystyle \underset{i,j,k}{}}S_{i,j,k}^{}S_{i,j,k+1}^++(\mathrm{H}.\mathrm{c}.)`$ where $`J`$ and $`J_{}`$ are respectively the intralayer and interlayer coupling constants. Under the JW transformation (11), the Hamiltonian (19) gives $`H_{XY}=`$ $`J{\displaystyle \underset{i,j,k}{}}c_{i,j,k}^{}c_{i+1,j,k}e^{i\pi _i\varphi _{i,j,k}}`$ (22) $`+J{\displaystyle \underset{i,j,k}{}}c_{i,j,k}^{}c_{i,j+1,k}`$ $`+J_{}{\displaystyle \underset{i,j,k}{}}c_{i,j,k}^{}c_{i,j,k+1}e^{i\pi _k\varphi _{i,j,k}}+(\mathrm{H}.\mathrm{c}).`$ Here $`_i\varphi _{i,j,k}=\varphi _{i+1,j,k}\varphi _{i,j,k}`$ and $`_k\varphi _{i,j,k}=\varphi _{i,j,k+1}\varphi _{i,j,k}`$ designate the discrete derivatives along the x and z-axes, respectively. This effective Hamiltonian describes the motion of spinless fermions coupled to a gauge field. The later field is given by $`A_i=_i\varphi _{i,j,k}`$ along the x-axis, $`A_j=_j\varphi _{i,j,k}`$ along the y-axis, and $`A_k=_k\varphi _{i,j,k}`$ along the z-axis. This leads to $`A_i=\pi [{\displaystyle \underset{\beta =0}{\overset{\mathrm{}}{}}}n_{i,\beta ,k}+{\displaystyle \underset{\beta =0}{\overset{j1}{}}}(n_{i+1,\beta ,k}n_{i,\beta ,k})],`$ (23) $`A_j=0,`$ (24) $`A_k=\pi [{\displaystyle \underset{\alpha =0}{\overset{i1}{}}}{\displaystyle \underset{\beta =0}{\overset{\mathrm{}}{}}}(n_{\alpha ,\beta ,k+1}n_{\alpha ,\beta ,k})`$ (25) $`+{\displaystyle \underset{\beta =0}{\overset{j1}{}}}(n_{i,\beta ,k+1}n_{i,\beta ,k})+{\displaystyle \underset{\alpha =0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{\beta =0}{\overset{\mathrm{}}{}}}n_{\alpha ,\beta ,k}].`$ (26) Note that the y-component of $`𝐀=(A_i,A_j,A_k)`$ is found to be identically zero. The effective magnetic field is obtained through the discrete curl of $`𝐀`$, that is $`𝐁(i,j,k)=\times 𝐀(i,j.k)`$, \[vectors $`𝐀`$ and $`𝐁`$ depend on the position $`(i,j,k)`$\]. The components of $`𝐁`$ which are given by $`B_i=\pi (n_{i,j,k+1}n_{i,j,k}),`$ (27) $`B_j=0,`$ (28) $`B_k=\pi (n_{i,j,k}n_{i+1,j,k}),`$ (29) are consistent with the results of the fluxes per plaquette in Eqs. (16) and (17), (the unit area being 1). Therefore the effective magnetic field to which the spinless fermions are coupled has two nonzero components only. In the second part of this paper we apply the above findings to the study of the Heisenberg bilayer antiferromagnet that has received a lot of attention in the past few years.- The use of the JW transformation and the bond-mean-field theory is advantageous in comparison with the Schwinger boson approach, modified spin-wave approaches, or the bond-operator mean-field theory, because no constraint on the number of particle per site is required for the JW spinless fermions. We would like to draw the attention to the bond-mean field approach which has so far given very good results in the case of the Heisenberg ladder, and coupled Heisenberg ladders. The bilayer antiferromagnet is a system that presents an order-disorder second order quantum transition at zero temperature, and may be of some relevance to some of the high temperature superconductors in their low-doping normal state. The Heisenberg model on a bilayer reads as follows $`H`$ $`=H_{XY}+J{\displaystyle \underset{i,j,k}{}}S_{i,j,k}^zS_{i+1,j,k}^z`$ (31) $`+J{\displaystyle \underset{i,j,k}{}}S_{i,j,k}^zS_{i,j+1,k}^z+J_{}{\displaystyle \underset{i,j,k}{}}S_{i,j,k}^zS_{i,j,k+1}^z.`$ Hamiltonian (31) is simplified using the approximation where the average flux per elementary plaquette is $`\pi `$ on the $`xy`$ and $`yz`$ planes, and zero on the $`xz`$ plane (as discussed above). To achieve this, we choose the following configuration: the phases are alternated $`\mathrm{}\pi 0\pi 0\mathrm{}`$ along the adjacent bonds on the y-axis, and zero on all bonds on the remaining axes. Note that the alternated phases could be put rather on the x-axis without changing the physical results because of gauge invariance. In the case of the bilayer system, the alternated phases cannot however be put on the z-axis because there is only one bond along that axis. Further simplification is done by using the bipartite character due to antiferromagnetic correlations. Then we decouple the interacting quartic terms by introducing the alternated magnetization parameter $`m_i=2S_{i,j,k}^z=2c_{i,j,k}^{}c_{i,j,k}1=(1)^im`$, and the bond parameters $`Q=c_{i,j,k}c_{i+1,j,k}^{}=c_{i,j,k}c_{i,j+1,k}^{}`$ within the planes and $`P=c_{i,j,k}c_{i,j,k+1}^{}`$ perpendicular to the planes. The Hamiltonian takes the following form in the reciprocal space: $`H`$ $`={\displaystyle \underset{𝐤}{}}(2J+J_{})m(e_𝐤^{}e_𝐤f_𝐤^{}f_𝐤)`$ (33) $`+{\displaystyle \underset{𝐤}{}}[iJ_1\mathrm{sin}k_x+\gamma (k_y,k_z)]e_𝐤^{}f_𝐤+\mathrm{H}.\mathrm{c}.`$ where $`\gamma (k_y,k_z)=J_1\mathrm{cos}k_y+J_1\mathrm{cos}k_z`$, $`e_𝐤=c_𝐤^A`$ (resp. $`f_𝐤=c_𝐤^B`$) is the fermion operator on the sublattice A (resp. B). The quantities $`J_1`$ and $`J_1`$ are given by $`J_1=J(1+2Q)`$ and $`J_1=J_{}(1+2P)`$. The energy spectrum takes the expression: $`E_\pm (k)=\pm \{(2J+`$ $`J_{})^2m^2`$ (35) $`+J_1^2\mathrm{sin}^2k_x+\gamma ^2(k_y,k_z)\}^{1/2}`$ leading to the ground-state energy per site: $`E_{GS}={\displaystyle \frac{2J+J_{}}{4}}m^2+2JQ^2+J_{}P^2{\displaystyle \frac{(d^3k)}{2}E_+}`$ (36) where by definition $`(d^3k)d𝐤/(2\pi )^3`$. Because of the periodic boundary conditions along the z-axis, $`J_{}`$ is counted twice. Thus we have to divide the value of $`J_{}`$ by 2 in all our equations, or multiply $`J_{}`$ by two in the final results wherever $`J_{}`$ appears. The parameters $`m`$, $`Q`$, and $`P`$ are calculated self-consistently through the equations: $`m={\displaystyle (d^3k)m(2J+J_{})/E_+}`$ (37) $`Q={\displaystyle (d^3k)\{J_1\mathrm{sin}^2k_x+\gamma (k_y,k_z)\mathrm{cos}k_y\}/4E_+}`$ (38) $`P={\displaystyle (d^3k)\gamma (k_y,k_z)\mathrm{cos}k_z/2E_+},`$ (39) obtained by minimizing the ground-state energy with respect to $`m`$, $`Q`$ and $`P`$. Two solutions for the magnetization $`m/2`$ are possible depending on the value of $`J_{}`$, Fig. 1. Let us first analyze the trivial solution $`m=0`$ of Eqs. (39). In this case the energy spectrum $`E_k=\pm \{J_1^2\mathrm{sin}^2k_x+\gamma ^2(k_y,k_z)\}^{1/2}`$ presents an energy gap at $`𝐤=(k_0,\pi ,0)`$ or $`(k_0,0,\pi )`$, ($`k_0=0,\pi `$): $`E_g=\|J_1J_1\|=\|J_{}(1+2P)J(1+2Q)\|,`$ (40) if $`J_1>J_1`$. This condition is deduced from $`\gamma (k_y,k_z)=0`$ which is essential for a vanishing gap. Numerically we find that the value above which an energy gap opens is $`J_{}=2\times 0.78J=1.56J`$, see Fig. 1. However, this does not mean that this is the critical value for the order-disorder second-order quantum transition. As for the finite $`m`$ solution, figure 1 shows that $`m0`$ for $`J_{}<J_c=2\times 2.1J=4.2J`$. The phase is ordered antiferromagnetically for $`J_{}<J_c`$, whereas it is a disordered quantum state for $`J_{}>J_c`$. The value $`J_c=4.2J`$ compares well to the result ($`4.48J`$) of the Schwinger boson approach and that ($`4.3`$) of the self-consistent spin-wave theory. But it disagrees with the Monte Carlo simulations result ($`2.5J`$) or the series expansion calculation ($`2.56J`$). Chubokov and Morr reported that while transverse spin-wave excitations are gapless (Goldstone modes) longitudinal spin-wave excitations are gapped for $`J_c>J_{}>0`$. Our result seems to indicate that a finite value of the interlayer coupling ($`J_{}=1.56J`$) would be needed for longitudinal fluctuations to become relevant. Note however that we have to consider gaussian fluctuations about our mean-field point to account for the spin-wave excitations in the ordered phase. We expect that fluctuations beyond the mean-field point will lead to a better estimate of $`J_c`$. The magnetization increases with $`J_{}`$, passes through a maximum, then decreases in agreement with previous results. As for the parameters $`P`$ and $`Q`$, when $`J_{}`$ increases $`P`$ saturates at $`0.5`$, whereas $`Q`$ decreases to 0, Fig. 1. For very large values of $`J_{}`$ this approach becomes less accurate. Perturbation expansion about the transverse dimers may be used. Finally, the result of the calculation of the ground-state energy, $`E_{GS}`$, is reported in Fig. 2 in order to support the validity of the present approach. To compare with the result of Weihong obtained using Ising and dimer expansions, $`\nu =E_{GS}/(4J+J_{})`$ is drawn as a function of $`\eta =J_{}/(J+J_{})`$. The maximum of $`\nu `$ at $`\eta =0.5`$ coincides with that reported by Weihong, and in general both results agree very well. In summary, we extended the definition of the JW transformation to three dimensions and gave a full analysis of its consequences in the case of the XY model. This model is found to be transformed to a system of spinless fermions that are coupled to a gauge field with two components only (one of the three components being zero). Accordingly, the flux per plaquette is nonzero for 2 elementary perpendicular plaquettes, but is zero for the third one. We gave a simple interpretation of the average flux being $`\pi `$ for those plaquettes where it is nonzero. These findings were used to study the Heisenberg bilayer antiferromagnet within the bond-mean field theory. We recalculated the order-disorder critical coupling and the ground-state energy. Very good agreement with previous results is found. Finally, the present three-dimensional JW transformation can be applied to models of three-dimensional quantum spin systems. For simplicity a cubic Bravais lattice was used in this work, but the transformation introduced here can be easily generalized to other Bravais lattices. B.B. aknowledges receipt of an undergraduate NSERC award. M.A. is grateful to Prof. Shegelski for his support and the useful discussions they had together, and to Prof. Hussein for his support. He thanks Prof. Thalmeier and Dr. Yuan for the helpful discussions they had during his stay at the Max-Planck Institute in Dresden. He is also grateful to this institute for the financial support.
warning/0007/hep-ph0007146.html	ar5iv	text	# Off-diagonal structure of neutrino mass matrix in see–saw mechanism and electron–muon–tau lepton universality ## I Restrictions on neutrino mixing angles Let us consider an off-diagonal Majorana mass matrix in $`(e,\mu ,\tau )`$ basis $$M=m_0\left(\begin{array}{ccc}\text{0}& a_{e\mu }& a_{e\tau }\\ a_{e\mu }& \text{0}& a_{\mu \tau }\\ a_{e\tau }& a_{\mu \tau }& \text{0}\end{array}\right)$$ (1) It is convenient to define neutrino mixing angles as follows $$\left(\begin{array}{c}\nu _e\hfill \\ \nu _\mu \hfill \\ \nu _\tau \hfill \end{array}\right)=U\left(\begin{array}{c}\nu _1\hfill \\ \nu _2\hfill \\ \nu _3\hfill \end{array}\right)$$ (2) where $$U=\left(\begin{array}{ccc}c_2c_3& c_2s_3& s_2e^{i\delta }\\ c_1s_3s_1s_2c_3e^{i\delta }& c_1c_3s_1s_2s_3e^{i\delta }& s_1c_2\\ s_1s_3c_1s_2c_3e^{i\delta }& s_1c_3c_1s_2s_3e^{i\delta }& c_1c_2\end{array}\right)$$ (3) with $`c_i=\mathrm{cos}\theta _i`$ and $`s_i=\mathrm{sin}\theta _i`$. We shall put$`\delta =0.`$ Due to the off-diagonal structure of the mass matrix (1), the following relations are derived : $`m_2`$ $`=`$ $`{\displaystyle \frac{\mathrm{cos}^2\theta _3\mathrm{tan}^2\theta _2}{\mathrm{sin}^2\theta _3\mathrm{tan}^2\theta _2}}m_1,\text{ }m_1+m_2+m_3=0`$ (4) $`\mathrm{cos}2\theta _1\mathrm{cos}2\theta _2\mathrm{cos}2\theta _3`$ $`=`$ $`{\displaystyle \frac{1}{2}}\mathrm{sin}2\theta _1\mathrm{sin}2\theta _2\mathrm{sin}2\theta _3\left(3\mathrm{cos}^2\theta _22\right)`$ (5) $`\mathrm{sin}2\theta _1\mathrm{cos}2\theta _2\mathrm{cos}2\theta _3{\displaystyle \frac{1}{2}}\mathrm{cos}2\theta _1\mathrm{sin}2\theta _2\mathrm{sin}2\theta _3\left(3\mathrm{cos}^2\theta _22\right)`$ $`=`$ $`a_2a_{\mu \tau }`$ (6) $`\mathrm{cos}2\theta _3\mathrm{sin}\theta _1\mathrm{sin}\theta _2\mathrm{cos}\theta _2+{\displaystyle \frac{1}{2}}\mathrm{cos}\theta _1\mathrm{sin}2\theta _3\mathrm{cos}\theta _2\left(3\mathrm{cos}^2\theta _22\right)`$ $`=`$ $`a_2a_{e\mu }`$ (7) $`\mathrm{cos}2\theta _3\mathrm{cos}\theta _1\mathrm{sin}\theta _2\mathrm{cos}\theta _2{\displaystyle \frac{1}{2}}\mathrm{sin}\theta _1\mathrm{cos}\theta _2\mathrm{sin}2\theta _3\left(3\mathrm{cos}^2\theta _22\right)`$ $`=`$ $`a_2a_{e\tau }`$ (8) where $$a_2=\frac{m_0}{m_2}\left(\mathrm{cos}^2\theta _3\mathrm{cos}^2\theta _2\mathrm{sin}^2\theta _2\right).$$ (9) We also give here the transition probabilities $`P\left(\nu _\mu \nu _\tau \right)`$ $`=`$ $`[{\displaystyle \frac{1}{4}}\mathrm{sin}^22\theta _1\mathrm{sin}^22\theta _3(1+\mathrm{sin}^2\theta _2)^2+\mathrm{sin}^22\theta _1\mathrm{sin}^2\theta _2`$ (13) $`+\mathrm{cos}2\theta _1\mathrm{sin}2\theta _1\mathrm{sin}2\theta _3\mathrm{cos}2\theta _3\mathrm{sin}\theta _2(1+\mathrm{sin}^2\theta _2)+\mathrm{cos}^22\theta _1\mathrm{sin}^22\theta _3\mathrm{sin}^2\theta _2\left]\mathrm{sin}^2\right({\displaystyle \frac{\mathrm{\Delta }m_{12}^2L}{4E}})`$ $`+\mathrm{sin}2\theta _1\mathrm{cos}^2\theta _2[(\mathrm{sin}2\theta _1\mathrm{cos}^2\theta _3\mathrm{sin}2\theta _1\mathrm{sin}^2\theta _2\mathrm{sin}^2\theta _3+\mathrm{sin}2\theta _3\mathrm{cos}2\theta _1\mathrm{sin}\theta _2)\mathrm{sin}^2\left({\displaystyle \frac{\mathrm{\Delta }m_{23}^2L}{4E}}\right)`$ $`+(\mathrm{sin}2\theta _1\mathrm{sin}^2\theta _3\mathrm{sin}2\theta _1\mathrm{sin}^2\theta _2\mathrm{cos}^2\theta _3\mathrm{sin}2\theta _3\mathrm{cos}2\theta _1\mathrm{sin}\theta _2)\mathrm{sin}^2\left({\displaystyle \frac{\mathrm{\Delta }m_{13}^2L}{4E}}\right)]`$ $`P\left(\nu _e\nu _\mu \right)`$ $`=`$ $`\left[\mathrm{sin}^22\theta _3\mathrm{cos}^2\theta _2\left(\mathrm{cos}^2\theta _1\mathrm{sin}^2\theta _1\mathrm{sin}^2\theta _2\right)+\mathrm{sin}2\theta _1\mathrm{sin}2\theta _3\mathrm{sin}\theta _2\mathrm{cos}^2\theta _2\right]\mathrm{sin}^2\left({\displaystyle \frac{\mathrm{\Delta }m_{12}^2L}{4E}}\right)`$ (16) $`+\mathrm{sin}2\theta _2\mathrm{sin}\theta _1[(\mathrm{cos}\theta _1\mathrm{cos}\theta _2\mathrm{sin}2\theta _3+\mathrm{sin}\theta _1\mathrm{sin}2\theta _2\mathrm{sin}^2\theta _3)\mathrm{sin}^2\left({\displaystyle \frac{\mathrm{\Delta }m_{23}^2L}{4E}}\right)`$ $`+(\mathrm{cos}\theta _1\mathrm{cos}\theta _2\mathrm{sin}2\theta _3+\mathrm{sin}\theta _1\mathrm{sin}2\theta _2\mathrm{cos}^2\theta _3)\mathrm{sin}^2\left({\displaystyle \frac{\mathrm{\Delta }m_{13}^2L}{4E}}\right)]`$ $`P\left(\nu _\mu \nu _\tau \right)`$ $`=`$ $`\left[\mathrm{sin}^22\theta _3\mathrm{cos}^2\theta _2\left(\mathrm{sin}^2\theta _1\mathrm{cos}^2\theta _1\mathrm{sin}^2\theta _2\right)\mathrm{sin}2\theta _1\mathrm{sin}2\theta _3\mathrm{cos}2\theta _3\mathrm{cos}^2\theta _2\mathrm{sin}\theta _2\right]\mathrm{sin}^2\left({\displaystyle \frac{\mathrm{\Delta }m_{12}^2L}{4E}}\right)`$ (19) $`+\mathrm{sin}2\theta _2\mathrm{cos}\theta _1[(\mathrm{sin}2\theta _3\mathrm{sin}\theta _1\mathrm{cos}\theta _2+\mathrm{cos}\theta _1\mathrm{sin}^2\theta _3\mathrm{sin}2\theta _2)\mathrm{sin}^2\left({\displaystyle \frac{\mathrm{\Delta }m_{23}^2L}{4E}}\right)`$ $`+(\mathrm{sin}2\theta _3\mathrm{sin}\theta _1\mathrm{cos}\theta _2+\mathrm{cos}\theta _1\mathrm{cos}^2\theta _3\mathrm{sin}2\theta _2)\mathrm{sin}^2\left({\displaystyle \frac{\mathrm{\Delta }m_{13}^2L}{4E}}\right)]`$ $$P\left(\nu _e\nu _e\right)=1\mathrm{cos}^4\theta _2\mathrm{sin}^22\theta _3\mathrm{sin}^2\left(\frac{\mathrm{\Delta }m_{12}^2L}{4E}\right)\mathrm{sin}^22\theta _2\mathrm{sin}^2\theta _3\mathrm{sin}^2\left(\frac{\mathrm{\Delta }m_{23}^2L}{4E}\right)\mathrm{sin}^22\theta _2\mathrm{cos}^2\theta _3\mathrm{sin}^2\left(\frac{\mathrm{\Delta }m_{13}^2L}{4E}\right)$$ (20) It may be noted that $$\mathrm{\Delta }m_{12}^2+\mathrm{\Delta }m_{23}^2+\mathrm{\Delta }m_{31}^2=m_2^2m_1^2+m_3^2m_2^2+m_1^2m_3^2=0$$ (21) ## II Extension of the standard model and neutrino mass matrix By a simple extension of the standard electroweak gauge group to $`GSU_L(2)\times U_e(1)\times U_\mu (1)\times U_\tau (1),`$ it was shown that the Majorana masses for light neutrinos are generated through diagrams shown in figure 1. Here $`\varphi ^{(i)}`$ and $`\mathrm{\Sigma }^{(i)}`$ are respectively three $`SU_L(2)`$ Higgs doublets and singlets with appropriate $`U_i(1)`$ quantum numbers; $`h`$’s and $`f`$’s are the corresponding Yukawa couplings. The symmetry is spontaneously broken by giving vacuum expectation values to Higgs bosons $`\varphi ^{(i)}`$ and $`\mathrm{\Sigma }^{(i)}`$: $`\varphi ^{(i)}=\frac{v_i}{\sqrt{2}}`$ and $`\mathrm{\Sigma }^{(i)}=\frac{\mathrm{\Lambda }_i}{\sqrt{2}}`$. For simplicity we shall take $`v_1=v_2=v_3=v`$ and $`\mathrm{\Lambda }_1=\mathrm{\Lambda }_2=\mathrm{\Lambda }_3=\mathrm{\Lambda }`$ (any difference can be absorbed in the corresponding Yukawa couplings $`h`$ and $`f)`$. We take $`\mathrm{\Lambda }v`$ so that $`X`$bosons which break the $`e\mu \tau `$ universality as well as the the Majorana mass term for heavy neutrinos $`N`$’s are superheavy. In order to simplify the calculation, we put $`f_{12}=f_{13}=f_{23}=f`$ (any differences can again be absorbed in $`h`$-couplings) and put $`f\mathrm{\Lambda }/\sqrt{2}=M_R`$. Thus finally we obtain the following off-diagonal mass matrix for light neutrinos $$M_\nu =\frac{v^2}{2M_R}\left(\begin{array}{ccc}0& h_1^{(2)}h_2^{(3)}& h_1^{(2)}h_3^{(1)}\\ h_1^{(2)}h_2^{(3)}& 0& h_2^{(3)}h_3^{(1)}\\ h_1^{(2)}h_3^{(1)}& h_2^{(3)}h_3^{(1)}& 0\end{array}\right)$$ (22) The Yukawa couplings here are arbitrary and different choices for them provide different predictions. We shall consider two choices, called $`A`$ and $`B,`$ with different mass hierarchies. For the choice $`A`$ we assume that the Yukawa couplings are proportional to the generation index of quarks $`h_1^{(2)}{\displaystyle \frac{v}{\sqrt{2}}}`$ $`=`$ $`{\displaystyle \frac{1}{K}}m_u`$ (23) $`h_2^{(3)}{\displaystyle \frac{v}{\sqrt{2}}}`$ $`=`$ $`{\displaystyle \frac{1}{K}}m_c`$ (24) $`h_3^{(1)}{\displaystyle \frac{v}{\sqrt{2}}}`$ $`=`$ $`{\displaystyle \frac{1}{K}}m_t`$ (25) where K is dimensioless parameter. Further we take $`m_u:m_c:m_t=\lambda ^6:\lambda ^4:1`$ as an order of magnitude relations. Then the mass matrix (22) can be written as Mν=m0( 0 λ6 λ2λ601λ210)subscript𝑀𝜈subscript𝑚0fragments 0 λ6 λ2fragmentsλ601fragmentsλ210M_{\nu}=m_{0}\left(\begin{tabular}[]{ccc}0&$\lambda^{6}$&$\lambda^{2}$\\ $\lambda^{6}$&0&1\\ $\lambda^{2}$&1&0\end{tabular}\right) (26) where $$m_0=\frac{\lambda ^4m_t^2}{K^2M_R}.$$ (27) In the first approximation, it has eigenvalues $`m_0(\pm 1,0)\left[m_2m_3\right]`$. The right-hand sides of Eqs. (6), (7) and (8) respectively become 1, $`\lambda ^6`$ and $`\lambda ^2.`$ Eqs. (7) and (8) can then only be satisfied if both $`\theta _2`$ and $`\theta _3`$ are small so that $`\mathrm{sin}\theta _{2,3}\theta _{2,3}`$ while Eqs. (5) and (6) are also then satisfied for $`\mathrm{sin}2\theta _11,`$ $`\mathrm{cos}2\theta _10,`$ $`\mathrm{sin}\theta _11/\sqrt{2},`$ $`\mathrm{cos}\theta _11/\sqrt{2}.`$ Writing Eqs. (6) and (7) in detail we have then $`{\displaystyle \frac{\theta _2}{\sqrt{2}}}+{\displaystyle \frac{1}{2}}{\displaystyle \frac{1}{\sqrt{2}}}\mathrm{sin}2\theta _3`$ $`=`$ $`\lambda ^6`$ (28) $`{\displaystyle \frac{\theta _2}{\sqrt{2}}}{\displaystyle \frac{1}{2}}{\displaystyle \frac{1}{\sqrt{2}}}\mathrm{sin}2\theta _3`$ $`=`$ $`\lambda ^2`$ (29) implying $`\mathrm{sin}2\theta _3`$ $``$ $`\sqrt{2}\lambda ^2\left(1\lambda ^4\right)`$ (30) $`\mathrm{sin}2\theta _2`$ $``$ $`2\theta _2=\sqrt{2}\lambda ^2\left(1+\lambda ^4\right)`$ (31) Finally from Eq. (4), $`m_21,`$ $`m_10`$ so that $`m_31`$. In fact diagonalization of the matrix (26) give $`m_3`$ $``$ $`m_0\left(\sqrt{1+\lambda ^4}+\lambda ^8\right)`$ (32) $`m_2`$ $``$ $`m_0\left(\sqrt{1+\lambda ^4}+\lambda ^8\right)`$ (33) $`m_1`$ $``$ $`2\lambda ^8m_0`$ (34) so that $`m_3\left\|m_2\right\|\left\|m_1\right\|`$ and $$\mathrm{\Delta }m_{12}^2=\mathrm{\Delta }m_{13}^2=m_0^2,\mathrm{\Delta }m_{23}^24m_0^2\lambda ^8.$$ (35) Finally from Eqs. (2) and (3) to leading orders in $`s_2`$ and $`s_3`$ $`\nu _1`$ $``$ $`\nu _es_3\left(c_1\nu _\mu s_1\nu _\tau \right)s_2\left(s_1\nu _\mu +c_1\nu _\tau \right)`$ (36) $`\nu _2`$ $``$ $`s_3\nu _e+c_3\left(c_1\nu _\mu s_1\nu _\tau \right)`$ (37) $`\nu _3`$ $``$ $`s_2\nu _e+c_2\left(s_1\nu _\mu +c_1\nu _\tau \right)`$ (38) showing that $`\nu _1`$ is primarily $`\nu _e`$ while $`\nu _2`$ and $`\nu _3`$ are primarily $`\left(c_1\nu _\mu s_1\nu _\tau \right)`$ and $`\left(s_1\nu _\mu +c_1\nu _\tau \right)`$ respectively. We now consider the choice $`B`$, where we assume $`h_1^{(2)}h_1^{(3)}h_2^{(3)}`$ and use the parametrization $`h_2^{(3)}`$ $`=`$ $`h\mathrm{cos}\theta ,h_1^{(3)}=h\mathrm{sin}\theta `$ (39) $`{\displaystyle \frac{h_1^{(3)}h_2^{(3)}}{h_1^{(2)}}}`$ $`=`$ $`h\sigma ,m_0={\displaystyle \frac{hh_1^{(2)}v^2}{2M_R}}`$ (40) where $`\sigma 1.`$ Then Mν=m0( 0cosθsinθcosθ0σsinθσ0)subscript𝑀𝜈subscript𝑚0 0cosθsinθcosθ0σsinθσ0M_{\nu}=m_{0}\left(\begin{tabular}[]{ccc}0&$\cos\theta$&$\sin\theta$\\ $\cos\theta$&0&$\sigma$\\ $\sin\theta$&$\sigma$&0\end{tabular}\right) (41) The diagonalization (neglecting $`\sigma ^2`$) gives $`m_{2,1}`$ $`=`$ $`m_0\left[\pm 1+{\displaystyle \frac{1}{2}}\sigma \mathrm{sin}2\theta \right]`$ (42) $`m_3`$ $`=`$ $`\left(m_1+m_2\right)=m_0\sigma \mathrm{sin}2\theta `$ (43) so that $`\left\|m_1\right\|\left\|m_2\right\|\left\|m_3\right\|`$ and $`\mathrm{\Delta }m_{12}^2=2\sigma \mathrm{sin}2\theta ,\mathrm{\Delta }m_{31}^2=\mathrm{\Delta }m_{32}^2=m_0^2`$ (44) We will take $`\theta _20`$ as before so that from Eq. (4), we must have $`\theta _3\frac{\pi }{4}`$ in order to have $`m_2=m_1.`$ Then from Eqs. (6), (7) and (8) in leading order, $`\theta =\theta _1.`$ In this case $`\nu _1`$ $``$ $`{\displaystyle \frac{1}{\sqrt{2}}}\left[\nu _e\left(\mathrm{cos}\theta _1\nu _\mu \mathrm{sin}\theta _1\nu _\tau \right)\right]`$ (45) $`\nu _2`$ $``$ $`{\displaystyle \frac{1}{\sqrt{2}}}\left[\nu _e+\left(\mathrm{cos}\theta _1\nu _\mu \mathrm{sin}\theta _1\nu _\tau \right)\right]`$ (46) $`\nu _3`$ $``$ $`\left(\mathrm{sin}\theta _1\nu _\mu +\mathrm{cos}\theta _1\nu _\tau \right)`$ (47) ## III Transition probabilities and conclusions For our choice $`A`$, $`\lambda `$ is expected to be of order $`0.22\mathrm{sin}\theta _c`$ ($`\theta _c`$ being the Cabibbo angle) so that from Eqs. (30), (31) and (35) $`\mathrm{sin}^22\theta _3`$ $``$ $`\mathrm{sin}^22\theta _22\lambda ^44.5\times 10^3`$ (48) $`\mathrm{\Delta }m_{23}^2`$ $`=`$ $`2\times 10^5m_0^2`$ (49) Thus with $$\mathrm{\Delta }m_{12}^2\mathrm{\Delta }m_{31}^2\mathrm{\Delta }m_{23}^2$$ (50) and neglecting terms of order $`s_3^4,s_2^2s_3^2,\mathrm{cos}2\theta _1s_2s_3,`$ we have from Eqs. (13)–(19) $`P\left(\nu _\mu \nu _\tau \right)\|_{\text{atm}}`$ $`=`$ $`\mathrm{sin}^22\theta _1\mathrm{cos}^2\theta _2\mathrm{cos}^2\theta _3\mathrm{sin}^2{\displaystyle \frac{\mathrm{\Delta }m_{23}^2R_a}{4E}}`$ (51) $``$ $`\mathrm{sin}^22\theta _1\mathrm{sin}^2{\displaystyle \frac{\mathrm{\Delta }m_{23}^2R_a}{4E}}`$ (52) $$P\left(\nu _e\nu _\mu \right)\|_{\text{LSND}}=\left[\mathrm{sin}^22\theta _3\mathrm{cos}^2\theta _1+\mathrm{sin}2\theta _1\mathrm{sin}\theta _2\mathrm{sin}2\theta _3+\mathrm{sin}2\theta _2\mathrm{sin}2\theta _3\mathrm{sin}\theta _1\mathrm{cos}\theta _1+\mathrm{sin}^22\theta _2\mathrm{sin}^2\theta _1\right]\mathrm{sin}^2\frac{\mathrm{\Delta }m_{12}^2R_{\text{LSND}}}{4E}$$ (53) and $`P\left(\nu _e\nu _\tau \right)`$ $`=`$ $`\left[\mathrm{sin}^2\theta _1\mathrm{sin}^22\theta _3\mathrm{sin}2\theta _1\mathrm{sin}\theta _2\mathrm{sin}2\theta _3\mathrm{sin}\theta _1\mathrm{cos}\theta _1\mathrm{sin}2\theta _2\mathrm{sin}2\theta _3+\mathrm{sin}^22\theta _2\mathrm{cos}^2\theta _1\right]\mathrm{sin}^2{\displaystyle \frac{\mathrm{\Delta }m_{12}^2R_s}{4E}}`$ (55) $`+\left[\mathrm{sin}\theta _1\mathrm{cos}\theta _1\mathrm{sin}2\theta _2\mathrm{sin}2\theta _3\right]\mathrm{sin}^2{\displaystyle \frac{\mathrm{\Delta }m_{23}^2R_s}{4E}}`$ Now with $`\theta _1\frac{\pi }{4}`$ and using Eq. (49) \[note that the coefficient of $`\mathrm{sin}^2\left(\mathrm{\Delta }m_{12}^2R_s/4E\right)`$ in Eq. (55) vanishes\], we have from Eqs. (53) and (55) $`P\left(\nu _e\nu _\mu \right)\|_{\text{LSND}}`$ $``$ $`\mathrm{sin}^22\theta _{\text{eff}}\mathrm{sin}^2{\displaystyle \frac{\mathrm{\Delta }m_{12}^2R_{\text{LSND}}}{4E}}`$ (56) $`P\left(\nu _e\nu _\tau \right)\|_{\text{solar}}`$ $``$ $`\mathrm{sin}^22\stackrel{~}{\theta }_{\text{eff}}\mathrm{sin}^2{\displaystyle \frac{\mathrm{\Delta }m_{23}^2R_s}{4E}}`$ (57) where $`\mathrm{sin}^22\theta _{\text{eff}}`$ $``$ $`2\left(4.5\right)\times 10^310^2`$ (58) $`\mathrm{sin}^22\stackrel{~}{\theta }_{\text{eff}}`$ $``$ $`{\displaystyle \frac{1}{2}}\left(4.5\right)\times 10^3=2.25\times 10^3`$ (59) Thus with $`m_0^20.3`$ eV$`^2,`$ $`\mathrm{\Delta }m_{23}4\lambda ^8m_0^22\times 10^5m_0^26\times 10^6`$ eV<sup>2</sup> the LSND data and solar neutrino oscillations are explained, the latter with SMA MSW solution. Finally with $`m_0\sqrt{0.3}`$ eV, $`\lambda ^4=4.5\times 10^3`$ and $`m_t175`$ GeV, we obtain from Eq. (27) $`K^2M_R10^{11}`$ GeV, giving the mass scale at which $`e\mu \tau `$ universality is broken and the scale associated with superheavy Majorana neutrinos. In the version of the model we have considered $`L_\tau L_\mu L_e`$ number is , however, conserved while in the particular version of the Zee model mentioned earlier as well as in our model $`B`$, it is the $`L_eL_\mu L_\tau `$ number which is conserved. We now consider the predictions of our version $`B`$ for which $$m_0^2=\mathrm{\Delta }m_{31}^2=\mathrm{\Delta }m_{32}^2\mathrm{\Delta }m_{12}^2$$ (60) and $`\theta _3\frac{\pi }{4},`$ $`\theta _20`$ so that neglecting $`\mathrm{sin}^22\theta _2`$ and $`\mathrm{cos}2\theta _3,`$ we obtain in the leading order from Eqs. (13) and (20) $`P\left(\nu _e\nu _\tau \right)\|_{\text{atm}}`$ $``$ $`\mathrm{sin}^22\theta _1\mathrm{sin}^2\left({\displaystyle \frac{\mathrm{\Delta }m_{32}^2R_a}{4E}}\right)`$ (61) $`P\left(\nu _e\nu _e\right)\|_{\text{solar}}`$ $``$ $`1\mathrm{sin}^22\theta _3\mathrm{sin}^2\left({\displaystyle \frac{\mathrm{\Delta }m_{12}^2R_s}{4E}}\right)`$ (62) Thus atmospheric neutrino experimental data is explained with $`\mathrm{\Delta }m_{32}^2m_0^210^3`$ eV<sup>2</sup> and $`\theta _1\frac{\pi }{4}`$ while the Eq. (62) is consistent with the large angle \[$`\mathrm{sin}^22\theta _31`$\] vacuum or MSW solution in solar neutrino experiments. Here $`m_3,`$ the mass of the lightest neutrino consisting mainly of $`\nu _\mu `$ and $`\nu _\tau `$ \[cf. Eq. (47)\] is given by $`m_3\sigma \mathrm{sin}2\theta _1m_0=\sigma m_0={\displaystyle \frac{\mathrm{\Delta }m_{12}^2}{2m_0}}={\displaystyle \frac{\mathrm{\Delta }m_{\text{solar}}^2}{2\sqrt{\mathrm{\Delta }m_{\text{atm}}^2}}}.`$ To conclude by considering a simple extention of the standard model in which $`\left(e\mu \tau \right)`$ universality is not conserved, we have presented a scenario within the framework of see–saw mechanism in which the neutrino mass matrix is strictly off-diagonal in the flavor basis. Further we have shown that a version of this scenario can accomodate the atmospheric $`\nu _\mu \nu _\tau `$ neutrino oscillations and large angle vacuum or MSW solution in solar neutrino experiments while another version is compatible with small angle MSW solution of solar neutrino oscillations and $`\nu _\mu \nu _e`$ oscillations claimed by the LSND collaboration.
warning/0007/math0007181.html	ar5iv	text	# A birational invariant for algebraic group actions ## 1. Introduction Let $`G`$ be an algebraic group and let $`X`$ be a smooth projective $`G`$-variety (i.e., an algebraic variety with a $`G`$-action) defined over an algebraically closed base field of characteristic zero. It is shown in \[RY<sub>1</sub>\] that for each finite abelian subgroup $`H`$, the presence of an $`H`$-fixed points is a birational invariant of $`X`$ as a $`G`$-variety. In other words, if $`X`$ and $`Y`$ are birationally isomorphic smooth projective $`G`$-varieties and $`H`$ is a finite abelian subgroup of $`G`$ then $`X^H\mathrm{}`$ iff $`Y^H\mathrm{}`$. (Note that only nontoral finite abelian subgroups $`H`$ are of interest in this setting. If $`H`$ lies in a torus $`T`$ of $`G`$ then the Borel Fixed Point Theorem says that $`X^H`$ can never be empty.) In \[RY<sub>1</sub>\]\[RY<sub>2</sub>\] and \[RY<sub>3</sub>\] we used this fact to study the geometry of $`G`$-varieties (their essential dimensions, splitting degrees, etc.) and properties of related algebraic objects (field extensions, division algebras, octonion algebras, etc.) In this paper we will associate (under additional assumptions on $`X`$, $`G`$ and $`H`$) a more subtle invariant $`i(X,x,H)`$ to a point $`xX^H`$; the precise definition is given in Section 4. Our main result about $`i(X,x,H)`$ is stated below. Recall that the rank of a finite abelian group $`H`$ is the minimal number of generators of $`H`$ (see Section 2) and that a $`G`$-variety $`X`$ is called generically free if $`\mathrm{Stab}(x)=\{1\}`$ for $`x`$ in general position in $`X`$ (see Section 3). ###### Theorem 1.1. Let $`G`$ be an algebraic group of dimension $`d`$, $`H`$ be a finite abelian subgroup of $`G`$ of rank $`r`$, and $`X`$, $`Y`$ be birationally isomorphic smooth projective irreducible generically free $`G`$-varieties of dimension $`d+r`$. Assume that $`\mathrm{Stab}(x)`$ is finite for every $`xX^H`$ and $`\mathrm{Stab}(y)`$ is finite for every $`yY^H`$. Then for every $`xX^H`$ there exists a $`yY^H`$ such that $`i(Y,y,H)=i(X,x,H)`$. Informally speaking, the presence of $`H`$-fixed points $`x`$ with a prescribed value $`i(X,x,H)`$ (on a suitable model) is a birational invariant of $`X`$ as a $`G`$-variety. Our proof of Theorem 1.1, presented in Section 6, relies on canonical resolution of singularities. Note that Theorem 1.1 remains valid even if $`X`$ and $`Y`$ are not assumed to be irreducible; see Remark 6.4. We give two applications of Theorem 1.1. ### A birational classification of linear representations For our first application recall that by the no-name lemma any two generically free linear representations of a given algebraic group $`G`$ are stably isomorphic as $`G`$-varieties; see, e.g., \[P, 1.5.3\]. Thus it is natural to try to classify such representations up to birational isomorphism. This problem was proposed by E. B. Vinberg \[PV<sub>2</sub>, pp. 494-496\]; see also \[P, 1.5.1\]. P. I. Katsylo has subsequently stated the following conjecture: ###### Conjecture 1.2. (\[K\]; see also \[P, 1.5.10\]) Let $`V`$ and $`W`$ be generically free linear representations of an algebraic group $`G`$. If $`dim(V)=dim(W)`$ then $`V`$ and $`W`$ are birationally isomorphic as $`G`$-varieties. In this paper we will establish the following birational classification of faithful linear representations of a diagonalizable group. Recall that every diagonalizable group $`G`$ can be uniquely written in the form $`G=𝔾_m(n_1)\times \mathrm{}\times 𝔾_m(n_r)`$ such that (1.1) $`𝔾_m(n_1)\mathrm{}𝔾_m(n_r)`$ $`\text{and each }n_i=0\text{ or }2,`$ where $`𝔾_m=k^{}=𝔾_m(0)`$ denotes the 1-dimensional torus, $`𝔾_m(n)/n`$ is the $`n`$-torsion subgroup of $`𝔾_m`$, and $`𝔾(a)𝔾(b)`$ iff either $`a`$ divides $`b`$ or $`a=b=0`$. Recall that a representation of a diagonalizable group is faithful iff it is generically free and that any such representation has dimension $`r`$. ###### Theorem 1.3. Let $`G`$ be a diagonalizable group, as in (1). (a) If $`dr+1`$ then any two faithful $`d`$-dimensional linear representations of $`H`$ are birationally equivalent. (b) If $`n_1=0`$ or $`2`$ then any two faithful $`r`$-dimensional linear representations of $`H`$ are birationally equivalent. (c) If $`n_13`$ then $`H`$ has exactly $`\varphi (n_1)/2`$ birational equivalence classes of faithful $`r`$-dimensional representations. Here $`\varphi `$ denotes the Euler $`\varphi `$-function. In particular, Conjecture 1.2 fails for $`G`$ if and only if $`n_1=5`$ or $`7`$. The birational equivalence classes of faithful linear representations of $`G`$ are explicitly described in Theorem 7.1. Later in Section 7 we will show that Conjecture 1.2 also fails for a class of nonabelian groups $`G`$. On the other hand, we remark that P. I. Katsylo \[K\] proved Conjecture 1.2 for $`G=\mathrm{SL}_2`$, $`G=\mathrm{PGL}_2`$, $`G=S_n`$ ($`n4`$) and that many interesting cases remain open, including $`G=\mathrm{S}_n`$ ($`n5`$) and $`G`$ = arbitrary connected semisimple group. ### Birational equivalence of quantum tori Our second application is based on the fact that a generically free $`\mathrm{PGL}_n`$-varieties $`X`$ with $`k(X)^{\mathrm{PGL}_n}=K`$ are in 1-1 correspondence with division algebras $`D`$ with center $`K`$; see, e.g., \[Se, X.5\] or \[RY<sub>2</sub>, Section 3\]. Thus Theorem 1.1 (with $`G=\mathrm{PGL}_n`$) will sometimes allow us to prove that certain division algebras are not isomorphic to each other over $`k`$. Let $`\omega _1,\mathrm{},\omega _r`$ be roots of unity and let $`R(\omega _1,\mathrm{},\omega _r)`$ be the $`k`$-algebra $`k\{x_1^{\pm 1},\mathrm{},x_{2r}^{\pm 1}\}`$, where $`x_{2i1}x_{2i}=\omega _ix_{2i}x_{2i1}`$ for $`i=1,\mathrm{},r`$ and all other pairs of variables commute. Denote the algebra of quotients of $`R(\omega _1,\mathrm{},\omega _r)`$ by $`Q(\omega _1,\mathrm{},\omega _r)`$. Note that $`Q(\omega _1,\mathrm{},\omega _r)`$ is obtained from $`R(\omega _1,\mathrm{},\omega _r)`$ by adjoining the inverses of all central elements and that $`Q(\omega _1,\mathrm{},\omega _r)`$ is a finite-dimensional division algebra (in fact, it is a tensor product of symbol algebras). M. Lorenz \[Lo<sub>1</sub>, Proposition 1.3\] showed that $`Q(\omega )`$ and $`Q(\omega ^{})`$ are isomorphic as $`k`$-algebras if and only if $`\omega ^{}=\omega ^{\pm 1}`$. In Section 8 we will give a geometric proof of the following variant of this result. ###### Theorem 1.4. Suppose $`\omega _i`$ is a primitive $`n_i`$th root of unity, $`n_i`$ divides $`n_{i+1}`$ for $`i=1,\mathrm{},r1`$, $`n_12`$, and $`(m_i,n_i)=1`$. Then $`Q(\omega _1,\mathrm{},\omega _r)`$ and $`Q(\omega _1^{m_1},\mathrm{},\omega _r^{m_r})`$ are isomorphic as $`k`$-algebras if and only if $`m_1\mathrm{}m_r\pm 1(modn_1)`$. The algebras $`Q(\omega _1,\mathrm{},\omega _r)`$ and $`Q(\omega _1^{m_1},\mathrm{},\omega _r^{m_r})`$ have a common center $`K=k(x_1^{n_1},\mathrm{},x_{2r}^{n_{2r}})`$. Note that these algebras may be $`k`$-isomorphic but not $`K`$-isomorphic, i.e., not Brauer equivalent. More precisely, $`Q(\omega _1,\mathrm{},\omega _r)`$ and $`Q(\omega _1^{m_1},\mathrm{},\omega _r^{m_r})`$ are Brauer equivalent iff $`m_i1(modn_i)`$ for every $`i=1,\mathrm{},r`$. It is natural to think of $`R(\omega _1,\mathrm{},\omega _r)`$ and $`Q(\omega _1,\mathrm{},\omega _r)`$ as, respectively, the “coordinate ring” and the “function field” of a quantum torus. Using this terminology, one may view Theorem 1.4 as a result about birational isomorphism classes of quantum tori. Finally we remark that M. Lorenz has communicated to us a proof of Theorem 1.4 based on the techniques of \[Lo<sub>1</sub>\] and \[Lo<sub>2</sub>\]. His argument works in arbitrary characteristic. ### Acknowledgements We are grateful to P. I. Katsylo, M. Lorenz and V. L. Popov fo helpful communications related to the subject matter of this paper. The second author warmly thanks Institute of Mathematics of Hebrew University for its hospitality during 1999/2000. ## 2. Linear algebra in abelian groups Recall that any finitely generated abelian group $`(A,+)`$ can be written in the form (2.1) $$\begin{array}{c}A/n_1\times \mathrm{}\times /n_r,\text{ where each }n_i=0\text{ or }2\text{ and}\\ \text{for every }i=1,\mathrm{},r1\text{ either }n_i=n_{i+1}=0\text{ or }n_i\text{ divides }n_{i+1}\text{.}\end{array}$$ Here $`r`$ and $`n_1,\mathrm{},n_r`$ are uniquely determined by the isomorphism type of $`A`$. We shall refer to the integer $`r`$ as the the rank of $`A`$; equivalently, the rank of $`A`$ equals the minimal possible number of generators of $`A`$. Recall that if $`B`$ is an abelian group then the dual group $`B^{}`$ is defined as $`\mathrm{Hom}(B,/)`$; we will often identify $`/`$ with the multiplicative group of roots of unity in $`k`$. The finitely generated group $`A`$ of (2.1) and the diagonalizable group $`G`$ of (1) are dual to each other. The rank of a diagonalizable group $`G`$ is defined as the rank of the finitely generated group $`G^{}`$ (in particular, the group $`G`$ of (1) has rank $`r`$). Note that this is consistent with the way we defined rank for a finitely generated group: indeed, if $`A`$ is both diagonalizable and finitely generated, i.e., is finite abelian, then $`A`$ and $`A^{}`$ are isomorphic, so that their ranks coincide. ### Skew-symmetric powers We will write $`^d(A)`$ for the $`d`$th skew-symmetric power of $`A`$, viewed as a $``$-module. The proof of the following lemma is elementary; we leave it as an exercise for the reader. ###### Lemma 2.1. Let $`A`$ be a finitely generated abelian group as in (2.1). Then (a) $`^r(A)/n_1`$. (b) $`^d(A)=(0)`$, if $`dr+1`$. ∎ ###### Definition 2.2. Let $`A`$ be a finite abelian group. Let $`\omega :A\times A/`$ be a $``$-bilinear form. As usual, we shall say that (a) $`\omega `$ is alternating if $`\omega (a,a)=0`$ for every $`aA`$, (b) $`\omega `$ is non-degenerate if $`\omega (a,)`$ is not identically zero for any $`aA\{0\}`$, (c) $`\omega `$ is symplectic if it is both alternating and non-degenerate. ###### Lemma 2.3. Let $`A`$ be a finite abelian group of rank $`r`$, $`\omega `$ be a symplectic form of $`A`$, and $`\psi `$ be an $`\omega `$-preserving automorphism $`AA`$. Then (a) $`^{2r}\psi `$ is the trivial automorphism of $`^{2r}(A)`$ and (b) $`^{2r}\psi ^{}`$ is the trivial automorphism of $`^{2r}(A^{})`$. ###### Proof. (a) It is well-known that $`A`$ can be written in the form $`A=A_0A_0^{}`$ such that (2.2) $$\omega ((a,a^{}),(b,b^{}))=a^{}(b)b^{}(a)$$ for any $`a,bA_0`$ and $`a^{},b^{}A_0^{}`$; see, e.g., \[TA, Theorem 4.1\]. Write $`A_0`$ as $`/n_1\times \mathrm{}\times /n_r`$, where $`n_i`$ divides $`n_{i+1}`$ for $`i=1,\mathrm{},r1`$ and $`n_12`$. Let $`e_iA_0`$ be a generator of the of the $`i`$th factor, and let $`f_iA_0^{}`$ be given by $`f_i:A_0/n_i/`$, where the first map the projection to the $`i`$th factor, and the second map takes $`e_i`$ to $`1/n_i`$. Then every $`aA`$ can be written in the form $`a=_{i=1}^n(\alpha _{2i1}e_i+\alpha _{2i}f_i)`$, where $`\alpha _{2i1},\alpha _{2i}`$, and (2.2) translates into $$\omega [\underset{i=1}{\overset{n}{}}(\alpha _{2i1}e_i+\alpha _{2i}f_i),\underset{i=1}{\overset{n}{}}(\beta _{2i1}e_i+\beta _{2i}f_i)]=\underset{i=1}{\overset{n}{}}\frac{1}{n_i}(\alpha _{2i}\beta _{2i1}\alpha _{2i1}\beta _{2i}).$$ Suppose $$\begin{array}{c}\psi (e_1)=c_{11}e_1+c_{12}f_1+\mathrm{}+c_{1,2r1}e_r+c_{1,2r1}f_r,\hfill \\ \psi (f_1)=c_{21}e_1+c_{22}f_1+\mathrm{}+c_{2,2r1}e_r+c_{2,2r}f_r,\hfill \\ \mathrm{}\hfill \\ \psi (e_r)=c_{2r1,1}e_1+c_{2r1,2}f_1+\mathrm{}+c_{2r1,2r1}e_r+c_{2r1,2r}f_r,\hfill \\ \psi (f_r)=c_{2r,1}e_1+c_{2r,2}f_1+\mathrm{}+c_{2r,2r1}e_r+c_{2r,2r}f_r,\hfill \end{array}$$ where $`C=(c_{ij})_{i,j=1,\mathrm{},2r}\mathrm{M}_n()`$. Since $`\lambda =e_1f_1\mathrm{}e_rf_r`$ generates $`^{2r}(A^{})/n_1`$ and $`[^{2r}\psi ](\lambda )=det(C)\lambda `$, it is sufficient to show that $`det(C)=1(modn_1)`$. The condition that $`\psi `$ preserves $`\omega `$ translates into $`CJC^t=J(mod1)`$, where $`C^t`$ is the transpose of $`C`$ and $$J=\left(\begin{array}{ccccc}0& 1/n_1& \mathrm{}& 0& 0\\ 1/n_1& 0& \mathrm{}& 0& 0\\ \mathrm{}& & \mathrm{}& & \mathrm{}\\ 0& 0& & 0& 1/n_r\\ 0& 0& & 1/n_r& 0\end{array}\right).$$ In other words, (2.3) $$CJC^t=J+N,$$ where $`N`$ is a skew-symmetric integral matrix. We shall deduce the desired equality, $`det(C)=1(modn_1)`$, by computing the Pfaffian on both sides of (2.3). On the one hand (2.4) $$\mathrm{Pf}(CJC^t)=det(C)\mathrm{Pf}(J)=(1)^rdet(C)\frac{1}{n_1\mathrm{}n_r};$$ see, e.g., \[Lang, XIV, 10, Theorem 7\]. On the other hand, suppose $`X=(x_{ij})`$, where $`x_{ji}=x_{ij}`$ for $`1i,j2r`$. Then $`\mathrm{Pf}(X)[x_{ij}\|\mathrm{\hspace{0.17em}1}i<j2r]`$ has degree 1 in every $`x_{ij}`$, where $`i<j`$. (Indeed, $`det(X)`$ has degree 2 in every $`x_{ij}`$, and $`\mathrm{Pf}(X)^2=det(X)`$.) Consequently, (2.5) $$\mathrm{Pf}(J+N)=\mathrm{Pf}(J)+\frac{z}{n_2\mathrm{}n_r}=(1)^r\frac{1}{n_1\mathrm{}n_r}+\frac{z}{n_2\mathrm{}n_r},$$ where $`z`$ is an integer. (Here we used the fact that $`n_i`$ divides $`n_{i+1}`$ for every $`i=1,\mathrm{},r1`$.) Putting (2.3), (2.4) and (2.5) together, we see that $`det(C)=1+(1)^rn_1z`$, i.e., $`det(C)=1(modn_1)`$, as claimed. (b) Let $`i:AA^{}`$ be the isomorphism $`ai_a`$, where $`i_a(b)=\omega (a,b)`$. It is easy to see that the automorphism $`\psi ^{}:A^{}A^{}`$ preserves the symplectic form $`\omega ^{}`$ on $`A^{}`$ given by $`\omega ^{}(a^{},b^{})=\omega (i^1a^{},i^1b^{})`$. Applying part (a) to $`\psi ^{}`$, we obtain the desired result. ∎ ### Elementary operations Let $`A`$ be an abelian group. We will say that two $`d`$-tuples $`(a_1,\mathrm{},a_d)`$ and $`(b_1,\mathrm{},b_d)A^d`$ are related by an elementary operation if $`(b_1,\mathrm{},b_d)=(a_1,\mathrm{},a_{i1},a_i+ma_j,a_{i+1},\mathrm{},a_d)`$ for some $`ij`$ and $`m`$. ###### Lemma 2.4. Let $`a_1,\mathrm{},a_d`$ and $`b_1,\mathrm{},b_d`$ be two sets of generators for an abelian group $`A`$. Then $`a=(a_1,\mathrm{},a_d)`$ and $`(b_1,\mathrm{},b_d)`$ are related by a sequence of elementary operations if and only if $`a_1\mathrm{}a_d=b_1\mathrm{}b_d`$ in $`^d(A)`$. ###### Proof. It is clear that if $`(a_1,\mathrm{},a_d)`$ and $`(b_1,\mathrm{},b_d)`$ are related by a sequence of elementary operations then $`a_1\mathrm{}a_d=b_1\mathrm{}b_d`$. We will prove the converse by induction on $`d`$. If $`d=1`$ there is nothing to prove, since $`^1(A)=A`$. For the induction step, assume $`d2`$ and $`A(/n_1)\times \mathrm{}\times (/n_r)`$ is as in (2.1) Here $`r=\mathrm{rank}(A)d`$, since we are assuming $`A`$ is generated by $`d`$ elements. A $`d`$-tuple of generators $`(a_1,\mathrm{},a_d)`$ of $`A`$ can now be represented by a $`d\times r`$-matrix $`a=(a_{ij})`$ whose $`i`$th row is $`a_i`$. Elements of the $`j`$th column of this matrix lie in $`/n_j`$. Elementary operation on such matrices amount to adding an integer multiple of the $`j`$th row to the $`i`$th row or some $`ij`$. Elementary operations allow us to perform the Euclidean algorithm in the last column of $`(a_{ij})`$. Since $`a_{1r},\mathrm{},a_{dr}`$ generate $`/n_r`$, after a sequence of elementary operations, we may assume that $$(a_{ij})=\left(\begin{array}{cccc}& & & 0\\ & X& & \mathrm{}\\ & & & 0\\ & \mathrm{}& & 1\end{array}\right),$$ where $`X`$ is a $`(d1)\times (r1)`$-matrix. Since the rows of $`(a_{ij})`$ generate $`A`$, the rows of $`X`$ generate $`\overline{A}=(/n_1)\times \mathrm{}\times (/n_{r1})`$. Thus after performing additional elementary operations, we may assume $$(a_{ij})=\left(\begin{array}{cccc}& & & 0\\ & X& & \mathrm{}\\ & & & 0\\ 0& \mathrm{}& 0& 1\end{array}\right),\text{and similarly}(b_{ij})=\left(\begin{array}{cccc}& & & 0\\ & Y& & \mathrm{}\\ & & & 0\\ 0& \mathrm{}& 0& 1\end{array}\right).$$ In other words, we may assume $`(a_1,\mathrm{},a_{d1})`$ and $`(b_1,\mathrm{},b_{d1})`$ are $`(d1)`$-tuples of generators in $`\overline{A}`$ and $`a_d=b_d`$ is the generator $`1/n_r`$. We claim that (2.6) $$a_1\mathrm{}a_{d1}=b_1\mathrm{}b_{d1}\text{in }^{(d1)}(\overline{A})\text{.}$$ Indeed, if $`dr+1`$, this is obvious, since $`^{d1}(\overline{A})=(0)`$; see Lemma 2.1(b). If $`d=r`$ then the isomorphism $`^d(A)/n_1`$ identifies $`a_1\mathrm{}a_d`$ with $`det(a_{ij})(modn_1)`$, and the isomorphism $`^{(d1)}(\overline{A})/n_1`$ identifies $`a_1\mathrm{}a_{d1}`$ with $`det(X)(modn_1)`$. Since $`a_1\mathrm{}a_d=b_1\mathrm{}b_d`$, we know that $`det(a_{ij})=det(b_{ij})(modn_1)`$; hence, $`det(X)=det(Y)(modn_1)`$, and (2.6) follows. Now by the induction assumption $`(a_1,\mathrm{},a_{d1})`$ and $`(b_1,\mathrm{},b_{d1})`$ are related by a sequence of elementary transformations. Since $`a_d=b_d`$, so are $`(a_1,\mathrm{},a_d)`$ and $`(b_1,\mathrm{},b_d)`$, as desired. ∎ ###### Corollary 2.5. Let $`a_1,\mathrm{},a_d`$ and $`b_1,\mathrm{},b_d`$ be two sets of generators for an abelian group $`A`$. Then the following conditions are equivalent. (a) There exists a matrix $`N=(n_{ij})\mathrm{GL}_d()`$ such that $`b_i=n_{i1}a_1+\mathrm{}+n_{id}a_d`$ for $`i=1,\mathrm{},d`$. (b) $`a_1\mathrm{}a_d=\pm b_1\mathrm{}b_d`$ in $`^d(A)`$. ###### Proof. The implication (a) $``$ (b) is obvious. To prove the converse, note that we may assume without loss of generality that $`a_1\mathrm{}a_d=b_1\mathrm{}b_d`$; indeed, if $`a_1\mathrm{}a_d=b_1\mathrm{}b_d`$ then we can simply replace $`(a_1,a_2,\mathrm{},a_d)`$ by $`(a_1,a_2,\mathrm{},a_d)`$. Now Lemma 2.4 says, $`(b_1,\mathrm{},b_d)`$ is obtained from $`(a_1,\mathrm{},a_d)`$ by a sequence of elementary operations. Each elementary operation relates the two sets of generators as in (a), with $`N`$ = elementary matrix $`\mathrm{SL}()`$. Multiplying these matrices we obtain the desired conclusion. ∎ ## 3. $`H`$-slices ###### Definition 3.1. Let $`G`$ be an algebraic group and $`X`$ be a $`G`$-variety. We will call a locally closed subvariety $`V`$ of $`X`$ a slice at $`xV`$ if $`X`$ and $`V`$ are smooth at $`x`$ and $`T_x(X)=T_x(Gx)T_x(V)`$. If, moreover, $`V`$ is invariant under the action of a subgroup $`H`$ of $`\mathrm{Stab}(x)`$, we will refer to $`V`$ as an $`H`$-slice. ###### Remark 3.2. Note that since $`V`$ is smooth at $`x`$, we may always replace $`V`$ by its (unique) irreducible component passing through $`x`$ and thus assume that $`V`$ is irreducible. ###### Example 3.3. Let $`G`$ be an algebraic group, $`H`$ be an algebraic subgroup of $`G`$ and $`V`$ be an $`H`$-variety. Recall that the homogeneous fiber product $`X=G_HX`$ is defined as the geometric quotient $`X=(G\times V)/H`$, where $`H`$ acts on $`G\times V`$ by $`h(g,v)=(gh^1,hv)`$ (This geometric quotient exists under mild assumptions on $`V`$; in particular, it exists whenever $`V`$ is quasi-projective; see \[PV<sub>1</sub>, Theorem 4.19\].) Note that $`X=G_HV`$ is naturally a $`G`$-variety, where $`G`$ acts by left multiplication on the first factor; the details of this construction can be found in \[PV<sub>1</sub>, Section 4.8\]. The points of $`X`$ are in 1-1 correspondence with $`H`$-orbits in $`G\times V`$; we shall denote the point $`xX`$ corresponding to the $`H`$-orbit of $`(g,v)`$ in $`G\times V`$ by $`x=[g,v]`$. Note that there is an $`H`$-equivariant map $`VX`$ given by $`v[1,v]`$; we shall denote the image of this map by $`\stackrel{~}{V}`$. With these notations, $`\stackrel{~}{V}`$ is an $`H`$-slice for $`X`$ at $`x=[1,v]`$ for every smooth point $`v`$ of $`V`$; see \[PV<sub>1</sub>, Proposition 4.22\]. ∎ The following three results, Lemma 3.4, Lemma 3.5 and Proposition 3.6, are immediate consequences of the Luna Slice Theorem (see \[Lu\] or \[PV<sub>1</sub>, Section 6\]), if we assume that $`X`$ is an affine variety and the orbit of $`x`$ is closed in $`X`$. The arguments we present below do not require these assumptions. ###### Lemma 3.4. Let $`G`$ be an algebraic group, $`X`$ be an irreducible $`G`$-variety, and $`V`$ be a slice at $`xX`$. Then $`GV`$ is dense in $`X`$. ###### Proof. Consider the map $`\varphi :G\times VX`$, given by $`\varphi (g,v)=gv`$. The differential $`d\varphi _{(1,x)}:T_1(G)\times T_x(V)T_x(X)`$ is onto, since its image contains both $`T_x(Gx)`$ and $`T_x(V)`$. Consequently, $`d\varphi `$ is onto at a general point of $`G\times V`$. Thus $`\varphi `$ is dominant, and the lemma follows. ∎ ###### Lemma 3.5. Let $`G`$ be an algebraic group, $`H`$ be a subgroup, $`X`$ be a $`G`$-variety, and $`x`$ is a smooth $`H`$-fixed point of $`X`$. If $`H`$ is reductive then $`X`$ has an $`H`$-slice at $`x`$. ###### Proof. Let $`_{x,X}`$ be the maximal ideal of the local ring of $`X`$ at $`x`$. Consider the natural $`H`$-equivariant linear maps $`_{x,X}T_x(X)^{}T_x(Gx)^{}`$. Since $`H`$ is reductive these maps have $`H`$-equivariant splittings as maps of $`k`$-vector spaces. Thus we can choose a local coordinate system $`u_1,\mathrm{},u_n`$ in $`_{x,X}`$ such that $`\mathrm{Span}_k(u_1,\mathrm{},u_d)`$ is an $`H`$-invariant $`k`$-vector subspace of $`_{x,X}`$ and $`u_1,\mathrm{},u_d`$ (restricted to $`Gx`$) form a local coordinate system in $`𝒪_{x,Gx}`$. (Here $`n=dim(X)`$ and $`d=dim(Gx)`$.) Note that $`u_1,\mathrm{},u_n`$ are regular functions in some $`H`$-invariant open neighborhood of $`x`$. In this neighborhood a slice $`V`$ with desired properties is given by $`u_1=\mathrm{}=u_d=0`$. ∎ Recall that a $`G`$-variety $`X`$ is called generically free if $`\mathrm{Stab}(x)=\{1\}`$ for $`x`$ in general position in $`X`$. ###### Proposition 3.6. Let $`G`$ be an algebraic group, $`H`$ be a reductive subgroup, $`X`$ be a generically free $`G`$-variety and $`x`$ be a smooth $`H`$-fixed point of $`X`$. Then $`H`$ acts faithfully on $`T_x(X)/T_x(Gx)`$. ###### Proof. Let $`X_0`$ be the unique component of $`X`$ passing through $`x`$ and $`G_0`$ be the subgroup of all elements of $`G`$ that preserve $`X_0`$. Note that $`G_0`$ has finite index in $`G`$ and $`HG_0`$. After replacing $`X`$ by $`X_0`$ and $`G`$ by $`G_0`$, we may assume $`X`$ is irreducible. We now argue by contradiction. Assume the kernel $`K`$ of the $`H`$-action on $`T_x(X)/T_x(Gx)`$ is non-trivial. Note that $`K`$ is a normal subgroup of $`H`$. Since $`H`$ is reductive, $`K`$ is not unipotent. Hence, we can find a non-identity element $`gK`$ of finite order. By Lemma 3.5 $`X`$ has an $`H`$-slice $`V`$ at $`x`$. Since $`T_x(V)T_x(X)/T_x(Gx)`$ as $`H`$-representations, $`g`$ acts trivially on $`T_x(V)`$. This implies that $`g`$ acts trivially on $`V`$; see, e.g., \[RY<sub>1</sub>, Lemma 4.2\]. On the other hand, by Lemma 3.4 $`GV`$ is dense in $`X`$; consequently, for every $`xX`$ in general position $`\mathrm{Stab}(x)`$ contains a conjugate of $`g`$. Thus the $`G`$-action on $`X`$ is not generically free, contradicting our assumption. ∎ ## 4. Definition and first properties of $`i(X,x,H)`$ Throughout this section we shall make the following assumptions: $$\begin{array}{ccc}G\hfill & & \text{algebraic group}\hfill \\ H\hfill & & \text{finite abelian subgroup of }G\text{ of rank }r\hfill \\ X\hfill & & G\text{-variety of dimension }dim(G)+r\hfill \\ x\hfill & & H\text{-fixed point of }X\text{ whose stabilizer is finite}\hfill \end{array}$$ ###### Definition 4.1. The $`H`$-representation on $`T_x(X)/T_x(Gx)`$ decomposes as a direct sum of $`r`$ characters $`\chi _1,\mathrm{},\chi _rH^{}`$. We define $$i(X,x,H)=\chi _1\mathrm{}\chi _r^r(H^{}).$$ Since the collection of characters $`\chi _1,\mathrm{},\chi _n`$ is well-defined, up to reordering, $`i(X,x,H)`$ is well-defined in $`^r(H^{})`$, up to multiplication by $`1`$. Thus, properly speaking, $`i(X,x,H)`$ should be viewed as an element of the factor set $`^r(H^{})/`$, where $`w_1w_2`$ iff $`w_1=\pm w_2`$. By abuse of notation we will sometimes write $`i(X,x,H)^r(H^{})`$; in such cases it should be understood that $`i(X,x,H)`$ is only defined up to sign. ###### Remark 4.2. It is clear from the definition that if $`V`$ is an $`H`$-slice for $`X`$ at $`x`$ then $`i(X,x,H)=i(V,x,H)`$. In particular, in the setting of Example 3.3, if $`V`$ is an $`r`$-dimensional $`H`$-variety, $`X=G_HV`$, and $`v`$ is a smooth $`H`$-fixed point of $`V`$ then $`i(X,[1,v],H)=i(\stackrel{~}{V},[1,v],H)=i(V,v,H)`$. ###### Remark 4.3. Let $`g`$ be an element of the normalizer $`N_G(H)`$ and let $`\varphi _g`$ be the automorphism of $`H`$ sending $`h`$ to $`ghg^1`$. Then it is easy to see that $`i(X,gx,H)=(^r\varphi _g^{})\left(i(X,x,H)\right)`$, where $`^r\varphi _g^{}`$ is the automorphism of $`^r(H^{})`$ induced by $`\varphi _g`$. ###### Example 4.4. Let $`G=H`$ be a finite abelian group of rank $`r`$, $`\chi _1,\mathrm{},\chi _r`$ be a generating set for $`H^{}`$, and $`V=𝔸^r`$ be a faithful $`r`$-dimensional linear representation of $`H`$, given by $$h:(v_1,\mathrm{},v_r)(\chi _1(h)v_1,\mathrm{},\chi _r(h)v_r).$$ Then the origin $`0_V`$ is the only $`H`$-fixed point of $`V`$, and Definition 4.1 immediately implies $`i(V,0_V,H)=\chi _1\mathrm{}\chi _r`$. The extended $`H`$-action on $`\overline{V}=^r`$, given by $$h(v_0:v_1:\mathrm{}:v_d)=(v_0:\chi _1(h)v_1:\mathrm{}:\chi _d(h)v_d),$$ has exactly $`r+1`$ fixed points: $$x_0=(1:0:\mathrm{}:0),\mathrm{},x_r=(0:\mathrm{}:0:1).$$ Note that $`x_0=0_V`$. We claim that, up to sign, $$i(\overline{V},x_j,H)=i(V,x_0,H)=\chi _1\mathrm{}\chi _r$$ for $`j=1,\mathrm{},r`$. To prove this claim, say for $`j=1`$, note that $`v_0/v_1,v_2/v_1,\mathrm{},v_r/v_1`$ form an affine coordinate system on $`\overline{V}`$ near $`x_1`$. The $`H`$-action is diagonal in these coordinates, and the representation of $`H`$ on $`T_{x_1}(\overline{V})`$ is the direct sum of the characters $`\chi _1^1,\chi _2\chi _1^1,\mathrm{},\chi _r\chi _1^1`$. Consequently, $`i(\overline{V},x_1,H)=\pm \chi _1^1\chi _2\chi _1^1\mathrm{}\chi _r\chi _1^1=\pm \chi _1\mathrm{}\chi _r`$, as claimed. ###### Lemma 4.5. Suppose $`X`$ is a generically free $`G`$-variety. Then $`i(X,x,H)`$ generates $`^r(H^{})`$ as a group. Note that the statement of the lemma makes sense, even though $`i(X,x,H)`$ is only defined up to sign: if $`a`$ generates $`^r(H^{})`$ then so does $`a`$. ###### Proof. By Proposition 3.6 $`H`$ acts faithfully on $`T_x(X)/T_x(Gx)`$. Hence, the characters $`\chi _1,\mathrm{},\chi _r`$ introduced in Definition 4.1 generate $`H^{}`$ as an abelian group, and the lemma follows. ∎ ## 5. $`i(X,x,H)`$ and birational morphisms The purpose of this section is to prove the following: ###### Theorem 5.1. Let $`G`$ be an algebraic group of dimension $`d`$, $`H`$ be a finite abelian subgroup of $`G`$ of rank $`r`$, $`f:XY`$ be birational morphism of irreducible generically free $`G`$-varieties of dimension $`d+r`$, $`x`$ is a smooth $`H`$-fixed point of $`X`$, $`y=f(x)`$ is a smooth point of $`Y`$, and $`\mathrm{Stab}(y)`$ is finite. Then $`i(X,x,H)=i(Y,y,H)`$. ### Case I: G = H As a first step towards proving Theorem 5.1, we will consider the case where $`G=H`$ is a finite abelian group. In this case Theorem 5.1 can be restated as follows. ###### Proposition 5.2. Let $`H`$ be a finite abelian group, and $`f:XY`$ be a birational morphism of irreducible generically free $`H`$-varieties of dimension $`r=\mathrm{rank}(H)`$. Assume that $`x`$ is a smooth $`H`$-fixed point of $`X`$ and $`y=f(x)`$ is a smooth point of $`Y`$. Then $`i(X,x,H)=i(Y,y,H)`$. Before proceeding with the proof of Proposition 5.2, we introduce some background material on the power series ring $`k[[u_1,\mathrm{},u_r]]`$. Given $`wk[[u_1,\mathrm{},u_r]]`$ we shall denote by $`\mathrm{lm}(w)`$ the lowest degree monomial in $`u_1,\mathrm{},u_r`$ which enters into $`w(u_1,\mathrm{},u_r)`$ with a nonzero coefficient. Here “lowest degree” refers to a fixed (lexicographic) monomial order $``$ given by, say, $`u_1\mathrm{}u_r`$. Suppose $`v_1,\mathrm{},v_m`$ lie in the maximal ideal of $`k[[u_1,\mathrm{},u_r]]`$, i.e, $`\mathrm{lm}(v_i)1`$ for any $`i=1,\mathrm{},m`$. Then we can substitute $`v_1,\mathrm{},v_m`$ into any power series $`pk[[z_1,\mathrm{},z_m]]`$; in other words, $`p(v_1,\mathrm{},v_m)`$ is a well-defined element of $`k[[u_1,\mathrm{},u_r]]`$. If $`p=Z`$ is a monomial in $`k[[z_1,\mathrm{},z_m]]`$ then clearly (5.1) $$\mathrm{lm}\left(Z(v_1,\mathrm{},v_m)\right)=Z(\mathrm{lm}(v_1),\mathrm{},\mathrm{lm}(v_m)).$$ We shall write $`<u_1,\mathrm{},u_r>`$ for the group of all monomials in $`u_1,\mathrm{},u_r`$ (here we allow negative exponents). ###### Lemma 5.3. Suppose $`v_1,\mathrm{},v_mk[[u_1,\mathrm{},u_r]]`$. If $`\mathrm{lm}(v_1),\mathrm{},\mathrm{lm}(v_m)`$ generate a rank $`m`$ subgroup $`\mathrm{\Lambda }`$ in $`<u_1,\mathrm{},u_r>^r`$ then $`\mathrm{lm}\left(p(v_1,\mathrm{},v_m)\right)\mathrm{\Lambda }`$ for any $`pk[[z_1,\mathrm{},z_m]]`$. Note that the conditions of the lemma imply $`mr`$; only the case $`m=r`$ will be used in the subsequent application. ###### Proof. Suppose $`p(z_1,\mathrm{},z_m)=a_ZZ`$, where $`Z`$ ranges over all monomials in $`z_1,\mathrm{},z_m`$ with non-negative exponents. By our assumption $`\mathrm{lm}(v_1),\mathrm{},\mathrm{lm}(v_m)`$ are (multiplicatively) linearly independent, i.e., $$Z(\mathrm{lm}(v_1),\mathrm{},\mathrm{lm}(v_m))Z^{}(\mathrm{lm}(v_1),\mathrm{},\mathrm{lm}(v_m))$$ for any two distinct monomials $`Z`$ and $`Z^{}`$. Suppose $`Z_0(\mathrm{lm}(v_1),\mathrm{},\mathrm{lm}(v_m))`$ has the smallest monomial order among all monomials (in $`u_1,\mathrm{},u_m`$) of the form $`Z(\mathrm{lm}(v_1),\mathrm{},\mathrm{lm}(v_m))`$, with $`c_Z0`$. Then (5.1) tells us that $$\mathrm{lm}(Z(v_1,\mathrm{},v_m))\mathrm{lm}(Z_0(v_1,\mathrm{},v_m))=Z_0(\mathrm{lm}(v_1),\mathrm{},\mathrm{lm}(v_m))$$ for any $`ZZ_0`$ with $`c_Z0`$. Thus $$\mathrm{lm}(p(v_1,\mathrm{},v_m))=\mathrm{lm}(Z_0(v_1,\mathrm{},v_m))=Z_0(\mathrm{lm}(v_1),\mathrm{},\mathrm{lm}(v_m))\mathrm{\Lambda },$$ as claimed. ∎ ###### Proof of Proposition 5.2. Diagonalizing the action of $`H`$ on the cotangent space $`T_x^{}(X)`$, we obtain a basis $`\overline{u}_1,\mathrm{},\overline{u}_rT_x^{}(X)`$ such that $`h\overline{u}_i=\chi _i(h)\overline{u}_i`$ for every $`hH`$; here $`\chi _1,\mathrm{},\chi _rH^{}`$. Since the $`k`$-linear map $`_{x,X}_{x,X}/_{x,X}^2=T_x^{}(X)`$ has an $`H`$-invariant $`k`$-linear splitting, we can find a local system of parameters $`u_1,\mathrm{},u_r_{x,X}`$ such that (5.2) $$hu_i=\chi _i(h)u_i$$ for every $`hH`$ and $`i=1,\mathrm{},r`$. Similarly, we can find a local coordinate system $`v_1,\mathrm{},v_r_{y,Y}`$ for $`Y`$ at $`y`$ and $`\eta _1,\mathrm{},\eta _rH^{}`$ such that (5.3) $$hv_i=\eta _i(h)v_i$$ for every $`hH`$ and $`i=1,\mathrm{},r`$. Clearly $`i(X,x,H)=\chi _1\mathrm{}\chi _r`$ and $`i(Y,y,H)=\eta _1\mathrm{}\eta _r`$. We shall identify the elements $`v_1,\mathrm{},v_r`$ with their images in $`𝒪_{x,X}`$ under $`f^{}:𝒪_{y,Y}𝒪_{x,X}`$. The $`H`$-action on $`𝒪_{x,X}`$ naturally extends to $`k[[u_1,\mathrm{},u_r]]`$; in view of (5.2) the leading term map $`w\mathrm{lm}(w)`$ is $`H`$-equivariant. Suppose $`\mathrm{lm}(v_i)=u_1^{a_{i1}}\mathrm{}u_r^{a_{ir}}`$ for some non-negative integers $`a_{ij}`$. Then (5.2) and (5.3) imply $`\eta _i=_j\chi _i^{a_{ij}}`$, and thus, up to sign, $$i(Y,y,H)=\eta _1\mathrm{}\eta _r=det(a_{ij})\chi _1\mathrm{}\chi _r=\pm det(a_{ij})i(X,x,H).$$ There are two conclusions we can draw from this formula. First of all, by Lemma 4.5 we know that both $`i(X,x,H)`$ and $`i(Y,y,H)`$ generate $`^r(H^{})`$; thus $`det(a_{ij})0`$. Secondly, in order to prove the proposition, it is sufficient to show that (5.4) $$det(a_{ij})=\pm 1\text{in }\text{.}$$ We now proceed with the proof of (5.4). Let $`<u_1,\mathrm{},u_r>`$ be the free abelian multiplicative group generated by $`u_1,\mathrm{},u_r`$. Since $`det(a_{ij})0`$, the leading monomials $`\mathrm{lm}(v_1),\mathrm{},\mathrm{lm}(v_r)`$ generate a (free abelian) subgroup $`\mathrm{\Lambda }`$ of rank $`r`$ in $`<u_1,\mathrm{},u_r>()^r`$; in other words, $`\mathrm{\Lambda }`$ has finite index in $`<u_1,\mathrm{},u_r>`$. On the other hand, (5.4) holds if and only if $`\mathrm{\Lambda }=<u_1,\mathrm{},u_r>`$. It is therefore sufficient to prove that $`u_i\mathrm{\Lambda }`$ for every $`i=1,\mathrm{},r`$. Since $`𝒪_{x,X}`$ and $`𝒪_{y,Y}`$ have the same field of fractions, each $`u_i`$ can be written as $`p/q`$, where $`p,q𝒪_{y,Y}\{0\}`$. Represent $`p`$ and $`q`$ by power series in $`v_1,\mathrm{},v_r`$. By Lemma 5.3 $`\mathrm{lm}(p(v_1,\mathrm{},v_r))`$ and $`\mathrm{lm}(q(v_1,\mathrm{},v_r))`$ lie in $`\mathrm{\Lambda }`$; thus taking the leading monomials on both sides of the equation $$q(v_1,\mathrm{},v_r)u_i=p(v_1,\mathrm{},v_r),$$ we conclude that $`u_i\mathrm{\Lambda }`$, as desired. ∎ ### Case II: G - arbitrary We are now ready to finish the proof of Theorem 5.1. The idea is to replace $`X`$ and $`Y`$ by suitable $`H`$-slices, then appeal to Proposition 5.2. Diagonalizing the $`H`$-action on $`T_y^{}(Gy)`$, we obtain a basis $`\overline{v}_1,\mathrm{},\overline{v}_d`$ such that $`h\overline{v}_i=\alpha _i(h)\overline{v}_i`$ for some characters $`\alpha _1,\mathrm{},\alpha _d`$ of $`H`$. Since the natural $`H`$-equivariant $`k`$-vector space maps $`_{y,Y}T_y^{}(Y)T_y^{}(Gy)`$ have $`H`$-equivariant splittings, we can lift $`\overline{v}_1,\mathrm{},\overline{v}_d`$ to $`v_1,\mathrm{},v_d_{y,Y}`$ such that such that $`hv_i=\alpha _i(h)v_i`$ for each $`hH`$. In other words, $`v_1,\mathrm{},v_d`$ form a local coordinate system for $`Gy`$ at $`y`$. Since both $`\mathrm{Stab}(x)`$ and $`\mathrm{Stab}(y)`$ are finite, $`df_x:T_x(Gx)T_y(Gy)`$ is an isomorphism. Thus $`f^{}(v_1),\mathrm{},f^{}(v_r)`$ form a local coordinate system for $`Gx`$ near $`x`$. Define $`WY`$ as the irreducible component of $`\{v_1=\mathrm{}=v_d=0\}`$ passing through $`y`$ and $`VX`$ as the irreducible component of $`\{f^{}(v_1)=\mathrm{}=f^{}(v_d)=0\}`$ passing through $`x`$. Then $`W`$ is an $`H`$-slice for $`Y`$ at $`y`$ and $`V`$ is an $`H`$-slice for $`X`$ at $`x`$. Clearly $`f(V)W`$, i.e., $`f_{\|V}:VW`$ is a well-defined morphism. We claim that $`f_{\|V}`$ is, in fact, a birational morphism. Theorem 5.1 follows from this claim because $$i(X,x,H)\stackrel{\text{Remark }\text{4.2}}{=}i(V,x,H)\stackrel{\text{Proposition }\text{5.2}}{=}i(W,y,H)\stackrel{\text{Remark }\text{4.2}}{=}i(Y,y,H).$$ To show that $`f_{\|V}`$ is dominant, assume, to the contrary, that $`dim(f(V))<r`$. Since $`X`$ is irreducible, $`GV`$ is dense in $`X`$ by Lemma 3.4 and thus $`dim(Y)=dim(f(X))=dim(f(GV))=dim(Gf(V))d+dim(f(V))<d+r`$, a contradiction. It remains to show that $`f_{\|V}`$ is generically 1-1 on closed points. Since $`f`$ is a birational morphism (i.e., has degree 1), there exists a dense $`G`$-invariant open subset $`Y_0`$ of $`Y`$ such that for every $`y_0Y_0`$, $`f^1(y_0)`$ is a single point in $`X`$. Since $`GW`$ is dense in $`Y`$ (by Lemma 3.4), $`Y_0W`$ is a dense open subset of $`W`$. Thus a general point of $`W`$ has exactly one preimage in $`X`$. On the other hand, a general point of $`W`$ has at least $`\mathrm{deg}(f_{\|V})1`$ preimages in $`X`$. This shows that $`\mathrm{deg}(f_{\|V})=1`$, i.e. $`f_{\|V}`$ is birational, as claimed. ∎ ## 6. Proof of Theorem 1.1 In this section we deduce Theorem 1.1 from Theorem 5.1. Our proof relies on canonical resolution of singularities. (We remark that canonical resolution of singularities is not used elsewhere in this paper.) We begin with two simple preliminary results. ###### Lemma 6.1. Let $`H`$ be an algebraic group, $`f:ZX`$ be a birational morphism of complete irreducible $`H`$-varieties and $`X_1`$ be an irreducible $`H`$-invariant codimension 1 subvariety of $`X`$. If $`X_1`$ passes through a normal point of $`X`$ then there exists an $`H`$-invariant irreducible codimension 1 subvariety $`Z_1Z`$ such that $`f_{\|Z_1}:Z_1X_1`$ is a birational morphism. ###### Proof. Since $`X_1`$ contains a normal point of $`X`$, the rational map $`f^1:XZ`$ is defined at a general point of $`X_1`$. Now $`Z_1`$ = the closure of $`f^1(X_1)`$ in $`Z`$, has the desired properties. ∎ ###### Proposition 6.2. Let $`H`$ be a diagonalizable group, $`\alpha :ZX`$ be a birational morphism of complete irreducible $`H`$-varieties and $`x`$ be a smooth $`H`$-fixed point of $`X`$. Then there exists an $`H`$-fixed point $`zZ`$ such that $`\alpha (z)=x`$. This result can be established by the argument used in the proof of \[RY<sub>1</sub>, Proposition A.4\], due to J. Kollár and E. Szabó. (In fact, if $`H`$ is a $`p`$-group then Proposition 6.2 follows from \[RY<sub>1</sub>, Proposition A.4\].) We give a simple self-contained proof below. ###### Proof. We argue by induction on $`dim(X)`$. The base case, $`dim(X)=0`$, is obvious. For the induction step, assume $`dim(X)1`$. We claim that there exists codimension 1 $`H`$-invariant irreducible subvariety $`X_1`$ such that $`x`$ is a smooth point of $`X_1`$. Arguing as we did at the beginning of the proof of Proposition 5.2, we see that there exists non-zero element $`u_{x,X}`$ such that for every $`hH`$, $`hu=\alpha (h)u`$, where $`\alpha `$ is a character of $`H`$. Then the (locally closed) subvariety $`\{u=0\}`$ of $`X`$ is $`H`$-invariant and smooth at $`x`$. Hence, its unique irreducible component passing through $`x`$ is also $`H`$-invariant, and we can define $`X_1`$ as the closure of this irreducible component. This proves the claim. Now by Lemma 6.1 there exists a codimension 1 irreducible $`H`$-invariant subvariety $`Z_1`$ of $`Z`$ such that $`\alpha _{\|Z_1}:Z_1X_1`$ is a birational morphism of $`H`$-varieties. Applying the induction assumption to this morphism, we construct $`zZ_1Z`$ with desired properties. ∎ We are now ready to complete the proof of Theorem 1.1. The idea is to construct a complete smooth model $`Z`$ of $`X`$ (or $`Y`$) that dominates them both, i.e., fits into a diagram $$\begin{array}{ccccc}& & Z& & \\ & \stackrel{\alpha }{}& & \stackrel{\beta }{}& \\ X& & & & Y.\end{array}$$ where $`\alpha `$ and $`\beta `$ are birational morphisms of $`G`$-varieties. If we find such a $`Z`$, Theorem 1.1 will easily follow. Indeed, by Proposition 6.2 there exists an $`H`$-fixed point $`zZ`$ such that $`\alpha (z)=x`$. Setting $`y=\beta (z)`$ and applying Theorem 5.1 to $`\alpha `$ and $`\beta `$, we conclude that $`i(X,x,H)=i(Z,z,H)=i(Y,y,H)`$, as desired. It remains to construct $`Z`$. Let $`WX\times Y`$ be the closure of the graph of a birational isomorphism $`f`$ between $`X`$ and $`Y`$. Then $`W`$ is a complete $`G`$-variety that dominates both $`X`$ and $`Y`$. In other words, $`W`$ satisfies all of our requirements for $`Z`$, with one exception: it may not be smooth. Let (6.1) $$\pi :Z=W_n\mathrm{@}>\pi _n>>\mathrm{}\mathrm{@}>\pi _1>>W_0=W,$$ be the canonical resolution of singularities of $`W`$, as in \[V, Theorem 7.6.1\] or \[BM, Theorem 1.6\]. Here $`Z`$ is smooth, and each $`\pi _i`$ is a blowup with a smooth center; since these centers are canonically chosen, they are $`G`$-invariant. Thus the $`G`$ action can be lifted to $`Z`$ so that $`\pi `$ is a birational morphism of complete $`G`$-varieties. The smooth complete $`G`$-variety $`Z`$ constructed in this way has the desired properties. ∎ ###### Remark 6.3. An alternative construction of $`Z`$ is given by the equivariant version of Hironaka’s theorem on elimination of points of indeterminacy (proved in \[RY<sub>4</sub>\]), which asserts that for every rational map $`f:XY`$ of $`G`$-varieties there exists a sequence of blowups $$\pi :Z=X_n\mathrm{@}>\pi _n>>\mathrm{}\mathrm{@}>\pi _1>>X_0$$ with smooth $`G`$-equivariant centers such that $`f\pi :ZY`$ is regular. The advantage of this approach is that it only uses Proposition 6.2 in the case where $`\alpha `$ is a single blowup with a smooth $`G`$-equivariant center (in which case the proof is immediate; see, e.g., \[RY<sub>1</sub>, Lemma 5.1\]). On the other hand, since the theorem on equivariant elimination of points of indeterminacy is itself deduced from canonical resolution of singularities in \[RY<sub>4</sub>\], we opted for a direct proof here. ###### Remark 6.4. A rational map $`f:XY`$ (respectively, a morphism $`f:XY`$) of reducible varieties is called a birational isomorphism (respectively, a birational morphism) if $`X`$ and $`Y`$ have irreducible decompositions $`X_1\mathrm{}X_n`$ and $`Y_1\mathrm{}Y_n`$ such that $`f`$ restricts to a birational isomorphism $`X_iY_i`$ (respectively, a birational morphism) for each $`i`$. With this definition, the irreducibility assumption in Proposition 6.2 can be removed. Indeed, we can reduce to the case where $`X`$ and $`Z`$ are irreducible by replacing them with suitable irreducible components. The irreducibility assumption in Theorem 1.1 (respectively, Theorem 5.1 and Proposition 5.2) can also be removed, if we assume $`dim_x(X)=d+r`$ (respectively, $`dim_x(X)=d+r`$ and $`dim_x(X)=r`$). Indeed, if $`X_1`$ is the (necessarily unique) irreducible component of $`X`$ containing $`x`$ and $`G_1=\{gG\|g(X_1)=X_1\}`$, then $`HG_1`$, $`[G:G_1]<\mathrm{}`$, and $`i(X,x,H)=i(X_1,x,H)`$, so that in each case we may replace $`X`$, $`Y`$, and $`G`$ by $`X_1`$, $`Y_1`$ and $`G_1`$. ###### Remark 6.5. One may ask if the condition that $`\mathrm{Stab}(x)`$ is finite for every $`xX^H`$ (and similarly for $`Y`$) of Theorem 1.1 is ever satisfied. Indeed, if $`H`$ is contained in a torus $`T`$ of $`G`$ then the answer is “no”, since $`\mathrm{Stab}(x)`$ is infinite for every $`xX^TX^H`$, and $`X^T\mathrm{}`$ by the Borel Fixed Point Theorem. On the other hand, if the centralizer $`C_G(H)`$ is finite, then we claim that every generically free $`G`$-variety $`X`$ has a birational model satisfying this condition. Indeed, by \[RY<sub>1</sub>, Theorem 1.1\] $`X`$ has a model with the property that the stabilizer of every point is of the form $`U>D`$, where $`D`$ is diagonalizable and $`U`$ is unipotent. Assume $`xX^H`$. By the Levi decomposition theorem, we may choose $`D`$ so that $`D`$ contains $`H`$. Now \[RY<sub>1</sub>, Lemma 7.3\] tells us that $`U=\{1\}`$. Thus $`\mathrm{Stab}(x)=DC_G(H)`$, which is finite, as claimed. Examples of pairs $`HG`$, where $`G`$ is a semi-simple algebraic group and $`H`$ is an abelian subgroups of $`G`$ whose centralizer is finite can be found in \[Gr\]; see also \[RY<sub>1</sub>, Section 8\]. ## 7. A birational classification of linear representations It is not difficult to see that Conjecture 1.2 fails if $`G`$ is a finite cyclic group of order $`n=5`$ or $`7`$. Indeed, let $`V_\omega `$ be the 1-dimensional representation of $`G`$ such that $`\sigma `$ acts on $`V_\omega `$ by the character $`x\omega x`$, where $`\omega `$ is a primitive $`n`$th root of unity. Since any birational automorphism of $`𝔸^1`$ lifts to a regular automorphism of $`^1`$, it is easy to see that $`V_\omega `$ is birationally isomorphic to $`V_\omega ^{}`$ iff $`\omega ^{}=\omega `$ or $`\omega ^{}=\omega ^1`$. (The two $`G`$-fixed points in $`^1`$ are preserved in the former case and interchanged in the latter.) If $`n=5`$ or $`7`$, we can find two primitive $`n`$th roots of unity $`\omega `$ and $`\omega ^{}`$ such that $`\omega ^{}\omega ^{\pm 1}`$, so that $`V_\omega `$ and $`V_\omega ^{}`$ are not birationally isomorphic. (P. I. Katsylo has informed us that this observation was independently made by E. A. Tevelev.) In this section we will classify faithful linear representations of diagonalizable group $`G`$ up to birational equivalence and show that Conjecture 1.2 fails for a number of groups, both abelian and non-abelian. The general flavor of the results we obtain will be similar to the situation described in the above paragraph but the arguments are more complicated due to the fact that we will be working with higher-dimensional varieties, rather than curves. ### Representations of diagonalizable groups Recall that every linear representation $`V`$ of $`G`$ decomposes as a sum of 1-dimensional character spaces; if the associated characters of $`G`$ are $`\chi _1,\mathrm{},\chi _d`$, we shall write $`V=\chi _1\mathrm{}\chi _d`$. ###### Theorem 7.1. Let $`G`$ be a diagonalizable group of rank $`r`$ and let $`V=\chi _1\mathrm{}\chi _d`$ and $`W=\eta _1\mathrm{}\eta _d`$ be faithful $`d`$-dimensional linear representations of $`G`$. (In particular, $`dr`$.) Then $`V`$ and $`W`$ are birationally isomorphic as $`G`$-varieties if and only if $`\chi _1\mathrm{}\chi _d=\pm \eta _1\mathrm{}\eta _d`$ in $`^d(G^{})`$. ###### Proof. Since $`G`$ acts faithfully on $`V`$ and $`W`$, we have (7.1) $$<\chi _1,\mathrm{},\chi _d>=<\eta _1,\mathrm{},\eta _d>=G^{}.$$ Assume $`\chi _1\mathrm{}\chi _r=\pm \eta _1\mathrm{}\eta _r`$. Then by Corollary 2.5 there exists an $`N=(n_{ij})\mathrm{GL}()`$ such that $`\eta _i=\chi _1^{n_{1i}}\mathrm{}\chi _d^{n_{di}}`$. The desired birational isomorphism $`VW`$ is can now be explicitly given, in the natural coordinates on $`V`$ and $`W`$, by $`(x_1,\mathrm{},x_d)(y_1,\mathrm{},y_d)`$, where $`y_i=x_1^{n_{1i}}\mathrm{}x_d^{n_{di}}`$. Conversely, suppose (7.2) $$\chi _1\mathrm{}\chi _d\pm \eta _1\mathrm{}\eta _d\text{in}^d(G^{})$$ We want to prove that $`V`$ and $`W`$ are not birationally isomorphic as $`G`$-varieties. Note that (7.2) is impossible if $`dr+1`$, since in this case $`^d(G^{})=(0)`$; see Lemma 2.1(b). Thus we will assume from now on that $`d=r=\mathrm{rank}(G)`$. We will consider three cases. Case 1: $`G=(G_m)^r`$ is a torus. In this case $`\chi _1\mathrm{}\chi _r`$ and $`\eta _1\mathrm{}\eta _r`$ are both generators of $`^r(G^{})=`$, so that (7.2) is impossible. Case 2: $`G`$ is a finite abelian group. The $`G`$-action on $`V=𝔸^r`$ (respectively $`W=𝔸^r`$) naturally extends to the projective space $`\overline{V}=^r`$ (respectively, $`\overline{W}=^r`$). Example 4.4 shows that for every $`G`$-fixed point $`x\overline{V}`$, $`i(\overline{V},x,G)=\pm \chi _1\mathrm{}\chi _r`$ and for every $`G`$-fixed point $`y\overline{W}`$, $`i(\overline{W},y,G)=\pm \eta _1\mathrm{}\eta _r`$. Thus in view of (7.2), Theorem 1.1 says that $`\overline{V}`$ and $`\overline{W}`$ (and, hence, $`V`$ and $`W`$) are not birationally isomorphic as $`G`$-varieties. Case 3: $`G`$ is a diagonalizable group but not a torus. Write $`G=𝔾_m(n_1)\times \mathrm{}\times 𝔾_m(n_r)`$, as in (2.1), with $`n_12`$. Let $`H=𝔾_m(n_1)^r=(/n_1)^r`$ be the $`n_1`$-torsion subgroup of $`G`$. It is sufficient to show that $`V`$ and $`W`$ are not birationally isomorphic as $`H`$-varieties; then they certainly cannot be birationally isomorphic as $`G`$-varieties. By Case 2, it is enough to show that (7.3) $$\chi _1^{}\mathrm{}\chi _r^{}\pm \eta _1^{}\mathrm{}\eta _r^{}\text{in}^r(H^{})$$ where $`\chi _i^{}`$ and $`\eta _j^{}`$ are the characters of $`H`$ obtained by restricting $`\chi _i`$ and $`\eta _j`$ from $`G`$ to $`H`$. Note that the inclusion $`\varphi :HG`$ induces a surjection $`\varphi ^{}:G^{}H^{}`$ of the dual group, which, in turn, induced a map of cyclic groups $`^r(\varphi ^{}):^r(G^{})^r(H^{})`$. Elementary group theory tells us that $`G^{}=(/n_1)\times \mathrm{}\times (/n_r)`$, $`H^{}=(/n_1)^r`$, $`\varphi ^{}:G^{}H^{}`$ is (componentwise) reduction modulo $`n_1`$, and $`^r(\varphi ^{})`$ is the identity map $`^r(G^{})=/n_1\stackrel{}{}/n_1=^r(H^{})`$. Applying $`^r(\varphi ^{})`$ to both sides of (7.2), we obtain (7.3), as desired. ∎ ### Proof of Theorem 1.3 Let $`G`$ be a diagonalizable group $`G`$ of rank $`r`$ of the form (1), and let $`V=\chi _1\mathrm{}\chi _d`$ and $`W=\eta _1\mathrm{}\eta _d`$ be faithful $`d`$-dimensional linear representations of $`G`$. (a) If $`dr+1`$ then $`^r(G^{})=(0)`$, so that $`\chi _1\mathrm{}\chi _d=0=\eta _1\mathrm{}\eta _d`$. Thus $`V`$ and $`W`$ are birationally isomorphic as $`G`$-varieties by Theorem 7.1. From now on we will assume $`d=r`$. Note that in this case both $`\chi _1\mathrm{}\chi _d`$ and $`\eta _1\mathrm{}\eta _d`$ are generators of $`^r(G^{})=/n_1`$. (b) Suppose $`n_1=2`$. Since $`^r(G^{})=/2`$ has only one generator, $`\chi _1\mathrm{}\chi _d=\eta _1\mathrm{}\eta _d`$. Thus $`V`$ and $`W`$ are birationally isomorphic as $`G`$-varieties by Theorem 7.1. Now assume $`n_1=0`$. Then $`^r(G^{})=`$. The only generators of this group are $`\pm 1`$; thus $`\chi _1\mathrm{}\chi _d=\pm \eta _1\mathrm{}\eta _d`$, and, once again, Theorem 7.1 tells us that $`V`$ and $`W`$ are birationally isomorphic. (c) Suppose $`n_13`$. By Theorem 7.1, birational isomorphism classes of $`r`$-dimensional linear representations of $`H`$ are in 1-1 correspondence generators of $`^r(H^{})/n_1`$ (as an additive group), modulo multiplication by $`1`$. Since $`n_r3`$, $`aa`$ for any generator $`a`$ of $`/n_1`$. Thus in this case the number of isomorphism classes of faithful $`r`$-dimensional $`H`$-representations is $`\varphi (n_r)/2`$, as claimed. ∎ ### Further counterexamples to Conjecture 1.2 Theorem 1.3 shows that Conjecture 1.2 fails for many diagonalizable groups. We will now see this conjecture fails for some non-abelian groups as well. ###### Proposition 7.2. Let $`n`$ and $`r`$ be a positive integers, $`P`$ be a subgroup of $`\mathrm{S}_r`$ and $`G=(/n)^r>P`$, where $`P`$ acts on $`(/n)^r`$ by permuting the factors. Assume there exists an $`m`$ such that $`(m,n)=1`$ and $`m^r\pm 1(modn)`$. Then there exist two birationally inequivalent $`r`$-dimensional representations of $`G`$. In particular, Conjecture 1.2 fails for this group. We remark that an integer $`m`$ satisfying the requirements of Proposition 7.2 always exists if the exponent of $`U_n`$ does not divide $`2r`$; here $`U_n`$ is the (multiplicative) group of units in $`/n`$. In particular, $`m`$ exists if there is a prime power $`p^e`$ such that $`p^e\|n`$ but $`\varphi (p^e)=(p1)p^{e1}2r`$. ###### Proof. Let $`\omega `$ be a primitive $`n`$th root of unity. We define the $`r`$-dimensional representations $`V`$ and $`W`$ of $`G`$ as follows: $`((a_1,\mathrm{},a_r),\sigma ):`$ $`(v_1,\mathrm{},v_r)`$ $`(\omega ^{a_1}v_{\sigma ^1(1)},\mathrm{},\omega ^{a_r}v_{\sigma ^1(r)})`$ and $`((a_1,\mathrm{},a_r),\sigma ):`$ $`(w_1,\mathrm{},w_r)`$ $`(\omega ^{ma_1}w_{\sigma ^1(1)},\mathrm{},\omega ^{ma_r}w_{\sigma ^1(r)}).`$ Here $`a_1,\mathrm{},a_n/n`$, $`\sigma P\mathrm{S}_r`$, $`(v_1,\mathrm{},v_r)V`$ and $`(w_1,\mathrm{},w_r)W`$. It is easy to see that $`V`$ and $`W`$ are, indeed, well-defined faithful $`r`$-dimensional representations of $`G`$. To prove the proposition it is sufficient to show that $`V`$ and $`W`$ are not birationally isomorphic as $`(/n)^r`$-varieties. Let $`\chi _i`$ be the character of $`(/n)^r`$ given by $`\chi (a_1,\mathrm{},a_r)=\omega ^{a_i}`$. Then, as $`/n`$-representations, $`V=\chi _1\mathrm{}\chi _r`$ and $`W=\chi _1^m\mathrm{}\chi _r^m`$. By our assumption $$\chi _1^m\mathrm{}\chi _r^m=m^r\chi _1\mathrm{}\chi _r\pm \chi _1\mathrm{}\chi _r$$ in $`^r(/n^{})/n`$. Thus Theorem 7.1 tells us that $`V`$ and $`W`$ are not isomorphic as $`(/n)^r`$-varieties (and hence, as $`G`$-varieties). ∎ ###### Remark 7.3. The same argument proves the following stronger result. Let $`n_1,n_2,\mathrm{},n_s`$, $`r_1,\mathrm{},r_s`$ be positive integers such that $`n_i`$ divides $`n_{i+1}`$ for $`i=2,\mathrm{},r`$, let $`P_i`$ be a subgroup of the symmetric group $`\mathrm{S}_{r_i}`$ and let $`G_i=(/n_i)^{r_i}>P_i`$. Assume there exist integers $`m_1,\mathrm{},m_s`$ such that $`(m_i,n_i)=1`$ and $`m_1^{r_1}\mathrm{}m_s^{r_s}\pm 1(modn_1)`$. Then $`G=(G_m)^a\times G_1\times \mathrm{}\times G_s`$ has two birationally inequivalent ($`a+r_1+\mathrm{}+r_s`$)-dimensional representations. In particular, Conjecture 1.2 fails for any $`G`$ of this form. ###### Remark 7.4. The proof of Proposition 7.2 shows that $`G=(/n)^r>P`$ has at least $`\|\pm U_n^r\|/2`$ birational isomorphism classes of $`r`$-dimensional representations. Here, as before, $`U_n`$ denotes the multiplicative group of units in the ring $`/n`$, and $`\pm U_n^r`$ denotes the subset of $`U_r`$ consisting of elements of the form $`\pm m^r`$, as $`m`$ ranges over $`U_n`$. A similar estimate can be given for the number of birational isomorphism classes of ($`a+r_1+\mathrm{}+r_s`$)-dimensional representations of the group $`G`$ in Remark 7.3. In particular, if $`n_1=\mathrm{}=n_s`$ and $`(r_1,\mathrm{},r_s)=1`$ then there are at least $`\varphi (n_1)/2`$ such classes. ## 8. Birational equivalence of quantum tori In this section we will use the invariant $`i(X,x,H)`$ to classify $`\mathrm{PGL}_n`$-varieties (and consequently central simple algebras) of a certain form. In particular, we will prove Theorem 1.4. ### Abelian subgroups of $`\mathrm{PGL}_n`$ Let $`A`$ be a finite abelian group of order $`n`$ and let $`V=k[A]`$ be the group ring of $`A`$. For $`aA`$ and $`\chi A^{}`$ define $`P_a,D_\chi \mathrm{GL}(V)`$ by $`P_a(b)=ab`$ and $`D_\chi (b)=\chi (b)b`$ for every $`bA`$. It is easy to see that $`D_\chi P_a=\chi (a)P_aD_\chi `$. Thus if $`p_a`$ and $`d_\chi `$ denote the elements of $`\mathrm{PGL}(V)`$ represented, respectively, by $`P_a`$ and $`D_\chi \mathrm{GL}(V)`$, then $`\varphi :A\times A^{}`$ $`\mathrm{PGL}(V)=\mathrm{PGL}_n`$ $`(a,\chi )`$ $`p_ad_\chi `$ defines an embedding of $`A\times A^{}`$ in $`\mathrm{PGL}_n`$. Let $`H`$ be an abelian subgroup of $`\mathrm{PGL}_n`$. Then $`H`$ is naturally equipped with an alternating bilinear form $`\omega _H:H\times H\mu _n`$ (cf. Definition 2.2(a)). Here $`\mu _n`$ is the group of $`n`$th roots of unity in $`k`$, identified with the center of $`\mathrm{SL}_n(k)`$, and $`\omega _H(a,b)=ABA^1B^1`$, where $`A`$ and $`B\mathrm{SL}_n`$ represent $`a`$ and $`b\mathrm{PGL}_n`$ respectively. ###### Lemma 8.1. Let $`A`$ be a finite abelian group of rank $`r`$, $`H=\varphi (A\times A^{})=\{p_ad_\chi \|aA,\chi A^{}\}`$ be the subgroup of $`\mathrm{PGL}_n`$ defined above. Then (a) the elements of $`P_aD_\chi `$ span $`\mathrm{M}_n`$ as a $`k`$-vector space, as $`a`$ ranges over $`A`$ and $`\chi `$ ranges over $`A^{}`$, and (b) the alternating bilinear form $`\omega _H`$ is symplectic (i.e., non-degenerate). Let $`g`$ be an element of the normalizer $`N_{\mathrm{PGL}_n}(H)`$, and $`\psi _g:HH`$ be conjugation by $`g`$. Then (c) $`\psi _g`$ preserves $`\omega _H`$, and (d) $`\psi _g`$ induces the identity map $`^{2r}(H^{})^{2r}(H^{})`$. ###### Proof. (a) See \[RY<sub>3</sub>, Lemma 3.2\]. (b) See \[RY<sub>2</sub>, Lemma 7.8\]. (c) Choose $`a`$ and $`bH\mathrm{PGL}_n`$ and lift them to $`A`$ and $`B\mathrm{SL}_n`$. Since $`ABA^1B^1`$ is a central element of $`\mathrm{SL}_n`$, we have $`\omega _H(\psi _g(a),\psi _g(b))`$ $`=\omega _H(gag^1,gbg^1)`$ $`=(gAg^1)(gBg^1)(gA^1g^1)(gB^1g^1)`$ $`=g(ABA^1B^1)g^1=ABA^1B^1=\omega _H(a,b),`$ as claimed. (d) Follows from (b), (c) and Lemma 2.3(b). ∎ ### $`\mathrm{PGL}_n`$-varieties ###### Proposition 8.2. Let $`A`$ be a finite abelian group of order $`n`$ and rank $`r`$ and let $`H=\varphi (A\times A^{})`$ be the subgroup of $`\mathrm{PGL}_n`$ defined above. Suppose $`V=\chi _1\mathrm{}\chi _{2r}`$ and $`W=\eta _1\mathrm{}\eta _{2r}`$ are faithful representations of $`H`$. Then the following are equivalent: (a) $`\chi _1\mathrm{}\chi _{2r}=\pm \eta _1\mathrm{}\eta _{2r}`$ in $`^{2r}(H^{})`$, (b) $`V`$ and $`W`$ are birationally isomorphic as $`H`$-varieties, (c) $`X=\mathrm{PGL}_n_HV`$ and $`Y=\mathrm{PGL}_n_HW`$ are birationally isomorphic as $`\mathrm{PGL}_n`$-varieties. Here $`\mathrm{PGL}_n_HV`$ and $`\mathrm{PGL}_n_HW`$ are homogeneous fiber products; see Example 3.3. ###### Proof. (a) and (b) are equivalent by Theorem 7.1. The implication (b) $``$ (c) is obvious. Thus we only need to show (c) $``$ (a). The idea of the proof is to appeal to Theorem 1.1. We begin by observing that $`X`$ and $`Y`$ naturally embed as dense open subsets in projective varieties $`\overline{X}=((\mathrm{M}_n)\times \overline{V})/H`$ and $`\overline{Y}=((\mathrm{M}_n)\times \overline{W})/H`$ respectively. Here $`\overline{V}=^{2r}`$ is the projective completion of $`V=𝔸^{2r}`$; $`\mathrm{PGL}_n`$ acts on $`(\mathrm{M}_n)\times \overline{V}`$ by left multiplication on the first factor; this action commutes with the $`H`$-action on $`(\mathrm{M}_n)\times \overline{V}`$ given by $`h:(x,y)(xh^1,hy)`$ and thus descends to the geometric quotient $`\overline{X}=((\mathrm{M}_n)\times \overline{V})/H`$. We shall denote the point $`x\overline{X}`$ corresponding to the orbit of $`(g,v)(\mathrm{M}_n)\times \overline{V}`$ by $`[g,v]`$. The $`H`$-variety $`\overline{W}`$ and the $`\mathrm{PGL}_n`$-variety $`\overline{Y}`$ are defined in a similar manner. Our goal is to show that (i) every $`H`$-fixed points of $`\overline{X}`$ is of the form $`x=[g,v]`$, where $`gN_{\mathrm{PGL}_n}(H)`$ and $`v`$ is an $`H`$-fixed point of $`\overline{V}`$, and for any such point $`x`$, (ii) $`\mathrm{Stab}(x)=H`$ and (iii) up to sign, $`i(X,x,H)=\chi _1\mathrm{}\chi _{2r}`$. These assertions, in combination with Theorem 1.1, will prove that if $`X`$ and $`Y`$ (and hence, $`\overline{X}`$ and $`\overline{Y}`$) are birationally isomorphic then $`\chi _1\mathrm{}\chi _{2r}=\pm \eta _1\mathrm{}\eta _{2r}`$, i.e., (c) $``$ (a). To prove (i), assume $`x=[g,v]`$ is an $`H`$-fixed point of $`\overline{X}`$ for some $`g(\mathrm{M}_n)`$ and $`v\overline{V}`$. This means that for every $`hH`$ there exists an $`h^{}H`$ such that $`(hg,v)=(gh^{},(h^{})^1v)`$ in $`(\mathrm{M}_n)\times \overline{V}`$. Equivalently, $`hg=gh^{}`$ and $`(h^{})^1v=v`$. Consider the vector space $`k^n`$ of $`(n\times 1)`$-row vectors. The multiplication by $`g`$ on the right yields a linear map $`k^nk^n`$; let $`\mathrm{RKer}(g)`$ be the kernel of this map. Note that since $`g(\mathrm{M}_n)`$, this linear map is only defined up to a nonzero constant multiple but $`\mathrm{RKer}(g)`$ is well-defined. The equality $`hg=gh^{}`$ implies that $`\mathrm{RKer}(g)`$ is an $`H`$-invariant subspace of $`k^n`$ with respect to the right action of $`H`$; again, as $`H\mathrm{PGL}_n`$, the right multipication by an element $`hH`$ is a linear map $`k^nk^n`$ defined up to a nonzero constant multiple, but the notion of $`H`$-invariance of a linear subspace of $`k^n`$ is well-defined. Now recall that by Lemma 8.1(a) the $`n^2`$ elements of the form $`P_aD_\chi `$ which represent the elements of $`H\mathrm{PGL}(V)=\mathrm{PGL}_n`$ in $`\mathrm{GL}(V)=\mathrm{GL}_n`$ span $`\mathrm{M}_n`$ as a $`k`$-vector space. Thus the only $`H`$-invariant subspaces of $`k^n`$ are the ones that are invariant under all of $`\mathrm{M}_n`$, namely $`k^n`$ and $`(0)`$. If $`\mathrm{RKer}(g)=k^n`$ then $`g`$ is the zero matrix, which is impossible for any $`g(\mathrm{M}_n)`$. Thus we conclude that $`\mathrm{RKer}(g)=(0)`$. This means that $`g`$ is nonsingular, i.e., $`g\mathrm{PGL}_n`$. Now we can rewrite $`hg=gh^{}`$ as $`g^1hg=h^{}H`$; this shows that $`gN_{\mathrm{PGL}_n}(H)`$. Moreover, as $`h`$ ranges over $`H`$, $`h^{}=g^1hg`$ also ranges over all of $`H`$. Thus the equality $`(h^{})^1v=v`$ implies that $`v`$ is an $`H`$-fixed point of $`\overline{V}`$. This proves (i). From now on let $`x=[g,v]`$ be an $`H`$-fixed point of $`\overline{X}`$, where $`gN_{\mathrm{PGL}_n}(H)`$ and $`v`$ is an $`H`$-fixed point of $`\overline{V}`$. To prove (ii), assume $`g^{}\mathrm{Stab}(x)`$, i.e., $`g^{}[g,v]=[g,v]`$. Then $`g^{}g=gh^{}`$ for some $`h^{}H^{}`$. Since $`gN_{\mathrm{PGL}_n}(H)`$, we conclude that $`g^{}=ghg^1H`$, as desired. To prove (iii), first note that $`i(\overline{X},[g,v],H)=^{2r}\psi _g^{}\left(i(\overline{X},[1,v],H)\right)`$, where $`\psi _g:HH`$ is conjugation by $`gN_{\mathrm{PGL}_n}(H)`$ and $`^{2r}(\psi _g^{})`$ is the automorphism of $`^{2r}(H^{})`$ induced by $`\psi _g`$; see Remark 4.3. By Lemma 8.1(d) $`^{2r}\psi _g^{}`$ is the identity automorphism. Thus $`i(\overline{X},[g,v],H)=i(\overline{X},[1,v],H)`$. On the other hand, by Remark 4.2 $`i(\overline{X},[1,v],H)=i(\overline{V},v,H)`$. Finally, recall that for any $`v\overline{V}^H`$, $`i(\overline{V},v,H)=i(V,0_V,H)=\chi _1\mathrm{}\chi _{2r}`$; see Example 4.4. In summary, $$i(\overline{X},[g,v],H)=i(\overline{X},[1,v],H)=i(\overline{V},v,H)=\chi _1\mathrm{}\chi _{2r},$$ as claimed. ∎ ###### Remark 8.3. Recall that the exceptional group $`E_8`$ has a unique nontoral subgroup isomorphic to $`(/5)^3`$ (up to conjugation). Denote this subgroup by $`H`$. Then, modifying the proof of Proposition 8.2, we can show the following: Let $`V=\chi _1\chi _2\chi _3`$ and $`W=\eta _1\eta _2\eta _3`$ be faithful 3-dimensional representations of $`H`$, where $`\chi _i`$ and $`\eta _j`$ are characters of $`H`$. Then the following are equivalent: (a) $`\chi _1\chi _2\chi _3=\pm \eta _1\eta _2\eta _3`$ in $`^3(H^{})/5`$, (b) $`V`$ and $`W`$ are birationally isomorphic as $`H`$-varieties, and (c) $`E_8_HV`$ and $`E_8_HW`$ are birationally isomorphic as $`E_8`$-varieties. In particular, there are exactly two birational isomorphism classes of $`E_8`$-varieties of the form $`E_8_HV`$, where $`V`$ is a faithful 3-dimensional representation of $`H`$: one corresponds to $`\pm 1`$, and the other to $`\pm 2`$ in $`/5^3(H^{})`$. ###### Remark 8.4. Note that the $`\mathrm{PGL}_n`$-varieties $`X=\mathrm{PGL}_n_HV`$ and $`Y=\mathrm{PGL}_n_HW`$ of Proposition 8.2, are stably isomorphic. In fact, $`X\times 𝔸^1Y\times 𝔸^1`$ because $`X\times 𝔸^1=\mathrm{PGL}_n_H(V\times 𝔸^1)`$, $`Y\times 𝔸^1=\mathrm{PGL}_n_H(W\times 𝔸^1)`$, and $`V\times 𝔸^1W\times 𝔸^1`$ as $`H`$-varieties by Theorem 7.1. For the same reason $`X\times 𝔸^1`$ and $`Y\times 𝔸^1`$ are isomorphic $`E_8`$-varieties, if $`X`$ and $`Y`$ are as in Remark 8.3. ### Proof of Theorem 1.4 Recall that birational isomorphism classes of generically free irreducible $`\mathrm{PGL}_n`$-varieties $`X`$ with $`k(X)^{\mathrm{PGL}_n}=K`$ are in 1-1 correspondence with central simple algebras of degree $`n`$ over $`K`$; see e.g., \[Se, X.5\] or \[RY<sub>2</sub>, Section 3\]. In particular, by \[RY<sub>3</sub>, Lemma 4.2\], the algebra $`Q(\omega _1,\mathrm{},\omega _r)`$ of Theorem 1.4 corresponds to the variety $`X=\mathrm{PGL}_n_HV`$, where $`V`$ is a faithful $`2r`$-dimensional representations of $`H=A\times A^{}`$ constructed as follows. Choose a set of generators $`a_1,\mathrm{},a_r`$ for $`A=/n_1\times \mathrm{}\times /n_r`$ and a “dual” set of generators $`\chi _1,\mathrm{},\chi _r`$ for $`A^{}`$ so that $$\chi _i(a_j)=\{\begin{array}{cc}1\hfill & \text{if }ij\hfill \\ \omega _i\hfill & \text{if }i=j,\hfill \end{array}$$ Now note that each $`aA`$ defines a character of $`H=A\times A^{}`$ by $`(b,\eta )\eta (a)`$. Similarly, each $`\chi A^{}`$ gives rise to a character $`H=A\times A^{}k^{}`$ via $`(b,\eta )\chi (b)`$; we shall denote these characters by $`c(a)`$ and $`c(\chi )`$ respectively. In these notations, $$V=c(a_1)\mathrm{}c(a_r)c(\chi _1)^1\mathrm{}c(\chi _r)^1.$$ see \[RY<sub>3</sub>, Proof of Lemma 4.2\]. Similarly the $`\mathrm{PGL}_n`$-variety associated to $`Q(\omega _1^{m_1},\mathrm{},\omega _r^{m_r})`$ is $`Y=\mathrm{PGL}_n_HW`$, where $`W=c(a_1^{})\mathrm{}c(a_r^{})c(\chi _1^{})^1\mathrm{}c(\chi _r^{})^1`$. Here $`a_1^{},\mathrm{},a_r^{}`$ are generators of $`A`$ and $`\chi _1^{},\mathrm{},\chi _r^{}`$ are generators of $`A^{}`$ such that $$\chi _i^{}(a_j^{})=\{\begin{array}{cc}1\hfill & \text{if }ij\hfill \\ \omega _i^{m_i}\hfill & \text{if }i=j.\hfill \end{array}$$ A natural choice for $`a_i^{}`$ and $`\chi _i^{}`$ is $`a_i^{}=a_i`$ and $`\chi _i^{}=\chi _i^{m_i}`$, so that $$W=c(a_1)\mathrm{}c(a_r)c(\chi _1)^{m_1}\mathrm{}c(\chi _r)^{m_r}.$$ As we mentioned above, $`Q(\omega _1,\mathrm{},\omega _r)`$ and $`Q(\omega _1^{m_1},\mathrm{},\omega _r^{m_r})`$ are isomorphic as $`k`$-algebras iff their associated $`\mathrm{PGL}_n`$-varieties, $`X=\mathrm{PGL}_n_HV`$ and $`Y=\mathrm{PGL}_n_HW`$, are birationally isomorphic. By Proposition 8.2 $`X`$ and $`Y`$ are birationally isomorphic iff $`c(a_1)\mathrm{}c(a_r)c(\chi _1)^1\mathrm{}c(\chi _r)^1=`$ $`\pm c(a_1)\mathrm{}c(a_r)c(\chi _1)^{m_1}\mathrm{}c(\chi _r)^{m_r}\text{in }^{2r}(H^{})/n_r\text{.}`$ The last condition is equivalent to $`m_1\mathrm{}m_r=\pm 1(modn_1)`$. ∎
warning/0007/cond-mat0007244.html	ar5iv	text	# Band gaps of the primary metallic carbon nanotubes ## Abstract Primary metallic, or small gap semiconducting nanotubes, are tubes with band gaps that arise solely from breaking the bond symmetry due to the curvature. We derive an analytic expression for these gaps by considering how a general symmetry breaking opens a gap in nanotubes with a well defined chiral wrapping vector. This approach provides a straightforward way to include all types of symmetry breaking effects, resulting in a simple unified gap equation as a function of chirality and deformations. Recently, individual single wall nanotubes (SWNT) exhibiting small band gaps of the order of 10 meV where observed for the first time. SWNT’s can be classified according to their electronic band gap into three groups: semiconductors, small gap semiconductors and real metals. A semiconducting gap arises when the graphite Fermi points are not allowed in the tube’s Brillouin zone, which is given by distinct quantization lines according to the tube’s circumferential boundary conditions. Such a gap is of the order of 1 eV and was predicted to scale with $`1/\mathrm{R}`$ where $`\mathrm{R}`$ is the tube’s radius, a prediction that was later verified experimentally by a Scanning Tunneling Microscopy measurement of the density of states. The graphite Fermi points lie on a quantization line if $`\mathrm{mod}(\frac{nm}{3})=0`$, where $`n`$ and $`m`$ are the two integers defining the tube’s chiral vector. Tubes that satisfy this condition are called primary metallic. In a real nanotube, however, effects of curvature and deformation break the nearest neighbor bond symmetry, resulting in a shift of the two distinct Fermi points of graphite which lie at the corners of the hexagonal first Brillouin zone where the bonding and antibonding bands are degenerate ($`K`$-points). This shift may open a gap depending on the position of the new $`K`$-points relative to the circumferential quantization line. These gaps are about the value of room temperature and thus coined small gap semiconductors. Much interest has been devoted to the study of these small gap semiconducting tubes, which has led to a good basic understanding. Already from numerical calculations it has been known for a while that only armchair tubes retain zero gap and therefore are truly metallic while chiral tubes open small gaps because of the intrinsic curvature. These gaps depend on the radius and the chiral angle and were found to be lower than room temperature for tubes with radius exceeding 6Å. Other numerical studies also examined the effect of structural deformation such as stress and twist on the electronic band gap of nanotubes where it was found that the armchair tubes open a gap under twist but are not affected by the uniaxial stress while an opposite response arises in zigzag tubes. Analytic progress was made by Kane et al where the effects of deformations were modeled as a perturbing vector potential in an effective mass Hamiltonian. Those calculations were able to verify the numerical findings for the gaps from the intrinsic curvature as well as twists. After the completion of this work we also became aware of a more recent analytic work by Yang et al where both twist and stress on the nanotubes were considered (but not the intrinsic curvature). From the previous works the gaps can in principle be calculated for almost any shape of nanotube either by using a numerical method or by determining the exact metric and curvature tensors. We now reconsider the effects of deformations on primary metallic nanotubes and formulate the problem in terms of a general symmetry breaking in the tight binding model. This gives a straightforward analysis of the effect which results in a surprisingly simple and useful formula for the gap. Our compact expression combines the effects from intrinsic curvature, twist and stress and is only a function of the chiral wrapping vector $`(n,m)`$ of the tube. This gives a quick and easy way to determine the gap and allows for a good insight in the physical effects as we will describe below. Our results are a direct consequence of any symmetry breaking effect in the tight binding model of the graphite sheet in carbon nanotubes. Our starting point is the observation that the energy separation between the bonding and antibonding bands according to the graphite tight binding scheme is $`2\gamma _{i=1}^3e^{i\stackrel{}{k}\stackrel{}{R}_i}`$, where $`\stackrel{}{R}_i`$ are the nearest neighbor bond vectors and $`\gamma `$ is the transfer integral which is the nearest neighbor Hamiltonian matrix element. Since at zero temperature the bonding band is occupied and the antibonding empty, the Fermi points $`\stackrel{}{K}_F`$ lie at the band crossings, which are, for the unperturbed graphite, the six corners of the hexagonal first Brillouin zone. If we now break the symmetry of graphite and allow different transfer integrals $`\gamma _i`$ depending on the direction of the bonds $`\stackrel{}{R}_i`$ we arrive at a more general equation for the $`k`$-vectors at which the bands cross $$\underset{i=1}{\overset{3}{}}\gamma _ie^{i\stackrel{}{k}\stackrel{}{R}_i}=0,$$ (1) which defines the points of zero gap in $`k`$-space $`\stackrel{}{K}_F^{}`$. For small changes $`\gamma _i=\gamma +\delta \gamma _i`$ we expect small shifts in the band crossing location $`\stackrel{}{K}_F^{}=\stackrel{}{K}_F+\mathrm{\Delta }\stackrel{}{k}`$. Since we are dealing with primary metallic tubes, we know that $`_{i=1}^3e^{i\stackrel{}{K}_F\stackrel{}{R}_i}=0`$ where $`\stackrel{}{K}_F`$ are the unperturbed Fermi points. Working in the nanotube coordinates $`(\widehat{c},\widehat{t})`$ adopted from Ref. where $`\widehat{c}`$ is the circumferential direction and $`\widehat{t}`$ stands for the translational direction along the axis, the bond vectors are $$\begin{array}{c}\stackrel{}{R}_1=\frac{a}{2c_h}[(n+m)\widehat{c}\frac{1}{\sqrt{3}}(nm)\widehat{t}]\hfill \\ \stackrel{}{R}_2=\frac{a}{2c_h}[m\widehat{c}+\frac{1}{\sqrt{3}}(2n+m)\widehat{t}]\hfill \\ \stackrel{}{R}_3=\frac{a}{2c_h}[n\widehat{c}\frac{1}{\sqrt{3}}(n+2m)\widehat{t}],\hfill \end{array}$$ (2) where $`a2.49\mathrm{\AA }`$ is the length of the honeycomb unit vector and $`c_h=\sqrt{n^2+nm+m^2}`$ is the circumference in units of $`a`$. Both inequivalent Fermi points in graphite give the same result when estimating the gap so it is sufficient to consider just one of them. For an unperturbed Fermi point we write $$\stackrel{}{K}_F=\frac{2\pi }{3ac_h}[(m+2n)\widehat{k}_c+m\sqrt{3}\widehat{k}_t],$$ (3) where $`\widehat{k}_c`$ and $`\widehat{k}_t`$ correspond to the $`k`$-vectors along the circumferential and translational directions, respectively. In order to get the gap, we now need to know the distance between the new Fermi point to the nearest quantization line at the quantized circumferential $`k_c`$-values. Since the quantization lines are parallel to $`\widehat{k}_t`$, this distance is given by $`\mathrm{\Delta }k_c`$ the shift along the circumferential direction $`\widehat{k}_c`$. Expanding Eq. (1) to linear order in the perturbations $`\delta \gamma _i\gamma _i\gamma `$ we find $$\mathrm{\Delta }k_c=\frac{1}{ac_h\gamma \sqrt{3}}[\delta \gamma _1(mn)+\delta \gamma _2(2n+m)\delta \gamma _3(n+2m)]$$ (4) The gap is then obtained by exploiting the fact that close to the Fermi point the dispersion relation is linear and isotropic $$E_g=\sqrt{3}a\gamma \|\mathrm{\Delta }k_c\|.$$ (5) We now want to determine the changes to the transfer integrals $`\delta \gamma _i`$ due to the curvature and deformations. To first order this change is proportional to the change of the bond length between two neighboring carbon atoms, but may also be created by a misalignment of two neighboring $`\pi `$ orbitals. In general we find that we can always express the change of the nearest neighbor transfer integrals in terms of a bond deviation matrix D $$\delta \gamma _i=\stackrel{}{R}_i\text{D}\stackrel{}{R}_i/R_i^2.$$ (6) This deviation matrix is useful for describing the effects of stress, twists, and curvature in a simple unified way as we will see below. All nanotubes have an intrinsic curvature which causes hybridization of the otherwise orthogonal $`\pi `$ and $`\sigma `$ orbitals. Since the $`\sigma `$ bands lie normally far from the Fermi energy, we only consider the hybridization effect on the $`\pi `$ band, which crosses the Fermi point in primary metallic tubes as long as the tube’s radius is $`R2.4\mathrm{\AA }`$. Following the calculations of Slater and Koster we can assume that the transfer integrals are proportional to cosine of the misalignment angle $`\varphi `$ between two neighboring $`\pi `$ orbitals. (A calculation that takes into account the full rehybridization of all orbitals will be considered in a future study.) Using $`\mathrm{cos}\varphi 1\frac{1}{8}\left(\frac{a_{cc}}{\mathrm{R}}\right)^2`$, we can immediately express the deformation tensor for the intrinsic curvature in the basis of $`\widehat{c}`$ and $`\widehat{t}`$ $$\text{D}^{\mathrm{curv}}=\left[\begin{array}{cc}\frac{\gamma }{8}\left(\frac{a_{cc}}{\mathrm{R}}\right)^2& 0\\ 0& 0\end{array}\right],$$ (7) where $`a_{cc}=a/\sqrt{3}`$ is the carbon-carbon bond length and $`\mathrm{R}=ac_h/2\pi `$ is the tube’s radius. From Eqs. (4-7) we find $$E_g=\frac{\gamma \pi ^2}{8c_h^5}(nm)(2n^2+5nm+2m^2).$$ (8) This formula agrees remarkably well with previous numerical studies if we choose $`\gamma =2.5`$eV and also agrees with the results of Ref. . Equation (8) gives the energy gap of all primary metallic nanotubes without any applied deformations. One observes that the armchair nanotubes $`m=n`$ are the only real metallic tubes, while the primary metallic zigzag tubes ($`m=0`$), open the highest gaps. Next we want to examine the effect of a general two dimensional structural deformation such as twists and stress on the gap in the primary metallic tubes. Our deformation can be written as $`\stackrel{}{R}(\text{I}+\text{S})\stackrel{}{R}`$, where $`\stackrel{}{R}`$ is any vector on the tube’s surface, I is the identity matrix and S is the deformation matrix $$\text{S}\left[\begin{array}{cc}ϵ_c& \xi \\ 0& ϵ_t\end{array}\right].$$ (9) Here $`ϵ_c`$ and $`ϵ_t`$ are uniaxial stresses along the circumferential and the translational directions and $`\xi `$ is the strain (nanotube twist). The bond deviation matrix is then given by $`\text{D}^{\mathrm{deform}}=\|\stackrel{}{R}\|b\text{S}`$, where $`b3.5`$ eV/Å is the linear change in the transfer integral with a change in the bond length $`\gamma _i\gamma _i+b\|\mathrm{\Delta }\stackrel{}{R}_i\|`$, and $`\|\stackrel{}{R}\|=a/\sqrt{3}`$ is the bond length. We now use $`\text{D}^{\mathrm{deform}}`$ to obtain the $`\delta \gamma _i`$ of Eq. (6) and as we did with the curvature, inserting in Eq. (4) and using the dispersion relation we find $`E_g`$ $`=`$ $`{\displaystyle \frac{ab}{4c_h^3}}\|\sqrt{3}(nm)(2n^2+5nm+2m^2)(ϵ_cϵ_t)`$ (10) $`9nm(n+m)\xi \|.`$ (11) This equation is the response to a two-dimensional linear deformation within the graphite sheet. We notice from Eq. (11) that in the presence of equal uniaxial stresses in both directions $`\widehat{c}`$ and $`\widehat{t}`$, a gap is not opened, as expected since the bonds would maintain their symmetry. The response of armchair and zigzag tubes is complimentary as noticed previously in the numerical studies, i.e. zigzag tubes have the maximum sensitivity for a uniaxial deformation ($`ϵ_c`$ or $`ϵ_t`$) but are insensitive to twists $`\xi `$, while the opposite is true for armchairs, reaffirming that a twist deformation is the only possible source for an energy gap in the armchair tube. A uniaxial stress $`ϵ_c`$ around the circumference corresponds to a global change of radius. This effect can come about as a time dependent deformation due to lattice vibrations in the breathing mode. Realistic static deformations on the other hand correspond to the intrinsic curvature, a stress along the tube $`ϵ_t`$ and a twist $`\xi `$. Therefore the total gap from static deformations is given by combining Eq. (8) and Eq. (11) with $`ϵ_c=0`$. The total band equation now reads $`E_g`$ $`=`$ $`\|({\displaystyle \frac{\gamma \pi ^2}{8c_h^5}}{\displaystyle \frac{ab\sqrt{3}}{4c_h^3}}ϵ_t)(nm)(2n^2+5nm+2m^2)`$ (12) $`{\displaystyle \frac{9ab}{4c_h^3}}nm(n+m)\xi \|`$ (13) which is the main result of our paper. In some cases it may be useful to express this formula in terms of the chiral angle $`\alpha `$ and the radius $`\mathrm{R}`$ of the tube, which gives $$E_g=\left\|\left(\frac{\gamma a^2}{16\mathrm{R}^2}\frac{ab\sqrt{3}}{2}ϵ_t\right)\mathrm{sin}3\alpha \frac{ab\sqrt{3}}{2}\xi \mathrm{cos}3\alpha \right\|.$$ (14) In this form our results are then consistent with previous calculations where twists and the intrinsic curvature (but not stress) have been considered as a function of the chiral angle. After the completion of this work, a paper deriving the change in the band gap due to deformations was published, which is also consistent with the angular dependence of the deformation part of Eq. (14) (i.e. without the intrinsic curvature). We see that Eq. (14) shows an interesting interplay between curvature and deformation effects as plotted in in Fig. 1. A very small twist can actually remove the gap due to the intrinsic curvature. For a given radius the gapless state is therefore moved towards tubes with higher chirality. In summary, we have presented a straightforward procedure to calculate the energy gap induced by a general broken bond symmetry. This leads to a simple analytic expression for the band gaps from both the intrinsic curvature and applied deformations, which provides a quick and reliable way to estimate the physical effects. These gaps have important consequences since they are generally of the same order as room temperature for most primary metallic SWNT. Only armchair tubes are generically gapless, but a very small twist induces a gap of the order of other small gap semiconducting nanotubes. Such a small twist, on the other hand, moves the gapless state to tubes with higher chirality. ###### Acknowledgements. We would like to thank Vitali Shumeiko for inspirational discussions. This research was supported in part by the Swedish Natural Science Research Council.
warning/0007/hep-ph0007232.html	ar5iv	text	# Real and Virtual Compton Scattering off the Nucleon ## 1 Introduction In the study of hadron structure, one of the main questions is how hadrons and nuclei are built from quarks and gluons and how one goes over from a description in terms of quark and gluon degrees of freedom to a description in terms of hadronic degrees of freedom. Nowadays, precision experiments at high energy have established Quantum Chromo Dynamics (QCD) as the gauge theory of strong interactions describing the dynamics between colored quarks and gluons. QCD exhibits the property that the interaction between the quarks becomes weak at very short distances, which is known as asymptotic freedom, and which allows us to use perturbation theory to describe high energy strong interaction phenomena. On the other hand at low energy, the QCD coupling constant grows, and quarks and gluons are confined into colorless mesons and baryons, which are the particles as seen through experiments. In this phase of hadronic matter, the underlying chiral symmetry of QCD, due to the nearly massless up, down and approximately also strange quarks, is spontaneously broken. To describe this regime directly from the QCD Lagrangian is a hard task which still defies a solution due to the strong coupling constant requiring non-perturbative methods. Probably the most promising and direct approach is the numerical solution of QCD through lattice calculations. For static hadronic properties, such as e.g. masses, much progress has already been made, but for more complicated hadron structure quantities, such as e.g. nucleon parton distributions, lattice calculations are still in their infancy. In absence of a full numerical solution of QCD, which would be able to describe the rich complexity of the hadronic many body systems from their underlying dynamics, a complementary strategy to refine our understanding of hadron structure is to perform new types of precision experiments in kinematical regimes at low energy, which require an inherent non-perturbative description. Besides, one may perform new types of experiments at high energies, in those kinematical regimes where factorization theorems have been proven. Such experiments will allow us to use an accurate perturbative QCD description of the reaction dynamics as a tool to extract new types of non-perturbative hadron structure information. In this quest at the intersection of particle and nuclear physics, the experiments performed with electromagnetic probes play an important role. A new generation of precision experiments has become possible with the advent of the new electron accelerators and in combination with high precision and large acceptance detectors. In particular, high precision Compton scattering experiments have become a reality in recent years. In Compton scattering, a real or virtual photon interacts with the nucleon and a real photon is emitted in the process. As this is a purely electromagnetic process, it yields small cross sections (compared to hadronic reactions), but on the other hand constitutes a clean probe of hadron structure. In this paper, a review will be given of very recent developments in the field of real and virtual Compton scattering off the nucleon. I shall discuss real and virtual Compton scattering at the same time and point out their complementarity and the differences in the extracted nucleon structure information. The emphasis is on those kinematical regimes where a fruitful interpretation is terms of nucleon structure observables has been shown to be possible. Virtual Compton scattering (VCS) off the nucleon has been reviewed before in Ref. Gui (98), which is referred to for most technical details. For the VCS part, the emphasis is on the progress over the past two years in the light of the first high precision VCS data in the threshold regime now available, and on the rapid development in the field of deeply virtual Compton scattering (DVCS) at large photon virtuality. In section 2, the real Compton scattering (RCS) process below pion threshold is discussed. In this regime, the RCS process can be interpreted as the global response of the nucleon to an applied electromagnetic field, which allows us to access global nucleon polarizabilities. A dispersion relation formalism is described, which was developed as a method to minimize the model uncertainty in the extraction of nucleon polarizabilites from both unpolarized and polarized RCS data at low energy. In section 3, the virtual Compton scattering (VCS) reaction at low energy is discussed. It consists of a generalization of RCS in which both energy and momentum of the virtual photon can be varied independently, which allows us to extract response functions, parametrized by the so-called generalized polarizabilities (GP’s) of the nucleon. A first dedicated VCS experiment was performed at the MAMI accelerator, and two combinations of those GP’s have been measured. Their values are compared with nucleon structure model predictions. Further experimental programs are underway at the major electron accelerators to measure the VCS observables. It is outlined how results of such experiments can be interpreted in terms of the nucleon GP’s, and in particular how polarization observables can separate the different GP’s. Besides the low energy region, RCS at high energy and large momentum transfer is a tool to access information on the partonic structure of the nucleon. In section 4, a leading order perturbative QCD calculation of RCS is described, which was developed to extract the valence quark wave function of the nucleon. Such an approach is compared with competing mechanisms, and it is pointed out how the planned experiments can shed light on this field. Section 5 discusses the recent developments in deeply virtual Compton scattering (DVCS) and associated meson electroproduction reactions at high energy, high photon virtuality $`Q^2`$ and small momentum transfer to the nucleon. It is shown how a unified theoretical description of those processes has recently emerged and how it gives access to new parton distributions, the so-called skewed parton distributions, which are generalizations of the usual parton distributions as known from inclusive deep inelastic lepton scattering. Leading order perturbative QCD calculations of DVCS and different meson electroproduction reactions, using an ansatz for the skewed parton distributions, are discussed in the kinematical regimes accessible at present or planned facilities. The corrections to those leading order QCD amplitudes and further open questions in this field are touched on briefly. In section 6, the calculation of the QED radiative corrections to the VCS process is discussed, which is indispensable to accurately extract the nucleon structure information from VCS experiments. Finally, the conclusions and perspectives are given in section 7. ## 2 Real Compton scattering (RCS) and nucleon polarizabilities ### 2.1 Introduction In the study of nucleon structure, real Compton scattering (RCS) at low energy is a precision tool to access global information on the nucleon ground state and its excitation spectrum. RCS off the nucleon is determined by 6 independent helicity amplitudes, which are functions of two variables, e.g. the Lorentz invariant variables $`\nu `$ (related to the $`lab`$ energy of the incident photon) and $`t`$ (related to the momentum transfer to the target). In the limit $`\nu 0`$, corresponding to wavelengths much larger than the nucleon size, the general structure of these amplitudes is governed by low energy theorems (LET) based on Lorentz invariance, gauge invariance and crossing symmetry. These theorems predict that the leading terms of an expansion in $`\nu `$ are determined by global properties of the nucleon, i.e. its charge, mass and anomalous magnetic moment. Furthermore, the internal structure shows up only at relative order $`\nu ^2`$ and can be parametrized in terms of the polarizabilities. In this way, there appear 6 polarizabilities for the nucleon, the electric and magnetic (scalar) polarizabilities $`\alpha `$ and $`\beta `$ respectively, familiar from classical physics, and 4 spin (vector) polarizabilities $`\gamma _1`$ to $`\gamma _4`$, originating from the spin 1/2 nature of the nucleon. These polarizabilities describe the response of the nucleon’s charge, magnetization, and spin distributions to an external quasistatic electromagnetic field. As such the polarizabilities are fundamental structure constants of the composite system. The differential cross section for RCS in the limit $`\nu 0`$ is given by the (model independent) Thomson term, as a consequence of the LET. The electric and magnetic polarizabilities then appear, in a low-energy expansion for RCS, as interference between the Thomson term and the subleading terms, i.e. as contribution of $`O(\nu ^2)`$ in the differential cross section. In this way, $`\alpha `$ and $`\beta `$ can in principle be separated by studying the RCS angular distributions. However, it has never been possible to isolate this term and thus to determine the polarizabilities in a model independent way. The obvious reason is that, for sufficiently small energies, say $`\nu 40`$ MeV, the structure effects are extremely small and hence the statistical errors for the polarizabilities large. Therefore, one has to go to larger energies, where the higher terms in the expansion become increasingly important and where also the spin polarizabilities come into play. To determine the nucleon polarizabilities from RCS observables, a reliable estimate of higher terms in the energy is therefore of utmost importance. To this end, actual experiments were usually analyzed in an unsubtracted dispersion relation formalism at fixed $`t`$ Lvo (97). Using such an analysis, the proton scalar polarizabilities were derived from Compton scattering data below pion threshold, with the results Mac (95) : $`\alpha `$ $`=`$ $`\left(12.1\pm \mathrm{\hspace{0.17em}0.8}\pm \mathrm{\hspace{0.17em}0.5}\right)\times \mathrm{\hspace{0.17em}10}^4\mathrm{fm}^3,`$ $`\beta `$ $`=`$ $`\left(2.1\mathrm{\hspace{0.17em}0.8}\mathrm{\hspace{0.17em}0.5}\right)\times \mathrm{\hspace{0.17em}10}^4\mathrm{fm}^3.`$ (1) Very recently, new RCS data on the proton below pion threshold have become available Olm (00). These data increase the available world data set substantially, and yield, in an unsubtracted DR formalism, the results : $`\alpha `$ $`=`$ $`\left(11.89\pm \mathrm{\hspace{0.17em}0.57}\right)\times \mathrm{\hspace{0.17em}10}^4\mathrm{fm}^3,`$ $`\beta `$ $`=`$ $`\left(1.17\pm \mathrm{\hspace{0.17em}0.75}\right)\times \mathrm{\hspace{0.17em}10}^4\mathrm{fm}^3.`$ (2) The sum of the scalar polarizabilities, which appears in the forward spin averaged Compton amplitude, can be determined directly from the total photoabsorption cross section by Baldin’s sum rule Bal (60), which yields a rather precise value : $`\alpha +\beta `$ $`=`$ $`\left(14.2\pm \mathrm{\hspace{0.17em}0.5}\right)\times \mathrm{\hspace{0.17em}10}^4\mathrm{fm}^3,`$ (3) $`=`$ $`\left(13.69\pm \mathrm{\hspace{0.17em}0.14}\right)\times \mathrm{\hspace{0.17em}10}^4\mathrm{fm}^3,`$ (4) with (3) from Ref. Dam (70) and (4) from Ref. Bab (98). Similarly, the proton forward spin polarizability can be evaluated by an integral over the difference of the absorption cross sections in states with helicity 3/2 and 1/2, $`\gamma _0=\gamma _1\gamma _22\gamma _4`$ $`=`$ $`1.34\times \mathrm{\hspace{0.17em}10}^4\mathrm{fm}^4,`$ (5) $`=`$ $`0.80\times \mathrm{\hspace{0.17em}10}^4\mathrm{fm}^4,`$ (6) with (5) from Ref. San (94) and (6) from Ref. Dre00a . While the predictions of Eqs. (5,6) rely on pion photoproduction multipoles, the helicity cross sections have now been directly determined at MAMI by scattering photons with circular polarizations on polarized protons Are (00). The contribution to $`\gamma _0`$ for the proton within the measured integration range (200 MeV $`\nu `$ 800 MeV) is Are (00) : $$\gamma _0\|_{\mathrm{\hspace{0.17em}200}\mathrm{MeV}}^{\mathrm{\hspace{0.17em}800}\mathrm{MeV}}=\left(1.68\pm 0.07\right)\times \mathrm{\hspace{0.17em}10}^4\mathrm{fm}^4.$$ (7) The contribution below 200 MeV can be estimated with the HDT pion photoproduction multipoles Han (98) to yield +1.0, and an estimate of the contribution above 800 MeV based on the SAID pion photoproduction multipoles Arn (96) yields -0.02, which results in a total value : $`\gamma _0`$ = -0.7 (here and in the following, all spin polarizabilities are given in units $`10^4\mathrm{fm}^4`$). Furthermore, unpolarized RCS data in the $`\mathrm{\Delta }`$(1232)-region were used to give - within a dispersion relation formalism - a first prediction for the so-called backward spin polarizability of the proton, i.e. the particular combination $`\gamma _\pi =\gamma _1+\gamma _2+2\gamma _4`$ entering the Compton spin-flip amplitude at $`\theta =180^{}`$ Ton (98) : $$\gamma _\pi =\left[27.1\pm \mathrm{\hspace{0.17em}2.2}(\mathrm{stat}+\mathrm{syst})\genfrac{}{}{0pt}{}{+2.8}{2.4}(\mathrm{model})\right]\times \mathrm{\hspace{0.17em}10}^4\mathrm{fm}^4.$$ (8) These values for the polarizabilities can be compared with our present day theoretical understanding from chiral perturbation theory (ChPT). A calculation to $`O(p^4)`$ in heavy baryon ChPT (HBChPT), where the expansion parameter $`p`$ is an external momentum or the quark mass, yields (here and in the following, $`\alpha `$ and $`\beta `$ are given in units $`10^4\mathrm{fm}^3`$) : $`\alpha =10.5\pm 2.0`$ and $`\beta =3.5\pm 3.6`$, the errors being due to 4 counter terms entering to that order, which were estimated by resonance saturation Ber (93). One of these counter terms describes the large paramagnetic contribution of the $`\mathrm{\Delta }`$(1232) resonance, which is partly cancelled by large diamagnetic contributions of pion-nucleon (N$`\pi `$)-loops. In view of the importance of the $`\mathrm{\Delta }`$ resonance, the calculation was also done by including the $`\mathrm{\Delta }`$ as a dynamical degree of freedom. This adds a further expansion parameter, the difference of the $`\mathrm{\Delta }`$ and nucleon masses (“$`ϵ`$ expansion”). A calculation to $`O(ϵ^3)`$ yielded $`\alpha `$ = 12.2 + 0 + 4.2 = 16.4 and $`\beta `$ = 1.2 + 7.2 + 0.7 = 9.1, the 3 separate terms referring to contributions of N$`\pi `$-loops (which is the $`O(p^3)`$ result), $`\mathrm{\Delta }`$-pole terms, and $`\mathrm{\Delta }\pi `$-loops Hem97a ; Hem (98). These $`O(ϵ^3)`$ predictions are clearly at variance with the data, in particular $`\alpha +\beta =25.5`$ is nearly twice the rather precise value determined from Baldin’s sum rule Eq. (4). The spin polarizabilities have also been calculated in HBChPT. The $`O(ϵ^3)`$ predictions for the proton are Hem (98) : $`\gamma _0=4.62.40.2+0=+2.0`$, and $`\gamma _\pi =4.6+2.40.243.5=36.7`$, the 4 separate contributions referring to N$`\pi `$-loops ($`O(p^3)`$ result), $`\mathrm{\Delta }`$-poles, $`\mathrm{\Delta }\pi `$-loops, and the triangle anomaly, in that order. It is obvious that the anomaly or $`\pi ^0`$-pole gives by far the most important contribution to $`\gamma _\pi `$, and that it would require surprisingly large higher order contributions to bring $`\gamma _\pi `$ close to the value of Eq. (8). Recently, the N$`\pi `$-loop contribution to the spin polarizabilities have been evaluated in HBChPT to $`O(p^4)`$ by several groups Ji (00); Vij (00); Gel (00). In Refs. Ji (00); Vij (00), the result for the proton is $`\gamma _0=+4.58.4`$, where the two contributions are the $`O(p^3)`$ and $`O(p^4)`$ N$`\pi `$-loop contributions, in this order. Based on the large $`O(p^4)`$ correction term, the authors in Ji (00); Vij (00) call the convergence of the chiral expansion into question. However in Ref. Gel (00), different results were obtained for the $`O(p^4)`$ N$`\pi `$-loop contributions to the 4 spin polarizabilities. It was argued that these differences are due to how one defines and extracts the $`O(p^4)`$ spin-dependent polarizabilities in chiral effective field theories. Following the procedure of Ref. Gel (00), which removes first all one-particle reducible contributions from the spin-dependent Compton amplitude, the resulting values for $`\gamma _0`$ and $`\gamma _\pi `$ of the proton are $`\gamma _0=+4.65.6=1.0`$, and $`\gamma _\pi =+4.61.2=+3.4`$ (without the $`\pi ^0`$-pole), the separate contributions being again the $`O(p^3)`$ and $`O(p^4)`$ N$`\pi `$-loop contributions respectively. For $`\gamma _0`$, a convergence of HBChPT at order $`O(p^4)`$ was not expected Gel (00), whereas the result for $`\gamma _\pi `$ \- when adding the $`\pi ^0`$-pole contribution - is not compatible with the estimate of Eq. (8) obtained by Ref. Ton (98). In order to refine our present understanding of the nucleon polarizabilities, a better understanding of the convergence of the HBChPT expansion is absolutely necessary, and it is to be hoped that a calculation to $`O(ϵ^4)`$ will clarify the status. On the other hand, it is also indispensable to minimize any model dependence in the extraction of the polarizabilities from the data. To this end, a fixed-$`t`$ subtracted dispersion relation (DR) formalism was developed in Ref. Dre00a for RCS off the nucleon at photon energies below 500 MeV, as a formalism to extract the nucleon polarizabilities with a minimum of model dependence as is described in the following. ### 2.2 Fixed-t subtracted dispersion relations for RCS To perform a dispersion theoretical analysis of Compton scattering, one has to calculate the 6 independent structure functions $`A_i(\nu ,t)`$, $`i=1,\mathrm{},6`$ (defined in Ref. Lvo (97)). They are functions of the usual Mandelstam variable $`t`$, and of $`\nu `$, defined in terms of the Mandelstam variables $`s`$ and $`u`$ as $`\nu =(su)/(4m_N)`$, with $`m_N`$ the nucleon mass. The invariant amplitudes $`A_i`$ are free of kinematical singularities and constraints, and because of the crossing symmetry they satisfy the relation $`A_i(\nu ,t)=A_i(\nu ,t)`$. Assuming further analyticity and an appropriate high-energy behavior, the amplitudes $`A_i`$ fulfill unsubtracted DR at fixed $`t`$ : $$\mathrm{R}eA_i(\nu ,t)=A_i^B(\nu ,t)+\frac{2}{\pi }𝒫_{\nu _{thr}}^+\mathrm{}𝑑\nu ^{}\frac{\nu ^{}\mathrm{I}m_sA_i(\nu ^{},t)}{\nu ^2\nu ^2},$$ (9) where $`A_i^B`$ are the Born (nucleon pole) contributions, and where $`\mathrm{I}m_sA_i`$ are the discontinuities across the $`s`$-channel cuts of the Compton process, starting from the threshold for pion production, $`\nu _{thr}`$. However, such unsubtracted DR require that at high energies ($`\nu \mathrm{}`$) the amplitudes $`\mathrm{I}m_sA_i(\nu ,t)`$ drop fast enough such that the integral of Eq. (9) is convergent and the contribution from the semi-circle at infinity can be neglected. For real Compton scattering, Regge theory predicts the following high-energy behavior for $`\nu \mathrm{}`$ and fixed $`t`$ Lvo (97) : $$A_{1,2}\nu ^{\alpha (t)},A_{3,5,6}\nu ^{\alpha (t)2},A_4\nu ^{\alpha (t)3},$$ (10) where $`\alpha (t)1`$ is the Regge trajectory. Due to the high-energy behavior of Eq. (10), the unsubtracted dispersion integral of Eq. (9) diverges for the amplitudes $`A_1`$ and $`A_2`$. In order to obtain useful results for these two amplitudes, L’vov et al. Lvo (97) proposed to close the contour of the integral in Eq. (9) by a semi-circle of finite radius $`\nu _{max}`$ in the complex plane (instead of the usually assumed infinite radius!), i.e. the real parts of $`A_1`$ and $`A_2`$ are calculated from the decomposition $$\mathrm{R}eA_i(\nu ,t)=A_i^B(\nu ,t)+A_i^{int}(\nu ,t)+A_i^{as}(\nu ,t),$$ (11) with $`A_i^{int}`$ the $`s`$-channel integral from pion threshold $`\nu _{thr}`$ to a finite upper limit $`\nu _{max}`$, and an ‘asymptotic contribution’ $`A_i^{as}`$ representing the contribution along the finite semi-circle of radius $`\nu _{max}`$ in the complex plane. In the actual calculations, the $`s`$-channel integral is typically evaluated up to a maximum photon energy of about $`1.5`$ GeV, for which the imaginary part of the amplitudes can be expressed through unitarity by meson photoproduction amplitudes (mainly 1$`\pi `$ and 2$`\pi `$ photoproduction) taken from experiment. All contributions from higher energies are then absorbed in the asymptotic terms $`A_i^{as}`$, which are replaced by a finite number of energy independent poles in the $`t`$ channel. In particular the asymptotic part of $`A_1`$ is parametrized by the exchange of a scalar particle in the $`t`$ channel, i.e. an effective “$`\sigma `$ meson” Lvo (97). In a similar way, the asymptotic part of $`A_2`$ is described by the $`\pi ^0`$ $`t`$-channel pole. This procedure is relatively safe for $`A_2`$ because of the dominance of the $`\pi ^0`$ pole or triangle anomaly, which is well established both experimentally and on general grounds as Wess-Zumino-Witten term. However, it introduces a considerable model-dependence in the case of $`A_1`$. It was therefore the aim of Ref. Dre00a to avoid the convergence problem of unsubtracted DR and the phenomenology necessary to determine the asymptotic contribution. To this end, it was proposed to consider DR’s at fixed $`t`$ that are once subtracted at $`\nu =0`$, $`\mathrm{R}eA_i(\nu ,t)`$ $`=`$ $`A_i^B(\nu ,t)+\left[A_i(0,t)A_i^B(0,t)\right]`$ (12) $`+`$ $`{\displaystyle \frac{2}{\pi }}\nu ^2𝒫{\displaystyle _{\nu _{thr}}^+\mathrm{}}𝑑\nu ^{}{\displaystyle \frac{\mathrm{I}m_sA_i(\nu ^{},t)}{\nu ^{}(\nu ^2\nu ^2)}}.`$ These subtracted DR should converge for all 6 invariant amplitudes due to the two additional powers of $`\nu ^{}`$ in the denominator, and they are essentially saturated by the $`\pi N`$ intermediate states. In other words, the lesser known contributions of two and more pions as well as higher continua are small and may be treated reliably by simple models. The price to pay for this alternative is the appearance of the subtraction functions $`A_i(\nu =0,t)`$, which have to be determined at some small (negative) value of $`t`$. This was achieved by setting up once-subtracted DR, this time in the variable $`t`$ Dre00a : $`A_i(0,t)A_i^B(0,t)=a_i+a_i^{tpole}`$ $`+{\displaystyle \frac{t}{\pi }}\left({\displaystyle _{(2m_\pi )^2}^+\mathrm{}}𝑑t^{}{\displaystyle _{\mathrm{}}^{2m_\pi ^24Mm_\pi }}𝑑t^{}\right){\displaystyle \frac{\mathrm{I}m_tA_i(0,t^{})}{t^{}(t^{}t)}},`$ (13) where the six coefficients $`a_iA_i(0,0)A_i^B(0,0)`$ are simply related to the six polarizabilities $`\alpha ,\beta ,\gamma _1,\gamma _2,\gamma _3,\gamma _4`$ (see Ref. Dre00a for details), and where $`a_i^{tpole}`$ represents, in the case of $`A_2`$, the contribution of the $`\pi ^0`$ pole in the $`t`$-channel. To evaluate the dispersion integrals, the imaginary part of the Compton amplitude due to the $`s`$-channel cuts in Eq. (12) is determined, through the unitarity relation, from the scattering amplitudes of photoproduction on the nucleon. Due to the energy denominator $`1/\nu ^{}(\nu ^2\nu ^2)`$ in the subtracted dispersion integrals, the most important contribution is from the $`\pi N`$ intermediate states, while mechanisms involving more pions or heavier mesons in the intermediate states are largely suppressed. In Ref. Dre00a , the $`\pi N`$ contribution was then evaluated using the pion photoproduction multipole amplitudes of Ref. Han (98) at photon energies below 500 MeV, and at the higher energies using the SAID multipoles (SP98K solution) Arn (96) as input. The multipion channels (in particular the $`\pi \pi N`$ channels) were approximated by the inelastic decay channels of the $`\pi N`$ resonances. It was found, however, that in the subtracted DR formalism, the sensitivity to the multipion channels is very small and that subtracted DR are essentially saturated at $`\nu `$ 0.4 GeV. The subtracted $`t`$-channel dispersion integral in Eq. (13) from $`4m_\pi ^2`$ to $`+\mathrm{}`$ is essentially saturated by the imaginary part of the $`t`$-channel amplitude $`\gamma \gamma N\overline{N}`$ due to $`\pi \pi `$ intermediate states. The dependence of the subtraction functions on momentum transfer $`t`$ can be calculated by including the experimental information on the $`t`$-channel process through $`\pi \pi `$ intermediate states as $`\gamma \gamma \pi \pi N\overline{N}`$. In Ref. Dre00a , a unitarized amplitude for the $`\gamma \gamma \pi \pi `$ subprocess was constructed, and a good description of the available data was found. This information is then combined with the $`\pi \pi N\overline{N}`$ amplitudes determined from dispersion theory by analytical continuation of $`\pi N`$ scattering. In this way, one avoids the uncertainties in Compton scattering associated with the two-pion continuum in the $`t`$ channel, usually modeled through the exchange of a somewhat fictitious $`\sigma `$ meson. The second integral in Eq. (13) extends from $`\mathrm{}`$ to $`2(m_\pi ^2+2Mm_\pi )0.56`$ GeV<sup>2</sup>. As we address Compton scattering for photon energies below about 500 MeV, the value of $`t`$ stays sufficiently small so that the denominator in the integral provides a rather large suppression, resulting in a small contribution. The contribution along the negative $`t`$-cut is estimated in the calculations Dre00b by saturation with $`\mathrm{\Delta }`$ intermediate states. Altogether the remaining uncertainties in the $`s`$\- and $`t`$\- channel subtracted integrals due to unknown high-energy contributions, are estimated to be less than 1%. As a consequence, this subtracted DR formalism provides a direct cross check between Compton scattering and one-pion photoproduction. Although all 6 subtraction constants $`a_1`$ to $`a_6`$ of Eq. (13) could be used as fit parameters in the present formalism, the fit was restricted to the parameters $`a_1`$ and $`a_2`$, or equivalently to $`\alpha \beta `$ and $`\gamma _\pi `$ in Dre00a . The subtraction constants $`a_4,a_5`$ and $`a_6`$ were calculated through an unsubtracted sum rule (Eq. (9) for $`\nu =t=0`$). The remaining subtraction constant $`a_3`$, related to $`\alpha +\beta `$ by $`\alpha +\beta =(a_3+a_6)/(2\pi )`$, was fixed through Baldin’s sum rule Bal (60), using the value $`\alpha +\beta =13.69\times 10^4`$ fm<sup>3</sup> Bab (98). ### 2.3 Results for RCS observables Since the polarizabilities enter as subtraction constants, the subtracted dispersion relations can be used to extract the nucleon polarizabilities from RCS data with a minimum of model dependence. The present formalism can be applied up to photon energies of about 500 MeV. Below pion production threshold, RCS data were usually analyzed to extract $`\alpha `$ and $`\beta `$. However, it was shown that the sensitivity to $`\gamma _\pi `$ is not at all negligible, especially at the backward angles and the higher energies, so that both $`\alpha \beta `$ and $`\gamma _\pi `$ should be fitted simultaneously Dre00a . RCS above pion threshold can serve as a complement to determine the polarizabilities, in particular the spin polarizabilities, and can provide valuable cross checks between Compton scattering and pion photoproduction, provided one can minimize the model uncertainties in the dispersion formalism. The three types of dispersion integrals of Eqs. (12) and (13) in the formalism outlined here are evaluated as described above. As a representative result obtained within the subtracted DR formalism, the RCS differential cross sections above pion threshold are shown in Fig. 1 at fixed $`\alpha \beta `$ = 10 (here and in the following in units of $`10^4`$ fm<sup>3</sup>), while $`\gamma _\pi `$ is varied between $`27`$ (here and in the following in units of $`10^4`$ fm<sup>4</sup>) and $`37`$ (for more details, see Refs. Dre00a ; Dre00b ). By comparing all available data above pion threshold, it was concluded Dre00a that there is no consistency between the pion photoproduction data from MAMI (entering through the dispersion integrals) and available Compton scattering data, in particular when comparing with the LEGS data, which were used in the extraction of Eq. (8) for $`\gamma _\pi `$. Therefore, new data in the $`\mathrm{\Delta }`$ region are called for, some of which have recently become available Wis (99). An analysis of those unpolarized data in a dispersion formalism favors also a much more negative value for $`\gamma _\pi `$ than extracted in Eq. (8). The fit performed by Ref. Wis (99) yields $`\alpha \beta `$ = 9.1 $`\pm `$ 1.7(stat + syst) $`\pm `$1.2(mod), when using a value of $`\gamma _\pi =37.6`$. Besides the existing information from unpolarized data, a full study of the spin (or vector) polarizabilities will however require double polarization experiments. It was in fact shown Dre00a that the scattering of polarized photons on polarized protons is very sensitive to $`\gamma _\pi `$, in particular in the backward hemisphere and at energies between threshold and the $`\mathrm{\Delta }`$ region. In addition, possible normalization problems can be avoided by measuring appropriate asymmetries. Therefore, future polarization experiments hold the promise to disentangle the scalar and vector polarizabilities of the nucleon with the help of the described subtracted DR formalism, and to further quantify the nucleon spin response in an external electromagnetic field. ### 2.4 Higher order polarizabilities of the proton As outlined above, the electric and magnetic polarizabilities arise as $`O(\nu ^2)`$ corrections to the lowest order (Thomson) scattering amplitude. One can then ask the question whether Compton scattering can also provide additional proton structure information via the use of higher-order polarizabilities. If one extends the analysis to consider spin-averaged $`O(\nu ^4)`$ terms (see Ref. Hol (00) for details), four new structures fulfill the requirements of gauge, P, and T invariance. Two new quantities, $`\alpha _{E\nu }^p`$ and $`\beta _{M\nu }^p`$, represent dispersive corrections to the lowest order static polarizabilities, $`\alpha `$ and $`\beta `$ respectively, and describe the response of the system to time-dependent fields. Two more quantities $`\alpha _{E2}^p`$, $`\beta _{M2}^p`$, represent quadrupole polarizabilities and measure the electric and magnetic quadrupole moments induced in a system by the presence of an applied field gradient. As to the experimental evaluation of such structure probes, it is, of course, in principle possible to extract them directly from Compton cross section measurements. However, it is already clear from the previous discussion of present data, that isolating such pieces from other terms which affect the cross section at energies above $``$ 100 MeV is virtually impossible since additional higher order effects soon become equally important. Thus an alternative procedure is required, which is made possible by the validity of dispersion relations. Within the subtracted DR formalism of Ref. Dre00a outlined above, those higher order terms in the expansion of the Compton amplitudes $`A_i`$ can be reasonably evaluated as in Hol (00). These higher order polarizabilities can be expressed in terms of appropriate derivatives of the RCS invariant amplitudes $`A_i`$ at $`t,\nu ^2=0`$, denoted by $`a_{i,t},a_{i,\nu }`$. In terms of subtracted DR’s, they take the form : $`a_{i,\nu }`$ $`=`$ $`{\displaystyle \frac{2}{\pi }}{\displaystyle _{\nu _{thr}}^{\mathrm{}}}𝑑\nu ^{}{\displaystyle \frac{\mathrm{Im}_sA_i(\nu ^{},t=0)}{\nu ^3}},`$ (14) $`a_{i,t}`$ $`=`$ $`{\displaystyle \frac{1}{\pi }}\left({\displaystyle _{4m_\pi ^2}^{\mathrm{}}}{\displaystyle _{\mathrm{}}^{4Mm_\pi 2m_\pi ^2}}𝑑t^{}{\displaystyle \frac{\mathrm{Im}_tA_i(0,t^{})}{t^2}}\right).`$ (15) The higher order polarizabilities are then obtained as linear combinations of the $`a_{i,t}`$ and $`a_{i,\nu }`$ (for details see Hol (00)). The subtracted DR in Eqs. (14, 15) were evaluated as described above and yield (all in units of $`10^4`$ fm<sup>5</sup>) : $`\alpha _{E\nu }^p`$ $`=`$ $`3.840.19+0.06,`$ $`\beta _{M\nu }^p`$ $`=`$ $`+9.29+0.150.07,`$ $`\alpha _{E2}^p`$ $`=`$ $`+29.310.100.17,`$ $`\beta _{M2}^p`$ $`=`$ $`24.33+0.100.34,`$ (16) where the second and third entries on the rhs of Eq. (2.4) estimate the uncertainties in the $`s`$\- and $`t`$-channel dispersion integrals. The values of Eq. (2.4) were then confronted in Hol (00) with the predictions of HBChPT at $`𝒪(p^3)`$ $`𝒪(p^3):`$ $`\alpha _{E\nu }^p=2.4,\beta _{M\nu }^p=3.7,`$ (17) $`\alpha _{E2}^p=22.1,\beta _{M2}^p=9.5.`$ By comparing Eq. (2.4) and (17) one finds that the size of $`\alpha _{E2}^p`$ is about right, while for both $`\beta _{M2}^p`$ and $`\beta _{M\nu }^p`$ the sign and order of magnitude is correct but additional contributions are called for. The most serious problem lies in the experimental determination of $`\alpha _{E\nu }`$ which is negative in contradistinction to the chiral prediction and to sum rule arguments which assert its positivity. To see if inclusion of $`\mathrm{\Delta }(1232)`$ degrees of freedom can help to resolve the above discrepancies, these quantities were also calculated in Hol (00) in HBChPT to $`O(ϵ^3)`$, with the result : $`𝒪(ϵ^3):`$ $`\alpha _{E\nu }^p=1.7,\beta _{M\nu }^p=7.5,`$ (18) $`\alpha _{E2}^p=26.2,\beta _{M2}^p=12.3.`$ Except for the sign problem with $`\alpha _{E\nu }^p`$ indicated above, which persists in the $`ϵ`$-expansion, the changes are generally helpful, although the magnetic quadrupole polarizability is still somewhat underpredicted. In Hol (00), the described analysis was also extended to higher order contributions $`O(\nu ^5)`$ to the proton spin polarizabilities, for which 8 new structures were found. A dispersive evaluation of those higher order spin polarizabilities showed a qualitative agreement with HBChPT $`O(ϵ^3)`$ predictions. Recently, an evaluation of the higher order polarizabilities of the proton in HBChPT to $`O(p^4)`$ has been reported Hem00b , providing an important new testing ground for the chiral predictions. It was found Hem00b that the $`O(p^4)`$ HBChPT result for the 4 quadrupole polarizabilities and the 8 spin polarizabilities at $`O(\nu ^5)`$ of the proton are in encouraging good agreement with the DR estimates of Ref. Hol (00). In summary, the subtracted DR formalism presented not only provides a formalism to extract the lowest order nucleon polarizabilities from present and forthcoming RCS data with a minimum of model dependence. It can also be used to obtain information about higher order polarizabilities of the proton, in this way providing a great deal of additional nucleon structure information. ## 3 Virtual Compton scattering (VCS) and generalized nucleon polarizabilities ### 3.1 Introduction The nucleon structure information obtained through RCS, as discussed in section 2, can be generalized by virtual Compton scattering (VCS) below pion threshold. VCS can be interpreted as electron scattering off a target polarized by the presence of constant electric and magnetic fields. To see how VCS generalizes the RCS process, it is useful to think of the analogy with the electromagnetic form factors. Their measurement through elastic electron-nucleon scattering reveals the spatial distribution of the charge and magnetization distributions of the target, whereas a real photon is only sensitive to the overall charge and magnetization of the target. The physics addressed with VCS is then the same as if one were performing an elastic electron scattering experiment on a target placed between the plates of a capacitor or between the poles of a magnet. In this way one studies the spatial distributions of the polarization densities of the target, by means of the generalized polarizabilities, which are functions of the square of the four-momentum, $`Q^2`$, transferred by the electron. Experimentally, the VCS process is accessed through the electroproduction of photons, and we consider in all of the following the reaction on a proton target, i.e. the reaction $`epep\gamma `$. One immediately sees a difference with regard to the RCS $`\gamma p\gamma p`$ reaction, because in the $`epep\gamma `$ reaction, the final photon can be emitted either by the proton, giving access to the VCS process, or by the electrons, which is referred to as the Bethe-Heitler (BH) process. The BH amplitude can be calculated exactly in QED, provided one knows the elastic form factors of the proton. Therefore it contains no new information on the structure. Unfortunately, light particles such as electrons radiate much more than the heavy proton. Therefore the BH process generally dominates or at least interferes strongly with the VCS process, and this may complicate the interpretation of the $`epep\gamma `$ reaction. The only way out of this problem is either to find kinematical regions where the BH process is suppressed or to have a very good theoretical control over the interference between the BH and the VCS amplitudes, as will be discussed below. Assuming that this problem is fixed, one can then proceed to extract the nucleon structure information from VCS. In doing so, care has to be taken to separate the trivial response of the target, due to its global charge and/or a global magnetic moment. Indeed, if we put a proton in an electric field, the first effect we observe is that it moves as a whole. Similarly, the magnetic field produces a precession of the magnetic moment. This problem is absent when one studies the polarizability of a macroscopic sample because it can be fixed in space by appropriate means, which is not possible for the proton. This absence of a restoring force explains why the trivial response due to the motion of charge and magnetic moment dominates over the response of the internal degrees of freedom. This is the physical origin of the low energy theorem (LET) Low (58) for VCS. All what is needed to calculate this part of the response, are the parameters which control the motion, that is the mass, the charge, and the magnetic moment. Once the motion is known, one can compute the amplitude for scattering an electron on this moving proton, the so-called Born amplitude. Having separated the trivial response, one can parametrize the structure part of interest in the VCS process through the so-called generalized polarizabilities (GP’s) as in Ref. Gui (95). ### 3.2 Definition of generalized polarizabilities The known BH + Born parts of the VCS amplitude at low energy start at order 1/$`\text{q}^{}`$ in an expansion in the outgoing photon energy $`\text{q}^{}`$. The LET Low (58) asserts that the non trivial part of the VCS amplitude, the so-called non-Born part (denoted by $`H_{NB}`$), begins at order $`\text{q}^{}`$. There is of course also a contribution of order $`\text{q}^{}`$ in the BH + Born amplitude, but this term is exactly known and therefore can be subtracted, at least in principle. So what is needed next is an adequate parametrization of $`H_{NB}`$. For this purpose, a multipole expansion (in the c.m. frame) was performed in Gui (95) in order to take advantage of angular momentum and parity conservation. The behaviour of the non-Born VCS amplitude $`H_{NB}`$ at low energy ($`\text{q}^{}0`$) but at arbitrary three-momentum q of the virtual photon, was then parametrized by 10 functions of q defined by : $$\text{Limit\hspace{0.17em} of}\frac{1}{\text{q}^{}}\frac{1}{\text{q}^L}H_{NB}^{(\rho ^{}1,\rho L)S}(\text{q}^{},\text{q})\text{when}\text{q}^{}0.$$ (19) In this notation, $`\rho `$ ($`\rho ^{}`$) refers to the electric (2), magnetic (1) or longitudinal (0) nature of the initial (final) photon, $`L`$ ($`L^{}=1`$) represents the angular momentum of the initial (final) photon, whereas $`S`$ differentiates between the spin-flip ($`S=1`$) and non spin-flip ($`S=0`$) character of the electromagnetic transition at the nucleon side. As the angular momentum of the outgoing photon is $`L^{}`$ = 1, this leads to 10 q-dependent GP’s, denoted generically by $`P^{(\rho ^{}L^{},\rho L)S}(\text{q})`$. By imposing the constraints due to nucleon crossing combined with charge conjugation invariance on the VCS amplitude, it was shown however in Dre (97); Dre98b that 4 of the GP’s can be eliminated. Thus only 6 GP’s , e.g. Gui (98) $`P^{(01,01)0}(\text{q}),P^{(11,11)0}(\text{q}),`$ $`P^{(01,01)1}(\text{q}),P^{(11,11)1}(\text{q}),P^{(11,02)1}(\text{q}),P^{(01,12)1}(\text{q}),`$ (20) are necessary to give the low energy behaviour of $`H_{NB}`$. In the limit $`\text{q}0`$ for the GP’s, one finds the following relations with the polarizabilities (in gaussian units) of RCS, as discussed in section 2 Dre98c : $`P^{(01,01)0}(0)={\displaystyle \frac{1}{\alpha _{em}}}\sqrt{{\displaystyle \frac{2}{3}}}\alpha ,`$ $`P^{(11,11)0}(0)={\displaystyle \frac{1}{\alpha _{em}}}\sqrt{{\displaystyle \frac{8}{3}}}\beta ,`$ $`P^{(01,12)1}(0)={\displaystyle \frac{1}{\alpha _{em}}}{\displaystyle \frac{\sqrt{2}}{3}}\gamma _3,`$ $`P^{(11,02)1}(0)={\displaystyle \frac{1}{\alpha _{em}}}{\displaystyle \frac{2\sqrt{2}}{3\sqrt{3}}}\left(\gamma _2+\gamma _4\right),`$ $`P^{(01,01)1}(0)=0,`$ $`P^{(11,11)1}(0)=0,`$ (21) where $`\alpha _{em}=1/137.036`$ is the QED fine structure constant. ### 3.3 VCS observables We next discuss how one can analyze $`epep\gamma `$ observables to extract the 6 GP’s of Eq. (20). The VCS unpolarized squared amplitude is denoted by $``$. Besides, one can consider VCS double polarization observables, which are denoted by $`\mathrm{\Delta }(h,i)`$ for an electron of helicity $`h`$, and which are defined as a difference of the squared amplitudes for recoil (or target) proton spin orientation in the direction and opposite to the axis $`i`$ ($`i=x,y,z`$) (see Ref. Vdh97a for details). In an expansion in $`\text{q}^{}`$, $``$ and $`\mathrm{\Delta }`$ take the form $`^{\mathrm{exp}}`$ $`=`$ $`{\displaystyle \frac{_2^{\mathrm{exp}}}{\text{q}_{}^{}{}_{}{}^{2}}}+{\displaystyle \frac{_1^{\mathrm{exp}}}{\text{q}^{}}}+_0^{\mathrm{exp}}+O(\text{q}^{}),`$ $`\mathrm{\Delta }^{\mathrm{exp}}`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Delta }_2^{\mathrm{exp}}}{\text{q}_{}^{}{}_{}{}^{2}}}+{\displaystyle \frac{\mathrm{\Delta }_1^{\mathrm{exp}}}{\text{q}^{}}}+\mathrm{\Delta }_0^{\mathrm{exp}}+O(\text{q}^{}).`$ (22) Due to the LET, the threshold coefficients $`_2`$, $`_1`$, $`\mathrm{\Delta }_2`$, $`\mathrm{\Delta }_1`$ are known. The information on the GP’s is contained in $`_0^{\mathrm{exp}}`$ and $`\mathrm{\Delta }_0^{\mathrm{exp}}`$. These coefficients contain a part which comes from the (BH+Born) amplitude and another one which is a linear combination of the GP’s with coefficients determined by the kinematics. The unpolarized observable $`_0^{\mathrm{exp}}`$ was obtained by Ref. Gui (95) in terms of 3 structure functions $`P_{LL}(\text{q})`$, $`P_{TT}(\text{q})`$, $`P_{LT}(\text{q})`$, which are linear combinations of the 6 GP’s, $`_0^{\mathrm{exp}}_0^{\mathrm{BH}+\mathrm{Born}}`$ $`=2K_2\{v_1[ϵP_{LL}(\text{q})P_{TT}(\text{q})]`$ (23) $`+(v_2{\displaystyle \frac{\stackrel{~}{q}_0}{\text{q}}}v_3)\sqrt{2\epsilon \left(1+\epsilon \right)}P_{LT}(\text{q})\},`$ where $`K_2,ϵ,\stackrel{~}{q}_0,v_1,v_2,v_3`$ are kinematical quantities (for details see Ref. Gui (98)). The three double-polarization observables $`\mathrm{\Delta }_0^{\mathrm{exp}}(h,i)`$ ($`i=x,y,z`$) were expressed by Vdh97a in terms of three new independent structure functions $`P_{LT}^z(\text{q})`$, $`P_{LT}^{}_{}{}^{}z(\text{q})`$, and $`P_{LT}^{{}_{}{}^{}}(\text{q})`$, which are also linear combinations of the 6 GP’s, $`\mathrm{\Delta }_0^{\mathrm{exp}}(h,z)\mathrm{\Delta }_0^{\mathrm{BH}+\mathrm{Born}}(h,z)`$ $`=4(2h)K_2\{v_1\sqrt{1\epsilon ^2}P_{TT}(\text{q})+v_2\sqrt{2\epsilon \left(1\epsilon \right)}P_{LT}^z(\text{q})`$ $`+v_3\sqrt{2\epsilon \left(1\epsilon \right)}P_{LT}^{}_{}{}^{}z(\text{q})\},`$ $`\mathrm{\Delta }_0^{\mathrm{exp}}(h,x)\mathrm{\Delta }_0^{\mathrm{BH}+\mathrm{Born}}(h,x)`$ $`=4(2h)K_2\{v_1^x\sqrt{2\epsilon \left(1\epsilon \right)}P_{LT}^{}(\text{q})+v_2^x\sqrt{1\epsilon ^2}P_{TT}^{}(\text{q})`$ $`+v_3^x\sqrt{1\epsilon ^2}P_{TT}^{{}_{}{}^{}}(\text{q})+v_4^x\sqrt{2\epsilon \left(1\epsilon \right)}P_{LT}^{{}_{}{}^{}}(\text{q})\},`$ $`\mathrm{\Delta }_0^{\mathrm{exp}}(h,y)\mathrm{\Delta }_0^{\mathrm{BH}+\mathrm{Born}}(h,y)`$ $`=4(2h)K_2\{v_1^y\sqrt{2\epsilon \left(1\epsilon \right)}P_{LT}^{}(\text{q})+v_2^y\sqrt{1\epsilon ^2}P_{TT}^{}(\text{q})`$ $`+v_3^y\sqrt{1\epsilon ^2}P_{TT}^{{}_{}{}^{}}(\text{q})+v_4^y\sqrt{2\epsilon \left(1\epsilon \right)}P_{LT}^{{}_{}{}^{}}(\text{q})\},`$ (24) where $`v_1^x,\mathrm{},v_4^x,v_1^y,\mathrm{},v_4^y`$ are kinematical coefficients. The other structure functions in Eq. (24) can be expressed in terms of $`P_{LL},P_{TT},P_{LT},P_{LT}^z,P_{LT}^{}_{}{}^{}z,P_{LT}^{{}_{}{}^{}}`$ Gui (98). Therefore, measuring those 6 structure functions amounts to determine the 6 independent GP’s. ### 3.4 Results for VCS observables below pion threshold In the previous sections, the observables of the $`epep\gamma `$ reaction below pion threshold were outlined, and it was shown how the nucleon structure effect can be parametrized in terms of 6 independent GP’s. To access the GP’s, the experimental strategy of VCS in the threshold region consists of two steps. First, one measures the VCS cross section at several values of the outgoing photon energy. At low energies, the measurement of the VCS observables provides a test of the LET. Once the LET is verified, the relative effect of the GP’s can be extracted using Eqs. (23,24). The predictions for the Bethe-Heitler (BH) and Born cross sections below pion threshold are shown in Fig. 2. The BH cross section has a characteristic angular shape and displays two “spikes”, which occur when the direction of the outgoing photon coincides with either the initial or final electron directions. In these regions, the cross section is completely dominated by the BH contributions. In order to determine the VCS contribution, one clearly has to minimize the BH contamination by detecting the photon in the half-plane opposite to the electron directions. The first dedicated VCS experiment has been performed at MAMI, and for the first time two combinations (see Eq. (23)) of GP’s have been determined at $`Q^2`$ = 0.33 GeV<sup>2</sup> and photon polarization $`ϵ`$ = 0.62 Roc (00), $`P_{LL}(Q^2){\displaystyle \frac{1}{ϵ}}P_{TT}(Q^2)=\left(23.7\pm 2.2\pm 0.6\pm 4.3\right)\mathrm{GeV}^2,`$ $`P_{LT}(Q^2)=\left(5.0\pm 0.8\pm 1.1\pm 1.4\right)\mathrm{GeV}^2.`$ (25) VCS experiments at higher $`Q^2`$ (1 - 2 GeV<sup>2</sup>) at JLab Bert (93), and at lower $`Q^2`$ at MIT-Bates Sha (97) have already been performed, and are under analysis at this time. The GP’s have been calculated in various approaches and nucleon structure models, ranging from constituent quark models Gui (95); Pas (98), a relativistic effective Lagrangian model Vdh (96), and the linear $`\sigma `$-model Met (96) to ChPT Hem97b ; Hem00a . The GP’s teach us about the interplay between nucleon-core excitations and pion-cloud effects, which are described differently in the various models. We focus here on the calculation of the GP’s in HBChPT to $`𝒪(p^3)`$, as it takes account of N$`\pi `$-loop contributions in a systematic way. The $`𝒪(p^3)`$ calculation yields for the two measured combinations at $`Q^2`$ = 0.33 GeV<sup>2</sup> and $`ϵ=0.62`$ of Eq. (25) the values Hem00a : $`𝒪(p^3):P_{LL}P_{TT}/ϵ`$ $`=`$ $`\mathrm{\hspace{0.33em}26.3}\mathrm{GeV}^2,`$ $`P_{LT}`$ $`=`$ $`5.7\mathrm{GeV}^2,`$ (26) which are in astonishing agreement with the experimentally determined values of Eq. (25). In particular, the $`𝒪(p^3)`$ ChPT calculation predicts quite large values for the spin GP’s. As for the case of the RCS polarizabilities, the importance of the $`𝒪(p^4)`$ corrections remains to be checked. If one wants to extract the different polarizabilities from experiment, and in particular in case of the spin polarizabilities, an unpolarized experiment is not sufficient as it gives access to 3 independent response functions only. To further separate the polarizabilities, one has to resort to double-polarization observables. Experimentally, at the existing high-duty-cycle electron facilities with polarized electron beams such as MAMI, MIT-Bates and JLab, double polarization VCS experiments can be performed by measuring the recoil polarization of the outgoing nucleon with a focal plane polarimeter. An experiment at MAMI has already been proposed d’Ho99b . In Fig. 3, the double-polarization asymmetry of Eq. (24), with proton recoil polarization either along the $`x`$\- or $`z`$\- directions, is shown at $`\mathrm{q}`$ = 600 MeV/c for in-plane kinematics. It is seen that the asymmetry yields a very large value (between 0.6 and 0.7) if the final proton is polarized parallel to the virtual photon. ### 3.5 Dispersion relation formalism for VCS At present, VCS experiments at low outgoing photon energies are analyzed in terms of low-energy expansions (LEXs) of Eq. (22). In the LEX, the non-Born response of the system to the quasi-constant electromagnetic field of the low energetic photon is proportional to the GP’s, as expressed in Eqs. (23,24). As the sensitivity of the VCS cross sections to the GP’s grows with the photon energy, it is advantageous to go to higher photon energies, provided one can keep the theoretical uncertainties under control when crossing the pion threshold. The situation can be compared to RCS, for which one uses a dispersion relation formalism as discussed in section 2, to extract the polarizabilities at energies above pion threshold, with generally larger effects on the observables. In order to go to higher energies and to check the validity of LEXs at these higher energies, a dispersion relation analysis for VCS has been developed very recently Pas (00); Vdh00b , which will allow to extract the GP’s from data over a larger energy range. The same formalism also provides for the first time a dispersive evaluation of 4 GP’s. To set up a dispersion formalism for the VCS process, one starts from the helicity amplitudes : $$T_{\lambda ^{}s^{};\lambda s}=e^2\epsilon _\mu (q,\lambda )\epsilon _\nu ^{{}_{}{}^{}}(q^{},\lambda ^{})\overline{u}(p^{},s^{})^{\mu \nu }u(p,s),$$ (27) with $`e`$ the electric charge, $`q`$ ($`q^{}`$) the four-vectors of the virtual (real) photon in the VCS process, and $`p`$ ($`p^{}`$) the four-momenta of the initial (final) nucleons respectively. The nucleon helicities are denoted by $`s,s^{}=\pm 1/2`$, and $`u,\overline{u}`$ are the nucleon spinors. The initial virtual photon has helicity $`\lambda =0,\pm 1`$ and polarization vector $`\epsilon _\mu `$, whereas the final real photon has helicity $`\lambda ^{}=\pm 1`$ and polarization vector $`\epsilon _\nu ^{^{}}`$. The VCS process is characterized by 12 independent helicity amplitudes $`T_{\lambda ^{}s^{};\lambda s}`$. The VCS tensor $`^{\mu \nu }`$ in Eq. (27) is then expanded into a basis of 12 independent gauge invariant tensors $`\rho _i^{\mu \nu }`$, $$^{\mu \nu }=\underset{i=1}{\overset{12}{}}F_i(Q^2,\nu ,t)\rho _i^{\mu \nu },$$ (28) as introduced in Ref. Dre (97) (starting from the amplitudes of Ref. Tar (75)). The amplitudes $`F_i`$ in Eq. (28) contain all nucleon structure information and are functions of 3 invariants for the VCS process : $`Q^2q^2`$, $`\nu =(su)/(4m_N)`$ which is odd under $`su`$ crossing, and $`t`$. The Mandelstam invariants $`s`$, $`t`$ and $`u`$ for VCS are defined by $`s=(q+p)^2`$, $`t=(qq^{})^2`$, and $`u=(qp^{})^2`$, with the constraint $`s+t+u=2m_N^2Q^2`$, and $`m_N`$ is the nucleon mass. Nucleon crossing combined with charge conjugation provides the following constraints on the amplitudes $`F_i`$ <sup>1</sup><sup>1</sup>1In Pas (00), 4 of the 12 amplitudes of Dre (97) were redefined by dividing them through $`\nu `$, such that all of them are even functions of $`\nu `$. This simplifies the formalism since only one type of dispersion integrals needs to be considered then. at arbitrary virtuality $`Q^2`$ : $$F_i(Q^2,\nu ,t)=F_i(Q^2,\nu ,t)(i=1,\mathrm{},12).$$ (29) With the choice of the tensor basis of Ref. Dre (97), the resulting non-Born amplitudes $`F_i^{NB}`$ ($`i`$ = 1,…,12) are free of all kinematical singularities and constraints. Assuming further analyticity and an appropriate high-energy behavior, the non-Born amplitudes $`F_i^{NB}(Q^2,\nu ,t)`$ fulfill unsubtracted dispersion relations (DR’s) with respect to the variable $`\nu `$ at fixed $`t`$ and fixed virtuality $`Q^2`$ : $$\mathrm{R}eF_i^{NB}(Q^2,\nu ,t)=\frac{2}{\pi }𝒫_{\nu _{thr}}^+\mathrm{}𝑑\nu ^{}\frac{\nu ^{}\mathrm{I}m_sF_i(Q^2,\nu ^{},t)}{\nu ^2\nu ^2},$$ (30) with $`\mathrm{I}m_sF_i`$ the discontinuities across the $`s`$-channel cuts of the VCS process. Since pion production is the first inelastic channel, $`\nu _{thr}=m_\pi +(m_\pi ^2+t/2+Q^2/2)/(2m_N)`$, where $`m_\pi `$ denotes the pion mass. The unsubtracted DR’s of Eq. (30) require that at sufficiently high energies ($`\nu \mathrm{}`$ at fixed $`t`$ and fixed $`Q^2`$) the amplitudes $`\mathrm{I}m_sF_i(Q^2,\nu ,t)`$ ($`i`$ = 1,…,12) drop fast enough such that the integrals are convergent and the contributions from the semi-circle at infinity can be neglected. It turns out that for two amplitudes, $`F_1`$ and $`F_5`$, an unsubtracted dispersion integral as in Eq. (30) does not exist Pas (00), whereas the other 10 amplitudes can be evaluated through unsubtracted dispersion integrals. This situation is similar as for RCS, where 2 of the 6 invariant amplitudes cannot be evaluated by unsubtracted dispersion relations either Lvo (97). The unsubtracted DR formalism for VCS also allows to predict 4 of the 6 GP’s. The appropriate limit in the definition of the GP’s is $`\text{q}^{}0`$ at finite q (see Eq. (19)), which corresponds in terms of VCS invariants to $`\nu 0`$ and $`tQ^2`$ at finite $`Q^2`$. One can therefore express the GP’s in terms of the VCS amplitudes $`F_i`$ at the point $`\nu =0`$, $`t=Q^2`$ at finite $`Q^2`$, denoted in the following as : $`\overline{F}_i(Q^2)F_i^{NB}(Q^2,\nu =0,t=Q^2)`$. The relations between the GP’s and the $`\overline{F}_i(Q^2)`$ can be found in Ref. Dre (97). From the high-energy behavior for the VCS invariant amplitudes, it follows that one can evaluate the $`\overline{F}_i`$ (for $`i`$ 1, 5) through the unsubtracted DR’s $$\overline{F}_i(Q^2)=\frac{2}{\pi }_{\nu _{thr}}^+\mathrm{}𝑑\nu ^{}\frac{\mathrm{I}m_sF_i(Q^2,\nu ^{},t=Q^2)}{\nu ^{}}.$$ (31) Unsubtracted DR’s for the GP’s will therefore hold for those combinations of GP’s that do not depend upon the amplitudes $`\overline{F}_1`$ and $`\overline{F}_5`$ <sup>2</sup><sup>2</sup>2$`\overline{F}_5`$ can appear however in the combination $`\overline{F}_5+4\overline{F}_{11}`$, in which the $`\pi ^0`$-pole drops out, and which has a high-energy behavior leading to a convergent integral (see Ref. Pas (00)).. Among the 6 GP’s, the following 4 combinations of GP’s were found in Ref. Pas (00) : $`P^{(01,01)0}+{\displaystyle \frac{1}{2}}P^{(11,11)0}={\displaystyle \frac{2}{\sqrt{3}}}\left({\displaystyle \frac{E+m_N}{E}}\right)^{1/2}m_N\stackrel{~}{q}_0`$ $`\times \left\{{\displaystyle \frac{\text{q}^2}{\stackrel{~}{q}_0^2}}\overline{F}_2+\left(2\overline{F}_6+\overline{F}_9\right)\overline{F}_{12}\right\},`$ (32) $`P^{(01,01)1}={\displaystyle \frac{1}{3\sqrt{2}}}\left({\displaystyle \frac{E+m_N}{E}}\right)^{1/2}\stackrel{~}{q}_0`$ $`\times \left\{\left(\overline{F}_5+\overline{F}_7+4\overline{F}_{11}\right)+4m_N\overline{F}_{12}\right\},`$ (33) $`P^{(01,12)1}{\displaystyle \frac{1}{\sqrt{2}\stackrel{~}{q}_0}}P^{(11,11)1}={\displaystyle \frac{1}{3}}\left({\displaystyle \frac{E+m_N}{E}}\right)^{1/2}{\displaystyle \frac{m_N\stackrel{~}{q}_0}{\text{q}^2}}`$ $`\times \left\{\left(\overline{F}_5+\overline{F}_7+4\overline{F}_{11}\right)+4m_N\left(2\overline{F}_6+\overline{F}_9\right)\right\},`$ (34) $`P^{(01,12)1}+{\displaystyle \frac{\sqrt{3}}{2}}P^{(11,02)1}={\displaystyle \frac{1}{6}}\left({\displaystyle \frac{E+m_N}{E}}\right)^{1/2}{\displaystyle \frac{\stackrel{~}{q}_0}{\text{q}^2}}`$ $`\times \left\{\stackrel{~}{q}_0\left(\overline{F}_5+\overline{F}_7+4\overline{F}_{11}\right)+8m_N^2\left(2\overline{F}_6+\overline{F}_9\right)\right\},`$ (35) where $`E=\sqrt{\text{q}^2+m_N^2}`$ denotes the initial proton c.m. energy, and $`\stackrel{~}{q}_0=m_NE`$ the virtual photon c.m. energy in the limit $`\text{q}^{}`$ = 0. Unfortunately, the 4 combinations of GP’s of Eqs. (32)-(35) can at present not yet be compared with the data. In particular, the only unpolarized experiment Roc (00) measured two structure functions which cannot be evaluated in an unsubtracted DR formalism, as they contain in addition to $`P^{(01,01)0}+1/2P^{(11,11)0}`$ of Eq. (32), which is proportional to $`\alpha +\beta `$ at $`Q^2`$ = 0, also the generalization of $`\alpha \beta `$. The 4 combinations of GP’s on the lhs of Eqs. (32)-(35) can then be evaluated by unsubtracted DR’s, from the dispersion integrals of Eq. (31) for the $`\overline{F}_i(Q^2)`$. To this end, the imaginary parts $`\mathrm{I}m_sF_i`$ in Eq. (31) have to be calculated by use of unitarity. For the VCS helicity amplitudes of Eq. (27) (denoted for short by $`T_{fi}`$), the unitarity equation reads : $$2\mathrm{Im}_sT_{fi}=\underset{X}{}(2\pi )^4\delta ^4(P_XP_i)T_{Xf}^{}T_{Xi},$$ (36) where the sum runs over all possible intermediate states $`X`$ that can be formed. In Ref. Pas (00), the dispersion integrals of Eq. (31) were saturated by the dominant contribution of the $`\pi N`$ intermediate states. For the pion photo- and electroproduction helicity amplitudes in the range $`Q^2`$ 0.5 GeV<sup>2</sup>, the phenomenological analysis of MAID Dre (99) was used, which contains both resonant and non-resonant pion production mechanisms. In Fig. 4, the results for the 4 combinations of GP’s of Eqs. (32)-(35) are shown in the DR formalism, and compared to the results of the $`O(p^3)`$ heavy-baryon chiral perturbation theory (HBChPT) Hem97b ; Hem00a and the linear $`\sigma `$-model Met (96). The $`\pi N`$ contribution to the sum $`P^{(01,01)0}`$ \+ $`1/2P^{(11,11)0}`$ gives only about 80% of the Baldin sum rule Bab (98), because of a non-negligible high-energy contribution (of heavier intermediate states) to the photoabsorption cross section entering the sum rule, which is not estimated here. On the other hand, for the 3 combinations of spin polarizabilities of Eqs. (33)-(35), the dispersive estimates with $`\pi N`$ states are expected to provide a rather reliable guidance. By comparing the DR results with those of HBChPT at $`O(p^3)`$, one remarks a rather good agreement for $`P^{(01,12)1}`$ \+ $`\sqrt{3}/2P^{(11,02)1}`$, whereas for the GP’s $`P^{(01,01)1}`$ and $`P^{(01,12)1}`$ \- $`1/(\sqrt{2}\stackrel{~}{q}_0)P^{(11,11)1}`$, the dispersive results drop much faster with $`Q^2`$. This trend is also seen in the relativistic linear $`\sigma `$-model, which takes account of some higher orders in the chiral expansion. It remains to be checked how the $`𝒪(p^4)`$ corrections in HBChPT change this comparison with the DR estimates. To complete the DR formalism for VCS, one further needs to construct the VCS amplitudes $`F_1`$ and $`F_5`$, for which the unsubtracted dispersion integrals of Eq. (30) do not converge. One strategy is to proceed in an analogous way as has been proposed in Ref. Lvo (97) in the case of RCS. The unsubtracted dispersion integrals for $`F_1`$ and $`F_5`$ are evaluated along the real $`\nu `$-axis in a finite range $`\nu _{max}\nu +\nu _{max}`$ (with $`\nu _{max}`$ 1.5 GeV). The integral along a semi-circle of finite radius $`\nu _{max}`$ in the complex $`\nu `$-plane is described by the asymptotic contribution $`F_i^{as}`$, which is parametrized by $`t`$-channel poles (e.g. for $`Q^2`$ = 0, $`F_1^{as}`$ corresponds to $`\sigma `$-exchange, and $`F_5^{as}`$ to $`\pi ^0`$-exchange). A full study of VCS observables within such a dispersion formalism, including a parametrization of the two asymptotic contributions, is presently underway Dre00c . This will yield a formalism to extract the nucleon GP’s over a larger range of energies from both unpolarized and polarized VCS data. ## 4 Compton scattering at large momentum transfer and the nucleon distribution amplitude ### 4.1 Introduction Besides the low energy region, as discussed in section 2, RCS will also provide access to information on the partonic structure of the nucleon at sufficiently large momentum transfer. This regime is defined by requiring that all three Mandelstam variables ($`s,t,u`$) be large with respect to a typical hadronic scale, say 1 GeV. In this case there is a prejudice (actually proven in the case of elastic electron scattering Ste (97)) that the amplitude factorizes in a soft non-perturbative part, the distribution amplitude, and a hard scattering kernel which is calculable from perturbative QCD (PQCD). Because of asymptotic freedom, the perturbative approach must be to some degree relevant to the hard scattering regime. However, since the binding of the quarks and gluons in the hadrons is a long distance, non-perturbative effect, the description of the reaction requires a consistent analysis of both large and small scales. When the reaction is hard enough, the relative velocities of the participating particles are nearly lightlike. Time dilatation increases the lifetime of the quantum configurations which build the hadron. As a result, the partonic content, as seen by the other particles, is frozen. Moreover, due to the apparent contraction of the hadron size, the time during which momentum can be exchanged is decreased. Therefore one expects a lack of coherence between the long-distance confining effects and the short distance reaction. This incoherence between the soft and hard physics is the origin of the factorization which is illustrated in Fig. 5. ### 4.2 Factorization and the nucleon distribution amplitude The calculation of the RCS amplitude at large momentum transfer, follows the Brodsky-Lepage formalism Bro (80), which leads to the factorized expression : $`T(\lambda ^{},h_N^{},\lambda ,h_N)=`$ $`{\displaystyle 𝑑x_i𝑑y_j\varphi _N^{}(y_j)T_H(\lambda ^{},h_N^{},y_j,\lambda ,h_N,x_i;s,t)\varphi _N(x_i)},`$ where ($`x_i,y_i`$) are the momentum fractions of the quarks in the initial and final nucleon respectively, $`T_H`$ is the hard scattering kernel and $`\varphi _N`$ is the distribution amplitude (DA). The evaluation of Eq. (4.2) requires a four-fold convolution integral since there are two constraint equations $`(x_1+x_2+x_3=1`$ and $`y_1+y_2+y_3=1)`$. In Eq. (4.2) a sufficiently large momentum transfer is assumed in order to neglect the transverse momentum dependence of the partons in the hard scattering amplitude $`T_H`$. In this limit, the integration over the transverse momenta $`\stackrel{}{k}_i`$ (where $`_i\stackrel{}{k}_i=\mathrm{𝟎}`$) acts only on the valence wavefunction $$\mathrm{\Psi }_V(x_1,x_2,x_3;\stackrel{}{k}_1,\stackrel{}{k}_2,\stackrel{}{k}_3),$$ (38) which is the amplitude of the three quark state in the Fock expansion of the proton: $`P>`$ $`=`$ $`\mathrm{\Psi }_Vqqq>+\mathrm{\Psi }_{q\overline{q}}qqq,q\overline{q}>+\mathrm{\Psi }_gqqq,g>`$ (39) $`+\mathrm{}`$ This valence wave function $`\mathrm{\Psi }_V`$ integrated up to a scale $`\mu `$ (which separates the soft and hard parts of the wavefunction) defines the DA which appears in Eq. (4.2) : $$\varphi _N(x_i,\mu )=^\mu d^2\stackrel{}{k}_i\mathrm{\Psi }_V(x_i;\stackrel{}{k}_i).$$ (40) For $`\mu `$ much larger than the average value of the transverse momentum in the proton, this function $`\varphi _N`$ depends only weakly on $`\mu `$ Bro (80) and this dependence can be neglected. The interest of the formalism is that the distribution amplitude is universal, that is independent of the particular reaction considered. Several distribution amplitudes have been modeled using QCD sum rules Che (84, 89); Kin (87). They have a characteristic shape and predict that in a proton, the $`u`$-quark with helicity along the proton helicity carries about 2/3 of its longitudinal momentum (see Fig. 6 ). For the computation of the hard scattering amplitude $`T_H`$ (black circle in Fig. 5), the leading order PQCD contribution corresponds to the exchange of the minimum number of gluons (two in the present case) between the three quarks. The number of diagrams grows rapidly with the number of elementary particles involved in the reaction (42 diagrams for the nucleon form factor, 336 diagrams in the case of real or virtual Compton scattering). Despite the large number of diagrams, the calculation of $`T_H`$ is a parameter free calculation once the scale $`\mathrm{\Lambda }_{QCD}`$ 200 MeV in the strong coupling $`\alpha _s(Q^2)`$ is given. Note that configurations with more than three valence quarks are a priori allowed but since this implies the exchange of more hard gluons, the corresponding contribution is suppressed by powers of $`1/t`$. There are two characteristic features of the Brodsky-Lepage model which are almost direct consequences of QCD: the dimensional counting rules Bro (73) and the conservation of hadronic helicities Bro (81). The latter feature implies that any helicity flip amplitude is zero and, hence, any single spin asymmetry too. The helicity sum rule is a consequence of utilizing the collinear approximation and of dealing with (almost) massless quarks which conserve their helicities when interacting with gluons. Whereas the dimensional counting rules are in reasonable agreement with experiment, the helicity sum rule seems to be violated even at moderately large momentum transfers. The prevailing opinion is that these phenomena cannot be explained in terms of perturbative QCD (see, for example, Ref. Siv (89)), but rather are generated by an interplay of perturbative and non-perturbative physics. An interesting aspect of real and virtual Compton scattering is that these are the simplest processes in which the integrals over the longitudinal momentum fractions yield imaginary parts. The reason is that, as in any scattering process, there are kinematical regions where internal quarks and gluons can go on their mass shell. The appearance of imaginary parts to leading order in $`\alpha _s`$ is a non-trivial prediction of PQCD, which should be tested experimentally. As discussed in Gui (98), the ($`e,e^{}\gamma `$) reaction with polarized incoming electrons seems to be a good candidate for this investigation. In contrast to the PQCD (or hard scattering) approach to RCS, it was argued in Refs. Rad98a ; Die (99) that wide angle Compton scattering at accessible energies is described by a competing mechanism, in which the large momentum transfer is absorbed on a single quark and shared by the overlap of high-momentum components in the soft wave function. This so-called soft-overlap mechanism gives a purely real amplitude, therefore displaying a different signature than the PQCD amplitude. The transition from such a soft-overlap mechanism to the perturbative, hard scattering approach when increasing the momentum transfer is an open question for a reaction such as wide angle Compton scattering. It is hoped that future experiments can shed light on this transition. ### 4.3 Results for RCS in PQCD The leading order PQCD prediction for RCS at large momentum transfer has been calculated several times in the literature Far (90, 91); Kro (91); Vdh97b ; Broo (00). The first step in such a calculation consists of evaluating the 336 diagrams entering the hard scattering amplitude $`T_H`$ for RCS. Next, the four-fold convolution integral of Eq. (4.2) has to be performed to obtain the Compton helicity amplitudes. The numerical integration requires some care because the quark and/or gluon propagators can go on-shell which leads to (integrable) singularities. The different numerical implementations of these singularities are probably the reason of the different results obtained in the literature. In Refs. Far (90, 91), the propagator singularities were integrated by taking a finite value for the imaginary part +i$`ϵ`$ of the propagator. The behavior of the result was then studied by decreasing the value of $`ϵ`$. To obtain convergence with a practical number of samples in the Monte Carlo integration performed in Far (90, 91), the smallest feasible value for $`ϵ`$ was $`ϵ0.005`$. In Ref. Kro (91), the propagator singularities were integrated by decomposing the propagators into a principal value (off-shell) part and an on-shell part. Both methods were implemented and compared in Ref. Vdh97b , and it was found that the +i$`ϵ`$ method yields differences of the order of 10% for every diagram as compared with the result of the principal value method. It is not surprising that, when summing hundreds of diagrams, an error of 10% on every diagram can easily be amplified due to the interference between the diagrams. To have confidence in the evaluation of the convolution of Eq. (4.2), the principal value integration method was compared in Ref. Vdh97b with a third independent method. This third method starts from the observation that the diagrams can be classified into four categories depending upon the number of propagators which can develop singularities : in the present case this number is 0, 1, 2 or 3. Besides the trivial case of zero singularities which can be integrated immediately, the diagrams with one or two propagator singularities can be integrated by performing a contour integration in the complex plane for one of the four integrations. For the most difficult case of three propagator singularities, it was found to be possible to evaluate the diagram by performing two contour integrations in the complex plane. In doing so, one achieves quite a fast convergence because the integrations along the real axis are replaced by integrations along semi-circles in the complex plane which are far from the propagator poles. This method was compared with the principal value integration method, and the same results were found up to 0.1% for each type of singularity Vdh97b . The principal value method was however found to converge much slower and is more complicated to implement, especially for the case with three singularities due to the coupled nature of the three principal value integrals. Comparing the results of Ref. Vdh97b with those of Ref. Kro (91), a rather good agreement was found for all helicity amplitudes, except for the helicity amplitude where both photon and proton helicities are positive, in which case both calculations differ strongly. Very recently, the PQCD calculation for RCS at large momentum transfer has been recalculated again in Ref. Broo (00), by also performing convolution integrals through contour integrations in the complex plane. The authors of Ref. Broo (00) also find a strong difference with the results of Kro (91) for the same helicity amplitude, where both photon and proton helicities are positive. Furthermore, in the angular region around $`90^o`$, where the PQCD formalism is supposed to be applicable, the authors of Ref. Broo (00) find a good agreement with the calculations of Ref. Vdh97b , keeping in mind that there is an overall normalization uncertainty in these PQCD calculations for RCS, associated with $`\alpha _s`$ and the valence quark wave function normalization. The remaining difference between the results of Refs. Vdh97b and those of Ref. Broo (00) seems to be isolated to a single helicity amplitudes and appears for backward scattering angles. We therefore limit ourselves in the following discussion to the results in the angular region around $`90^o`$ where the calculations of Refs. Vdh97b and Broo (00) are in good agreement, and which is the most relevant region for the PQCD calculation as it corresponds to the largest momentum transfer for a given photon energy. In Figs. 7 and 8, the PQCD calculations for RCS are shown for several model DA’s denoted as CZ Che (84), COZ Che (89), KS Kin (87), and the asymptotic DA. The highest energy data which exist for real Compton scattering were taken around 5 GeV and are shown in Fig. 7. Although the energy at which these experiments were performed is probably too low to justify a PQCD calculation, the comparison with these data is nevertheless shown in Fig. 7 for illustrative purposes. The normalization of the calculations shown at these very low scales corresponds to using a frozen coupling constant, with $`\alpha _s0.5`$. One first notices that the hard scattering amplitude for RCS has the $`s`$-dependence ($`Ts^2`$) which leads to the QCD scaling laws Bro (73), that is $`\frac{d\sigma }{dt}s^6`$ for Compton scattering or VCS. The unpolarized real Compton differential cross section (multiplied by the scaling factor $`s^6`$) is shown in Fig. 7 as function of the photon c.m. angle. It is observed that the result with the asymptotic DA ($`120x_1x_2x_3`$) is more than one decade below the results obtained with the amplitudes KS, COZ, and CZ, motivated by QCD sum rules. The results with KS, COZ and CZ show a similar characteristic angular dependence which is asymmetric around 90<sup>o</sup>. Note that in the forward and backward directions, which are dominated by diffractive mechanisms, a PQCD calculation is not reliable. Comparing the results obtained with KS, COZ and CZ, one notices that although these DA’s have nearly the same lowest moments, they lead to differences of a factor of two in the Compton scattering cross section. Consequently, this observable is sensitive enough to distinguish between various distribution amplitudes, provided, of course, one is in the regime where the hard scattering mechanism dominates. In Fig. 8, the polarized Compton cross sections are shown for the two helicity states of the photon and for a target proton with positive helicity. One remarks that for all DA’s there is a marked difference both in magnitude and angular dependence between the cross sections for the two photon helicities. Consequently, the resulting photon asymmetry $`\mathrm{\Sigma }`$, defined as $$\mathrm{\Sigma }_{}=\frac{\frac{d\sigma }{dt}(,\lambda =1)\frac{d\sigma }{dt}(,\lambda =1)}{\frac{d\sigma }{dt}(,\lambda =1)+\frac{d\sigma }{dt}(,\lambda =1)},$$ (41) where $`\lambda `$ is the helicity of the incoming photon and $``$ denotes a positive hadron helicity, changes sign for the DA’s KS, COZ, and CZ for different values of $`\mathrm{\Theta }_{\mathrm{c}.\mathrm{m}.}`$ as shown in Fig. 8. It is seen that the asymptotic DA on the other hand yields a very large, negative asymmetry around $`90^o`$. Therefore, it was suggested in Ref. Vdh97b that the photon asymmetry might be a particularly useful observable to distinguish between nucleon distribution amplitudes. The predicted sensitivity of the asymmetry to the nucleon DA can be used in the extraction of a DA from Compton scattering data in the scaling region. In Ref. Vdh97b , a procedure was outlined to extract a DA from Compton data in a model independent way by first expanding the DA in a set of basis functions and then using the angular information of the cross sections to fit the expansion coefficients. It was seen that the precision for these coefficients is greatly improved when one measures both unpolarized cross sections and photon asymmetries. A first dedicated experiment to measure the RCS differential cross section and the asymmetry of Eq. (41) for $`\mathrm{\Theta }_{\mathrm{c}.\mathrm{m}.}`$ around $`90^o`$, and for a real photon energy of 6 GeV, is planned at JLab Woj (99). In particular, it will be interesting to see if one approaches the PQCD result at these “lower” energies, and to study the interplay with soft-overlap type contributions for RCS as proposed in Refs. Rad98a ; Die (99). RCS experiments using a real photon energy in the 15 GeV range, might be feasible e.g. at the HERA ring in the foreseeable future d’Ho (96); D’An (97) and might open up prospects to study the nucleon valence wave function in a direct way. ## 5 Deeply virtual Compton scattering and skewed parton distributions ### 5.1 Introduction Much of the internal structure of the nucleon has been revealed during the last two decades through the inclusive scattering of high energy leptons on the nucleon in the Bjorken -or “Deep Inelastic Scattering” (DIS)- regime (where the photon virtuality $`Q^2`$ is very large, and $`x_B=Q^2/2p.q`$ finite). Unpolarized DIS experiments have mapped out the quark and gluon distributions in the nucleon, while polarized DIS experiments have shown that only a small fraction of the nucleon spin is carried by the quarks. This has stimulated new investigations to understand the nucleon spin. With the advent of the new generation of high-energy, high-luminosity lepton accelerators combined with large acceptance spectrometers, a wide variety of exclusive processes in the Bjorken regime are considered as experimentally accessible. In recent years, a unified theoretical description of such processes has emerged through a formalism introducing a new type of parton distributions, commonly denoted as skewed parton distributions (SPD’s) Ji (97); Rad96a . These SPD’s are generalizations of the parton distributions measured in DIS. It has been shown that these SPD’s, which parametrize the structure of the nucleon, allow one to describe, in leading order perturbative QCD (PQCD), various exclusive processes in the near forward direction, where the momentum transfer to the nucleon is small. Such non-forward processes were already considered in the literature a longer time ago, see e.g. Wat (82); Bar (82); Dit (88); Jai (93); Mul (94). The most promising of these non-forward hard exclusive processes are deeply virtual Compton scattering (DVCS) and longitudinal electroproduction of vector or pseudoscalar mesons at large $`Q^2`$. ### 5.2 Definitions and modelizations of skewed parton distributions The leading order PQCD diagrams for DVCS and hard meson electroproduction are of the type as shown in Fig. 9. The hard scale in Fig. 9 is the photon virtuality $`Q^2`$, which should be large (of the order of several GeV<sup>2</sup>), so as to be in the Bjorken regime. It has been proven Ji (97); Rad96a that the leading order DVCS amplitude in the forward direction can be factorized in a hard scattering part (which is exactly calculable in PQCD) and a soft, nonperturbative nucleon structure part as illustrated on the left panel of Fig. 9. The nucleon structure information can be parametrized, at leading order, in terms of four (quark helicity conserving) generalized structure functions. These functions are the SPD’s denoted by $`H,\stackrel{~}{H},E,\stackrel{~}{E}`$ which depend upon three variables : $`x`$, $`\xi `$ and $`t`$. The light-cone momentum <sup>3</sup><sup>3</sup>3using the definition $`a^\pm 1/\sqrt{2}(a^0\pm a^3)`$ for the light-cone components fraction $`x`$ is defined by $`k^+=xP^+`$, where $`k`$ is the quark loop momentum and $`P`$ is the average nucleon momentum ($`P=(p+p^{})/2`$, where $`p(p^{})`$ are the initial (final) nucleon four-momenta respectively). The skewedness variable $`\xi `$ is defined by $`\mathrm{\Delta }^+=2\xi P^+`$, where $`\mathrm{\Delta }=p^{}p`$ is the overall momentum transfer in the process, and where $`2\xi x_B/(1x_B/2)`$ in the Bjorken limit. Furthermore, the third variable entering the SPD’s is given by the Mandelstam invariant $`t=\mathrm{\Delta }^2`$, being the total squared momentum transfer to the nucleon. In a frame where the virtual photon momentum $`q^\mu `$ and the average nucleon momentum $`P^\mu `$ are collinear along the $`z`$-axis and in opposite direction, one can parametrize the non-perturbative object in the lower blobs of Fig. 9 as : $`{\displaystyle \frac{P^+}{2\pi }}{\displaystyle 𝑑y^{}e^{ixP^+y^{}}p^{^{}}\|\overline{\psi }_\beta (y/2)\psi _\alpha (y/2)\|p}\|_{y^+=\stackrel{}{y}_{}=0}`$ (42) $`=`$ $`{\displaystyle \frac{1}{4}}\{(\gamma ^{})_{\alpha \beta }[H^q(x,\xi ,t)\overline{N}(p^{^{}})\gamma ^+N(p)`$ $`+E^q(x,\xi ,t)\overline{N}(p^{^{}})i\sigma ^{+\kappa }{\displaystyle \frac{\mathrm{\Delta }_\kappa }{2m_N}}N(p)]`$ $`+(\gamma _5\gamma ^{})_{\alpha \beta }[\stackrel{~}{H}^q(x,\xi ,t)\overline{N}(p^{^{}})\gamma ^+\gamma _5N(p)`$ $`+\stackrel{~}{E}^q(x,\xi ,t)\overline{N}(p^{^{}})\gamma _5{\displaystyle \frac{\mathrm{\Delta }^+}{2m_N}}N(p)]\},`$ where $`\psi `$ is the quark field, $`N`$ the nucleon spinor and $`m_N`$ the nucleon mass. The lhs of Eq. (42) can be interpreted as a Fourier integral along the light-cone distance $`y^{}`$ of a quark-quark correlation function, representing the process where a quark is taken out of the initial nucleon (with momentum $`p`$) at the space-time point $`y/2`$, and is put back in the final nucleon (with momentum $`p^{}`$) at the space-time point $`y/2`$. This process takes place at equal light-cone time ($`y^+=0`$) and at zero transverse separation ($`\stackrel{}{y}_{}=0`$) between the quarks. The resulting one-dimensional Fourier integral along the light-cone distance $`y^{}`$ is with respect to the quark light-cone momentum $`xP^+`$. The rhs of Eq. (42) parametrizes this non-perturbative object in terms of four SPD’s, according to whether they correspond to a vector operator $`(\gamma ^{})_{\alpha \beta }`$ or an axial-vector operator $`(\gamma _5\gamma ^{})_{\alpha \beta }`$ at the quark level. The vector operator corresponds at the nucleon side to a vector transition (parametrized by the function $`H^q`$, for a quark of flavor $`q`$) and a tensor transition (parametrized by the function $`E^q`$). The axial-vector operator corresponds at the nucleon side to an axial-vector transition (function $`\stackrel{~}{H}^q`$) and a pseudoscalar transition (function $`\stackrel{~}{E}^q`$). In Fig. 9, the variable $`x`$ runs from -1 to 1. Therefore, the momentum fractions ($`x+\xi `$ or $`x\xi `$) of the active quarks can either be positive or negative. Since positive (negative) momentum fractions correspond to quarks (antiquarks), it has been noted in Rad96a that in this way, one can identify two regions for the SPD’s : when $`x>\xi `$ both partons represent quarks, whereas for $`x<\xi `$ both partons represent antiquarks. In these regions, the SPD’s are the generalizations of the usual parton distributions from DIS. Actually, in the forward direction, the SPD’s $`H`$ and $`\stackrel{~}{H}`$ reduce to the quark density distribution $`q(x)`$ and quark helicity distribution $`\mathrm{\Delta }q(x)`$ respectively, obtained from DIS : $$H^q(x,0,0)=q(x),\stackrel{~}{H}^q(x,0,0)=\mathrm{\Delta }q(x).$$ (43) The functions $`E`$ and $`\stackrel{~}{E}`$ are not measurable through DIS because the associated tensors in Eq. (42) vanish in the forward limit ($`\mathrm{\Delta }0`$). Therefore, $`E`$ and $`\stackrel{~}{E}`$ are new leading twist functions, which are accessible through the hard exclusive electroproduction reactions, discussed in the following. In the region $`\xi <x<\xi `$, one parton connected to the lower blob in Fig. 9 represents a quark and the other one an antiquark. In this region, the SPD’s behave like a meson distribution amplitude and contain completely new information about nucleon structure, because the region $`\xi <x<\xi `$ is absent in DIS, which corresponds to the limit $`\xi 0`$. Besides coinciding with the quark distributions at vanishing momentum transfer, the skewed parton distributions have interesting links with other nucleon structure quantities. The first moments of the SPD’s are related to the elastic form factors (FF) of the nucleon through model independent sum rules. By integrating Eq. (42) over $`x`$, one obtains the following relations for one quark flavor : $`{\displaystyle _1^{+1}}𝑑xH^q(x,\xi ,t)`$ $`=`$ $`F_1^q(t),`$ $`{\displaystyle _1^{+1}}𝑑xE^q(x,\xi ,t)`$ $`=`$ $`F_2^q(t),`$ $`{\displaystyle _1^{+1}}𝑑x\stackrel{~}{H}^q(x,\xi ,t)`$ $`=`$ $`g_A^q(t),`$ $`{\displaystyle _1^{+1}}𝑑x\stackrel{~}{E}^q(x,\xi ,t)`$ $`=`$ $`h_A^q(t).`$ (44) The elastic FF for one quark flavor on the rhs of Eqs. (44) are related to the physical ones (restricting oneself to $`u,d`$ and $`s`$ quark flavors) as : $$F_1^u=\mathrm{\hspace{0.17em}2}F_1^p+F_1^n+F_1^s,F_1^d=\mathrm{\hspace{0.17em}2}F_1^n+F_1^p+F_1^s,$$ (45) where $`F_1^p`$ and $`F_1^n`$ are the usual proton and neutron Dirac FF respectively, and where $`F_1^s`$ is the strangeness form factor. Relations similar to Eq. (45) hold for the Pauli FF $`F_2^q`$. For the axial vector FF one uses the isospin decomposition : $$g_A^u=\frac{1}{2}g_A+\frac{1}{2}g_A^0,g_A^d=\frac{1}{2}g_A+\frac{1}{2}g_A^0,$$ (46) where $`g_A(g_A^0)`$ are the isovector (isoscalar) axial FF respectively. Similar relations exist for $`h_A`$. The isovector axial form factor $`g_A`$ is known from experiment, with $`g_A(0)1.267`$. The induced pseudoscalar form factor $`h_A`$ contains an important pion pole contribution, through the partial conservation of the axial current (PCAC). A lot of the recent interest and activity in this field has been triggered by the observation of Ji (97) that the SPD’s may shed a new light on the “spin-puzzle”. Starting from a (color) gauge-invariant decomposition of the nucleon spin : $`1/2=J_q+J_g`$ , where $`J_q`$ and $`J_g`$ are the total quark and gluon angular momentum respectively, it was shown in Ji (97) that the second moment of the unpolarized SPD’s at $`t=0`$ gives $$J_q=\frac{1}{2}_1^{+1}𝑑xx\left[H^q(x,\xi ,t=0)+E^q(x,\xi ,t=0)\right],$$ (47) and this relation is independent of $`\xi `$. The quark angular momentum $`J_q`$ decomposes as : $`J_q=\mathrm{\Delta }\mathrm{\Sigma }/2+L_q`$ , where $`\mathrm{\Delta }\mathrm{\Sigma }/2`$ and $`L_q`$ are the quark spin and orbital angular momentum respectively. As $`\mathrm{\Delta }\mathrm{\Sigma }`$ is measured through polarized DIS experiments, a measurement of the sum rule of Eq. (47) in terms of the SPD’s, provides a model independent way to determine the quark orbital contribution $`L_q`$ to the nucleon spin. Ultimately, one wants to extract the SPD’s from data, but in order to evaluate the electroproduction observables, and to study their sensitivity to the new physics, one needs an educated guess for the SPD’s. In Ref. Vdh (99), a model for the SPD’s was constructed using a $`\xi `$-dependent product ansatz (for the double distributions introduced in Ref. Rad98b ) of a quark distribution and an asymptotic “meson-like” distribution amplitude (see Ref. Vdh (99) for more details). For the quark distributions, the MRST98 parametrization Mar (98) is used as input. The $`t`$-dependence of the model for the SPD’s is given by the corresponding FF (Dirac form factor for $`H`$, axial form factor for $`\stackrel{~}{H}`$), so that the first moments of the SPD’s are satisfied by construction. As an example, the $`d`$-quark SPD (formerly also denoted as off-forward parton distribution (OFPD)), using the above described ansatz, is shown in Fig. 10. One observes from Fig. 10 the transition from a quark distribution ($`\xi =0`$) to a meson distribution amplitude ($`\xi =1`$). Model calculation of the SPD’s are currently possible within the QCD chiral models for intermediate $`x_B`$. In particular, a calculation Pet (98) in the chiral quark soliton model (see Ref. Chr (96) for a review) found a strong dependence of the SPD’s on $`\xi `$ and fast “crossovers” at $`\|x\|=\xi `$. Such behavior is related to the fact that the SPD’s in the region $`\xi <x<\xi `$ have properties of meson distribution amplitudes. In particular for the SPD $`H`$, this can be seen as being due to a scalar-isoscalar two-pion exchange contribution Pol (99), indicating that the SPD’s are qualitatively a richer source of nucleon structure information than ordinary parton distributions. One may expect that eventually it will be possible to calculate SPD’s for intermediate $`x_B`$ using lattice QCD. ### 5.3 Leading order amplitudes and observables for DVCS and hard meson electroproduction The leading order (L.O.) DVCS amplitude in the forward direction is given Ji (97) by the handbag diagram shown on the left panel of Fig. 9 (the crossed diagram which is not shown is also understood). A formal factorization proof for DVCS has been given in Refs. Ji98a ; Col (99). To calculate the DVCS amplitude in the Bjorken regime, it is natural to express the momenta in the process ($`q^\mu `$ of the virtual photon, $`q^\mu `$ of the real photon, and $`P^\mu `$ denoting the average nucleon momentum) in terms of the lightlike vectors $$\stackrel{~}{p}^\mu =\frac{P^+}{\sqrt{2}}(1,0,0,1),n^\mu =\frac{1}{P^+\sqrt{2}}(1,0,0,1).$$ (48) Using the parametrization of Eq. (42) for the bilocal quark operator, the L.O. DVCS tensor $`H_{L.O.DVCS}^{\mu \nu }`$ (defined e.g. in Gui (98)) follows from the handbag diagrams as : $`H_{L.O.DVCS}^{\mu \nu }`$ (49) $`=`$ $`{\displaystyle \frac{1}{2}}\left[\stackrel{~}{p}^\mu n^\nu +\stackrel{~}{p}^\nu n^\mu g^{\mu \nu }\right]`$ $`\times {\displaystyle _1^{+1}}dx[{\displaystyle \frac{1}{x\xi +iϵ}}+{\displaystyle \frac{1}{x+\xi iϵ}}]`$ $`\times [H_{DVCS}^p(x,\xi ,t)\overline{N}(p^{^{}})\gamma .nN(p)`$ $`+E_{DVCS}^p(x,\xi ,t)\overline{N}(p^{^{}})i\sigma ^{\kappa \lambda }{\displaystyle \frac{n_\kappa \mathrm{\Delta }_\lambda }{2m_N}}N(p)]`$ $`+`$ $`{\displaystyle \frac{1}{2}}\left[i\epsilon ^{\mu \nu \kappa \lambda }\stackrel{~}{p}_\kappa n_\lambda \right]{\displaystyle _1^{+1}}𝑑x\left[{\displaystyle \frac{1}{x\xi +iϵ}}{\displaystyle \frac{1}{x+\xi iϵ}}\right]`$ $`\times [\stackrel{~}{H}_{DVCS}^p(x,\xi ,t)\overline{N}(p^{^{}})\gamma .n\gamma _5N(p)`$ $`+\stackrel{~}{E}_{DVCS}^p(x,\xi ,t)\overline{N}(p^{^{}})\gamma _5{\displaystyle \frac{\mathrm{\Delta }n}{2m_N}}N(p)].`$ On the rhs of the DVCS tensor of Eq. (49), the SPD’s $`H,\stackrel{~}{H},E,\stackrel{~}{E}`$ enter in a convolution integral over the quark momentum fraction $`x`$. This is a qualitative difference compared with the DIS amplitude, where one is only sensitive (through the optical theorem) to the imaginary part of the forward double virtual Compton amplitude. We refer to Ref. Gui (98) for details and for the formalism to calculate DVCS observables starting from the DVCS tensor of Eq. (49). The leading order DVCS amplitude corresponding to Eq. (49), is exactly gauge invariant with respect to the virtual photon, i.e. $`q_\nu H_{L.O.DVCS}^{\mu \nu }=0`$. However, electromagnetic gauge invariance is violated by the real photon except in the forward direction. This violation of gauge invariance is a higher twist (twist-3) effect compared to the leading order term $`H_{L.O.DVCS}^{\mu \nu }`$. Since $`q_\mu ^{^{}}H_{L.O.DVCS}^{\mu \nu }\mathrm{\Delta }_{}`$, an improved DVCS amplitude linear in $`\mathrm{\Delta }_{}`$ has been proposed in Ref. Gui (98) to restore gauge invariance (in the nonforward direction) in a heuristic way : $$H_{DVCS}^{\mu \nu }=H_{L.O.DVCS}^{\mu \nu }+\frac{\stackrel{~}{p}^\mu }{\left(\stackrel{~}{p}q^{^{}}\right)}\left(\mathrm{\Delta }_{}\right)_\lambda H_{L.O.DVCS}^{\lambda \nu },$$ (50) leading to a correction term to the L.O. DVCS amplitude of order $`O\left(\mathrm{\Delta }_{}/Q\right)`$. Very recently, the gauge invariance of the DVCS amplitude was addressed in much more detail in several works Ani (00); Pen (00); Bel00b . It was found that the twist-3 terms which restore gauge invariance (to twist-4 accuracy) involve two contributions. First there are terms proportional to the twist-2 SPD’s of Eq. (49), which were found to completely coincide with the improved DVCS amplitude of Eq. (50). In addition, there are terms which are characterized by new ‘transverse’ SPD’s (see Refs. Ani (00); Pen (00); Bel00b for details). These latter functions are suppressed by one power $`1/Q`$ compared with the contribution of the twist-2 SPD’s in DVCS cross sections, and could in principle be separated by measuring DVCS observables over a sufficiently large $`Q^2`$ range (see e.g. Die (97) for tests of the handbag approximation to DVCS). In view of current DVCS experiments which are performed or planned at $`Q^2`$ in the few GeV<sup>2</sup> range only, the numerical importance of those additional contributions remains to be investigated. Besides the DVCS process, a factorization proof was also given for the L.O. meson electroproduction amplitudes in the Bjorken regime Col (97); Rad96b , which is illustrated on the right panel of Fig. 9. This factorization theorem only applies when the virtual photon is longitudinally polarized. In the valence region, the L.O. amplitude $`^L`$ for meson production by a longitudinal photon consists of evaluating Fig. 9 (right panel, where only one of the four L.O. diagrams is shown) with the one-gluon exchange diagrams as hard scattering kernel. In this way, the L.O. expressions for $`\rho _L^0`$ (longitudinally polarized vector meson) and $`\pi ^0`$ electroproduction were calculated in Vdh (98) (see also Ref. Man (98)) as : $`_{\rho _L^0}^L`$ $`=`$ $`ie{\displaystyle \frac{4}{9}}{\displaystyle \frac{1}{Q}}\left[{\displaystyle _0^1}𝑑z{\displaystyle \frac{\mathrm{\Phi }_\rho (z)}{z}}\right]`$ (51) $`\times `$ $`{\displaystyle \frac{1}{2}}{\displaystyle _1^{+1}}𝑑x\left[{\displaystyle \frac{1}{x\xi +iϵ}}+{\displaystyle \frac{1}{x+\xi iϵ}}\right]`$ $`\times `$ $`(4\pi \alpha _s)\{H_{\rho _L^0}^p(x,\xi ,t)\overline{N}(p^{^{}})\gamma .nN(p)`$ $`+E_{\rho _L^0}^p(x,\xi ,t)\overline{N}(p^{^{}})i\sigma ^{\kappa \lambda }{\displaystyle \frac{n_\kappa \mathrm{\Delta }_\lambda }{2m_N}}N(p)\},`$ $`_{\pi ^0}^L`$ $`=`$ $`ie{\displaystyle \frac{4}{9}}{\displaystyle \frac{1}{Q}}\left[{\displaystyle _0^1}𝑑z{\displaystyle \frac{\mathrm{\Phi }_\pi (z)}{z}}\right]`$ (52) $`\times `$ $`{\displaystyle \frac{1}{2}}{\displaystyle _1^{+1}}𝑑x\left[{\displaystyle \frac{1}{x\xi +iϵ}}+{\displaystyle \frac{1}{x+\xi iϵ}}\right]`$ $`\times `$ $`(4\pi \alpha _s)\{\stackrel{~}{H}_{\pi ^0}^p(x,\xi ,t)\overline{N}(p^{^{}})\gamma .n\gamma _5N(p)`$ $`+\stackrel{~}{E}_{\pi ^0}^p(x,\xi ,t)\overline{N}(p^{^{}})\gamma _5{\displaystyle \frac{\mathrm{\Delta }n}{2m_N}}N(p)\},`$ where $`\alpha _s`$ is the QCD coupling constant. Because the quark helicity is conserved in the hard scattering process, one finds the interesting result that the vector meson electroproduction amplitude depends only on the unpolarized SPD’s $`H`$ and $`E`$, whereas the pseudoscalar meson electroproduction amplitudes depend only on the polarized SPD’s $`\stackrel{~}{H}`$ and $`\stackrel{~}{E}`$. In contrast, the DVCS amplitude of Eq. (49) depends on both the unpolarized and polarized SPD’s. Another difference from DVCS, is the fact that the meson electroproduction amplitudes require additional non-perturbative input from the meson distribution amplitudes $`\mathrm{\Phi }_\rho (z)`$ and $`\mathrm{\Phi }_\pi (z)`$ respectively, for which the asymptotic forms are taken in the calculations. From Eqs. (51,52), one furthermore sees that the L.O. longitudinal amplitudes for meson electroproduction behave as $`1/Q`$. At large $`Q^2`$, fixed $`x_B`$ and fixed $`t`$, this leads to a $`1/Q^6`$ behavior for the longitudinal cross section $`d\sigma _L/dt`$, which provides an experimental signature (scaling) of the leading order mechanism. Expressions analogous to Eqs. (51, 52) have also been worked out for the charged meson channels $`\rho ^\pm ,\pi ^\pm `$ as well as for the $`\omega ,\varphi `$ and $`\eta `$ channels (see Refs. Fra (99); Man99a ; Man99b ; Vdh (99) for details). According to the considered reaction, the SPD’s enter in different combinations due to the charges and isospin factors. For DVCS on the proton, the combination is $$H_{DVCS}^p(x,\xi ,t)=\frac{4}{9}H^u+\frac{1}{9}H^d+\frac{1}{9}H^s,$$ (53) and similarly for $`\stackrel{~}{H}`$, $`E`$ and $`\stackrel{~}{E}`$. For electroproduction of $`\rho ^0`$ and $`\pi ^0`$ on the proton, the isospin structure yields the combination $`H_{\rho ^0}^p(x,\xi ,t)`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}\left\{{\displaystyle \frac{2}{3}}H^u+{\displaystyle \frac{1}{3}}H^d\right\},`$ (54) $`\stackrel{~}{H}_{\pi ^0}^p(x,\xi ,t)`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}\left\{{\displaystyle \frac{2}{3}}\stackrel{~}{H}^u+{\displaystyle \frac{1}{3}}\stackrel{~}{H}^d\right\},`$ (55) and similar for $`E`$ and $`\stackrel{~}{E}`$. Corresponding relations for the $`\rho ^\pm ,\omega ,\varphi ,\pi ^\pm `$ and $`\eta `$ channels can be found in Refs. Fra (99); Man99a ; Man99b ; Vdh (99). Therefore, the measurements of the different meson production channels are sensitive to different combinations of the same universal SPD’s, and allow us to perform a flavor separation of the SPD’s, provided one is able to deconvolute the SPD’s from the leading order amplitudes. Some representative results for DVCS and meson electroproduction observables using the $`\xi `$-dependent ansatz for the SPD’s, are shown in the following. More detailed results can be found in Refs. Vdh (98); Gui (98); Vdh (99). Before considering the extraction of the SPD’s from electroproduction data, it is compulsory to demonstrate that the scaling regime has been reached. In Fig. 11, the forward longitudinal electroproduction cross sections are shown as a function of $`Q^2`$ and the L.O. predictions are compared for different mesons. The L.O. amplitude for longitudinal electroproduction of mesons was seen to behave as $`1/Q`$, leading to a $`1/Q^6`$ scaling behavior for $`d\sigma _L/dt`$. By comparing the different vector meson channels in Fig. 11, one sees that the $`\rho _L^0`$ channel yields the largest cross section. The $`\omega _L`$ channel in the valence region ($`x_B`$ 0.3) is about a factor of 5 smaller than the $`\rho _L^0`$ channel, which is to be compared with the ratio at small $`x_B`$ (in the diffractive regime) where $`\rho ^0`$ : $`\omega `$ = 9 : 1. The $`\rho _L^+`$ channel, which is sensitive to the isovector combination of the unpolarized SPD’s, yields a cross section comparable to the $`\omega _L`$ channel. The $`\rho _L^+`$ channel is interesting as there is no competing diffractive contribution, and therefore allows to test directly the quark SPD’s. The three vector meson channels ($`\rho _L^0`$, $`\rho _L^+`$, $`\omega _L`$) are highly complementary in order to perform a flavor separation of the unpolarized SPD’s $`H^u`$ and $`H^d`$. A dedicated experiment is planned at JLab at 6 GeV in the near future Guid (98) to investigate the onset of the scaling behavior for $`\rho _L^0`$ electroproduction in the valence region ($`Q^23.5`$ GeV<sup>2</sup>, $`x_B0.3`$). For the pseudoscalar mesons which involve the polarized SPD’s, one remarks in Fig. 11 the prominent contribution of the charged pion pole to the $`\pi ^+`$ cross section. For the contribution proportional to the SPD $`\stackrel{~}{H}`$, it is also seen that the $`\pi ^0`$ channel is about a factor of 5 below the $`\pi ^+`$ channel due to isospin factors. In the $`\pi ^0`$ channel, the $`u`$\- and $`d`$-quark polarized SPD’s enter with the same sign, whereas in the $`\pi ^+`$ channel, they enter with opposite signs. As the polarized SPD’s are constructed here from the corresponding polarized parton distributions, the difference between the predictions for the $`\pi ^0`$ and $`\pi ^+`$ channels results from the fact that the polarized $`d`$-quark distribution is opposite in sign to the polarized $`u`$-quark distribution. For the $`\eta `$ channel, the ansatz for the SPD $`\stackrel{~}{H}`$ based on the polarized quark distributions yields a prediction comparable to the $`\pi ^0`$ cross section. For the $`\gamma `$ leptoproduction in the few GeV beam energy range, the cross section is dominated by the Bethe-Heitler (BH) process (see Ref. Gui (98)). However, it was suggested in Ref. Gui (98) that an exploration of DVCS might be possible if the beam is polarized. The electron single spin asymmetry (SSA) does not vanish out of plane due to the interference between the purely real BH process and the imaginary part of the DVCS amplitude. Because the SSA measures the imaginary part of the DVCS amplitude, it is directly proportional to a linear combination of the SPD’s along the line $`x=\xi `$. In fact, the SSA maps out an ‘envelope’ function, e.g. $`H(x=\xi ,\xi ,t)`$, as shown e.g. in Fig. 13 for the valence down quark SPD in the ansatz corresponding with Fig. 10. In Fig. 13, it is shown that the SSA yields a sizeable asymmetry for JLab kinematics, and displays a sensitivity to the $`\xi `$-dependent shape of the SPD’s. An experiment to measure the SSA for DVCS has very recently been proposed at JLab at 6 GeV Bert (99). The SSA for DVCS is at present also measured at HERMES Arm (99). Going up in energy, the increasing virtual photon flux factor boosts the DVCS part of the $`\gamma `$ leptoproduction cross section, making it more important compared to the BH contribution. This provides a nice opportunity for COMPASS at 200 GeV beam energy, where experiments have been proposed for DVCS d’Ho99a and meson electroproduction Poc (99). ### 5.4 Extension to hard exclusive electroproduction of decuplet baryons In the previous section, the main focus were the reactions $`\gamma ^{}+N\gamma +N^{}`$ and $`\gamma _L^{}+NM+N^{}`$ with $`M`$ a meson, and where $`N^{}`$ is an octet baryon. One of the intriguing questions of medium-energy QCD dynamics is the differences and similarities in the structure of baryons belonging to the different $`SU(3)_f`$ multiplets. In particular, a naive constituent quark model predicts that they are similar, while there are suggestions that due to a strong attraction between the quarks in the spin-isospin zero channel, diquark correlations should be important in the baryon octet but not in the decuplet Sch (98). At the same time the chiral models suggest that in the limit of a large number of colors (large $`N_c`$) of QCD, which is known to be a useful guideline, nucleons and $`\mathrm{\Delta }`$ isobars are different rotational excitations of the same soliton Adk (83); Dia (88). For these studies, the potential of the process $`\gamma _L^{}+N\pi +\mathrm{\Delta }`$ as well as the DVCS process $`\gamma ^{}+N\gamma +\mathrm{\Delta }`$, was explored in Ref. Fra (00). In addition, the study of the processes with production of decuplet baryons has also a practical usefulness, because in the experiments with low resolution in the mass of the recoiling system ($`\mathrm{\Delta }M`$ 300 MeV for HERMES in the current set-up), the estimates of $`\mathrm{\Delta }`$ production are necessary to extract the $`NN`$ SPD’s from such data. In Ref. Fra (00), a new set of SPD’s were introduced for the axial $`N\mathrm{\Delta }`$ (isovector) transition, denoted as $`C_i^{(3)}`$, which enter into $`\pi \mathrm{\Delta }`$ electroproduction : $`{\displaystyle \frac{P^+}{2\pi }}{\displaystyle 𝑑y^{}e^{ixP^+y^{}}\mathrm{\Delta }^+\|\overline{\psi }(y/2)n/\gamma ^5\psi (y/2)\|N}\|_{y^+=\stackrel{}{y}_{}=0}`$ $`=\overline{\psi }^\beta (p^{})[C_1^{(3)}(x,\xi ,t)n_\beta `$ $`+C_2^{(3)}(x,\xi ,t){\displaystyle \frac{\mathrm{\Delta }_\beta (n\mathrm{\Delta })}{m_N^2}}+\mathrm{}]N(p),`$ (56) where the same notations are used as before, and where $`\psi ^\beta (p^{})`$ is the Rarita-Schwinger spinor for the $`\mathrm{\Delta }`$ isobar. In Eq. (56), the ellipses denote other contributions which are suppressed at large $`N_c`$. For the $`N\mathrm{\Delta }`$ DVCS process, besides the axial SPD’s, also vector SPD’s enter which were also defined in Ref. Fra (00). The observation that in the large $`N_c`$ limit, the nucleon and $`\mathrm{\Delta }`$ are rotational excitations of the same classical soliton object, allows us to derive a number of relations between $`NN`$ and $`N\mathrm{\Delta }`$ SPD’s. For $`C_1^{(3)}`$ and $`C_2^{(3)}`$, these have the form Fra (99) : $`C_1^{(3)}(x,\xi ,t)`$ $`=`$ $`\sqrt{3}\stackrel{~}{H}^{(3)}(x,\xi ,t),`$ (57) $`C_2^{(3)}(x,\xi ,t)`$ $`=`$ $`\sqrt{3}/4\stackrel{~}{E}^{(3)}(x,\xi ,t),`$ (58) in terms of the (isovector) $`NN`$ SPD $`\stackrel{~}{H}^{(3)}=\stackrel{~}{H}^u\stackrel{~}{H}^d`$, and analogously for $`\stackrel{~}{E}^{(3)}`$. Using the large $`N_c`$ relations of Eq. (57), one can easily derive the relations between the different cross sections for charged pion production as $`\sigma _L^{\gamma ^{}p\pi ^+n}:\sigma _L^{\gamma ^{}p\pi ^+\mathrm{\Delta }^0}:\sigma _L^{\gamma ^{}p\pi ^{}\mathrm{\Delta }^{++}}:\sigma _L^{\gamma ^{}n\pi ^{}p}1:0.5:1.25:0.8`$. Besides the cross section $`\sigma _L`$, the second observable involving only longitudinal amplitudes and being a leading order observable for hard exclusive meson electroproduction, is the single spin asymmetry, for a proton target polarized perpendicular to the reaction plane (or the equivalent recoil polarization observable) Fra (99). These transverse spin asymmetries for $`\pi ^+n`$ and $`\pi ^+\mathrm{\Delta }^0`$ are shown in Fig. 14. It is obvious from Fig. 14, that large transverse spin asymmetries are predicted for these processes, related to the peculiar feature of chiral QCD. As a consequence, investigations of these processes can provide unique tests of the soliton type approach to baryon structure. The spin asymmetries are likely to be less sensitive to higher twist effects and hence can be explored already using the HERMES detector and JLab at higher energies. Furthermore, the study of these processes would allow one to make a more reliable separation of the $`\pi `$ pole contribution in the electroproduction of pions, which is mandatory for the measurement of the pion elastic form factor at higher $`Q^2`$. ### 5.5 Power corrections to the leading order amplitudes When measuring hard electroproduction reactions in the region $`Q^2120`$ GeV<sup>2</sup>, there arises the question of the importance of power corrections to the leading order amplitudes, i.e. how fast does one approach the scaling regime predicted by the L.O. amplitudes. One source of power corrections is evident from the structure of the matrix element of Eq. (42) which defines the SPD’s, where the quarks are taken at zero transverse separation. This amounts to neglect, at leading order, the quark’s transverse momentum compared with its large longitudinal (+ component) momentum. A first estimate of these corrections due to the quark’s intrinsic transverse momentum has been obtained in Ref. Vdh (99), which is referred to for details. This correction is known to be important for a successful description at the lower $`Q^2`$ values of the $`\pi ^0\gamma ^{}\gamma `$ transition FF, for which data exist in the range $`Q^2110`$ GeV<sup>2</sup>. For the pion elastic FF in the transition region before asymptotia is reached, the power corrections due to both the transverse momentum dependence and the soft overlap mechanism (i.e. the process which does not proceed through one-gluon exchange) are quantitatively important. The result for the pion elastic FF is shown in Fig. 15, where it is seen that the leading order PQCD result is approached only at very large $`Q^2`$. The correction including the transverse momentum dependence gives a substantial suppression at lower $`Q^2`$ (about a factor of two around $`Q^2`$ 5 GeV<sup>2</sup>). At these lower $`Q^2`$ values, the inclusion of the transverse momentum dependence renders the PQCD calculation internally consistent. These form factors were taken as a guide in Ref. Vdh (99) to estimate the corrections due to the parton’s intrinsic transverse momentum dependence in the DVCS and hard meson electroproduction amplitudes. Although experimental data for $`\rho _L^0`$ electroproduction at larger $`Q^2`$ do not yet exist in the valence region ($`x_B`$ 0.3), the reaction $`\gamma _L^{}p\rho _L^0p`$ has been measured at smaller values of $`x_B`$. Therefore, in Fig. 16 the calculations are compared to those data, in order to study how the valence region is approached, in which one is sensitive to the quark SPD’s. For the purpose of this discussion, we call the mechanism proceeding through the quark SPD’s (i.e. the right panel of Fig. 9), the Quark Exchange Mechanism (QEM). Besides the QEM, $`\rho ^0`$ electroproduction at large $`Q^2`$ and small $`x_B`$ proceeds predominantly through a perturbative two-gluon exchange mechanism (PTGEM) as studied in Ref. Fra (96). To compare to the data at intermediate $`Q^2`$, the power corrections due to the parton’s intrinsic transverse momentum dependence were implemented in both mechanisms (see Ref. Vdh (99) for details), which gives a significant reduction at the lower $`Q^2`$. The results are compared with the data in Fig. 16, showing that the PTGEM explains well the fast increase of the cross section at high c.m. energy ($`W`$), but substantially underestimates the data at the lower energies. This is where the QEM is expected to contribute since $`x_B`$ is then in the valence region. The results including the QEM describe the change of behavior of the data at lower $`W`$ quite nicely. Recently, $`\rho _L^0`$ data have been obtained by the HERMES Collaboration for $`Q^2`$ up to 5 GeV<sup>2</sup> and around $`W`$ 5 GeV Air (00). These data show a clear dominance of the QEM in the intermediate $`W`$ range as predicted in Vdh (98, 99). The model ansatz for the SPD’s of Ref. Vdh (99) gives a fairly good agreement with these longitudinal $`\rho ^0`$ electroproduction data Air (00). ### 5.6 Perspectives and outlook In order to extract SPD’s from forthcoming data, the $`Q^2`$ evolution of the SPD’s has already been worked out at the next-to-leading order level Bel (98), which shows that the $`Q^2`$ evolution of the SPD’s interpolates between the evolution of the parton distributions and the evolution of the distribution amplitudes. Also radiative corrections to the coefficient functions have been calculated recently in next-to-leading order Bel00a . A major open theoretical question in this field is how the SPD’s can be deconvoluted from the leading order amplitudes. Suitable parametrizations of the SPD’s, incorporating all physical constraints, might be one avenue to tackle this problem. In absence of a solution to this problem, one has to resort to model calcuations or educated guesses for the SPD’s in order to compare with the data. On the experimental side, it is clear that new and accurate data are needed for various exclusive channels at large $`Q^2`$ in the valence region, where the quark exchange mechanism dominates. Several experiments are being performed or are planned or proposed at JLab Guid (98); Bert (99, 00), HERMES and COMPASS Poc (99); d’Ho99a . Looking somewhat further into the future, the measurement of hard exclusive reactions will be one of the central themes for the planned upgrade of JLab to 12 GeV Bur (00). A facility with high luminosity combined with an intermediate energy of around 25 GeV, such as e.g. the ELFE project Bur (99), will be a dedicated facility to measure these exclusive reactions at high momentum transfer and to map out these new SPD’s in detail. Although such exclusive experiments at large $`Q^2`$ are quite demanding, the fundamental interest of the SPD’s justifies an effort towards their experimental determination. ## 6 QED radiative corrections to virtual Compton scattering ### 6.1 Introduction We have discussed in section 3 how VCS below pion production threshold, allows us to access generalized polarizabilities of the proton. Furthermore, we have seen in section 5 that VCS in the Bjorken regime determines generalized parton distributions of the nucleon. In both regimes, experiments are either being done or planned for the near future. In order to extract the nucleon structure information of interest from the $`epep\gamma `$ reaction, especially in those kinematical situations where the Bethe-Heitler process is not negligible, it is indispensable to have a very good understanding of the QED radiative corrections to the $`epep\gamma `$ reaction. The $`epep\gamma `$ reaction is particular in comparison with other electron scattering reactions, because the photon can be emitted from both the proton side (this is the VCS process which contains the nucleon structure information of interest) or from one of the electrons (which is the Bethe-Heitler process). The radiative corrections obtained from the Bethe-Heitler process differ formally from the case of elastic electron scattering. ### 6.2 Results for the QED radiative corrections to VCS The first order QED radiative corrections to the $`epep\gamma `$ reaction were calculated in Ref. Vdh00a . The one-loop virtual radiative corrections have been evaluated by a combined analytical-numerical method. Several tests were performed to cross-check the numerical method used. It was shown in Ref. Vdh00a how all IR divergences cancel when adding the soft-photon emission processes. Furthermore, a fully numerical method was presented for the photon emission processes where the photon energy is not very small compared with the electron energies, which makes up the radiative tail. Fig. 17 shows as representative result the effect of the radiative corrections on the VCS differential cross section for the MAMI VCS experiment Roc (00) at an outgoing photon energy of $`\mathrm{q}^{}`$ = 111.5 MeV (we refer to Vdh00a for all details and more results). It is seen that the total effect of the radiative corrections in these kinematics is a reduction of the BH+Born cross section of the order of 20%. The effect of the radiative corrections was also confirmed by the experimental results at very low energy of the outgoing photon ($`\mathrm{q}^{}`$ = 33 MeV), where the effect of the GP’s is negligible. From the difference between the radiatively corrected data and the BH + Born result, the two values of Eq. (25) for the combinations of the proton’s GP’s at $`Q^2`$ 0.33 GeV<sup>2</sup> have been extracted in Ref. Roc (00). In Ref. Vdh00a , calculations of the VCS radiative corrections were also given for unpolarized and polarized VCS observables both at low energies and in the Bjorken regime. The results will be an important tool for the analysis of present and forthcoming VCS experiments, in order to extract the nucleon structure information from these experiments. ## 7 Conclusions and outlook It has been discussed how the real and virtual Compton scattering in different kinematical regimes provide new tools to extract nucleon structure information. It has been seen that for RCS at low energy, new accurate data have become available which not only allow the extraction of scalar polarizabilities of the proton, but also start to explore the spin polarizabilities of the nucleon. Those spin polarizabilities have been calculated recently to $`O(p^4)`$ in HBChPT. A fixed-t dispersion relation formalism was developed to extract the nucleon polarizabilities with a minimum of model dependence from both unpolarized and polarized RCS data. The DR formalism was also used to obtain information on new higher order polarizabilities of the proton, providing new nucleon structure observables and a new testing ground for the chiral calculations. The VCS reaction at low photon energy maps out the spatial distribution of the polarization densities of the proton, through generalized polarizabilities. Over the last few years, the VCS has become a mature field and a first experiment at MAMI at low energy has been successfully completed. In order to extract GP’s from VCS data over a larger range of energies, a dispersion relation formalism is underway, providing a new tool to analyze such data. The DR formalism provides already results for 4 of the 6 GP’s, which can be confronted with chiral predictions. The RCS reaction at high energy and large momentum transfer is a tool to access information on the partonic structure of the nucleon. The PQCD predictions for wide angle real Compton scattering (90<sup>o</sup> in the c.m.) show a strikingly different behavior than competing soft-overlap type mechanisms, and forthcoming experiments can teach us about the interplay of both mechanism at accessible energies. The VCS reaction in the Bjorken regime and associated hard electroproduction reactions give access to new, generalized (skewed) parton distributions. The study of SPD’s has opened up a whole new field in the study of nucleon structure. These observables unify two different fields as they interpolate between purely inclusive quantities (parton distributions) on the one hand and between simple exclusive quantities (such as form factors) on the other hand. Besides the SPD’s $`H`$ and $`\stackrel{~}{H}`$, which reduce in the DIS limit to the quark distribution and quark helicity distribution respectively, there are two entirely new leading twist SPD’s ($`E`$ and $`\stackrel{~}{E}`$), which cannot be accessed in DIS. The non-perturbative information contained in the SPD’s is rather rich as they are functions of 3 different variables. In particular, the skewedness variable $`\xi `$ leads to different regions where one is sensitive either to quark distribution type information or to meson distribution amplitude information in the nucleon. Through a sum rule, two of these SPD’s ($`H`$ and $`E`$) determine the quark orbital angular momentum contribution to the nucleon spin. An educated guess was shown for these SPD’s, which was used to estimate the leading order DVCS amplitude. Furthermore, the leading order meson electroproduction amplitudes were discussed and compared to the available data. In particular it was seen that in the intermediate $`W`$ range (valence region), a dominance of the handbag mechanism is predicted for $`\rho ^0`$ electroproduction, which seems to be well confirmed by recent HERMES data. Furthermore, the extension of the formalism of the SPD’s to the $`N\mathrm{\Delta }`$ transition was discussed. The large $`N_c`$ limit allows to relate the $`N\mathrm{\Delta }`$ SPD’s to the $`NN`$ SPD’s. The transverse spin asymmetry was discussed as a promising observable, likely to be less sensitive to higher twist effects. It is easy to foresee that the fields of real and virtual Compton scattering will show important activities in the near future both on the theoretical and experimental sides, and that an attempt to review them is very timely. It is hoped however, that the works initiated and discussed in this paper will stimulate further efforts on the theoretical and experimental sides to extend our knowledge of nucleon structure in new directions. ## Acknowledgements It is a pleasure for me to thank my collaborators who participated in an important way in the different works referred to in this paper : N. D’Hose, D. Drechsel, L.L. Frankfurt, J.M. Friedrich, M. Gorchtein, P.A.M. Guichon, M. Guidal, B. Holstein, D. Lhuillier, D. Marchand, A. Metz, B. Pasquini, M.V. Polyakov, M. Strikman, L. Van Hoorebeke, and J. Van de Wiele. I also like to thank the many experimental colleagues working in these fields for numerous and very useful discussions. Furthermore, I would like to thank D. Drechsel also for a careful reading of the text. This work was supported by the Deutsche Forschungsgemeinschaft (SFB443).
warning/0007/math0007184.html	ar5iv	text	# Untitled Document 3-Sasakian Geometry, Nilpotent Orbits, and Exceptional Quotients Charles P. Boyer Krzysztof Galicki and Paolo Piccinni During the preparation of this work the first two authors were supported by NSF grant DMS-9970904. The third author was supported by MURST and CNR. Abstract. Using 3-Sasakian reduction techniques we obtain infinite families of new 3-Sasakian manifolds $`(p_1,p_2,p_3)`$ and $`(p_1,p_2,p_3,p_4)`$ in dimension 11 and 15 respectively. The metric cone on $`(p_1,p_2,p_3)`$ is a generalization of the Kronheimer hyperkähler metric on the regular maximal nilpotent orbit of $`\text{s}\text{l}(3,\mathrm{})`$ whereas the cone on $`(p_1,p_2,p_3,p_4)`$ generalizes the hyperkähler metric on the 16-dimensional orbit of $`\text{s}\text{o}(6,\mathrm{})`$. These are first examples of 3-Sasakian metrics which are neither homogeneous nor toric. In addition we consider some further $`U\left(1\right)`$-reductions of $`(p_1,p_2,p_3)`$. These yield examples of non-toric 3-Sasakian orbifold metrics in dimensions 7. As a result we obtain explicit families $`𝒪\left(\mathrm{\Theta }\right)`$ of compact self-dual positive scalar curvature Einstein metrics with orbifold singularities and with only one Killing vector field. Introduction In 1990 Kronheimer showed that the co-adjoint orbits in the complex Lie algebra $`\text{g}^{\mathrm{}}`$ of both semi-simple and nilpotent elements are hyperkähler \[Kr1, Kr2\]. In particular these orbits carry Ricci-flat metrics. Later, using Kronheimer’s result, Swann showed that the hyperkähler structure on the nilpotent orbits is very special \[Sw\]. Such orbits admit an action of $`\mathrm{}^{}`$ with the orbit space being a compact quaternionic Kähler orbifold of positive scalar curvature. Another way of expressing this result is to say that the nilpotent orbits are metric cones $`C(𝒮)`$ on compact 3-Sasakian orbifolds \[BGM1\]. The mimimal nilpotent orbit is easily seen to be a metric cone $`C(G/K)`$, where $`G/K`$ is a simply-connected 3-Sasakian homogeneous space of \[BGM1\]. Quaternionic geometry of the regular maximal nilpotent orbit of $`\text{s}\text{l}(3,\mathrm{})`$ was investigated by Kobak and Swann in great detail \[KS1\]. This orbit is 12-dimensional and, in the language of 3-Sasakian geometry, it is a cone $`𝒩=C(𝒮)`$ on the 3-Sasakian orbifold $`𝒮=\mathrm{}_3\backslash G_2/Sp(1)`$. Here $`𝒮`$ is simply an $`\mathrm{}_3`$ quotient of the 3-Sasakian homogeneous space associated to the exceptional Lie group $`G_2`$, where $`\mathrm{}_3`$ is the center of $`SU(3)G_2`$. This is a typical example of a bi-quotient which has orbifold singularities and the singular locus can be easily identified with a homogeneous 3-Sasakian 7-manifold $`SU(3)/U(1)`$. Furthermore, Kobak and Swann \[KS1\] show that this particular nilpotent orbit can be obtained as a hyperkähler quotient of another nilpotent orbit by a circle action. In the language of 3-Sasakian geometry it is simply a hyperkähler reduction of the metric cone $`C\left({\scriptscriptstyle \frac{SO(7)}{SO(3)\times Sp(1)}}\right)`$ associated to the action of the diagonal $`U(1)U(3)SO(7)`$. As $`C\left({\scriptscriptstyle \frac{SO(7)}{SO(3)\times Sp(1)}}\right)`$ can be realized as an $`Sp(1)`$ reduction of the flat space $`\mathrm{}^7=C(S^{27})`$ \[G\] we have the following hyperkähler (3-Sasakian \[BGM1\], quaternionic Kähler \[GL\]) quotients: $$\begin{array}{ccccccccc}C(S^{27})& & \stackrel{Sp\left(1\right)}{}& & C(\frac{SO(7)}{SO(3)\times Sp(1)})& & \stackrel{U\left(1\right)}{}& & C(\mathrm{}_3\backslash G_2/Sp(1))\\ & & & & & & & & \\ & & & & & & & & \\ & & & & & & & & \\ S^{27}& & \stackrel{Sp\left(1\right)}{}& & \frac{SO(7)}{SO(3)\times Sp(1)}& & \stackrel{U\left(1\right)}{}& & \mathrm{}_3\backslash G_2/Sp(1)\\ & & & & & & & & \\ & & & & & & & & \\ & & & & & & & & \\ \mathrm{}\mathrm{}^6& & \stackrel{Sp\left(1\right)}{}& & \frac{SO(7)}{SO(4)\times SO(3)}& & \stackrel{U\left(1\right)}{}& & \mathrm{}_3\backslash G_2/SO(4).\end{array}$$ $`(0.1)`$ In a later paper Kobak and Swann show that any classical nilpotent orbit can be obtained as a hyperkähler quotient of a flat space of an appropriate dimension so that the above diagram is only an example \[KS2\]. The middle horizontal line of Diagram 0.1 is absent from the discussion in \[KS1\]as the importance of 3-Sasakian geometry in this context was realized later \[BGM1\]. The starting point of this paper is to revisit the Kobak-Swann construction in the context of the associated 3-Sasakian geometry. We are going to examine the quotient construction of the orbifold fibration $`\mathrm{}_3\backslash G_2/Sp(1)\mathrm{}_3\backslash G_2/SO(4)`$ showing that it admits interesting generalizations. More precisely, the Kobak-Swann quotient can be “deformed” by introducing weights in much the same way toric 3-Sasakian manifolds with $`b_2=1`$ can all be obtained by “deforming” the classical homogeneous example of the fibration $`𝒮(\mathrm{𝟏})=SU(n)/S(U(n2)\times U(1))\mathrm{Gr}_2(\mathrm{}^n)`$ \[BGM1\]. On the other hand our new quotients are quite different from the construction of $`𝒮(𝐩)`$ considered in \[BGM1\] as they cannot be interpreted as bi-quotients. In Section 2 we prove Theorem A: Let $`𝐩=(p_1,p_2,p_3)\mathrm{}^3`$. For any such non-zero $`𝐩`$ one can define an isometric action of $`Sp(1)\times U(1)_𝐩Sp(7)`$ on $`S^{27}`$ with the following property: If $`0<p_1<p_2<p_3`$ are pairwise relatively prime and $`\mathrm{gcd}(p_1\pm p_2,p_1\pm p_3)=1`$ then the reductions $`(𝐩)𝒵(𝐩)𝒪(𝐩)`$ of $`S^{27}\mathrm{}\mathrm{}^{13}\mathrm{}\mathrm{}^6`$ give a compact smooth 11-dimensional 3-Sasakian manifold $`(𝐩)`$ together with the (orbifold) leaf spaces of its fundamental foliations $`𝒵(𝐩)`$ and $`𝒪(𝐩)`$. Furthermore, the reduced space $`(1,1,1)`$ is a 3-Sasakian orbifold $`\mathrm{}_3\backslash G_2/Sp(1)`$. Analysis of the symmetry structure of all the quotients together with the associated foliations gives Theorem B: The manifold $`(𝐩)`$ of Theorem A is not toric. The corresponding leaf spaces $`𝒵(𝐩)`$ and $`𝒪(𝐩)`$ are compact Riemannian orbifolds with inhomogeneous Einstein metrics of positive scalar curvature. In Section 3 we investigate whether $`(p_1,p_2,p_3)`$ admits further reduction by an isometric circle action. More generally we consider $`Sp(1)\times S^1\times S^1`$ actions on the 27-sphere and ask if one can get any smooth quotients. Surprisingly, no smooth examples can be found but orbifold quotients exist in profusion. These are interesting, since, due to Theorem B, they are necessarily non-toric (more precisely of cohomogeneity 3). More importantly, they yield new explicit self-dual Einstein metrics of positive scalar curvature with only orbifold singularities and with one-dimensional group of isometries. We get Theorem C: Let $`\mathrm{\Theta }_{2\times 3}(\mathrm{})`$ be any integral $`2\times 3`$ matrix such that each of its three $`2\times 2`$ minor determinants does not vanish. In addition suppose that the sum of the all minor determinants is nonvanishing, and none of them is equal to the some of the other two. For any such $`\mathrm{\Theta }`$ there exists a compact 4-dimensional orbifold $`𝒪(\mathrm{\Theta })`$ which admits a self-dual Einstein metric of positive scalar curvature with a one-dimensional group of isometries. Moreover, this metric can be constructed explicitly as a quaternionic Kähler reduction of the real Grassmannian $`Gr_4(\mathrm{}^7)`$ by an isometric action of the 2-torus $`T_\mathrm{\Theta }^2`$ defined by $`\mathrm{\Theta }.`$ The first examples of positive self-dual Einstein metrics on orbifolds were obtained in \[GL\]. Later such metrics were considered by Hitchin \[Hi1, Hi2\]. Hitchin’s examples have large group of isometries. More recently many new orbifold metrics with $`T^2`$-symmetry group were constructed in \[BGMR\]. The examples presented here are perhaps the first self-dual Einstein metrics with only one Killing vector field. We are not aware of any self-dual Einstein metrics which have only discrete isometries. In Section 4 we consider the obvious higher-dimensional extension of the problem. Not surprisingly, once again the new examples involve hyperkähler geometry of a nilpotent variety. This time it is the 16-dimensional nilpotent orbit of $`\text{s}\text{o}(6,\mathrm{})`$ which is the Swann bundle over the quaternionic Kähler orbifold $`Gr_4(\mathrm{}^7)/\mathrm{}_2`$. In \[KS2\] it is shown how all classical nilpotent orbits can be obtained as hyperkähler reductions from flat spaces (typically in more than one way). This orbit can be constructed as a $`Sp(1)\times U(1)`$ reduction of $`\mathrm{}^8`$ \[KS3\] and thus it appears as part of a diagram similar to the one in (0.1). In this context our construction is a systematic study of the general $`U(1)`$ actions which, at the 3-Sasakian level, produce smooth metrics. We prove the following analogue of the Theorem A: Theorem D: Let $`𝐩=(p_1,p_2,p_3,p_4)\mathrm{}^4`$. For any such non-zero $`𝐩`$ one can define an action of $`Sp(1)\times U(1)_𝐩Sp(8)`$ on $`S^{31}`$ with the following property: If $`0p_1<p_2<p_3<p_4`$, any triple $`p_i<p_j<p_k`$ satisfies $`\mathrm{gcd}(p_i,p_j,p_k)=1`$, and $`\mathrm{gcd}(p_i\pm p_j,p_i\pm p_k)=1`$ then the the reductions $`(𝐩)𝒵(𝐩)𝒪(𝐩)`$ of $`S^{31}\mathrm{}\mathrm{}^{15}\mathrm{}\mathrm{}^7`$ give a compact smooth 15-dimensional 3-Sasakian manifold $`(𝐩)`$ together with the (orbifold) leaf spaces of its fundamental foliations $`𝒵(𝐩)`$ and $`𝒪(𝐩)`$. Furthermore, the reduced space $`(1,1,1,1)`$ is a 3-Sasakian orbifold $`\mathrm{}_2\backslash \mathrm{Spin}(7)/\mathrm{Spin}(4)`$. These 15-dimensional quotients are “deformations” of the standard homogeneous 3-Sasakian structure on $`SO(7)/SO(3)\times Sp(1)`$ which projects to the Wolf space $`\mathrm{Gr}_4(\mathrm{}^7)`$ in the quaternionic Kähler base. In higher dimensions the orbifold $`\mathrm{Gr}_4(\mathrm{}^{n+3})/\mathrm{}_2`$ can equally be obtained as a quaternionic Kähler reduction of $`\mathrm{}\mathrm{}^{n+4}`$ by $`U(1)\times Sp(1)`$. Our construction shows that the $`n=4`$ case, from the standpoint of 3-Sasakian geometry, is somewhat exceptional: In higher dimension we can only get orbifold metrics. The orbifold bundles $`^{15}(𝐩)𝒵^{15}(𝐩)𝒪^{15}(𝐩)`$ can be viewed as singular analogues of $`SO(7)/SO(3)\times Sp(1)SO(7)/SO(3)\times U(2)\mathrm{Gr}_4(\mathrm{}^7).`$ This makes use of the well-known isometry between $`\mathrm{Gr}_4(\mathrm{}^7)`$ and the space $`\mathrm{Spin}(7)/(Sp(1)\times Sp(1)\times Sp(1))/\mathrm{}_2`$ of the Cayley 4-planes in $`\mathrm{}^8`$. Again, standard analysis of the symmetry structure of all the quotients together with the associated foliations gives Theorem B: The manifold $`(𝐩)`$ of Theorem D is not toric. The corresponding leaf spaces $`𝒵(𝐩)`$ and $`𝒪(𝐩)`$ are compact Riemannian orbifolds with inhomogeneous Einstein metrics of positive scalar curvature. In Section 5 we give what topological information is available to us. In particular, we show that as long as $`𝐩`$ satifies the conditions of Theorems A and D (actually this hypothesis can be weakened) the rational cohomology of the corresponding $`(𝐩)`$ is independent of $`𝐩.`$ Finally, in section 6 we briefly mention the construction of hypercomples structures on circle bundles over our new 3-Sasakian manifolds (orbifolds). Acknowledgements: The second named author would like to thank Università di Roma “La Sapienza”, C.N.R, M.P.I-Bonn, and I.H.E.S-Bures sur Yvette for hospitality and support. Parts of this paper were written during his visits there. We would also like to thank Mike Buchner and Andrew Swann for comments and discussion. 1. Quotient construction of 3-Sasakian structure on $`\mathrm{}_3\backslash G_2/Sp(1)`$ There are two homogeneous Sasakian-Einstein geometries that are naturally associated with the exceptional Lie group $`G_2`$ \[BG1, BG2\]. They both come from the classical Lie group isomorphism between $`SO(4)G_2`$ and $`Sp(1)_{}Sp(1)_+G_2`$. The two $`Sp(1)_\pm `$ subgroups are very different. One of them has index 1 in $`G_2`$ and the other one has index 3. Consequently, the quotients are not of the same homotopy type as can be seen from the exact sequence in homotopy for the fibration $$Sp(1)_\pm G_2G_2/Sp(1)_\pm .$$ In particular, the two spaces can be distinguished by their third homotopy groups being trivial in one case and $`\mathrm{}_3`$ in the other. One of these quotients, which we shall denote by $`G_2/Sp(1)_{}`$ is diffeomorphic to the real Stiefel manifold $`V_{7,2}(\mathrm{})=SO(7)/SO(5)`$ of 2-frames in $`\mathrm{}^7`$ \[HL\]. As $`V_{7,2}(\mathrm{})`$ is 4-connected $`\pi _3(G_2/Sp(1)_{})=0`$. The other quotient denoted here by $`G_2/Sp(1)_+`$ is one of the 11-dimensional 3-Sasakian homogeneous spaces and $`\pi _3(G_2/Sp(1)_+)=\mathrm{}_3`$. $`G_2/Sp(1)_+`$ fibers as a circle bundle over a generalized flag $`𝒵=G_2/U(2)_+`$, which in turn is well-known to be the twistor space of the exceptional 8-dimensional Wolf space $`G_2/SO(4)`$. The second homogeneous Sasakian-Einstein manifold is a circle bundle over the complex flag $`G_2/U(2)_{}`$ which can be identified with the complex quadric in the 6-dimensional complex projective space $`\mathrm{}\mathrm{}^6`$ or, equivalently, the real Grassmannian $`\mathrm{Gr}_2(\mathrm{}^7)=SO(7)/SO(2)\times SO(5)`$ of oriented 2-planes in $`\mathrm{}^7`$. Both Sasakian-Einstein metrics are well-known and have been studied in the context of homogeneous Einstein geometries \[BG1\]. We have the following diagram of Riemannian submersions: $$\begin{array}{ccccccc}& & & G_2& & & \\ & & & & & & \\ & & & & & & \\ & & & & & & \\ & \frac{G_2}{Sp(1)_+}& & & & \frac{G_2}{Sp(1)_{}}V_{7,2}(\mathrm{})& \\ & & & & & & \\ & & & & & & \\ & & & & & & \\ & 𝒵=\frac{G_2}{U(2)_+}& & & & \frac{G_2}{U(2)_{}}\mathrm{Gr}_2(\mathrm{}^7)& \\ & & & & & & \\ & & & & & & \\ & & & & & & \\ & & & \frac{G_2}{SO(4)}& & & \end{array}$$ $`1.1`$ Poon and Salamon \[PS\] proved that $`G_2/SO(4)`$ is one of the three possible models of positive quaternionic Kähler manifolds in dimension 8. Later the geometry of $`G_2/SO(4)`$ was examined by Kobak and Swann \[KS1\] who proved the following remarkable theorem: Theorem 1.2: The quaternionic Kähler manifold $`\mathrm{Gr}_4(\mathrm{}^7)`$ admits an action of $`U(1)`$ such that the quaternionic Kähler quotient is a compact quaternionic Kähler orbifold $`𝒪=𝒪_r\mathrm{}\mathrm{}(2)=G_2/(SO(4)\times \mathrm{}_3)`$. We will first re-examine the Kobak-Swann construction from the point of view of the 3-Sasakian geometry of a certain $`SO(3)`$ V-bundle over $`𝒪`$ (or, equivalently, the hyperkähler geometry of the regular nilpotent orbit of $`\text{s}\text{l}(3,\mathrm{})`$). In particular, we have the following: Theorem 1.3: The 3-Sasakian homogeneous manifold $`SO(7)/SO(3)\times Sp(1)`$ admits an action of $`U(1)`$ such that the 3-Sasakian quotient is a compact 3-Sasakian orbifold $`=_rSU(3)/U(1)=\mathrm{}_3\backslash G_2/Sp(1)_+`$. Theorem 1.3 is a straightforward translation of Theorem 1.2 into the language of 3-Sasakian geometry and we could leave it at that. However, we will outline a constructive proof of this result as our description of the corresponding quotient differs slightly from the one given in \[KS1\]. One can think of the homogeneous 3-Sasakian manifold $`SO(7)/SO(3)\times Sp(1)`$ as the 3-Sasakian reduction $`S^{4n1}///Sp(1)`$ as follows \[G, BGM1\]: Let $`𝐮=(u_1,\mathrm{},u_7)S^{27}`$. Consider the $`Sp(1)`$ action given by multiplication by unit quaternion $`\lambda Sp(1)`$ on the left that is $$\phi _\lambda (𝐮)=\lambda 𝐮.$$ $`1.4`$ In the $`\{i,j,k\}`$ basis the 3-Sasakian moment maps for this action read: $$\mu _i(𝐮)=\underset{\alpha =1}{\overset{7}{}}\overline{u}_\alpha iu_\alpha ,\mu _j(𝐮)=\underset{\alpha =1}{\overset{7}{}}\overline{u}_\alpha ju_\alpha ,\mu _k(𝐮)=\underset{\alpha =1}{\overset{7}{}}\overline{u}_\alpha ku_\alpha .$$ $`1.5`$ Then, the common zero-locus of the moment maps $$N=\{𝐮S^{4n1}:\mu _i(𝐮)=\mu _j(𝐮)=\mu _k(𝐮)=0\}$$ $`1.6`$ is the Stiefel manifold $`NSO(7)/SO(3)=V_{7,4}(\mathrm{})`$ of the orthonormal 4-frames in $`\mathrm{}^7`$ and the corresponding 3-Sasakian quotient $`𝒮=N/Sp(1)`$ is Konishi’s $`\mathrm{}\mathrm{}^3`$-bundle over the real Grassmannian of oriented 4-planes in $`\mathrm{}^7`$. We can combine Theorem 1.3 with this description to get Corollary 1.7: The 3-Sasakian sphere $`S^{27}`$ admits an action of $`U(1)\times Sp(1)`$ such that the 3-Sasakian quotient is a compact 3-Sasakian orbifold $`=_rSU(3)/U(1)=\mathrm{}_3\backslash G_2/Sp(1)_+`$. We now turn to the explicit description of the $`U(1)`$ quotient. Consider the following subgroups of the group of 3-Sasakian isometries of the 27-sphere: $$Sp(7)SO(7)1\times SO(6)1\times U(3),$$ $`1.8`$ where $`U(1)U(3)`$ is the central subgroup. Explicitly, we shall write $`f:[0,2\pi )SO(7)`$ $$f(t)=\left(\begin{array}{cccc}1& 0& 0& 0\\ 0& A(t)& 0& 0\\ 0& 0& A(t)& 0\\ 0& 0& 0& A(t)\end{array}\right)SO(7),$$ $`1.9`$ where $$A(t)=\left(\begin{array}{cc}\mathrm{cos}t& \mathrm{sin}t\\ \mathrm{sin}t& \mathrm{cos}t\end{array}\right)$$ $`1.10`$ are the real rotations in $`\mathrm{}^2`$. The homomorphism $`f(t)`$ yields a circle action on $`S^{27}`$ or, equivalently, (after performing the “$`Sp(1)`$-reduction” first) on the homogeneous 3-Sasakian manifold $`SO(7)/SO(3)\times Sp(1)`$ via left multiplication $`f(t)𝐮`$ and the associated 3-Sasakian moment map can be written as $$\nu (𝐮)=\underset{\alpha =1,2,3}{}(\overline{u}_{2\alpha }u_{2\alpha +1}\overline{u}_{2\alpha +1}u_{2\alpha }).$$ $`1.11`$ Note here that $`\nu (𝐮)`$ does not depend on the $`u_1`$ quaternionic coordinate. Definition 1.12: Let us define the zero level set of this new moment map intersected with $`N`$, that is $$N_\nu N\nu ^1(0)$$ $`1.12`$ First we observe, following Kobak and Swann \[KS1\] that Lemma 1.13: The manifold $`N_\nu `$ can be identified with $`U(1)G_2=(S^1\times G_2)/\mathrm{}_3`$ where $`U(1)G_2=\mathrm{}_3.`$ Proof: The argument is similar here to the one used by Kobak and Swann in \[KS1\] and it is based on the Proposition 1.10 of \[HL\]. First, using the basis $`\{i,j,k\}`$ of unit imaginary quaternions, we write $`u_\alpha =u_\alpha ^0+iu_\alpha ^1+ju_\alpha ^2+ku_\alpha ^3`$ and introduce the $`4\times 7`$ real matrix $$\mathrm{𝔸}=\left(\begin{array}{ccccccc}u_1^0& u_2^0& u_3^0& u_4^0& u_5^0& u_6^0& u_7^0\\ u_1^1& u_2^1& u_3^1& u_4^1& u_5^1& u_6^1& u_7^1\\ u_1^2& u_2^2& u_3^2& u_4^2& u_5^2& u_6^2& u_7^2\\ u_1^3& u_2^3& u_3^3& u_4^3& u_5^3& u_6^3& u_7^3\end{array}\right)\left(\begin{array}{c}f^0\\ f^1\\ f^2\\ f^3\end{array}\right),$$ $`1.14`$ where to make the connection with the notation in \[KS1\] we also think of the rows of $`\mathrm{𝔸}`$ as purely imaginary octonions $`\mathrm{Im}(\mathrm{𝕆}).`$ In the standard basis of $`\mathrm{Im}(\mathrm{𝕆})`$ we write $`f^a=u_1^ai+u_2^aj+u_3^ak+u_4^ae+u_5^aie+u_6^aje+u_7^ake.`$ Let $`\varphi (a,b,c)=<ab,c>`$ denote the 3-form defining the associative calibration \[HL\] on $`\mathrm{Im}(\mathrm{𝕆})`$ where $`<,>`$ denotes the standard Euclidean inner product. Then writing $`\nu =\nu _1i+\nu _2j+\nu _3k`$ a straightforward computation shows that for $`a=1,2,3`$ and $`ϵ^{abc}`$ the totally antisymmetric tensor satisfying $`ϵ^{123}=1`$ $$\nu _a=2<f^0f^a+ϵ^{abc}f^bf^c,i>,$$ $`1.15`$ where the summation convention on repeated indices is used. Now, $`𝐮N`$ if and only if the rows $`\{f^0,f^1,f^2,f^3\}`$ of $`\mathrm{𝔸}`$ form an orthonormal frame in $`\mathrm{}^7\mathrm{Im}(\mathrm{𝕆}),`$ and one can identify $`G_2`$ with a special kind of oriented orthonormal 4-frame, namely those which are co-associative. This means that the 3-plane that is orthogonal to the 4-plane defined by the frame $`\{f^0,f^1,f^2,f^3\}`$ is spanned by an associative subalgebra of $`\mathrm{Im}(\mathrm{𝕆}).`$ Then one shows that these special 4-frames satisfy the $`U(1)`$-moment map equation $`\nu (𝐮)=0`$ and, hence, $`U(1)G_2N\nu ^1(0)`$. As, $`U(1)G_2=\mathrm{}_3`$ it is enough to show that by acting with $`U(1)`$ one gets the whole $`N\nu ^1(0)`$. The argument is similar to the one presented in \[KS1\]. (See \[KS1\] Lemma 5.1 and the discussion that follows.) Now, Theorems 1.2 and 1.3 and Corollary 1.7 all follow from the above lemma as we get the quotient $$=\frac{N_\nu }{U(1)\times Sp(1)}\frac{U(1)G_2}{U(1)\times Sp(1)}\mathrm{}_3\backslash G_2/Sp(1).$$ $`1.16`$ remark 1.17: The $`U(1)\times Sp(1)`$ action on the level set $`N_\nu `$ is not locally free. If we divide by $`Sp(1)`$ first and consider the $`U(1)`$ action on the orbit space $`N_\nu /Sp(1)`$ this circle action is quasi-free. This means that there are only two kinds orbits: regular orbits with the trivial isotropy group and singular orbits (points) where the isotropy group is the whole $`U(1)`$. In such cases the quotient space is often an orbifold (or even a smooth manifold). The stratification of the Theorem 2 is precisely with respect to the orbit types as will be seen in the next section. 2. Generalizations of the Kobak-Swann Quotient In this section we will consider the simplest possible family of quotients which generalize the construction of Kobak and Swann via an introduction of weights. Instead of considering the central $`U(1)U(3)`$ in 1.9 we can consider and arbitrary circle subgroup of the maximal torus $`U(1)T^3U(3)`$. Again, to be more specific, we have the following inclusions: $$Sp(7)SO(7)1\times SO(6)1\times SO(2)\times SO(2)\times SO(2).$$ $`2.1`$ We can consider arbitrary circle subgroups of the last 3-torus. Let $`𝐩=(p_1,p_2,p_3)\mathrm{}^3`$ and define the following homomorphism $$f_𝐩(t)=\left(\begin{array}{cccc}1& 0& 0& 0\\ 0& A(p_1t)& 0& 0\\ 0& 0& A(p_2t)& 0\\ 0& 0& 0& A(p_3t)\end{array}\right)SO(7),$$ $`2.2`$ where $$A(p_it)=\left(\begin{array}{cc}\mathrm{cos}(p_it)& \mathrm{sin}(p_it)\\ \mathrm{sin}(p_it)& \mathrm{cos}(p_it)\end{array}\right)SO(2),i=1,2,3$$ $`2.3`$ are 2-dimensional real rotations. Note that for $`p_1=p_2=p_3=1`$ we recover the example of the previous section. The homomorphism $`f_𝐩(t)`$ yields a circle action on the homogeneous 3-Sasakian manifold $`SO(7)/SO(3)\times Sp(1)`$ via left multiplication $`f_𝐩(t)𝐮`$ and the moment map can now be written as $$\nu _𝐩(u_1,\mathrm{},u_7)=\underset{\alpha =1,2,3}{}p_\alpha (\overline{u}_{2\alpha }u_{2\alpha +1}\overline{u}_{2\alpha +1}u_{2\alpha }).$$ $`2.4`$ Observe that without loss of generality we can assume all weights to be non-negative as $`p_i`$ can be changed to $`p_i`$ by renaming the quaternions in the associated pair $`(u_{2i},u_{2i+1})`$. We begin analysis of this quotient by considering the level set of the moment map Definition 2.5: Let $`N_\nu (𝐩)S^{27}`$ be the level set of the 3-Sasakian moment map of the $`Sp(1)\times U(1)_𝐩`$-action, i.e., $`N_\nu (𝐩)N\{\nu _𝐩^1(0)\}.`$ We want to consider a stratification of the level set $`N_\nu (𝐩)`$ that will allow us to analyze the quotient space $$(𝐩)=\frac{N_\nu (𝐩)V_{7,4}(\mathrm{})\{\nu _𝐩^1(0)\}}{Sp(1)\times U(1)_𝐩}.$$ $`2.6`$ Since $`N_\nu (𝐩)`$ is a submanifold of the Stiefel manifold $`V_{7,4}(\mathrm{})`$ at most 3 quaternionic coordinates can vanish on $`N_\nu (𝐩).`$ So setting various quaternionic coordinates equal to zero determines a stratification of $`N_\nu (𝐩)`$ in which the strata of minimal dimension play an important role. We call these strata vertices although, as we shall see, they each have two connected components. Lemma 2.7: Let $`0<p_1<p_2<p_3`$. At a vertex neither $`u_1`$ nor any of the three pairs of quaternions $`(u_{2i},u_{2i+1})`$, $`i=1,2,3`$ can vanish. Thus, there are precisely eight vertices and they are all diffeomorphic to $`O(4).`$ Proof: Every vertex has precisely 3 quaternionic coordinates vanishing, so the Stiefel manifold becomes $`V_{4,4}(\mathrm{})=O(4).`$ Let $`V`$ be a vertex. Then $`V=V_{ijkl}`$ can be represented by a matrix of the form $$\mathrm{𝔹}=\left(\begin{array}{cccc}u_i^0& u_i^1& u_i^2& u_i^3\\ u_j^0& u_j^1& u_j^2& u_j^3\\ u_k^0& u_k^1& u_k^2& u_k^3\\ u_l^0& u_l^1& u_l^2& u_l^3\end{array}\right),\mathrm{𝔹}^t\mathrm{𝔹}=\mathrm{𝔹}\mathrm{𝔹}^t=\frac{1}{4}\mathrm{𝕀}_4,$$ $`2.8`$ where $`1i<j<k<l7`$. Here we have written a quaternionic coordinate as $`u_i=u_i^0+iu_i^1+ju_i^2+ku_i^3.`$ Suppose $`i>1`$, i.e., $`u_1=0`$ on on $`V_{ijkl}`$. Then there are two possibilities: (1) the quadruple $`i<j<k<l`$ contains only one quaternionic pair $`(2\alpha ,2\alpha +1)`$ or (2) it contains two such pairs. Let us examine the first possibility. Without loss of generality we can take $`i=2,j=3,k=4,l=6`$. With this choice the vanishing of the $`U(1)`$ moment map 2.4 now becomes $$\mathrm{Im}(\overline{u}_2u_3)=0.$$ One can easily check that orthogonality of the basis forces $$\mathrm{Re}(\overline{u}_2u_3)=u_2^1u_3^1+u_2^2u_3^2+u_2^3u_3^3+u_2^4u_3^4=0.$$ But this implies $`\overline{u}_2u_3=0`$ which forces either $`u_2`$ or $`u_3`$ to vanish, giving a contradiction. Now assume the second case, that is, that $`i<j<k<l`$ consists of two quaternionic pairs $`(2\alpha ,2\alpha +1)`$. Again, when we take $`i=2,j=3,k=4,l=5`$, the orthogonality of the vectors in the associated $`\mathrm{𝔹}`$-matrix forces $$\begin{array}{cc}\hfill \mathrm{Re}(\overline{u}_2u_3)& =u_2^1u_3^1+u_2^2u_3^2+u_2^3u_3^3+u_2^4u_3^4=0,\hfill \\ \hfill \mathrm{Re}(\overline{u}_4u_5)& =u_4^1u_5^1+u_4^2u_5^2+u_4^3u_5^3+u_4^4u_5^4=0.\hfill \end{array}$$ $`2.9`$ But orthogonality also implies $$\|u_2\|^2=\|u_3\|^2=\|u_4\|^2=\|u_5\|^2=\frac{1}{4}.$$ Then we have $$\|\mathrm{Im}(\overline{u}_2u_3)\|^2=\|\mathrm{Im}(\overline{u}_4u_5)\|^2=\frac{1}{16},$$ and this contradicts the $`U(1)`$ moment map equation $$2p_1\mathrm{Im}(\overline{u}_2u_3)+2p_2\mathrm{Im}(\overline{u}_4u_5)=0$$ if $`p_1p_2`$. Repeating the argument for the other choices of $`1<i<j<k<l`$ gives the result under the hypothesis $`p_1<p_2<p_3.`$ Thus, at a vertex we must have $`u_10.`$ This proves the “neither” part of the statement. Now assume that $`i=1<j<k<l`$ and suppose either $`j<k`$ or $`k<l`$ is a quaternionic pair. We can take $`j=2,k=3`$ and let $`l`$ be arbitrary. Then the orthogonality of the corresponding $`\mathrm{𝔹}`$-matrix again forces $$\mathrm{Re}(\overline{u}_2u_3)=0.$$ But again the $`U(1)`$ moment map constraint is $$\mathrm{Im}(\overline{u}_2u_3)=0$$ giving a contradiction. This proves the “nor” part of the lemma. It is now clear that the vertices must automatically satisfy the $`U(1)`$ moment map constraint $`\nu _𝐩(𝐮)=0;`$ hence, they are all diffeomorphic to $`O(4).`$ Moreover, a simple counting shows that there are precisely eight vertices. Our analysis suggests the importance of the following strata: $$\begin{array}{cc}\hfill S_0& =\{𝐮N_\nu (𝐩)\|u_1=0\},\hfill \\ \hfill S_1& =\{𝐮N_\nu (𝐩)\|u_10\},\hfill \\ \hfill S_2& =\{𝐮N_\nu (𝐩)\|\text{some quaternionic pair }(u_{2i},u_{2i+1})\text{ vanishes}\},\hfill \\ \hfill S_3& =\{𝐮N_\nu (𝐩)\|\text{no quaternionic pair }(u_{2i},u_{2i+1})\text{ vanishes}\}.\hfill \end{array}$$ $`2.10`$ Then Lemma 2.7 easily implies that Corollary 2.11: Let $`0<p_1<p_2<p_3`$. Then (i) $`S_0S_1=S_2S_3=N_\nu (𝐩).`$ (ii) $`S_0S_1=S_2S_3=\mathrm{}.`$ (iii) $`S_0S_2=\mathrm{}.`$ (iv) $`S_2S_1.`$ (v) $`S_0S_3.`$ Notice that (iii) fails if $`p_i=p_j`$ for some $`ij.`$ In particular it fails for the level set $`N_\nu `$ of 1.12 in the previous section, and this is the reason that the quotient $``$ of 1.16 is not smooth. We now are ready to give necessary conditions to guarantee a smooth quotient. Lemma 2.12: Let $`𝐩=(p_1,p_2,p_3)(\mathrm{}_+)^3`$ be pairwise relatively prime. Then the isotropy group of the $`Sp(1)\times U(1)_𝐩`$ action at every point of $`S_1`$ is the identity. Proof: The action of $`Sp(1)\times U(1)_𝐩`$ on $`\mathrm{}^7`$ is the diagonal action of $`Sp(1)`$ by quaternionic multiplication by a unit quaternion $`\lambda `$ on the left, and the matrix multiplication $`𝐮f_𝐩(t)𝐮`$ for the $`U(1)_𝐩`$ action. These two actions clearly commute. Since $`u_10`$ we immediately get that $`\lambda =1`$. Consider the set where a quaternionic pair $`(u_6,u_7)=(0,0)`$. Then the fixed point equation becomes $$A(p_1t)=A(p_2t)=\mathrm{𝕀}_2,$$ $`2.13`$ which has only the trivial solution provided that $`\mathrm{gcd}(p_1,p_2)=1`$. Setting the other two quaternionic pairs to be zero gives $`\mathrm{gcd}(p_1,p_3)=\mathrm{gcd}(p_2,p_3)=1`$. As one cannot set more than one quaternionic pair equal to $`(0,0)`$ the lemma is proved. Lemma 2.14: Let $`𝐩=(p_1,p_2,p_3)(\mathrm{}_+)^3`$ satisfy the four conditions $`\mathrm{gcd}(p_1\pm p_2,p_1\pm p_3)=1.`$ Then the isotropy group of the $`Sp(1)\times U(1)_𝐩`$ action at every point of $`S_0`$ is the identity. Proof: To determine the conditions for fixed points of the action we consider the following equations $$A(p_it)\left(\begin{array}{c}u_{2i}\\ u_{2i+1}\end{array}\right)=\left(\begin{array}{cc}a_i& b_i\\ b_i& a_i\end{array}\right)\left(\begin{array}{c}u_{2i}\\ u_{2i+1}\end{array}\right)=\lambda \left(\begin{array}{c}u_{2i}\\ u_{2i+1},\end{array}\right)i=1,2,3$$ for $`\lambda Sp(1)`$ and $`t[0,2\pi ).`$ For each $`i=1,2,3`$ this reads $$a_iu_{2i}+b_iu_{2i+1}=\lambda u_{2i},$$ $$b_iu_{2i}+a_iu_{2i+1}=\lambda u_{2i+1}.$$ For each $`i=1,2,3`$ we multiply the first equation from the right by $`\overline{u}_{2i}`$ and the second by $`\overline{u}_{2i+1}`$ to get $$a_i\|u_{2i}\|^2+b_iu_{2i+1}\overline{u}_{2i}=\lambda \|u_{2i}\|^2,$$ $$b_iu_{2i}\overline{u}_{2i+1}+a_i\|u_{2i+1}\|^2=\lambda \|u_{2i+1}\|^2.$$ $`2.15`$ By adding these two equations we get $$a_i(\|u_{2i}\|^2+\|u_{2i+1}\|^2)+b_i(u_{2i+1}\overline{u}_{2i}u_{2i}\overline{u}_{2i+1})=\lambda (\|u_{2i}\|^2+\|u_{2i+1}\|^2),i=1,2,3.$$ By (iii) of Lemma 2.10 the term multiplying $`\lambda `$ on the right hand side of this equation never vanishes. This gives for each $`i=1,2,3`$ $$\begin{array}{cc}\hfill \mathrm{Re}(\lambda )& =a_i,\hfill \\ \hfill \mathrm{Im}(\lambda )& =b_i\frac{u_{2i+1}\overline{u}_{2i}u_{2i}\overline{u}_{2i+1}}{\|u_{2i}\|^2+\|u_{2i+1}\|^2}.\hfill \end{array}$$ $`2.16`$ The first of these equations gives $$a_1=a_2=a_3,$$ $`2.17`$ and combining this with $`a_i^2+b_i^2=1`$ implies $$b_1=\pm b_2=\pm b_3.$$ $`2.18`$ Let us write $`\tau =e^{it}`$. Then 2.17 and 2.18 give $$\tau ^{p_1}=\tau ^{\pm p_2},\tau ^{p_1}=\tau ^{\pm p_3}.$$ $`2.19`$ These have only trivial solutions if and only if $`\mathrm{gcd}(p_1\pm p_2,p_1\pm p_3)=1.`$ It is convenient to make the following: Definition 2.20: Let $`𝐩=(p_1,p_2,p_3)\mathrm{}^3`$. We say that the weight vector $`𝐩`$ is admissible if $`0<p_1<p_2<p_3`$, $`\mathrm{gcd}(p_i,p_j)=1`$ for all $`i<j`$, and $`\mathrm{gcd}(p_1\pm p_2,p_1\pm p_3)=1.`$ It now follows immediately from Lemmas 2.12, 2.14 and Definition 2.20 that Theorem 2.21: The $`Sp(1)\times U(1)_𝐩`$ action on $`N_\nu (𝐩)`$ is free if and only if $`𝐩\mathrm{}_+^3`$ is admissible. Note that there are infinitely many admissible weight vectors. For example we can take $`𝐩=(2k1,2k,2k+1)`$, where $`k\mathrm{}_+`$. Thus, there are infinite families of smooth quotients $`(𝐩)`$ and infinite families of the associated triples $`(𝐩)𝒵(𝐩)𝒪(𝐩)`$ with their (orbifold) Einstein metrics. Theorem A now follows from the Theorem 2.21 and various theorems concerning 3-Sasakian (complex contact, quaternionic Kähler) reductions \[BGM1, GL\]. The last statement of Theorem A follows from Corollary 1.7. We briefly return to the $`𝐩=(1,1,1)`$ case of the previous section. As already observed we get singularities here because the corresponding action is not even locally free on the level set of the moment map. But it is easy to see that it is quasi-free. That is there are only two types of isotropy subgroups in the circle: the identity and the whole group. In such a situation Dancer and Swann \[DS\] observed that 3-Sasakian quotients, stratify as the union of 3-Sasakian manifolds \[DS\]. This is true in this case in particular, but as we saw in the previous section the two strata nicely fit together and one gets a compact 3-Sasakian orbifold. Let us explicitly describe the singular part $`_1(1,1,1).`$ Note that if $`p_1=p_2=p_3=1`$ then $`a_1=a_2=a_3=a`$ and $`b_1=b_2=b_3=b`$ and we can add equations (2.15) to get $$a+b\rho =\lambda ,$$ $`2.22`$ where now $$\rho =\underset{i=1,2,3}{}(u_{2i+1}\overline{u}_{2i}u_{2i}\overline{u}_{2i+1}).$$ $`2.23`$ When $`u_10`$ one does not get any fixed points of the action. But when $`u_1=0`$ for any imaginary unit $`\rho `$ there is a $`U(1)Sp(1)\times U(1)`$ subgroup $$(\mathrm{cos}t+\rho \mathrm{sin}t,A(t))Sp(1)\times U(1),$$ which acts trivially on the following set $$\{𝐮N_\nu \|u_{2i+1}=\rho u_{2i}\}.$$ $`2.24`$ In this case all 4 moment map equations reduce to the same one and it reads $$\underset{i=1,2,3}{}\overline{u}_{2i}\rho u_{2i}=0.$$ For any fixed $`\rho `$ we can recognize this set as the complex Stiefel manifold $`U(3)/U(1)`$ and it follows that the singular stratum $`_1`$ is precisely the quotient $`SU(3)/U(1)=𝒮(1,1,1).`$ The geometry of the smooth families $`(𝐩)`$ is rather interesting. First we observe that these spaces cannot be toric. This can be seen in several different ways, for example by careful analysis of the associated foliations. One can also generalize the analysis of \[BGM2\] to show that the only isometries of the level set of the moment map $`N_\nu (𝐩)S^{27}\mathrm{}^{28}`$ can come from the restriction of the isometries of the Euclidean space $`\mathrm{}^{28}`$. From this we conclude Theorem 2.25: Let $`𝐩`$ be admissible so that $`(𝐩)`$ is a smooth compact 3-Sasakian 11-manifold. Then the Lie algebra $`\mathrm{Isom}^0((𝐩),g(𝐩))`$ of the group of 3-Sasakian isometries of $`(𝐩)`$ is isomorphic to $`\mathrm{}^2\text{s}\text{p}(1)`$. In particular, all such quotients are non-toric. This proves part of Theorem C of the introduction which relates to the 11-dimensional quotients. Finally, observe that $`(𝐩)`$ contains a special 7-manifold which is embedded as a 3-Sasakian submanifold. Define $$S_0(𝐩)=\frac{N_\nu (𝐩)\{u_1=0\}}{Sp(1)\times U(1)_𝐩}.$$ $`2.26`$ One can see that $`S_0(𝐩)`$ is a submanifold and, as it is itself a 3-Sasakian reduction, it must be a 3-Sasakian submanifold. One can even identify this space. Observe that the classical group isomorphism $`\mathrm{Spin}(6)SU(4)`$ implies that the $`Sp(1)`$ quotient yields the homogeneous 3-Sasakian 11-dimensional manifold $`𝒮(1,1,1,1)`$. Hence, $`S_0(𝐩)`$ is either a $`U(1)_𝐩`$-reduction of $`𝒮(1,1,1,1)`$ or, equivalently, a $`T^2`$-reduction of $`S^{15}`$. Hence, there exists an admissible integer weight matrix $`\mathrm{\Omega }𝒜_{2\times 4}(\mathrm{})`$ (see \[BGMR\]) such that $`S_0(𝐩)𝒮(\mathrm{\Omega })`$. Hence $`S_0(𝐩)`$ is toric with second Betti number equal to 2. 3. Further Reductions of $`(p_1,p_2,p_3)`$ by a Circle In this section we will examine reductions of $`(p_1,p_2,p_3)`$ by isometric circle actions. More generally we shall consider an arbitrary 2-torus subgroup of the maximal torus $`T^2T^3SO(7)`$. Let $$\mathrm{\Theta }=\left(\begin{array}{ccc}p_1& p_2& p_3\\ q_1& q_2& q_3\end{array}\right)_{2\times 3}(\mathrm{})$$ $`3.1`$ be an arbitrary integral $`2\times 3`$ matrix and define the homomorphism $`f_\mathrm{\Theta }:T^2SO(7)`$ $$f_\mathrm{\Theta }(t,s)=\left(\begin{array}{cccc}1& 0& 0& 0\\ 0& A(p_1t+q_1s)& 0& 0\\ 0& 0& A(p_2t+q_2s)& 0\\ 0& 0& 0& A(p_3t+q_3s)\end{array}\right)SO(7),$$ $`3.2`$ where $`A`$ is the $`SO(2)`$ rotations as in 2.3. Note that if $`𝐩=(p_1,p_2,p_3)`$ were admissible then we would be considering arbitrary isometric circle actions on the quotient $`(𝐩)`$; however, so far we assume nothing about $`\mathrm{\Theta }`$. The homomorphism $`f_\mathrm{\Theta }(t,s)`$ yields a 2-torus action on the homogeneous 3-Sasakian manifold $`SO(7)/SO(3)\times Sp(1)`$ via left multiplication $`f_\mathrm{\Theta }(t,s)𝐮`$ and the moment map can now be written as $$\nu _\mathrm{\Theta }(u_1,\mathrm{},u_7)=\left(\begin{array}{c}_{\alpha =1,2,3}p_\alpha (\overline{u}_{2\alpha }u_{2\alpha +1}\overline{u}_{2\alpha +1}u_{2\alpha })\\ _{\alpha =1,2,3}q_\alpha (\overline{u}_{2\alpha }u_{2\alpha +1}\overline{u}_{2\alpha +1}u_{2\alpha })\end{array}\right)\mathrm{}^2\text{s}\text{p}(1).$$ $`3.3`$ We begin our analysis of this quotient by considering the level set of the moment map. Definition 3.4: Let $`N_\nu (\mathrm{\Theta })S^{27}`$ denote the level set of the 3-Sasakian moment map of the $`Sp(1)\times T_\mathrm{\Theta }^2`$-action, i.e., $`N_\nu (\mathrm{\Theta })N\{\nu _\mathrm{\Theta }^1(\mathrm{𝟎})\}`$ and let $$(\mathrm{\Theta })=N_\nu (\mathrm{\Theta })/Sp(1)\times T_\mathrm{\Theta }^2.$$ $`3.5`$ We want to determine for which $`\mathrm{\Theta }_{2\times 3}(\mathrm{})`$ the 7-dimensional quotient $`(\mathrm{\Theta })`$ is an orbifold and, if possible which weight matrices yield smooth quotients. Let us define $$\mathrm{\Delta }_{ij}=\mathrm{\Delta }_{ij}(\mathrm{\Theta })=\mathrm{det}\left(\begin{array}{cc}p_i& p_j\\ q_i& q_j\end{array}\right),1i<j3,$$ $`3.6`$ the three minor determinants of $`\mathrm{\Theta }`$. Lemma 3.7: The action of $`Sp(1)\times T_\mathrm{\Theta }^2`$ on $`N_\nu (\mathrm{\Theta })\{u_10\}`$ is (i) locally free if and only if $`\mathrm{\Delta }_{ij}(\mathrm{\Theta })0,1i<j3`$, (ii) free if and only if $`\|\mathrm{\Delta }_{ij}(\mathrm{\Theta })\|=1,1i<j3`$. Proof: Since $`u_10`$ we must have $`\lambda =1`$ and, hence, it is enough to consider the $`T_\mathrm{\Theta }^2`$-action. Now, suppose that a quaternionic pair, say $`(u_6,u_7)`$ vanishes. Then the fixed point equation reads: $$A(p_it+q_is)\left(\begin{array}{c}u_{2i}\\ u_{2i+1}\end{array}\right)=\left(\begin{array}{c}u_{2i}\\ u_{2i+1}\end{array}\right),i=1,2$$ $`3.8`$ for $`t,s[0,2\pi ),`$ or, equivalently, $$A(p_it+q_is)=\mathrm{𝕀}_2,i=1,2.$$ $`3.9`$ Let $`\tau =e^{it}`$ and $`\rho =e^{is}`$. Then we can rewrite 3.9 as $$\tau ^{p_i}\rho ^{q_i}=1,i=1,2.$$ $`3.10`$ This has only discrete solutions provided $`\mathrm{\Delta }_{12}(\mathrm{\Theta })0`$. Furthermore, the isotropy group at all such points will be trivial provided $`\mathrm{\Delta }_{12}(\mathrm{\Theta })=\pm 1`$. This proves the lemma. Note that the second condition is already very restrictive as it says that any $`2\times 2`$ submatrix of $`\mathrm{\Theta }`$ must be an element of $`PSL(2,\mathrm{})`$. However, there are many matrices which satisfy both conditions, for example $$\mathrm{\Theta }_1=\left(\begin{array}{ccc}1& 0& 1\\ 0& 1& 1\end{array}\right),\mathrm{\Theta }_2=\left(\begin{array}{ccc}9& 2& 7\\ 40& 9& 31\end{array}\right).$$ It remains to analyze the fixed point equations on $`N_\nu (\mathrm{\Theta })\{u_1=0\}`$. We now prove Lemma 3.11: The action of $`Sp(1)\times T_\mathrm{\Theta }^2`$ on $`N_\nu (\mathrm{\Theta })`$ is locally free if and only if $`\mathrm{\Delta }_{ij}(\mathrm{\Theta })0,1i<j3`$ and $$\mathrm{}_{}^{}=\mathrm{det}\left(\begin{array}{cc}p_1p_2& q_1q_2\\ p_1p_3& q_1q_3\end{array}\right)0$$ $`3.12`$ Furthermore, there is no weight matrix $`\mathrm{\Theta }_{2\times 3}(\mathrm{})`$ for which the actions is free. Proof: First let us clarify that what we mean in 3.12 is that four determinants $`\mathrm{}_+^+`$, $`\mathrm{}_+^{}`$, $`\mathrm{}_{}^+`$, and $`\mathrm{}_{}^{}`$ must vanish (in any row we can choose either upper or lower signs). Since the $`\mathrm{\Delta }_{ij}(\mathrm{\Theta })0`$, by the previous Lemma, we know that the action is locally free on the $`u_10`$ part. Hence, it is enough to consider $`u_1=0`$. Here the analysis is similar to the one presented in Lemma 2.14 and it is entirely based on the fact that no quaternionic pair can vanish. As a result we get the following analogue of the fixed point equations 2.19: $$\tau ^{p_1}\rho ^{q_1}=\left(\tau ^{p_2}\rho ^{q_2}\right)^{\pm 1},\tau ^{p_1}\rho ^{q_1}=\left(\tau ^{p_3}\rho ^{q_3}\right)^{\pm 1}.$$ $`3.13`$ We can rewrite these as $$\tau ^{p_1p_2}=\rho ^{q_1\pm q_2},\tau ^{p_1p_3}=\rho ^{q_1\pm q_3}.$$ $`3.14`$ These are four systems of two equations in $`(\tau ,\rho )`$ variables. We want all four of them to have at most discrete solutions. This requires that the four determinants $$\mathrm{det}\left(\begin{array}{cc}p_1p_2& q_1\pm q_2\\ p_1p_3& q_1\pm q_3\end{array}\right)=\mathrm{det}\left(\begin{array}{cc}p_1p_2& q_1q_2\\ p_1p_3& q_1q_3\end{array}\right)$$ $`3.15`$ do not vanish and gives 3.12. The fact that orbifold singularities are always present requires more subtle analysis. To have smooth quotients we must assume $$\mathrm{\Delta }_{ij}(\mathrm{\Theta })=\pm 1,1i<j3,$$ $`3.16`$ one one hand, and $$(\mathrm{}_{}^{})=\left\|\begin{array}{cc}p_1p_2& q_1q_2\\ p_1p_3& q_1q_3\end{array}\right\|=\pm 1$$ $`3.17`$ on the other. A simple computation relates all of these four determinants to the 3 minor determinants $`\mathrm{\Delta }_{ij}(\mathrm{\Theta })`$ and we get $$\mathrm{}_{}^{}=\mathrm{det}\left(\begin{array}{cc}p_1p_2& q_1q_2\\ p_1p_3& q_1q_3\end{array}\right)=\mathrm{\Delta }_{12}(\mathrm{\Theta })+\mathrm{\Delta }_{23}(\mathrm{\Theta })\mathrm{\Delta }_{13}(\mathrm{\Theta }),$$ $$\mathrm{}_+^{}=\mathrm{det}\left(\begin{array}{cc}p_1p_2& q_1q_2\\ p_1+p_3& q_1+q_3\end{array}\right)=\mathrm{\Delta }_{12}(\mathrm{\Theta })\mathrm{\Delta }_{23}(\mathrm{\Theta })+\mathrm{\Delta }_{13}(\mathrm{\Theta }),$$ $$\mathrm{}_{}^+=\mathrm{det}\left(\begin{array}{cc}p_1+p_2& q_1+q_2\\ p_1p_3& q_1q_3\end{array}\right)=\mathrm{\Delta }_{12}(\mathrm{\Theta })\mathrm{\Delta }_{23}(\mathrm{\Theta })\mathrm{\Delta }_{13}(\mathrm{\Theta }),$$ $$\mathrm{}_+^+=\mathrm{det}\left(\begin{array}{cc}p_1+p_2& q_1+q_2\\ p_1+p_3& q_1+q_3\end{array}\right)=\mathrm{\Delta }_{12}(\mathrm{\Theta })+\mathrm{\Delta }_{23}(\mathrm{\Theta })+\mathrm{\Delta }_{13}(\mathrm{\Theta }).$$ $`3.18`$ Now, because of 3.16, all three minor determinants must be $`\pm 1`$. This gives 8 possible combinations of the values of $`\mathrm{\Delta }_{ij}(\mathrm{\Theta })`$. It is trivial to check that for any one out of these eight at least two determinants $`\mathrm{}_{}^{}`$ will be equal to $`\pm 3`$ (the other six all being equal to $`\pm 1`$). Hence, even if we choose $`\mathrm{\Theta }`$ so that 3.16 holds, the quotient will necessarily have orbifold singularities of type $`\mathrm{}_3`$. This concludes the proof of Lemma 3.11 Using the calculation in the proof of the above lemma we restate condition 3.12 to get Theorem 3.19: The action of $`Sp(1)\times T_\mathrm{\Theta }^2`$ on $`N_\nu (\mathrm{\Theta })`$ is locally free if and only if (1) all their determinants $`\mathrm{\Delta }_{12}(\mathrm{\Theta }),\mathrm{\Delta }_{23}(\mathrm{\Theta }),\mathrm{\Delta }_{13}(\mathrm{\Theta })`$ do not vanish, and (2) their sum does not vanish, and (3) none of the determinants is equal to the sum of the other two. In such a case the quotient $`(\mathrm{\Theta })`$ is a compact 7-dimensional 3-Sasakian orbifold. Furthermore, there is no weight matrix $`\mathrm{\Theta }`$ for which $`(\mathrm{\Theta })`$ is a smooth manifold. Now, Theorem C of the introduction follows from Theorem 3.19, and the fact that the fundamental 3-dimensional foliation $`(\mathrm{\Theta })𝒪(\mathrm{\Theta })`$, in the case $`(\mathrm{\Theta })`$ is a compact orbifold, yields a compact self-dual Einstein orbifold with a positive scalar curvature (orbifold) metric as the space of leaves. The fact that this metric has only one Killing vector field follows from an appropriate generalization of Theorem B. Remark 3.20: Note that both $`(\mathrm{\Theta })`$ and $`𝒪(\mathrm{\Theta })`$ depend only on the three minor determinants $`\mathrm{\Delta }_{12}(\mathrm{\Theta }),\mathrm{\Delta }_{23}(\mathrm{\Theta }),\mathrm{\Delta }_{13}(\mathrm{\Theta })`$ rather that on $`\mathrm{\Theta }`$ itself. Different weight matrices can certainly lead to equivalent quotients. One could compute the self-dual Einstein metrics $`g(\mathrm{\Theta })`$ on $`𝒪(\mathrm{\Theta })`$ explicitly. Locally we can change variables so that $$\left(\begin{array}{cccc}1& 0& 0& 0\\ 0& A(p_1t+q_1s)& 0& 0\\ 0& 0& A(p_2t+q_2s)& 0\\ 0& 0& 0& A(p_3t+q_3s)\end{array}\right)=\left(\begin{array}{cccc}1& 0& 0& 0\\ 0& A(\lambda _1)& 0& 0\\ 0& 0& A(\lambda _2)& 0\\ 0& 0& 0& A(a\lambda _1+b\lambda _2)\end{array}\right),$$ where $`a=\frac{\mathrm{\Delta }_{23}(\mathrm{\Theta })}{\mathrm{\Delta }_{12}(\mathrm{\Theta })}`$ and $`b=\frac{\mathrm{\Delta }_{13}(\mathrm{\Theta })}{\mathrm{\Delta }_{12}(\mathrm{\Theta })}`$ are now non-zero rational numbers. Such a choice simplifies both the 2-torus action as well as the moment map equations. The self-dual Einstein quotient metric in question (up to scale) will depend on these two parameters $`g(\mathrm{\Theta })=g(a,b)`$. 4. 15-Dimensional Examples At first, it may appear that there are obvious higher dimensional analogues of our construction. However, a simple parity argument shows that such actions do not yield 3-Sasakian metrics other than in dimensions 7, 11, and 15. One can have orbifold metrics in any 3-Sasakian dimension. On the other hand the existence of these smooth quotient in dimension 11 and 15 is closely related to the geometry of $`G_2`$ and $`\mathrm{Spin}(7)`$. We consider two separate cases of $`U(1)_𝐩SO(2k+1)`$ and $`U(1)_𝐩SO(2k)`$. In the $`SO(2k+1)`$ case it suffices to take $`k=4`$. Then $`𝐩=(p_1,p_2,p_3,p_4)`$ and the new $`U(1)_𝐩`$ action is defined by adding an $`SO(2)`$ matrix $`A(p_4t)`$ and let it act by rotating the additional quaternionic coordinates $`(u_8,u_9)`$. To get a smooth quotient any triple $`p_i<p_j<p_k`$ would have to be admissible according to the definition 2.31, as is easily seen from the analysis given in Section 2. However, this is clearly impossible as admissibility implies that each triple contains two odd and one even number. In the $`SO(2k)`$ case we have already seen in the previous section that $`k=3`$ leads to the toric 7-manifolds $`𝒮(\mathrm{\Omega })`$ with $`b_2=2`$. Now we will show that, in fact, $`k=4`$ leads to smooth quotients. Let $`𝐩=(p_1,p_2,p_3,p_4)\mathrm{}^4`$ and define the following homomorphism $$f_𝐩(t)=\left(\begin{array}{cccc}A(p_1t)& 0& 0& 0\\ 0& A(p_2t)& 0& 0\\ 0& 0& A(p_3t)& 0\\ & 0& 0& A(p_4t)\end{array}\right)SO(8),$$ $`4.1`$ where $$A(p_it)=\left(\begin{array}{cc}\mathrm{cos}(p_it)& \mathrm{sin}(p_it)\\ \mathrm{sin}(p_it)& \mathrm{cos}(p_it)\end{array}\right)SO(2),i=1,2,3,4$$ $`4.2`$ are 2-dimensional real rotations. As before we can choose all weights to be non-negative. Further note that at most one of the weights can vanish. We prove the following Lemma 4.3: Let $`(p_1,p_2,p_3,p_4)\mathrm{}^4`$. Then the action of $`U(1)_𝐩\times Sp(1)`$ on $`N_\nu (𝐩)`$ is free if and only if $`0p_1<p_2<p_3<p_4`$, $`\mathrm{gcd}(p_i,p_j,p_k)=1`$ and $`\mathrm{gcd}(p_i\pm p_j,p_i\pm p_k)=1`$ for any triple in $`𝐩`$. Proof: First, let us assume that all of the weights are non-negative. Consider a triple, say $`(p_1,p_2,p_3)`$. Set $`u_7=u_8=0.`$ Then the analysis of the previous section shows that the three must be distinct and that we must have $`\mathrm{gcd}(p_1\pm p_2,p_1\pm p_2)=1`$. However, we no longer need the three weights to be pairwise relatively prime as one cannot set two of the quaternionic pairs $`(u_{2i1},u_{2i}),i=1,2,3,4`$ equal to $`(0,0)`$ at the same time. One such quaternionic pair can vanish; hence, we need $`\mathrm{gcd}(p_1,p_2,p_3)=1`$ to get a free action. The analysis in the case when $`p_1=0`$ is similar. Then one sees that the triple $`(p_2,p_3,p_4)`$ has to be admissible in the sense of the Definition 2.31. But that is what Lemma 2.9 says in this case. It is immediately clear that one cannot extend this construction for $`k>4`$ without admitting orbifold singularities in the quotient. Theorem B follows from the above lemma except for its last statement. When $`𝐩=(1,1,1,1)`$ the quotient is singular. It easy to see that the action is quasi-regular just as in the 11-dimensional case. The fact that the two strata fit together giving the 3-Sasakian orbifold $`\mathrm{}_2\backslash \mathrm{Spin}(7)/\mathrm{Spin}(4)`$ can be seen by first identifying the zero level set $`N_\nu (1,1,1,1)`$ with the $`U(1)\left(\mathrm{Spin}(7)/Sp(1)\right)`$ and observing that $`U(1)\mathrm{Spin}(7)U(1)SU(4)\mathrm{}_2`$ \[BGOP\]. Actually, $`𝐩=(1,1,1,1)`$ is not quite the reduction considered by Kobak and Swann \[KS3\]. Instead they use $`𝐩=(0,0,0,1)`$ for the $`U(1)`$-action. The latter is easily seen to give $`\mathrm{Gr}_4(\mathrm{}^7)/\mathrm{}_2`$ as the quaternionic Kähler quotient. This is a simple consequence of the fact that a ”zero momentum” hyperkähler reduction of $`\mathrm{}^2`$ by any circle action is isometric to $`\mathrm{}/\mathrm{\Gamma }`$ ($`\mathrm{\Gamma }=\mathrm{}_2`$ for the standard action). The identification of the cases $`𝐩=(1,1,1,1)`$ and $`𝐩=(0,0,0,1)`$ owes to the isomorphism between $`\text{s}\text{o}(6)`$ and $`\text{s}\text{u}(4)`$. The fact that these quotients cannot be toric follows from the same type of argument as the one used for 11-dimensional quotients. The Lie algebra of the of the group of the 3-Sasakian isometries is $`\mathrm{}^3\text{s}\text{p}(1)`$. The smooth 15-dimensional manifolds $`(𝐩)`$ for $`p_1=0`$ and $`p_1>0`$ are geometrically different. In the first case $`(𝐩)`$ contains two copies of the 11-dimensional 3-Sasakian manifold $`(p_2,p_3,p_4)`$ which intersect in the 7-dimensional toric 3-Sasakian submanifold $`𝒮(\mathrm{\Omega }(p_2,p_3,p_4))`$. When $`p_1>0`$ then $`(𝐩)`$ does not have any obvious 11-dimensional 3-Sasakian submanifolds. However, we do get 4 disjoint toric 7-dimensional 3-Sasakian submanifolds $`𝒮(\mathrm{\Omega }(p_1,\mathrm{},\widehat{p}_i,\mathrm{},p_4))`$ by setting one of the quaternionic pairs $`(u_{2i1},u_{2i})=(0,0)`$. We can also consider non-zero momentum $`\xi \text{s}\text{p}(1)`$-deformations of the hyperkähler metric on the cones $`C((𝐩))`$. (Up to scale one can set $`\xi `$ to be any imaginary quaternion, so we really have just one parameter family). In some sense they are all deformations of Kronheimer’s hyperkähler metrics on the two nilpotent orbits in $`\text{s}\text{l}(3,\mathrm{})`$ and $`\text{s}\text{o}(6,\mathrm{})`$. Unfortunately such hyperkähler quotients are rarely free from orbifold singularities. In fact, in the 12-dimensional case we never get any complete metrics. In 16 dimensions we do get a complete metric only when $`𝐩=(1,1,1,1)`$ (or equivalently $`𝐩=(0,0,0,1)`$). The metric is $`SU(4)`$-invariant and it gives the Kronheimer metric on the 16-dimensional nilpotent orbit of $`\text{s}\text{o}(6,\mathrm{})`$ in the $`\xi 0`$ scaling limit. More generally, cohomogeneity 2 metrics were studied by Kobak and Swann \[KS4\]. As each classical nilpotent orbit is a hyperkähler reduction of some quaternionic vector space \[KS2\] it is tempting to undertake a more systematic study of the following problem: Which of the nilpotent orbits can give rise to compact 3-Sasakian manifolds? Certainly, any time a quotient involves some $`U(1)`$-factor one can introduce weights. However, as demonstrated here, requiring smoothness often puts very severe restrictions on the weights. We plan to address some of these questions in a future work \[BGOP\]. Remark 4.4: Note that all of the quotients considered in this paper are examples of toric reductions of the 3-Sasakian homogeneous space associated to the real Grassmanian $`Gr_4(\mathrm{}^n)`$. This space is $`SO(n)/SO(n4)\times Sp(1)`$ and it has $`SO(n)`$ as the group of isometries preserving the 3-Sasakian structure. Consider the maximal torus $`T^lSO(n)`$. Then the relevant question is: Which subgroups $`T^mT^l`$ yield smooth 3-Sasakian quotients? Sections 2 and 3 give the complete analysis of the $`n=2k+1=7`$ case. In Section 2 $`m=1`$ and in Section 3 $`m=2`$ which exhaust all interesting possibilities. This section considers $`n=2k=8`$ case with $`m=1`$. One can see that $`m>1`$ does not yield any smooth quotient but we leave the analysis of this to a future work where we plan to give a complete answer in the most general case of arbitrary $`(m,n)`$ \[BGP\]. 5. Comments on the Topology of $`(𝐩)`$ and Related Spaces In this section we denote by $`(𝐩)`$ either one of the 11 or 15 dimensional 3-Sasakian manifolds discussed in sections 2 and 4 with $`𝐩`$ admissible. Actually as discussed below $`(𝐩)`$ will denote a component of the manifolds discussed previously. It would be interesting to know the topology of our quotients, most importantly $`\pi _1((𝐩))`$ and $`H_2((𝐩),\mathrm{})`$. For this one needs to understand the topology of the level set of the moment map $`N_\nu (𝐩).`$ Of course, we do know that $`\pi _1((𝐩))`$ is finite and that the odd Betti numbers of $`(𝐩)`$ vanish up to the middle dimension \[GS\]. However, beyond this not much explicit topological information can be obtained. For example, so far we have been unable to determine whether $`(𝐩)`$ and $`N_\nu (𝐩)`$ are even connected. This presents no real problem as we shall always mean by these spaces connected components such that $`N_\nu (𝐩)`$ is a $`S^1\times Sp(1)`$ bundle over $`(𝐩).`$ Generally, the determination of the topology of an intersection of real quadrics such as $`N_\nu (𝐩)`$ is quite complicated. The analysis in previous work \[BGM1\] and \[BGMR\] relied heavily on very specialized information. In the former case the level sets in question were diffeomorphic to certain Stiefel manifolds whose topology is completely understood, whereas in the later case there was a large symmetry group whose quotient was a two dimensional space whose topology could be analysized. In the present case we meet with no such good fortune. Nevertheless, we are able to obtain a small amount of general information about the topology of $`N_\nu (𝐩).`$ Our first result is that for $`𝐩`$ and $`𝐩^{}`$ admissible the level sets $`N_\nu (𝐩)`$ and $`N_\nu (𝐩^{})`$ are diffeomorphic. This does not hold generally as the case of $`N_\nu (1,1,1)`$ shows. In this case the Jacobian matrix drops rank making the quotient singular. Lemma 5.1: For $`𝐩`$ and $`𝐩^{}`$ admissible, the level sets $`N_\nu (𝐩)`$ and $`N_\nu (𝐩^{})`$ are diffeomorphic. Proof: For $`𝐩`$ admissible there are no fixed points of the $`U(2)`$ action so by the 3-Sasakian version of a well known result in symplectic geometry, zero is a regular value of the $`U(2)`$ moment map $$\stackrel{~}{\nu }_𝐩:S^{27}\mathrm{}^{12}=\mathrm{}^3\times \mathrm{}^3\times \mathrm{}^3\times \mathrm{}^3$$ defined by $`\stackrel{~}{\nu }_𝐩=(\nu _𝐩,\mu _i,\mu _j,\mu _k).`$ Now the level sets $`N_\nu (𝐩)`$ are defined for all $`𝐩\mathrm{}^3.`$ Moreover, for any non-zero $`\rho \mathrm{}`$ we see that level sets of $`\stackrel{~}{\nu }_{\rho 𝐩}`$ and $`\stackrel{~}{\nu }_𝐩`$ coincide, that is $`N_\nu (\rho 𝐩)=N_\nu (𝐩).`$ Thus, by scaling we can choose $`\rho `$ such that $`\stackrel{~}{\nu }_{\rho 𝐩^{}}`$ is in an $`ϵ`$-neighborhood of $`\stackrel{~}{\nu }_{\rho 𝐩}`$ in $`C^{\mathrm{}}(S^{27},\mathrm{}^{12})`$ with the $`C^{\mathrm{}}`$ compact-open topology. Since zero is a regular value of both $`\stackrel{~}{\nu }_{\rho 𝐩}`$ and $`\stackrel{~}{\nu }_{\rho 𝐩^{}},`$ it follows from a well known theorem (cf. \[BCR\], 14.1.1) that $`N_\nu (𝐩)`$ and $`N_\nu (𝐩^{})`$ are diffeotopic, that is there is a one parameter family of diffeomorphisms $`\varphi _t:S^{27}S^{27}`$ parameterized by the unit interval such that $`\varphi _0`$ is the identity and $`\varphi _1`$ takes $`N_\nu (𝐩^{})`$ diffeomorphically to $`N_\nu (𝐩).`$ This immediately implies that the homotopy groups as well as the cohomology rings of $`N_\nu (𝐩)`$ and $`N_\nu (𝐩^{})`$ are isomorphic. Now once and for all we shall choose a component of $`N_\nu (𝐩)`$ and the corresponding component of $`(𝐩)`$ so that $`N_\nu (𝐩)`$ is an $`S^1\times Sp(1)`$ bundle over $`(𝐩)`$ with both base space and total space connected. We also denote by $`(𝐩)`$ the circle bundle over $`(𝐩)`$ that coincides with the quotient of the corresponding component of $`N_\nu (𝐩)`$ by the $`Sp(1)`$ action. We now study the latter as a fibration, namely $`S^3N_\nu (𝐩)\genfrac{}{}{0pt}{}{\pi }{}(𝐩).`$ Note that since $`\varphi _1`$ in Lemma 5.1 is not necessarily a bundle map, we cannot claim that the manifolds $`(𝐩)`$ and $`(𝐩^{})`$ are diffeomorphic. Nevertheless, we can obtain some useful information about their cohomology and homotopy groups. The manifolds $`(𝐩)`$ are of interest in their own right since as discussed briefly in the next section they admit hypercomplex structures. Consider the following commutative diagram of Gysin sequences with $`\mathrm{}`$ coefficients: $$\begin{array}{ccccc}H^{r+3}(N_\nu (𝐩))& H^r((𝐩))\genfrac{}{}{0pt}{}{\chi }{}& H^{r+4}((𝐩))\genfrac{}{}{0pt}{}{\pi ^{}}{}& H^{r+4}(N_\nu (𝐩))& H^{r+1}((𝐩))\\ \text{}& & & \text{}& \\ H^{r+3}(N_\nu (𝐩^{}))& H^r((𝐩^{}))\genfrac{}{}{0pt}{}{\chi }{}& H^{r+4}((𝐩^{}))\genfrac{}{}{0pt}{}{\pi ^{}}{}& H^{r+4}(N_\nu (𝐩^{}))& H^{r+1}((𝐩^{})),\end{array}$$ $`5.2`$ where the two vertical arrows are isomorphisms by Lemma 5.1, and $`\chi `$ denotes cupping by the Euler class of the bundle. The idea is to construct, as best as possible, the missing vertical arrows and relate the cohomology of $`(𝐩)`$ and $`(𝐩^{})`$ by the Five Lemma. First we notice that setting $`r=3`$ and $`2`$ and using Lemma 5.1 gives isomorphisms $$H^1((𝐩),\mathrm{})H^1(N_\nu (𝐩),\mathrm{})H^1(N_\nu (𝐩^{}),\mathrm{})H^1((𝐩^{}),\mathrm{})$$ $$H^2((𝐩),\mathrm{})H^2(N_\nu (𝐩),\mathrm{})H^2(N_\nu (𝐩^{}),\mathrm{})H^2((𝐩^{}),\mathrm{})$$ Next by setting $`r=1`$ we have 0H3((𝐩))πH3(Nν(𝐩))H0((𝐩)) = ϕ 0H3((𝐩))(π)H3(Nν(𝐩))H0((𝐩))\matrix{0&\hbox to18.0pt{\rightarrowfill}&H^{3}({\cal H}({\bf p}))\raise 4.0pt\hbox{$\pi^{}\atop\hbox to18.0pt{\rightarrowfill}$}&H^{3}(N_{\nu}({\bf p}))&\hbox to18.0pt{\rightarrowfill}&H^{0}({\cal H}({\bf p}))\cr\phantom{\hbox{$\scriptstyle{=}$}}\left\downarrow\vbox{\vskip 15.0pt\hbox{$\scriptstyle{=}$}}\right.&&&\phantom{\hbox{$\scriptstyle{\phi^{}}$}}\left\downarrow\vbox{\vskip 15.0pt\hbox{$\scriptstyle{\phi^{}}$}}\right.&&\phantom{\hbox{$\scriptstyle{\approx}$}}\left\downarrow\vbox{\vskip 15.0pt\hbox{$\scriptstyle{\approx}$}}\right.\cr 0&\hbox to18.0pt{\rightarrowfill}&H^{3}({\cal H}({\bf p}^{\prime}))\raise 4.0pt\hbox{$(\pi^{\prime})^{}\atop\hbox to18.0pt{\rightarrowfill}$}&H^{3}(N_{\nu}({\bf p}^{\prime}))&\hbox to18.0pt{\rightarrowfill}&H^{0}({\cal H}({\bf p}^{\prime}))} $`5.3`$ Since both $`\pi ^{}`$ and $`(\pi ^{})^{}`$ are injective, $`\varphi ^{}`$ is an isomorphism, and the diagram is exact and commutative, we can define the missing vertical arrow by $`\psi =((\pi ^{})^{})^1\varphi ^{}\pi ^{}`$ and it is an isomorphism. We thus have $$H^3((𝐩),\mathrm{})H^3((𝐩^{}),\mathrm{}).$$ Now generally we cannot construct the missing vertical maps; however, we can construct them if the groups are free. We thus change to rational coefficients $`\mathrm{}.`$ Consider now diagram 5.2 for $`r=0.`$ We can now fill in the second and last columns with vertical arrows that are isomorphisms. Now considering the diagram with rational coefficients, we can split the middle groups, for all admissable $`𝐩,`$ as $$H^4((𝐩),\mathrm{})\text{im}(\chi )\text{coker}(\chi ).$$ Choosing bases for these groups we can define the middle vertical map simply by sending a basis element in $`\text{im}(\chi )H^4((𝐩),\mathrm{})`$ to a basis element in $`\text{im}(\chi )H^4((𝐩^{}),\mathrm{}),`$ and a basis element in $`\text{coker}(\chi )H^4((𝐩),\mathrm{})`$ to a basis element in $`\text{coker}(\chi )H^4((𝐩^{}),\mathrm{}).`$ That these subspaces have the same dimension making this possible follows from exactness and commutativity of the diagram. It then follows from the Five Lemma that this middle arrow is an isomorphism. This argument is general and using a simple induction we arrive at Theorem 5.4: For admissible $`𝐩`$ and $`𝐩^{}`$ we have (i) $`H^r((𝐩),\mathrm{})H^r((𝐩^{}),\mathrm{})`$ for $`r=0,1,2,3.`$ (ii) $`b_r((𝐩))=b_r((𝐩^{}))`$ for all $`r.`$ In particular, $`(𝐩)`$ and $`(𝐩^{})`$ have isomorphic rational cohomology groups. Here $`b_r(M)`$ denotes the rth Betti number of $`M.`$ Similar arguments can be used for the homotopy groups, and combining these results with well known facts about circle bundles over 3-Sasakian manifolds \[BGM3\] we obtain Proposition 5.5: For $`𝐩`$ and $`𝐩^{}`$ admissible, we have isomorphisms: (i) $`\pi _i((𝐩))\pi _i(N_\nu (𝐩))\pi _i((𝐩^{}))`$ for $`i=0,1,2.`$ (ii) $`H^1((𝐩),\mathrm{})H^1(N_\nu (𝐩),\mathrm{})H^1((𝐩^{}),\mathrm{})0\text{or}\mathrm{}.`$ (iii) $`H^2((𝐩),\mathrm{})H^2(N_\nu (𝐩),\mathrm{}).`$ Next we consider our manifolds of primary interest, namely $`(𝐩).`$ First, it is easy to see the following relations: Proposition 5.6: For admissible $`𝐩`$ we have: (i) $`\pi _i((𝐩))\pi _i((𝐩))`$ for all $`i>2.`$ (ii) $`b_2(N_\nu (𝐩))=\{\begin{array}{cc}b_2((𝐩))1\hfill & \text{if }b_1(N_\nu (𝐩))=0;\hfill \\ b_2((𝐩))\hfill & \text{if }b_1(N_\nu (𝐩))=1.\hfill \end{array}`$ Next we have the analogue of Theorem 5.4 for our $`3`$-Sasakian manifolds $`(𝐩).`$ Theorem 5.7: For admissible $`𝐩`$ and $`𝐩^{}`$ we have (i) $`H^r((𝐩),\mathrm{})H^r((𝐩^{}),\mathrm{})`$ for $`r=0,1,2,3.`$ (ii) $`b_r((𝐩))=b_r((𝐩^{}))`$ for all $`r.`$ In particular, $`(𝐩)`$ and $`(𝐩^{})`$ have isomorphic rational cohomology groups. Proof: The proof of this Theorem is analagous to the proof of Theorem 5.4 with diagram 5.2 replaced by the following commutative diagram of Gysin sequences: $$\begin{array}{ccccc}H^{r+1}((𝐩))& H^r((𝐩))\genfrac{}{}{0pt}{}{\chi }{}& H^{r+2}((𝐩))& H^{r+2}((𝐩))& H^{r+1}((𝐩))\\ \text{}& & & \text{}& \\ H^{r+1}((𝐩^{}))& H^r((𝐩^{}))\genfrac{}{}{0pt}{}{\chi }{}& H^{r+2}((𝐩^{}))& H^{r+2}((𝐩^{}))& H^{r+1}((𝐩^{})),\end{array}$$ where we have used Theorem 5.4 and its proof to construct the isomorphisms indicated by the vertical arrows. Remark 5.8: Actually we can weaken the hypothesis that $`𝐩`$ be admissible by noting any of the results of this section concerning rational cohomology hold in the cases when $`(𝐩)`$ is an orbifold obtained as the quotient by a locally free action. It follows from the analysis in section 2, that the action is locally free precisely when the components of $`𝐩`$ are all distinct, and in this case the level sets $`(𝐩)`$ are smooth manifolds; hence, Theorem 5.4 and Proposition 5.5 hold in this case as well. It is interesting to note that in the case that $`(𝐩)`$ is an orbifold, but not a smooth manifold, the smooth manifold $`(𝐩)`$ cannot be the trivial V-bundle. The above remarks apply equally as well to the 7-dimensional orbifolds $`(\mathrm{\Theta })`$ constructed in section 3 with the condition for a locally free action being that all the minor determinants $`\mathrm{\Delta }_{ij}(\mathrm{\Theta })`$ are nonvanishing. Finally we briefly discuss the two singular cases $$(1,1,1)\mathrm{}_3\backslash G_2/Sp(1)\text{and}(1,1,1,1)\mathrm{}_2\backslash \text{Spin}(7)/\text{Spin}(4).$$ Since these are biquotients of Lie groups the topology is more accessible. In particular, their rational cohomology is that of the corresponding 3-Sasakian homogeneous space, $`G_2/Sp(1)`$ and $`\text{Spin}(7)/\text{Spin}(4),`$ respectively, which is well known \[GS,BG2\]. Thus, $`(1,1,1)`$ has the rational cohomology of $`S^{11},`$ whereas $`(1,1,1,1)`$ has the rational cohomology of $`S^4\times S^{11}.`$ In both cases $`b_2`$ vanishes, and we do not expect this in the non-singular cases. 6. Hypercomplex Structures on Circle Bundles over $`(𝐩)`$ According to the general theory described in \[BGM3\] 3-Sasakian manifolds (orbifolds) give rise to hypercomplex structures on circle bundles over them. In this short section we give new hypercomplex structures in dimensions 12 and 16 constructed as circle V-bundles over the 3-Sasakian orbifolds $`(𝐩)`$ constructed in sections 2 and 4. Of course, there is the trivial bundle $`(𝐩)\times S^1`$ over $`(𝐩)`$ which always admits a locally conformally hyperkähler structure, but here we concentrate on the level sets $`(𝐩).`$ As discussed in Remark 5.8 these level sets will be smooth manifolds as long as $`0<p_1<p_2<p_3`$ in the 12 dimensional case and $`0p_1<p_2<p_3<p_4`$ in the 16 dimensional case. We now have from our previous results \[BGM3\]: Theorem 6.1: Let $`𝐩`$ have components satisfying the inequalities above, then $`(𝐩)`$ is a compact hypercomplex manifold of dimension 12 or 16. Furthermore, the connected component of the Lie group of hypercomplex automorphisms is $`T^3`$ in the 12 dimensional case and $`T^4`$ in the 16 dimensional case. The last statement of Theorem 6.1 implies that these hypercomplex structures are distinct from any of those known previously. We do not know whether for different $`𝐩`$ the manifolds $`(𝐩)`$ are diffeomorphic or not; however, arguments similar to those in \[BGM2\] show that the hypercomplex structures are distinct. In fact, each smooth manifold $`(𝐩)`$ has a real one parameter family of distinct hypercomplex structures on them given by sending $`𝐩\lambda 𝐩`$ for any real $`\lambda >0.`$ We can also construct hypercomplex structures on the total space $`(\mathrm{\Theta })`$ of circle V-bundles over the 7-dimensional 3-Sasakian orbifolds $`(\mathrm{\Theta })`$ constructed in Section 3. In this case as in \[BGM3\] there should be gcd conditions on the entries of the matrix $`\mathrm{\Theta }`$ that gaurentee that $`(\mathrm{\Theta })`$ be a smooth manifold. These then give new hypercomplex manifolds in dimension 8 with a two-dimensional group of hypercomplex automorphisms. Bibliography \[BCR\] J. Bochnak, M. Coste, and M.-F. Roy, Géométrie algébrique réele, Springer-Verlag, Berlin, 1987. \[BG1\] C.P. Boyer and K. Galicki, On Sasakian-Einstein geometry, mathDG/9811098 to appear in Int. J. Math. \[BG2\] C.P. Boyer and K. Galicki, 3-Sasakian Manifolds, Surveys in Differential Geometry, Volume VI, Essays on Einstein Manifolds, A supplement to the Journal of Differential Geometry, pp. 123-184, C. LeBrun and M. Wang, Eds., International Press, Cambridge 1999. \[BGM1\] C.P. Boyer, K. Galicki, and B.M. Mann, The geometry and topology of 3-Sasakian manifolds, J. reine angew. Math., 455 (1994), 183-220. \[BGM2\] C.P. Boyer, K. Galicki, and B.M. Mann, Hypercomplex structures on Stiefel manifolds, Ann. Global Anal. Geom. 14 (1996), 81-105. \[BGM3\] C.P. Boyer, K. Galicki, and B.M. Mann, Hypercomplex structures from 3-Sasakian structures, J. reine angew. Math., 501 (1998), 115-141. \[BGMR\] C.P. Boyer, K. Galicki, B.M. Mann, and E. Rees, Compact 3-Sasakian 7-Manifolds with Arbitrary Second Betti Number, Invent. Math. 131 (1998), 321-344. \[BGP\] C.P. Boyer, K. Galicki, and P. Piccinni, Torus actions on $`Gr_4(\mathrm{}^n)`$ and 3-Sasakian manifolds, in preparation. \[BGOP\] C.P. Boyer, K. Galicki, L. Ornea, and P. Piccinni, Geometry of Exceptional Quotients, in preparation. \[DS\] A. Dancer and A. Swann, The Geometry of Singular Quaternionic Kähler Quotients, Int. J. Math., 8 (1997), 595-610. \[G\] K. Galicki, A generalization of the momentum mapping construction for quaternionic Kähler manifolds, Commun. Math. Phys 108 (1987), 117-138. \[GL\] K. Galicki and B. H. Lawson, Jr., Quaternionic Reduction and Quaternionic Orbifolds, Math. Ann., 282 (1988), 1-21. \[GS\] K. Galicki and S. Salamon, On Betti numbers of 3-Sasakian manifolds, Geom. Ded. 63 (1996), 45-68. \[Hi1\] N. J. Hitchin, A new family of Einstein metrics. Manifolds and geometry (Pisa, 1993), 190–222, Sympos. Math., XXXVI, Cambridge Univ. Press, Cambridge, 1996. \[Hi2\] N. J. Hitchin, Twistor spaces, Einstein metrics and isomonodromic deformations, J. Differential Geom. 42 (1995), 30–112. \[HL\] R. Harvey and B. H. Lawson, Jr., Calibrated geometries, Acta Math. 148 (1982), 47–157. \[Kr1\] P. Kronheimer, A hyperkähler structure on coadjoint orbits of a semisimple complex group, J. Lond. Math. Soc., II. Ser. 42 (1990), 193-208. \[Kr2\] P. Kronheimer, Instantons and geometry of the nilpotent variety, J. Differential Geom. 32 (1990), 473-490. \[KS1\] P. Kobak and A. Swann, Quaternionic geometry of a nilpotent variety, Math. Ann. 297 (1993), 747-764. \[KS2\] P. Kobak and A. Swann, Classical nilpotent orbits as hyperkähler quotients, Internat. J. Math. 7 (1996), 193-210. \[KS3\] P. Kobak and A. Swann, Exceptional hyperkähler reductions, Twistor Newsletter 44 (1998), 23-26. \[KS4\] P. Kobak and A. Swann, Hyperkähler Potentials in Cohomogeneity Two, math.DG/0001024. \[OP\] L. Ornea and P. Piccinni, On some moment maps and induced Hopf bundles on the quaternionic projective plane, mathDG/0001066, to appear in Int. J. Math. \[PS\] Y. S. Poon and S. Salamon, Eight-dimensional quaternionic Kähler manifolds with positive scalar curvature, J. Differential Geom. 33 (1990), 363-378. \[Sw\] A. F. Swann, Hyperkähler and quaternionic Kähler geometry, Math. Ann., 289 (1991), 421-450. Department of Mathematics and Statistics July 2000 University of New Mexico Albuquerque, NM 87131 email: cboyer@math.unm.edu, galicki@math.unm.edu Università degli Studi di Roma, “La Sapienza” Piazzale Aldo Moro 2 I-00185 Roma, Italia email: piccinni@mat.uniroma1.it
warning/0007/astro-ph0007262.html	ar5iv	text	# The combination of ground-based astrometric compilation catalogues with the HIPPARCOS Catalogue ## 1 Introduction In Paper I (Wielen et al. 1999b), we have shown that the combination of the data of the HIPPARCOS astrometric satellite (ESA 1997) with ground-based results (such as the FK5) is able to provide for many stars individual proper motions which are significantly more accurate than the HIPPARCOS proper motions themselves. The method has been already successfully applied in the construction of the FK6, the Sixth Catalogue of Fundamental Stars (Part I of the FK6: Wielen et al. 1999d, Part III of the FK6: Wielen et al. 2000a). The method of combination presented in Paper I is, however, strictly valid for single stars only. In the FK6, we call this procedure therefore the ‘single-star mode’. In reality most of the stars are members of binaries or of multiple systems. If the duplicity of an individual object is already definitely known, either from ground-based investigations or from HIPPARCOS observations, then the method of combination has to be changed properly in order to obtain meaningful results. We call such procedures ‘special solutions’. Paper III of this series of papers will discuss the special solutions for visual binaries and other types of double stars. The present paper (Paper II) describes an appropriate method of the combination of HIPPARCOS results with ground-based observations for ‘apparently single stars’. Even if we have removed from such a sample of apparently single stars all the objects with known duplicity, there remains beside the truly single-stars a large number of hitherto undetected astrometric binaries. The measured photo-center of an unresolved astrometric binary moves on the sky on a wavy curve, in contrast to the linear motion of single stars (Fig. 1). In such a case, an ‘instantaneously’ measured proper motion deviates from a ‘mean’ proper motion, averaged over a long interval of time. The proper motions provided by HIPPARCOS (ESA 1997) are essentially such instantaneously measured quantities, since they are derived from positional measurements spread over about three years only. In contrast, the proper motions given in the FK5 (Fricke et al. 1988, 1991) are long-term averages over up to about 200 years. We have called the difference between the instantaneous proper-motion and the mean one the ‘cosmic error’ of the instantaneous proper motion (Wielen 1995a, b, 1997, Wielen et al. 1997). In some cases, the individual cosmic error is so large that the duplicity of an apparently single star is strongly indicated by this fact. We have called such objects ‘$`\mathrm{\Delta }\mu `$ binaries’ (Wielen et al. 1999a). In most cases, however, the cosmic error in the HIPPARCOS proper motion of a hitherto undetected astrometric binary is not individually significant, but can be shown to exist only statistically in a larger sample of apparently single stars. Nevertheless, an appropriate combination method should not neglect the statistical consequences of these cosmic errors. Our comparison of HIPPARCOS proper motions with ground-based results has shown that, at least for brighter stars, the average cosmic error in a HIPPARCOS proper motion is often larger than the HIPPARCOS measuring error, typically by a factor of three (Wielen 1995a, b, Wielen et al. 1997, 1998, 1999c). Wielen (1997, henceforth called Paper WPSA) has developed a coherent scheme of ‘statistical astrometry’ for handling the effects of cosmic errors in high-precision astrometry. In the following sections, we shall apply the concepts of statistical astrometry to the problem of combining the HIPPARCOS results with ground-based measurements for a sample of apparently single stars. The main results will be solutions which we call the ‘long-term prediction mode (LTP)’ and the ‘short-term prediction mode (STP)’. Our results have already been applied for the LTP and STP solutions given in Part I and Part III of the FK6 (Wielen et al. 1999d, 2000a). It is clear that in principle the most desirable solution for our problem would be to treat each star individually and fully correctly. This would mean (1) for truly single stars: to use the ‘single-star mode (SI)’ described in Paper I, and (2) for binaries: to apply individual orbital corrections, as e.g. done for Polaris by Wielen et al. (2000b). However, since for ‘apparently single stars’ the true nature (single or double) of the individual objects is unknown, we have to rely on statistical methods in order to handle such a sample of stars properly. The treatment of a sample of apparently single stars by our statistical procedures gives on average the best astrometric prediction, and it provides the most realistic error budget for such a sample. In this sense, our statistical treatment is certainly much more appropriate than to ignore the binary nature of a large fraction of a sample of apparently single stars altogether. In Paper I, we have already pointed out that the older observations carry a high weight in the combination of ground-based measurements with HIPPARCOS data, and that therefore the GC (Boss et al. 1937) should be also considered here because of its high number of rather well-measured stars. We call the result of the combination of the GC with HIPPARCOS the combination catalogue GC+HIP. ## 2 Basic concepts and equations ### 2.1 Concepts With respect to the two catalogues which should be combined into a new one, we follow closely the situation described in Sect. 2 of Paper I. We are using also as far as possible the nomenclature of that section. We assume that two astrometric compilation catalogues are available, identified by the indices 1 (e.g. for the FK5) and 2 (e.g. for HIPPARCOS). For the combined catalogue (e.g. the FK6), we use the index C, usually now supplemented by an additional subindex which identifies the special mode of the solution (e.g. LTP or STP). Each of the two basic catalogues ($`i=1,2`$) provides for the stars a position $`x_i(T_i)`$ at a central epoch $`T_i`$ and a proper motion $`\mu _i`$ for two angular coordinates (e.g. right ascension $`\alpha _{}=\alpha \mathrm{cos}\delta `$ and declination $`\delta `$). The mean measuring errors of $`x_i(T_i)`$ and $`\mu _i`$ are denoted by $`\epsilon _{x,i}`$ and $`\epsilon _{\mu ,i}`$. Usually one of the catalogues (e.g. $`i=1`$) has first to be reduced to the astrometrical system of the other catalogue (e.g. the FK5 to the HIPPARCOS system). In this case, $`x_1`$ and $`\mu _1`$ denote the already systematically corrected quantities, and their mean errors $`\epsilon _{x,1}`$ and $`\epsilon _{\mu ,1}`$ include the uncertainty of the systematic corrections. For the determination of the systematic differences between two catalogues, we use methods developed for the construction of the FK5 (Bien et al. 1978). With respect to the principles of statistical astrometry, we make now the important additional assumptions that Catalogue 1 gives ‘mean’ quantities for $`x`$ and $`\mu `$, averaged over a long interval of time, while Catalogue 2 provides ‘instantaneously’ measured values of $`x`$ and $`\mu `$. If we apply our scheme to a combination of the FK5 with HIPPARCOS, both assumptions are fulfilled to a high degree of approximation: The FK5 is based on ground-based observations spread over about two centuries, while the HIPPARCOS results are obtained from measurements made during a short interval of time, about three years only. In the terminology of statistical astrometry, our assumptions mean that we suppose that the FK5 is free from cosmic errors. The cosmic errors in the HIPPARCOS positions and proper motions are denoted by $`c_x=(\xi (0))^{1/2}`$ and $`c_\mu =(\eta (0))^{1/2}`$. The correlation functions $`\xi (\mathrm{\Delta }t)`$, $`\eta (\mathrm{\Delta }t)`$, and $`\zeta (\mathrm{\Delta }t)`$ are explained in WPSA (especially Sect. 3). Numerical values for $`c_x(p)`$ and $`c_\mu (p)`$ as functions of the parallax $`p`$ will be provided in Sect. 5. ### 2.2 Basic equations In the combination of two astrometric catalogues of which at least one contains instantaneously measured data (affected by cosmic errors), the predicted position $`x_\mathrm{p}(t)`$ at an arbitrary epoch $`t`$ should not be anymore a linear function of time. According to the principles of statistical astrometry, the ‘best’ prediction $`x_\mathrm{p}(t)`$ for the true instantaneous position $`x(t)`$ is given by the non-linear expression (WPSA, Eq. (82), with a slight change in nomenclature): $`x_\mathrm{p}(t)`$ $`=`$ $`x_1(T_1)+\mu _1(tT_1)`$ (1) $`+`$ $`\gamma (t)\left(x_2(T_2)x_1(T_1)\mu _1(T_2T_1)\right)`$ $`+`$ $`\beta (t)\left(\mu _2(T_2)\mu _1\right)(tT_2).`$ The quantities $`\gamma (t)`$ and $`\beta (t)`$ are functions of time $`t`$ to be determined. For computational reasons, we often use $$B(t)=\beta (t)(tT_2)$$ (2) instead of $`\beta (t)`$. The functions $`\gamma (t)`$ and $`\beta (t)`$ are determined from the condition that the mean error $`\epsilon _{x,\mathrm{p}}(t)`$ of the predicted positions $`x_\mathrm{p}(t)`$ should be a minimum for every value of $`t`$, averaged over the ensemble: $$\epsilon _{x,\mathrm{p}}^2(t)=<\left(x_\mathrm{p}(t)x(t)\right)^2>=min..$$ (3) The operator $`<q>`$ means, as in WPSA (1997), the average of the quantity $`q`$ over the ensemble of similar stars. For an individual star, $`<q>`$ can be interpreted as the ‘expectation value’ of the quantity $`q`$ for this star. In this sense, we consider the ensemble averages as statistical predictions for the mean behaviour of a typical individual member of the ensemble. Inserting Eq. (1) into Eq. (3), and using the scheme of statistical astrometry, we obtain for $`\epsilon _{x,\mathrm{p}}^2(t)`$: $`\epsilon _{x,\mathrm{p}}^2(t)=<\left(\left(x_\mathrm{p}(t)x_\mathrm{m}(t)\right)\left(x(t)x_\mathrm{m}(t)\right)\right)^2>`$ (4) $`=<((x_1(T_1)x_\mathrm{m}(T_1))(1\gamma (t))`$ $`+(\mu _1\mu _\mathrm{m})\left(tT_1\gamma (t)(T_2T_1)\beta (t)(tT_2)\right)`$ $`+\gamma (t)\left((x_2(T_2)x(T_2))+(x(T_2)x_\mathrm{m}(T_2))\right)`$ $`+\beta (t)(tT_2)\left((\mu _2(T_2)\mu (T_2))+(\mu (T_2)\mu _\mathrm{m})\right)`$ $`(x(t)x_\mathrm{m}(t)))^2>`$ $`=\left(1+(\gamma (t))^2\right)\xi (0)2\gamma (t)\xi (tT_2)`$ $`+\left(\beta (t)(tT_2)\right)^2\eta (0)2\beta (t)(tT_2)\zeta (tT_2)`$ $`+\epsilon _{x,1}^2\left(1\gamma (t)\right)^2+\epsilon _{x,2}^2\left(\gamma (t)\right)^2`$ $`+\epsilon _{\mu ,1}^2\left(tT_1\gamma (t)(T_2T_1)\beta (t)(tT_2)\right)^2`$ $`+\epsilon _{\mu ,2}^2(\beta (t)(tT_2))^2.`$ The quantities $`x_\mathrm{m}(t)`$ and $`\mu _\mathrm{m}=\dot{x}_\mathrm{m}`$ are the true mean position and the true mean proper motion of the star (in the sense of statistical astrometry). $`x_\mathrm{m}`$ is a linear function of $`t`$, e.g. $$x_\mathrm{m}(t)=x_\mathrm{m}(T_1)+\mu _\mathrm{m}(tT_1).$$ (5) In deriving Eq. (4) we have used the following relations: According to our assumption that Catalogue 1 provides mean quantities, $`x_1(T_1)`$ and $`\mu _1`$ differ from $`x_\mathrm{m}(T_1)`$ and $`\mu _\mathrm{m}`$ by their measuring errors only: $`<(x_1(T_1)x_\mathrm{m}(T_1))^2>`$ $`=`$ $`\epsilon _{x,1}^2,`$ (6) $`<(\mu _1\mu _\mathrm{m})^2>`$ $`=`$ $`\epsilon _{\mu ,1}^2.`$ (7) The quantities $`x(t)`$ and $`\mu (t)`$ are the true instantaneous position and proper motion of the star at epoch $`t`$. Since Catalogue 2 is assumed to provide instantaneous quantities, $`x_2(T_2)`$ and $`\mu _2(T_2)`$ differ from $`x(T_2)`$ and $`\mu (T_2)`$ also by their measuring errors only: $`<(x_2(T_2)x(T_2))^2>`$ $`=`$ $`\epsilon _{x,2}^2,`$ (8) $`<(\mu _2(T_2)\mu (T_2))^2>`$ $`=`$ $`\epsilon _{\mu ,2}^2.`$ (9) The correlation functions $`\xi (\mathrm{\Delta }t),\eta (\mathrm{\Delta }t)`$, and $`\zeta (\mathrm{\Delta }t)`$ (see WPSA) occur because of: $$<(x(T_2)x_\mathrm{m}(T_2))^2>=<(x(t)x_\mathrm{m}(t))^2>=\xi (0),$$ (10) $$<(x(T_2)x_\mathrm{m}(T_2))(x(t)x_\mathrm{m}(t))>=\xi (tT_2),$$ (11) $$<(\mu (T_2)\mu _\mathrm{m})^2>=\eta (0),$$ (12) $$<(x(t)x_\mathrm{m}(t))(\mu (T_2)\mu _\mathrm{m})>=\zeta (tT_2).$$ (13) We have further to remember $$<(x(T_2)x_\mathrm{m}(T_2))(\mu (T_2)\mu _\mathrm{m})>=\zeta (0)=0.$$ (14) The products of the other terms occuring in the second line of Eq. (4) all vanish, because the measuring errors are not correlated with the cosmic errors, for example $$<(x_1(T_1)x_\mathrm{m}(T_1))(x(t)x_\mathrm{m}(t))>=0,$$ (15) and because the measuring errors of $`x_i(T_i)`$ and $`\mu _i(T_i)`$ are not correlated at the central epoch $`T_i`$ (by definition of the central epoch), for example: $$<(x_1(T_1)x_\mathrm{m}(T_1))(\mu _1\mu _\mathrm{m})>=0,$$ (16) and finally because the measuring errors of $`x_1`$ and $`\mu _1`$ are not correlated with those of $`x_2`$ and $`\mu _2`$, for example $$<(x_1(T_1)x_\mathrm{m}(T_1))(\mu _2(T_2)\mu (T_2))>=0.$$ (17) The conditions for the minimum of $`\epsilon _{x,\mathrm{p}}^2(t)`$ with respect to $`\gamma `$ and $`\beta `$ are $$\frac{\epsilon _{x,\mathrm{p}}^2}{\gamma }=0,$$ (18) and $$\frac{\epsilon _{x,\mathrm{p}}^2}{\beta }=0.$$ (19) Carrying out these procedures and using the final part of Eq. (4) for $`\epsilon _{x,\mathrm{p}}^2(t)`$, we obtain the following two equations of condition for the functions $`\gamma (t)`$ and $`\beta (t)`$ which give the ‘best’ values of $`\epsilon _{x,\mathrm{p}}(t)`$: $`\gamma (t)(\xi (0)+\epsilon _{x,1}^2+\epsilon _{x,2}^2+\epsilon _{\mu ,1}^2(T_2T_1)^2)`$ $`+\beta (t)(tT_2)\epsilon _{\mu ,1}^2(T_2T_1)`$ $`=\xi (tT_2)+\epsilon _{x,1}^2+\epsilon _{\mu ,1}^2(tT_1)(T_2T_1),`$ $`\gamma (t)\epsilon _{\mu ,1}^2(T_2T_1)`$ $`+\beta (t)(tT_2)(\eta (0)+\epsilon _{\mu ,1}^2+\epsilon _{\mu ,2}^2)`$ (21) $`=\zeta (tT_2)+\epsilon _{\mu ,1}^2(tT_1).`$ The Eqs. (20) and (21) correspond to the Eqs. (86) and (87) in WPSA. Solving the Eqs. (20) and (21) for the unknowns $`\gamma (t)`$ and $`\beta (t)`$ we obtain: $`\gamma (t)=((\xi (tT_2)+\epsilon _{x,1}^2)(\eta (0)+\epsilon _{\mu ,1}^2+\epsilon _{\mu ,2}^2)`$ $`+\epsilon _{\mu ,1}^2(T_2T_1)((\eta (0)+\epsilon _{\mu ,2}^2)(tT_1)\zeta (tT_2)))/N,`$ (22) $`B(t)`$ $`=`$ $`\beta (t)(tT_2)`$ (23) $`=`$ $`(\zeta (tT_2)(\xi (0)+\epsilon _{x,1}^2+\epsilon _{x,2}^2+\epsilon _{\mu ,1}^2(T_2T_1)^2)`$ $`+\epsilon _{\mu ,1}^2(tT_1)\left(\xi (0)+\epsilon _{x,1}^2+\epsilon _{x,2}^2\right)`$ $`\epsilon _{\mu ,1}^2(T_2T_1)(\xi (tT_2)+\epsilon _{x,1}^2))/N,`$ with the auxiliary quantity $`N`$: $`N=(\xi (0)+\epsilon _{x,1}^2+\epsilon _{x,2}^2)(\eta (0)+\epsilon _{\mu ,1}^2+\epsilon _{\mu ,2}^2)`$ (24) $`+\epsilon _{\mu ,1}^2(T_2T_1)^2(\eta (0)+\epsilon _{\mu ,2}^2).`$ For deriving the functions $`\gamma (t)`$ and $`\beta (t)`$, we have up to now implicitely assumed that all the stars in the ensemble have the same measuring errors. This is, of course, not strictly true in reality. However, the ensemble averages are actually neccessary for handling the cosmic errors only, but not for treating the measuring errors. Hence we shall use all the equations derived above by inserting the individual measuring errors if we treat individual stars of the ensemble. Formally we may imagine to handle subsamples of stars in which the stars have the overall behaviour with respect to the cosmic errors, but in which the common measuring errors are equal to those of the individual star under consideration. It is a more severe problem that we do not make use of our knowledge of how large the individual cosmic errors in $`x`$ and $`\mu `$ are for a given individual object (e.g. $`\mu _2\mu _1`$). As discussed in WPSA, this would require ‘conditioned’ correlation functions. Since this information is presently not available, we are treating here the consequences of the cosmic errors on the level of ensemble averages only. Inserting these results for $`\gamma (t)`$ and $`\beta (t)`$ from Eqs. (22) and (23) into Eqs. (1) and (4), we derive the prediction $`x_\mathrm{p}(t)`$ for the instantaneous position $`x(t)`$ of the star and the mean error $`\epsilon _{x,\mathrm{p}}(t)`$ of this prediction. The prediction $`x_\mathrm{p}(t)`$ is a non-linear function of $`t`$, because the correlation functions $`\xi (tT_2)`$ and $`\zeta (tT_2)`$, which occur in the formulae for $`\gamma (t)`$ and $`\beta (t)`$, are non-linear functions. A typical run of $`x_\mathrm{p}(t)`$ is shown in Fig. 2. In order to illustrate the properties of our prediction $`x_\mathrm{p}(t)`$, we consider in the following Subsections 2.3 and 2.4 two limiting cases in which we either neglect the measuring errors or the cosmic errors. ### 2.3 Measuring errors neglected If we neglect all the measuring errors and set $`\epsilon _{x,1}(T_1)=\epsilon _{x,2}(T_2)=\epsilon _{\mu ,1}=\epsilon _{\mu ,2}=0`$, then we obtain for $`\gamma (t)`$ and $`\beta (t)`$ $`\gamma _{\mathrm{nme}}(t)`$ $`=`$ $`\xi (tT_2)/\xi (0),`$ (25) $`B_{\mathrm{nme}}(t)`$ $`=`$ $`\beta _{\mathrm{nme}}(t)(tT_2)=\zeta (tT_2)/\eta (0).`$ (26) These results were already derived and discussed in WPSA (Sect. 4.2.4 and Fig. 10). The corresponding prediction $`x_{\mathrm{p},\mathrm{nme}}(t)`$, shown in Fig. 2, passes through the point $`x_2(T_2)`$ with the slope $`\mu _2(T_2)`$. Therefore, the measured instantaneous position and proper motion at epoch $`T_2`$ are exactly reproduced by $`x_{\mathrm{p},\mathrm{nme}}(t)`$. For $`t\pm \mathrm{}`$, the prediction $`x_{\mathrm{p},\mathrm{nme}}(t)`$ approaches asymptotically the mean position $`x_\mathrm{m}(t)=x_1(t)=x_1(T_1)+\mu _1(tT_1)`$ of the star. The uncertainty $`\epsilon _{x,\mathrm{p},\mathrm{nme}}(t)`$ of the prediction is given by $`\epsilon _{x,\mathrm{p},\mathrm{nme}}^2(t)={\displaystyle \frac{(\xi (0))^2(\xi (tT_2))^2}{\xi (0)}}{\displaystyle \frac{(\zeta (tT_2))^2}{\eta (0)}}.`$ (27) We see that $`\epsilon _{x,\mathrm{p},\mathrm{nme}}`$ is zero at the epoch $`T_2`$ and approaches $`c_x=(\xi (0))^{1/2}`$ for $`t\pm \mathrm{}`$, since in this limit $`x_\mathrm{p}(t)`$ is equal to $`x_\mathrm{m}(t)`$, and $`x_\mathrm{m}(t)`$ differs from the instantaneous position $`x(t)`$ by the cosmic error $`c_x`$ on average. ### 2.4 Cosmic errors neglected If we neglect the cosmic errors and set $`\xi =\eta =\zeta =0`$, then we obtain $`\gamma _{\mathrm{nce}}(t)=(\epsilon _{x,1}^2(\epsilon _{\mu ,1}^2+\epsilon _{\mu ,2}^2)`$ (28) $`+\epsilon _{\mu ,1}^2\epsilon _{\mu ,2}^2(T_2T_1)(tT_1))/N_{\mathrm{nce}},`$ $`B_{\mathrm{nce}}(t)=\beta _{\mathrm{nce}}(t)(tT_2)=(\epsilon _{x,1}^2\epsilon _{\mu ,1}^2(T_2T_1)`$ (29) $`+(\epsilon _{x,1}^2+\epsilon _{x,2}^2)\epsilon _{\mu ,1}^2(tT_1))/N_{\mathrm{nce}},`$ with $$N_{\mathrm{nce}}=(\epsilon _{x,1}^2+\epsilon _{x,2}^2)(\epsilon _{\mu ,1}^2+\epsilon _{\mu ,2}^2)+\epsilon _{\mu ,1}^2\epsilon _{\mu ,2}^2(T_2T_1)^2.$$ (30) Equation (29) illuminates one of the advantages of introducing $`B(t)=\beta (t)(tT_2)`$ as a substitute for $`\beta (t)`$. While $`B_{\mathrm{nce}}(T_2)`$, and in general $`B(T_2)`$ for non-zero measuring errors, remains finite, the quantity $`\beta _{\mathrm{nce}}`$, and in general $`\beta `$, tends towards infinity for $`tT_2`$. Only in some degenerated cases, such as $`\epsilon _{\mu ,1}=0`$ or $`\epsilon _{x,2}=0`$ or $`\epsilon _{x,1}\mathrm{}`$, the quantity $`\beta _{\mathrm{nce}}`$, and in general $`\beta `$, remains finite for $`tT_2`$. If we insert $`\gamma _{\mathrm{nce}}(t)`$ and $`B_{\mathrm{nce}}(t)`$ into Eq. (1), the corresponding prediction $`x_{\mathrm{p},\mathrm{nce}}(t)`$ for the position of the star at an epoch $`t`$ is now a strictly linear function of $`t`$, since $`\gamma _{\mathrm{nce}}(t)`$ and $`B_{\mathrm{nce}}(t)`$ are linear in $`t`$. In order to facilitate the understanding of the behaviour of the prediction $`x_{\mathrm{p},\mathrm{nce}}(t)`$, we rewrite $`x_{\mathrm{p},\mathrm{nce}}`$ by using the auxiliary quantities $`T_{\mathrm{C},\mathrm{nce}},x_{\mathrm{C},\mathrm{nce}},\mu _{\mathrm{C},\mathrm{nce}}`$. If we insert Eqs. (28) and (29) into Eq. (1), we obtain after some lengthy algebra: $`x_{\mathrm{p},\mathrm{nce}}(t)`$ $`=`$ $`x_{\mathrm{C},\mathrm{nce}}(T_{\mathrm{C},\mathrm{nce}})+\mu _{\mathrm{C},\mathrm{nce}}(tT_{\mathrm{C},\mathrm{nce}}),`$ (31) with the auxiliary quantities $`T_{\mathrm{C},\mathrm{nce}}`$ $`=`$ $`{\displaystyle \frac{w_{x,1}T_1+w_{x,2}T_2}{w_{x,1}+w_{x,2}}},`$ (32) $`x_{\mathrm{C},\mathrm{nce}}(T_{\mathrm{C},\mathrm{nce}})`$ $`=`$ $`{\displaystyle \frac{w_{x,1}x_1(T_1)+w_{x,2}x_2(T_2)}{w_{x,1}+w_{x,2}}},`$ (33) $`\mu _{\mathrm{C},\mathrm{nce}}`$ $`=`$ $`{\displaystyle \frac{w_{\mu ,1}\mu _1+w_{\mu ,2}\mu _2+w_{\mu ,0}\mu _0}{w_{\mu ,1}+w_{\mu ,2}+w_{\mu ,0}}},`$ (34) $`\mu _0`$ $`=`$ $`(x_2(T_2)x_1(T_1))/(T_2T_1).`$ (35) The weights $`w`$ are given by $`w_{x,1}`$ $`=`$ $`{\displaystyle \frac{1}{\epsilon _{x,1}^2}},`$ (36) $`w_{x,2}`$ $`=`$ $`{\displaystyle \frac{1}{\epsilon _{x,2}^2}},`$ (37) $`w_{\mu ,1}`$ $`=`$ $`{\displaystyle \frac{1}{\epsilon _{\mu ,1}^2}},`$ (38) $`w_{\mu ,2}`$ $`=`$ $`{\displaystyle \frac{1}{\epsilon _{\mu ,2}^2}},`$ (39) $`w_{\mu ,0}`$ $`=`$ $`{\displaystyle \frac{1}{\epsilon _{\mu ,0}^2}}={\displaystyle \frac{(T_2T_1)^2}{\epsilon _{x,1}^2+\epsilon _{x,2}^2}}.`$ (40) Inserting $`\gamma _{\mathrm{nce}}(t)`$ and $`B_{\mathrm{nce}}(t)`$ into Eq. (4), we obtain the mean error $`\epsilon _{x,\mathrm{p},\mathrm{nce}}(t)`$ of $`x_{\mathrm{p},\mathrm{nce}}`$. Using the form of $`x_{\mathrm{p},\mathrm{nce}}`$ as given in Eq. (31) and the auxiliary quantities with the index C, we find $$\epsilon _{x,\mathrm{p},\mathrm{nce}}^2(t)=\epsilon _{x,\mathrm{C},\mathrm{nce}}^2(T_{\mathrm{C},\mathrm{nce}})+\epsilon _{\mu ,\mathrm{C},\mathrm{nce}}^2(tT_{\mathrm{C},\mathrm{nce}})^2,$$ (41) with $$\epsilon _{x,\mathrm{C},\mathrm{nce}}^2(T_{\mathrm{C},\mathrm{nce}})=\frac{1}{w_{x,1}+w_{x,2}},$$ (42) $$\epsilon _{\mu ,\mathrm{C},\mathrm{nce}}^2=\frac{1}{w_{\mu ,1}+w_{\mu ,2}+w_{\mu ,0}}.$$ (43) A comparison of the Eqs. (31)–(42) with the analytic version of the single-star-mode solution of Paper I (Sect. 2, especially Eqs. (19), (23), (30)–(32)) shows that the prediction $`x_{\mathrm{p},\mathrm{nce}}(t)`$ and its mean error $`\epsilon _{x,\mathrm{p},\mathrm{nce}}(t)`$ are strictly identical with the single-star-mode solution $`x_{\mathrm{SI}}(t)`$ and its mean error $`\epsilon _{x,\mathrm{SI}}(t)`$. This result is very pleasing, since it proves the internal consistency of our scheme: In the limit of vanishing cosmic errors, the prediction $`x_\mathrm{p}(t)`$ according to Eq. (1) is asymptotically approaching the single-star mode solution $`x_{\mathrm{SI}}(t)`$. This result is not apriori self-evident, since our definitions for the ‘best’ solution for predicting $`x(t)`$ differ at least formally in Paper I and in this Paper II (i.e., Eqs. (5) and (35) of Paper I versus Eq. (3) of Paper II). ### 2.5 Motivation for introducing the long-term and short-term prediction The general solution of our combination problem is given in Sect. 2.2. for all epochs $`t`$. The solution $`x_\mathrm{p}(t)`$ is a non-linear function of $`t`$, and requires the knowledge of the correlation functions $`\xi (\mathrm{\Delta }t)`$ and $`\zeta (\mathrm{\Delta }t)`$ as functions of the epoch difference $`\mathrm{\Delta }t`$. At present, we do not have a well-established knowledge about the run of $`\xi (\mathrm{\Delta }t)`$ and $`\zeta (\mathrm{\Delta }t)`$. Only the cosmic errors $`c_\mu =(\eta (0))^{1/2}`$ and $`c_x=(\xi (0))^{1/2}`$ can be empirically determined from the comparison of the FK5 with HIPPARCOS, assuming that the FK5 is giving ‘mean’ quantities. Even if we would know the run of $`\xi (\mathrm{\Delta }t)`$ and $`\zeta (\mathrm{\Delta }t)`$ as a function of $`\mathrm{\Delta }t`$, the non-linearity of $`x_\mathrm{p}(t)`$ would demand a table of $`x_\mathrm{p}(t)`$ for a sequence of epochs, e.g. for each year, if the user should not have the burden to do the full calculation himself by running a program. We propose the following solution: The general solution for $`x_\mathrm{p}(t)`$ allows rather easily to obtain two limiting solutions for $`\mathrm{\Delta }t=tT_2\pm \mathrm{}`$ and for $`\mathrm{\Delta }t0`$. We call the solution for $`\mathrm{\Delta }t\pm \mathrm{}`$ the ‘long-term prediction (LTP)’, and the solution for $`\mathrm{\Delta }t0`$ the ‘short-term prediction (STP)’ around the epoch $`T_2`$. Both the LTP and STP solutions are linear in $`t`$. They can be therefore given in the usual astrometric style, i.e. as a position at a central epoch and a proper motion, together with their mean errors. The details on the LTP and STP solutions $`x_{\mathrm{LTP}}(t)`$ and $`x_{\mathrm{STP}}(t)`$ are given in the Sects. 3 and 4. The general solution $`x_\mathrm{p}(t)`$ is a smooth transition from short-term prediction $`x_{\mathrm{STP}}(t)`$ (for epochs around $`T_2`$) to the long-term prediction $`x_{\mathrm{LTP}}(t)`$ for epochs $`t`$ far away from $`T_2`$. In Sect. 5, we will discuss a convenient (but only approximately valid) method to carry out this transition, if we know the run of $`\zeta (\mathrm{\Delta }t)`$. This gives at least a rough indication on the process of transition as a function of the epoch difference $`\mathrm{\Delta }t`$, even if $`\zeta (\mathrm{\Delta }t)`$ is not well-established. If, in the future, $`\zeta (\mathrm{\Delta }t)`$ should be better determined, then our method would allow rather conveniently the (approximate) determination of $`x_\mathrm{p}(t)`$ also for epochs inbetween the validity ranges of $`x_{\mathrm{STP}}(t)`$ and $`x_{\mathrm{LTP}}(t)`$. ## 3 Long-term prediction (LTP) We consider in this section the limiting case of the general solution $`x_\mathrm{p}(t)`$ for $`\|tT_2\|\pm \mathrm{}`$. This ‘long-term prediction’ $`x_{\mathrm{LTP}}(t)`$ is valid for epochs not too close to the epoch $`T_2`$ of the instantaneous Catalogue 2 (i.e. here the HIPPARCOS Catalogue with $`T_21991.25`$). We assume that the epoch difference $`\mathrm{\Delta }t=tT_2`$ is so large that the correlation functions $`\xi (\mathrm{\Delta }t)`$ and $`\zeta (\mathrm{\Delta }t)`$ both vanish. Setting $`\xi (tT_2)=0`$ and $`\zeta (tT_2)=0`$, we obtain from the general Eqs. (22) and (23) for the LTP solution $`\gamma _{\mathrm{LTP}}(t)`$ $`=`$ $`(\epsilon _{x,1}^2(\epsilon _{\mu ,1}^2+[\epsilon _{\mu ,2}^2+\eta (0)])`$ (44) $`+\epsilon _{\mu ,1}^2[\epsilon _{\mu ,2}^2+\eta (0)](T_2T_1)(tT_1))/N,`$ $`B_{\mathrm{LTP}}(t)`$ $`=`$ $`\beta _{\mathrm{LTP}}(t)\left(tT_2\right)`$ (45) $`=`$ $`(\epsilon _{x,1}^2\epsilon _{\mu ,1}^2(T_2T_1)`$ $`+(\epsilon _{x,1}^2+[\epsilon _{x,2}^2+\xi (0)])\epsilon _{\mu ,1}^2(tT_1))/N,`$ where $`N`$ is still given by Eq. (24). If we compare the Eqs. (43), (44), and (24) for $`\gamma _{\mathrm{LTP}}`$, $`B_{\mathrm{LTP}}`$, and $`N=N_{\mathrm{LTP}}`$ with the corresponding Eqs. (28)–(30) for $`\gamma _{\mathrm{nce}}`$, $`B_{\mathrm{nce}}`$, and $`N_{\mathrm{nce}}`$, we find that they are identical if we replace in the equations for the nce solution the quantity $`\epsilon _{x,2}^2`$ by $$\epsilon _{x,2,\mathrm{LTP}}^2=\epsilon _{x,2}^2+\xi (0)=\epsilon _{x,2}^2+c_x^2,$$ (46) and $`\epsilon _{\mu ,2}^2`$ by $$\epsilon _{\mu ,2,\mathrm{LTP}}^2=\epsilon _{\mu ,2}^2+\eta (0)=\epsilon _{\mu ,2}^2+c_\mu ^2.$$ (47) This is very plausible, since the instantaneously measured quantities $`x_2(T_2)`$ and $`\mu _2(T_2)`$ are affected by the cosmic errors $`c_x`$ and $`c_\mu `$. If we add these cosmic errors to the corresponding measuring errors $`\epsilon _{x,2}`$ and $`\epsilon _{\mu ,2}`$, we obtain ‘apparent’ measuring errors $`\epsilon _{x,2,\mathrm{LTP}}`$ and $`\epsilon _{\mu ,2,\mathrm{LTP}}`$. Since the cosmic errors are not correlated with the measuring errors, the summation has to be done quadratically. Using this finding we obtain for the long-term prediction $`x_{\mathrm{LTP}}(t)`$ $`=`$ $`x_{\mathrm{LTP}}(T_{\mathrm{LTP}})+\mu _{\mathrm{LTP}}(tT_{\mathrm{LTP}}),`$ (48) with $`T_{\mathrm{LTP}}`$ $`=`$ $`{\displaystyle \frac{w_{x,1}T_1+w_{x,2,\mathrm{LTP}}T_2}{w_{x,1}+w_{x,2,\mathrm{LTP}}}},`$ (49) $`x_{\mathrm{LTP}}(T_{\mathrm{LTP}})`$ $`=`$ $`{\displaystyle \frac{w_{x,1}x_1(T_1)+w_{x,2,\mathrm{LTP}}x_2(T_2)}{w_{x,1}+w_{x,2,\mathrm{LTP}}}},`$ (50) $`\mu _{\mathrm{LTP}}`$ $`=`$ $`{\displaystyle \frac{w_{\mu ,1}\mu _1+w_{\mu ,2,\mathrm{LTP}}\mu _2+w_{\mu ,0,\mathrm{LTP}}\mu _0}{w_{\mu ,1}+w_{\mu ,2,\mathrm{LTP}}+w_{\mu ,0,\mathrm{LTP}}}}.`$ (51) The weights $`w`$ are given by Eqs. (36) and (38), and by $`w_{x,2,\mathrm{LTP}}`$ $`=`$ $`{\displaystyle \frac{1}{\epsilon _{x,2}^2+c_x^2}},`$ (52) $`w_{\mu ,2,\mathrm{LTP}}`$ $`=`$ $`{\displaystyle \frac{1}{\epsilon _{\mu ,2}^2+c_\mu ^2}},`$ (53) $`w_{\mu ,0,\mathrm{LTP}}`$ $`=`$ $`{\displaystyle \frac{(T_2T_1)^2}{\epsilon _{x,1}^2+\epsilon _{x,2}^2+c_x^2}}.`$ (54) The long-term prediction $`x_{\mathrm{LTP}}(t)`$ has in fact two conceptionally different properties: (1) As described above, it is the limit of the general prediction $`x_\mathrm{p}(t)`$ for the true instantaneous position $`x(t)`$ for $`t\pm \mathrm{}`$. (2) On the other hand, $`x_{\mathrm{LTP}}(t)`$ is for all epochs $`t`$ the best prediction for the mean position $`x_\mathrm{m}(t)`$ of the object. This means especially that $`\mu _{\mathrm{LTP}}`$ is the best estimate of the center-of-mass velocity of the object. The two different concepts produce two different error estimates $`\epsilon _{x,\mathrm{LTP}}(t)`$ for $`x_{\mathrm{LTP}}(t)`$. If we consider $`x_{\mathrm{LTP}}`$ as the prediction for the mean position $`x_\mathrm{m}(t)`$ then the mean error $`\epsilon _{x,\mathrm{LTP},\mathrm{m}}(t)`$ is given by: $$\epsilon _{x,\mathrm{LTP},\mathrm{m}}^2(t)=\epsilon _{x,\mathrm{LTP},\mathrm{m}}^2(T_{\mathrm{LTP}})+\epsilon _{\mu ,\mathrm{LTP},\mathrm{m}}^2(tT_{\mathrm{LTP}})^2,$$ (55) with $$\epsilon _{x,\mathrm{LTP},\mathrm{m}}^2(T_{\mathrm{LTP}})=\frac{1}{w_{x,1}^2+w_{x,2,\mathrm{LTP}}^2},$$ (56) $$\epsilon _{\mu ,\mathrm{LTP},\mathrm{m}}^2=\frac{1}{w_{\mu ,1}^2+w_{\mu ,2,\mathrm{LTP}}^2+w_{\mu ,0,\mathrm{LTP}}^2}.$$ (57) If we consider $`x_{\mathrm{LTP}}`$ as the prediction for the instantaneous position $`x(t)`$ for large values of $`\|\mathrm{\Delta }t\|=\|tT_2\|`$, then the uncertainty $`\epsilon _{x,\mathrm{LTP},\mathrm{ins}}(t)`$ of $`x_{\mathrm{LTP}}(t)`$ is given by $$\epsilon _{x,\mathrm{LTP},\mathrm{ins}}^2(t)=\epsilon _{x,\mathrm{LTP},\mathrm{m}}^2(t)+c_x^2.$$ (58) Similarly, the mean error $`\epsilon _{\mu ,\mathrm{LTP},\mathrm{ins}}`$ of the predicted instantaneous proper motion for large $`\mathrm{\Delta }t`$ is given by $$\epsilon _{\mu ,\mathrm{LTP},\mathrm{ins}}^2=\epsilon _{\mu ,\mathrm{LTP},\mathrm{m}}^2+c_\mu ^2.$$ (59) The equations above describe what we have called the ‘analytic’ approach in Paper I. For the ‘numerical’ approach, we can take over for the LTP solution the formulae given in Paper I for the single-star mode (SI) with the following changes: (1) We should not redetermine the parallax $`p`$ of the star. The determination of $`p`$ by HIPPARCOS requires instantaneous values of $`x`$ and $`\mu `$, not the ‘mean’ LTP values. Formally, we set all correlations between $`p`$ and the other quantities (position, proper motion) equal to zero. (2) In the variance-covariance matrix D we replace in the diagonal line the HIPPARCOS values of $`D_{\mathrm{H},11}=\epsilon _{\alpha ,\mathrm{H}}^2(T_\mathrm{H})`$ by $$D_{\mathrm{H},11,\mathrm{LTP}}=\epsilon _{\alpha ,\mathrm{H}}^2(T_\mathrm{H})+c_x^2,$$ (60) and $`\epsilon _{\delta ,\mathrm{H}}^2`$ by $`\epsilon _{\delta ,\mathrm{H}}^2+c_x^2`$, $`\epsilon _{\mu ,\alpha ,\mathrm{H}}^2`$ by $`\epsilon _{\mu ,\alpha }^2+c_\mu ^2`$, and $`\epsilon _{\mu ,\delta ,\mathrm{H}}^2`$ by $`\epsilon _{\mu ,\delta ,\mathrm{H}}^2+c_\mu ^2`$, respectively. No changes are made in the non-diagonal elements (except for decoupling the parallax as described above), i.e. all the covariances remain in the LTP as they were in the SI mode, e.g. $$D_{\mathrm{H},12,\mathrm{LTP}}=D_{\mathrm{H},12}=\epsilon _{\alpha ,\mathrm{H}}(T_\mathrm{H})\epsilon _{\delta ,\mathrm{H}}(T_\mathrm{H})\rho _{\alpha \delta ,\mathrm{H}}(T_\mathrm{H}).$$ (61) ## 4 Short-term prediction (STP) In this section, we consider the other limiting case of the general solution $`x_\mathrm{p}(t)`$, namely the limit for $`\mathrm{\Delta }t=tT_20`$. This ‘short-term prediction’ $`x_{\mathrm{STP}}(t)`$ is valid for epochs close to $`T_2`$ (in the case of using HIPPARCOS: $`T_21991.25`$). For the case $`\mathrm{\Delta }t0`$, we use for the correlation function $`\xi (\mathrm{\Delta }t)`$ and $`\zeta (\mathrm{\Delta }t)`$ Taylor series in $`\mathrm{\Delta }t`$, and keep only terms which are linear in $`\mathrm{\Delta }t`$. From the Eqs. (54) and (55) of WPSA, we obtain for small values of $`\mathrm{\Delta }t=tT_2`$ $`\xi (\mathrm{\Delta }t)`$ $``$ $`\xi (0),`$ (62) $`\zeta (\mathrm{\Delta }t)`$ $``$ $`\eta (0)\mathrm{\Delta }t.`$ (63) In Eq. (63), we have made use of the differential relation (28) of WPSA. Although HIPPARCOS values are already averaged over about 3 years of observation, the use of Eq. (63) is justified according to numerical investigations carried out by M. Biermann (1996, unpublished, see also WPSA, Sect. 3.4). Inserting Eqs. (62) and (63) into the Eqs. (22) and (23), we find for the short-term prediction $`\gamma _{\mathrm{STP}}(t)=((\xi (0)+\epsilon _{x,1}^2)(\eta (0)+\epsilon _{\mu ,1}^2+\epsilon _{\mu ,2}^2)`$ (64) $`+\epsilon _{\mu ,1}^2(T_2T_1)^2(\eta (0)+\epsilon _{\mu ,2}^2)`$ $`+\epsilon _{\mu ,1}^2\epsilon _{\mu ,2}^2(T_2T_1)(tT_2))/N,`$ $`B_{\mathrm{STP}}(t)=\beta _{\mathrm{STP}}(t)(tT_2)`$ (65) $`=(\epsilon _{\mu ,1}^2\epsilon _{x,2}^2(T_2T_1)`$ $`+(\eta (0)(\xi (0)+\epsilon _{x,1}^2+\epsilon _{x,2}^2+\epsilon _{\mu ,1}^2(T_2T_1)^2)`$ $`+\epsilon _{\mu ,1}^2(\xi (0)+\epsilon _{x,1}^2+\epsilon _{x,2}^2))(tT_2))/N,`$ where $`N`$ is still given by Eq. (24). Inserting $`\gamma _{\mathrm{STP}}(t)`$ and $`B_{\mathrm{STP}}(t)`$ into Eq. (1), we obtain $$x_{\mathrm{STP}}(t)=x_{\mathrm{STP}}(T_2)+\mu _{\mathrm{STP}}(tT_2),$$ (66) which is a linear function of the epoch $`t`$. We do not use here a ‘central’ epoch $`T_{\mathrm{STP}}`$, since our Taylor series for $`x_{\mathrm{STP}}(t)`$ is centered around $`t=T_2`$. The proper motion $`\mu _{\mathrm{STP}}`$ can be derived as the sum of the coefficients in front of $`t`$ in the Eqs. (1), (63), and (64): $`\mu _{\mathrm{STP}}=\mu _1`$ $`+\left(\left(x_2(T_2)x_1(T_1)\mu _1(T_2T_1)\right)\epsilon _{\mu ,1}^2\epsilon _{\mu ,2}^2(T_2T_1)\right)/N`$ $`+(\mu _2\mu _1)(\eta (0)(\xi (0)+\epsilon _{x,1}^2+\epsilon _{x,2}^2+\epsilon _{\mu ,1}^2(T_2T_1)^2)`$ (67) $`+\epsilon _{\mu ,1}^2(\xi (0)+\epsilon _{x,1}^2+\epsilon _{x,2}^2))/N.`$ After some algebra, $`\mu _{\mathrm{STP}}`$ can be rewritten as $$\mu _{\mathrm{STP}}=\frac{w_{\mu ,10,\mathrm{STP}}\mu _{10}+w_{\mu ,2}\mu _2}{w_{\mu ,10,\mathrm{STP}}+w_{\mu ,2}}.$$ (68) The ‘combined’ mean proper motion $`\mu _{10}`$ is derived from $`\mu _1`$ and $`\mu _0`$ by $$\mu _{10}=\frac{w_{\mu ,1}\mu _1+w_{\mu ,0,\mathrm{LTP}}\mu _0}{w_{\mu ,1}+w_{\mu ,0,\mathrm{LTP}}}.$$ (69) $`\mu _0,w_{\mu ,1}`$, and $`w_{\mu ,2}`$ are given by the corresponding equations in Sect. 2.4, and $`w_{\mu ,0,\mathrm{LTP}}`$ by Eq. (53). The weight $`w_{\mu ,10,\mathrm{STP}}`$ of $`\mu _{10}`$ in the STP solution is given by $$w_{\mu ,10,\mathrm{STP}}=\frac{1}{(1/w_{\mu ,10,\mathrm{LTP}})+c_\mu ^2}=\frac{1}{\epsilon _{\mu ,10,\mathrm{LTP}}^2+c_\mu ^2},$$ (70) with $$w_{\mu ,10,\mathrm{LTP}}=w_{\mu ,1}+w_{\mu ,0,\mathrm{LTP}}=\frac{1}{\epsilon _{\mu ,10,\mathrm{LTP}}^2}.$$ (71) The meaning of Eqs. (68)-(71) is the following: $`\mu _{\mathrm{STP}}`$ is the weighted mean of $`\mu _2`$ and the combined mean proper motion $`\mu _{10}`$, where the cosmic error $`c_\mu `$ in $`\mu _{10}`$ has to be taken into account in $`w_{\mu ,10,\mathrm{STP}}`$. The combined mean proper motion $`\mu _{10}`$ is itself a weighted mean of $`\mu _1`$ and $`\mu _0`$, where for $`\mu _0`$ the cosmic error $`c_x`$ in $`x_2(T_2)`$ has to be included into $`w_{\mu ,0,\mathrm{LTP}}`$. The exact expression for $`x_{\mathrm{STP}}(T_2)`$ is given by $`x_{\mathrm{STP}}(T_2)=x_2(T_2)`$ $`\epsilon _{x,2}^2((\eta (0)+\epsilon _{\mu ,1}^2+\epsilon _{\mu ,2}^2)(x_2(T_2)x_1(T_1)\mu _1(T_2T_1))`$ $`\epsilon _{\mu ,1}^2(T_2T_1)(\mu _2(T_2)\mu _1))/N.`$ (72) For all practical purposes we can, however, neglect the last term in Eq. (72), which is proportional to $`\epsilon _{x,2}^2`$. The measuring accuracy $`\epsilon _{x,2}`$ for the HIPPARCOS positions $`x_2(T_2)`$ at $`T_21991.25`$ is so high, relative to the other measuring errors and cosmic errors, that this term is usually of the order of 0.01 mas only. We use instead of Eq. (72) the very good approximation $$x_{\mathrm{STP}}(T_2)=x_2(T_2).$$ (73) The mean errors of $`x_{\mathrm{STP}}(T_2)`$ and of $`\mu _{\mathrm{STP}}`$ are given by $$\epsilon _{x,\mathrm{STP}}^2(T_2)=\epsilon _{x,2}^2(T_2),$$ (74) $$\epsilon _{\mu ,\mathrm{STP}}^2=\frac{1}{w_{\mu ,10,\mathrm{STP}}+w_{\mu ,2}}.$$ (75) The full uncertainty $`\epsilon _{x,\mathrm{STP}}(t)`$ of $`x_{\mathrm{STP}}(t)`$ is given by $`\epsilon _{x,\mathrm{STP}}^2(t)=\epsilon _{x,\mathrm{STP}}^2(T_2)+\epsilon _{\mu ,\mathrm{STP}}^2(tT_2)^2`$ (76) $`+{\displaystyle \frac{1}{4}}\xi _0^{(IV)}(tT_2)^4.`$ In Eq. (76), we have neglected the very small correlation between $`x_{\mathrm{STP}}(T_2)=x_2(T_2)`$ and $`\mu _{\mathrm{STP}}`$ which is caused by the use of $`x_2(T_2)`$ in deriving $`\mu _0`$ which in turn enters into $`\mu _{\mathrm{STP}}`$. The last term in Eq. (76) is the statistical uncertainty of a prediction based on instantaneous data at $`T_2`$ for small epoch differences $`\mathrm{\Delta }t`$ (see Eq. (57) of WPSA). $`\xi _0^{(IV)}`$ is the fourth derivative of the correlation function $`\xi (\mathrm{\Delta }t)`$ with respect to $`\mathrm{\Delta }t`$ at $`\mathrm{\Delta }t=0`$. In the numerical approach for deriving the short-term prediction, we modify the procedure for the single-star mode presented in Paper I more strongly than for the LTP solution. As ‘observations’ b we use first b$`{}_{\mathrm{H}}{}^{}(T_\mathrm{H})`$, as in Paper I, the parallax $`p_\mathrm{H}`$ inclusive. The corresponding part D<sub>H</sub> of the variance-covariance matrix D remains also unchanged. The second part of b, which we now call b<sub>m</sub> is given by the $`\alpha _{}`$ and $`\delta `$ components of the combined mean proper motion $`\mu _{10}`$. The quantity $`\mu _{10}`$ is obtained in a preparatory step from Eq. (69), and its mean error from Eq. (70) as $`(w_{\mu ,10,\mathrm{STP}})^{1/2}`$. Each component of $`\mu _{10}`$, i.e. of b<sub>m</sub>, is considered not to be correlated with any other component of b. The numerical approach for STP produces values for $`x_{\mathrm{STP}}`$ and $`\mu _{\mathrm{STP}}`$ (in $`\alpha _{}`$ and $`\delta `$), a new parallax $`p_{\mathrm{STP}}`$, and the corresponding variance-covariance matrix. In presenting the results in printed form, we use again central epochs $`T_{\mathrm{STP}}`$ (different for $`\alpha _{}`$ and $`\delta `$), at which $`x_{\mathrm{STP}}`$ and $`\mu _{\mathrm{STP}}`$ are uncorrelated. However, $`T_{\mathrm{STP}}`$ differs usually only very slightly from the individual central epochs of the basic HIPPARCOS data. ## 5 Transition from the short-term prediction to the long-term prediction ### 5.1 Transition in position As discussed in Sect. 2.4 and illustrated in Fig. 2, the general solution $`x_\mathrm{p}(t)`$ is a smooth transition from the short-term prediction $`x_{\mathrm{STP}}(t)`$ for epochs $`t`$ close to $`T_2`$ to the long-term prediction $`x_{\mathrm{LTP}}(t)`$ for $`t\pm \mathrm{}`$. We are now asking for the ‘transition function’ $`\beta _{\mathrm{trans}}(t)`$ which describes this transition in $`x`$: $$x_\mathrm{p}(t)=(1\beta _{\mathrm{trans}}(t))x_{\mathrm{LTP}}(t)+\beta _{\mathrm{trans}}(t)x_{\mathrm{STP}}(t).$$ (77) Formally we can always solve Eq. (77) for $`\beta _{\mathrm{trans}}(t)`$, using our former results for $`x_\mathrm{p}(t),x_{\mathrm{LTP}}(t),`$ and $`x_{\mathrm{STP}}(t)`$. However, the resulting transition function $`\beta _{\mathrm{trans}}(t)`$ is then very complicated and depends unfortunately not only on the correlation functions, but also explicitely on the measured values of $`x_1,x_2,\mu _1,\mu _2`$, and on their measuring errors. Such a result is not very suitable for a practical application. There does exist, however, an approximate treatment for the transition function $`\beta _{\mathrm{trans}}(t)`$ which gives a very simple and easily applicable result, and which nevertheless describes the transition quantitatively rather accurately. The basic idea is the observation that in real applications the mean position $`x_1(T_1)`$ enters into the final result $`x_\mathrm{p}(t)`$ mainly through the proper motion $`\mu _0`$. This is caused by the small error of the HIPPARCOS position $`x_2(T_2)`$ with respect to the error of the mean (FK5 or GC) position $`x_1(T_1)`$. Only in cases of a large cosmic error $`c_x(p)`$ in $`x_2(T_2)`$, our approximation becomes less accurate. We therefore consider the transition function $`\beta _{\mathrm{trans}}(t)`$ for the limiting case in which $`\epsilon _{x,1}`$ tends towards infinity while $`\epsilon _{\mu ,0}`$ remains finite (equal to its actual value). The latter can be enforced by setting $`T_2T_1=\epsilon _{x,1}/\epsilon _{\mu ,0}`$. This means that we let go $`T_1\mathrm{}`$ and $`\epsilon _{x,1}+\mathrm{}`$ in such a way that $`\epsilon _{\mu ,0}`$ remains constant. If we use this special case as an approximation, we derive after some lengthy algebra the following rather simple expression for the transition function: $$\beta _{\mathrm{trans}}(t)=\frac{\zeta (tT_2)}{\eta (0)(tT_2)}.$$ (78) A similar function has already been derived as Eq. (70) in Sect. 4.2.4 of WPSA. The function $`\beta _{\mathrm{trans}}`$ has the welcome property that it does neither depend on the measured values of $`x`$ and $`\mu `$ nor on the mean errors of a given star. These values are fully absorbed in the individual solutions $`x_{\mathrm{STP}}(t)`$ and $`x_{\mathrm{LTP}}(t)`$. Hence $`\beta _{\mathrm{trans}}(t)`$ is the same function for all the objects. We should remark here that the apparently more complicated form of the transition function according to Eq. (75) of WPSA is caused by the fact that we have used in WPSA a slightly different definition of the long and short-term prediction, namely $`x_2(t)`$ instead of $`x_{\mathrm{STP}}(t)`$ and $`x_1(t)`$ instead of $`x_{\mathrm{LTP}}(t)`$. The function $`\beta _{\mathrm{trans}}(t)`$ has the desired properties in the limits $`\mathrm{\Delta }t0`$ or $`\mathrm{}`$. For $`t=T_2`$, we have $`\beta _{\mathrm{trans}}(T_2)=1`$, because of $`\zeta (\mathrm{\Delta }t)\eta (0)\mathrm{\Delta }t`$ for small values of $`\mathrm{\Delta }t`$. For $`t\mathrm{},\zeta (tT_2)`$ vanishes and hence $`\beta _{\mathrm{trans}}(\mathrm{})=0`$. An example for the full run of $`\beta _{\mathrm{trans}}(tT_2)`$, from 1 to 0, is illustrated by the full curve shown in Fig. 8 of WPSA. Table 1 gives $`\beta _{\mathrm{trans}}(tT_2)`$ for a few values of $`\|tT_2\|`$, using the example given in Sect. 3.6 of WPSA for the correlation function $`\zeta (\mathrm{\Delta }t)`$. At $`\|tT_2\|`$ 5-6 years, the general solution $`x_\mathrm{p}(t)`$ is about half-way between $`x_{\mathrm{STP}}(t)`$ and $`x_{\mathrm{LTP}}(t)`$. From both the Fig. 8 of WPSA and Table 1, we get an indication for the range of applicability of the short and long-term prediction. The short-term prediction $`x_{\mathrm{STP}}(t)`$ has a rather limited range of applicability, namely a few years around $`T_2`$ only. The long-term prediction $`x_{\mathrm{LTP}}(t)`$ is a good approximation for the general solution $`x_\mathrm{p}(t)`$ for epoch differences $`\|tT_2\|`$ which are larger than about 10 years. Hence the transition from STP to LTP is rather rapid, at least in our example. We should remark here that (by chance) the transition function $`\beta _{\mathrm{trans}}(t)`$ given by Eq. (78) is even strictly valid (exact) for the example used for the correlation functions in Sect. 3.6 of WPSA and adopted in Table 1. For more general runs of the correlation functions, $`\beta _{\mathrm{trans}}`$ can become quite large at epochs around the ‘crossing time’ at which $`x_{\mathrm{LTP}}=x_{\mathrm{STP}}`$, if $`x_{\mathrm{LTP}}`$ and $`x_{\mathrm{STP}}`$ are significantly different from $`x_\mathrm{p}`$ at that time. ### 5.2 Transition in proper motion The most accurate prediction $`\mu _\mathrm{p}(t)`$ for the instantaneous proper motion $`\mu (t)`$ at an arbitrary epoch $`t`$ is formally given by $$\mu _\mathrm{p}(t)=\dot{x}_\mathrm{p}(t)=(1\nu _{\mathrm{trans}}(t))\mu _{\mathrm{LTP}}+\nu _{\mathrm{trans}}(t)\mu _{\mathrm{STP}}.$$ (79) The transition function $`\nu _{\mathrm{trans}}(t)`$ for the proper motion $`\mu `$ is, similar to $`\beta _{\mathrm{trans}}`$, in general a rather complicated function, which depends on the measured values of $`\mu _1,\mu _2,`$ and $`\mu _0`$. If we adopt the same approximation as used in Sect. 5.1 for obtaining Eq. (78) for $`\beta _{\mathrm{trans}}`$, namely $`\epsilon _{x,1}\mathrm{}`$ while keeping $`\epsilon _{\mu ,0}`$ constant, we derive for $`\nu _{\mathrm{trans}}(t)`$ also a very simple expression: $$\nu _{\mathrm{trans}}(t)=\beta _{\mathrm{trans}}(t)+\dot{\beta }_{\mathrm{trans}}(t)(tT_{\mathrm{STP}})=\frac{\eta (tT_{\mathrm{STP}})}{\eta (0)}.$$ (80) This transition function $`\nu _{\mathrm{trans}}(t)`$ for $`\mu (t)`$ has the same proper limiting values as $`\beta _{\mathrm{trans}}(t)`$: For $`t\mathrm{}`$, we have $`\nu _{\mathrm{trans}}(\mathrm{})=0`$ and hence $`\mu _\mathrm{p}(\mathrm{})=\mu _{\mathrm{LTP}}`$, and for $`t=T_{\mathrm{STP}}`$ we obtain $`\nu _{\mathrm{trans}}(T_{\mathrm{STP}})=1`$ and $`\mu _\mathrm{p}(T_{\mathrm{STP}})=\mu _{\mathrm{STP}}`$ . In Table 1 we list $`\nu _{\mathrm{trans}}(\mathrm{\Delta }t)`$ for a few values of $`\|\mathrm{\Delta }t\|=\|tT_{\mathrm{STP}}\|`$, using Eq. (80) and the simple example for $`\eta (\mathrm{\Delta }t)`$ given in Sect. 3.6 of WPSA. The transition functions $`\beta _{\mathrm{trans}}(t)`$ and $`\nu _{\mathrm{trans}}(t)`$ according to Eqs. (78) and (80) have another nice property: If we use the example given in Sect. 3.6 of WPSA, then the function $`\beta _{\mathrm{trans}}`$ depends only on $`tT_2`$, but not on the individual cosmic error $`c_\mu (p)=(\eta (0,p))^{1/2}`$, since $`\eta (0)`$ occurs also as a factor in $`\zeta (tT_2)`$ and cancels out in Eq. (78). Similarly, the transition function $`\nu _{\mathrm{trans}}(t)`$ is a ‘scaled’ function and depends on $`tT_2`$ only, but not on $`\eta (0,p)`$. ## 6 Cosmic errors For deriving the solutions in the LTP and STP mode, we need to know the cosmic errors $`c_\mu =(\eta (0))^{1/2}`$ and $`c_x=(\xi (0))^{1/2}`$. The cosmic errors depend strongly on the distance $`r`$ of the star from the Sun, or equivalently on the stellar parallax $`p`$. For 1202 ‘apparently single stars’ from the basic FK5, we have obtained the functions $`c_\mu (p)`$ and $`c_x(p)`$ empirically by using groups of stars in various distance intervals. The data can be represented by the following fit functions: $`c_\mu (p)`$ $`=(\eta (0,p))^{1/2}`$ $`=\left({\displaystyle \frac{C_1p}{(C_2^2+p^2)^{1/2}}}\right)^{1/2},`$ (81) $`c_x(p)`$ $`=(\xi (0,p))^{1/2}`$ $`=C_3c_\mu (p),`$ (82) with $`C_1`$ $`=`$ $`9.30(\mathrm{mas}/\mathrm{year})^2,`$ (83) $`C_2`$ $`=`$ $`22.14\mathrm{mas},`$ (84) $`C_3`$ $`=`$ $`5.93\mathrm{years}.`$ (85) These versions of $`c_\mu (p)`$ and $`c_x(p)`$ have been used for the FK6 (Wielen et al. 1999d, 2000a). A table for $`c_\mu `$ and $`c_x`$ as a function of $`p`$ or $`r`$, based on Eqs. (81)-(85), is given on page 12 of Wielen et al. (1999d). In the future we hope to determine also the dependence of the cosmic errors on the absolute magnitude (or mass) of the stars for a given parallax. A comparison of the results for the cosmic errors based on the FK5 stars with those derived from the (on average fainter) GC stars seems to indicate a rather weak dependence on the brightness of the stars (Wielen et al. 1998). ## 7 An example: $`\alpha `$ Ari In order to illustrate our combination method for all the three modes (SI, LTP, STP), we give in Table 2 the results for one individual star. As in Paper I (Sect. 5), we give the positions $`x(t)`$ and the proper motions $`\mu (t)`$ always relative to the HIPPARCOS solution as $`\mathrm{\Delta }x(t)=x(t)x_\mathrm{H}(t)`$ and $`\mathrm{\Delta }\mu (t)=\mu (t)\mu _\mathrm{H}(t)`$, in order to save printing space and to make the comparison of the results easier. The ground-based data are always reduced to the HIPPARCOS system. The results of the single-star mode presented here in Paper II differ somewhat from those of Paper I, because we now adopt slightly improved systematic differences FK5-HIP and GC-HIP. For a valid comparison of the results of the three different modes it is necessary to use exactly the same basic input data. Table 2 shows that the short-term prediction is usually quite close to the HIPPARCOS solution. On the other hand, the long-term prediction differs most strongly from the HIPPARCOS solution, since the HIPPARCOS values are entering with a lower weight into the LTP mode than in the SI mode, because of the cosmic errors in the HIPPARCOS data. The central epoch $`T_{\mathrm{LTP}}`$ is the only one which is usually significantly earlier than $`T_\mathrm{H}1991.25`$, and the mean error of the central position $`x_{\mathrm{LTP}}(T_{\mathrm{LTP}})`$ is typically only slightly smaller than the cosmic error $`c_x(p)`$. ## 8 Error budget In Table 3 we present the error budget of proper motions in the three different modes (SI, LTP, STP) for two samples of basic FK5 stars. The mean errors $`\epsilon _\mu `$ given in Table 3 refer to one ‘mean’ coordinate component. It is obtained as an rms average over $`\epsilon _{\mu ,\alpha }`$ and $`\epsilon _{\mu ,\delta }`$, and over all the stars in the corresponding sample. The error budget for the 1535 basic FK5 stars is slightly fictious, since this sample contains double stars for which the FK6 provides in reality ‘special’ solutions instead of the ‘direct’ combination solutions discussed in this paper. Nevertheless, the results for this sample provide a valid indication for the overall accuracy of our combination method in the direct modes SI, LTP, and STP. The sample of 1202 basic FK5 stars contains ‘apparently single objects’ only. Most of these stars (878 objects) have direct solutions in the FK6. The error budget for these 878 basic FK5 stars in Part I of the FK6 is given in Wielen et al. (1999d). The error budgets for 3272 additional fundamental stars with direct solutions in the three modes are presented in Part III of the FK6 (Wielen et al. 2000a). Table 4 gives the error budget for the combination of the GC (Boss et al. 1937) with HIPPARCOS. The sample of GC stars are the ‘full sample’ of 29 717 GC stars observed by HIPPARCOS, and the ‘subsample’ of 11 737 GC stars with linear HIPPARCOS standard solutions. From Table 4 it is obvious that the original proper motions $`\mu _{\mathrm{GC}}`$ of the GC do not contribute significantly to the GC+HIP on average. However, for some brighter and well-observed stars, the accuracy of $`\mu _{\mathrm{GC}}`$ is much better than the rms value of $`\epsilon _{\mu ,\mathrm{GC}}`$ seen in Table 4. Furthermore, the proper motion $`\mu _{0(\mathrm{GC})\mathrm{H}}`$, derived from the central positions of the GC and the HIPPARCOS Catalogue, has usually a quite reasonable accuracy for most of the GC stars. The typical gain in accuracy in the proper motions derived in the long-term prediction mode, relative to the original HIPPARCOS proper motions, is a factor of 4.6 for the basic fundamental stars in the FK6 = FK5+HIP, and a factor of 1.8 for the 11 773 GC stars in the GC+HIP. This improvement in the LTP mode with respect to HIPPARCOS is a consequence of the cosmic errors in the instantaneously measured HIPPARCOS proper motions. In contrast to the LTP mode, the short-term predictions (STP mode) do not differ so much from the HIPPARCOS solutions. However, in most cases we are more interested in the long-term averaged proper motion (LTP) or in the single-mode result (SI), where the gain in accuracy is quite significant. ## 9 Problems and applications The long-term predictions and the short-term predictions derived in the former sections are statistically the best astrometric solution for a sample of ‘apparently single-stars’. However, an unsatisfactory property of these statistically valid solutions is the fact that we are not able to make proper use of the available information on the individual behaviour of the stars. For example, we use the overall cosmic errors $`c_\mu (p)`$ and $`c_x(p)`$, no matter whether the object is a $`\mathrm{\Delta }\mu `$ binary or a single-star candidate (Wielen et al. 1999a). In principle, for each individual star, one would have to use ‘conditioned correlation functions’ which are based on the individually observed differences between the instantaneous measurements and the mean data (e.g. on $`\mu _{\mathrm{FK5}}\mu _{\mathrm{HIP}}`$). Unfortunately, the conditioned correlation functions, i.e. ‘conditioned cosmic errors’ in particular, are not available at present. We should also point out that presently the cosmic errors $`c_x`$ in position are much more uncertain than the cosmic errors $`c_\mu `$ in proper motion. While the typical values of $`c_\mu `$ are larger than the measuring errors $`\epsilon _{\mu ,\mathrm{FK5}}`$, $`\epsilon _{\mu ,0}`$, and $`\epsilon _{\mu ,\mathrm{HIP}}`$, the typical values of $`c_x`$ are nearly lost in the measuring errors of the ground-based data. Furthermore, our knowledge about the actual form of the correlation functions $`\xi (\mathrm{\Delta }t)`$, $`\eta (\mathrm{\Delta }t)`$, $`\zeta (\mathrm{\Delta }t)`$ for $`\mathrm{\Delta }t=\|\mathrm{\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}0}`$ is presently still quite rudimentary. The simple example for the correlation functions presented in Sect. 3.6 of WPSA is mainly given for illustrating the general behaviour of these functions. The application of this example to real data should be done very cautiously. Hence the real transition from the STP solution to the LTP solution (discussed in Sect. 5) is quantitatively not well-determined. Which of the three solutions offered (SI, LTP, STP) should be used in real applications ? For single-star candidates, the single-star mode should be adopted, although some of the single-star candidates may nevertheless be binaries. If the user is handling a sample of ‘apparently single stars’ (which usually contains single-star candidates, $`\mathrm{\Delta }\mu `$ binaries, and intermediate cases), then the LTP or STP solutions are recommended, depending on the corresponding epoch difference $`\mathrm{\Delta }t=tT_\mathrm{H}`$ (with $`T_\mathrm{H}1991`$). The problems discussed above are especially severe for objects detected as $`\mathrm{\Delta }\mu `$ binaries (Wielen et al. 1999a). For $`\mathrm{\Delta }\mu `$ binaries the difference between the instantaneous proper motion ($`\mu _{\mathrm{HIP}}`$) and the mean one (e.g., $`\mu _{\mathrm{FK5}}`$ or $`\mu _0`$) is sometimes much larger than the cosmic error $`c_\mu `$ expected on average. In such a case the weight $`w_{\mu ,2,\mathrm{LTP}}`$ (Eq. (53)) of the HIPPARCOS proper motion ($`\mu _{\mathrm{HIP}}=\mu _2`$) is higher than appropriate for this individual object. The derived mean proper motion $`\mu _{\mathrm{LTP}}`$ is then biased towards the HIPPARCOS proper motion, and the derived mean error $`\epsilon _{\mu ,\mathrm{LTP}}`$ is too small. In the case of extreme $`\mathrm{\Delta }\mu `$ binaries, it is better to adopt a properly weighted mean of $`\mu _1`$ (e.g. $`\mu _{\mathrm{FK5}}`$) and $`\mu _0`$ as a prediction for the mean proper motion $`\mu _\mathrm{m}`$. The main problem in such a procedure is the unknown individual value of the cosmic error in the HIPPARCOS position, which is higher than the value of $`c_x`$ expected on average. This value enters into the weight of $`\mu _0`$, and hence into the predicted value of $`\mu _\mathrm{m}`$. We shall discuss this problem in more detail in a subsequent paper. In any case, the LTP solutions (and, of course, the SI solutions) for $`\mathrm{\Delta }\mu `$ binaries are inherently the least accurate ones among the class of direct solutions, because of the disturbing double-star nature of these objects. ## 10 Summary and outlook In this Paper II, we have derived and discussed an appropriate method to combine a ground-based astrometric catalogue (such as the FK5 or GC) with the HIPPARCOS Catalogue, taking cosmic errors (due to undetected binaries) in the quasi-instantaneously measured data into account. The method leads to long-term predictions (LTP mode) and to short-term predictions (STP mode), which are the limiting cases of the general solution. The general solution is a smooth transition from the STP to the LTP mode. The case of single stars with no cosmic errors was already treated in Paper I (SI mode). In a subsequent paper, we shall present ‘special solutions’ for known double stars. x
warning/0007/cond-mat0007391.html	ar5iv	text	# Phenomenological model for magnetotransport in a multi-orbital system \[ ## Abstract By means of the Boltzmann equation, we have calculated some magnetotransport quantities for the layered multi-orbital compound Sr<sub>2</sub>RuO<sub>4</sub>. The Hall coefficient, the magnetoresistance and the in-plane resistivity have been determined taking into account the Fermi surface curvature and different time collisions for the electrons in the $`t_{2g}`$ bands. A consistent explanation of the experimental results has been obtained assuming different relaxation rates for the in-plane transport with and without an applied magnetic field, respectively. \] The layered perovskite oxide Sr<sub>2</sub>RuO<sub>4</sub> has attracted considerable experimental and theoretical attention since the recent discovery of superconductivity in this compound. Quantitative similarities between the Fermi liquid in <sup>3</sup>He and that in Sr<sub>2</sub>RuO<sub>4</sub> hint at the possibility of p-wave superconducting pairing. A growing body of experimental evidence, including results obtained from muon spin relaxation, NMR 1/T<sub>1</sub> and Knight shift, neutron scattering, impurity effect, and specific heat measurements, has shown that the pairing symmetry is unconventional, most likely p-wave. This compound shows also interesting normal state transport properties characterized by multi-bands effects, and provides an ideal opportunity to investigate the crossover from standard to non standard conduction processes. The in plane zero-field resistivity exhibits a Fermi-liquid like behaviour up to T$`{}_{FL}{}^{}`$25 K; above this temperature $`\rho _{ab}`$ rises monotonically with a curvature which is much weaker than T<sup>2</sup>. Moreover, around 1300 K the in-plane mean free path falls smoothly to less than $`1\AA `$ with no sign of resistivity saturation at the Mott-Ioffe-Regel limit, clearly indicating an unconventional high temperature conducting mechanism. $`\rho _c`$ reaches a maximum at 130 K then slowly decreases with increasing temperature. Below $``$ T<sub>FL</sub> the out-of -plane resistivity has a quadratic temperature dependence indicating that a crossover to a Fermi-liquid state takes place below 25K prior to the superconducting transition. Thus, below T<sub>FL</sub> Sr<sub>2</sub>RuO<sub>4</sub> has a resistivity with a quadratic temperature dependence in all directions and it is extremely anisotropic, with $`\rho _c/`$ $`\rho _{ab}>`$500, a ratio that is observed in most of the cuprate materials. The Hall coefficient exhibits a complicated temperature behavior: below 1 K it becomes almost temperature independent assuming a value of -1.15$`\times `$10<sup>-10</sup>m<sup>3</sup>/C; it has a strong temperature dependence below 25 K then, changes sign at temperatures approximately above $`30`$ K and shows a return to negative values above approximately $`130`$ K. The $`c`$axis magnetoresistance is large and positive and varies linearly with the applied magnetic field; with the increase of the temperature it falls sharply becoming negative above $`75`$ K. The in-plane magnetoresistance is positive and large at low temperatures, then decreases as T is raised up to $`80`$ K. The aim of the present work is to show, by means of the Boltzmann equation, that an overall agreement for the in-plane magnetotransport quantities can be reproduced in a temperature range up to 150 K, where the relaxation time approximation is applicable, i.e. the mean free path is greater than the lattice spacings. In order to determine the magnetotransport quantities, one has to calculate the conductivity tensor up to second order and this in turn implies the knowledge of the energy band spectrum as well as the explicit expression for the relaxation rate. The information of the electronic structure is given in Ref. for the three bands crossing the Fermi level and the loss of coherence of the conduction electrons is simulated by assuming an ad hoc temperature behavior for the scattering rates for electrons belonging to different bands. More precisely, the key assumption is that the temperature behaviour of the scattering rate of the carriers is modified by the presence of an applied magnetic field. We argue that the relaxation rate produced by spin fluctuations is suppressed in presence of an external magnetic field implying that the main contribution comes only from the usual electron-electron, impurity, or phonon- mediated scattering, which sum up to be $``$ T<sup>2</sup> at low temperatures. This effect applies in determining the Hall- and magneto- resistance. On the other hand, in the absence of magnetic field, according to the Matthiessen rule, the induced spin-fluctuation scattering rate adds to the previous relaxation mechanisms in determining the total resistivity of the system. The presence of an additional scattering rate is supported by the experimental evidence that the Sr<sub>2</sub>RuO<sub>4</sub> is close to magnetic instabilities which are the source of enhanced spin fluctuations at different $`q`$ points of the Brillouin zone, thus giving rise to large scattering amplitudes between charge carriers and the spin fluctuations themselves. <sup>17</sup>O NMR measurements probe spin correlations in Ru d<sub>xy</sub> and d<sub>xz,yz</sub> orbitals separately and show that only $`\chi _{xy}`$ increases monotonically with decreasing temperature down to about 40K, following a Curie-like behaviour, then turns over and tends to level below T<sub>FL</sub>, implying that the spin correlations in the d<sub>xy</sub> band are predominantly ferromagnetic in origin. Furthermore, by comparing $`{}_{}{}^{101}1/T_1T`$ at the Ru site and $`{}_{}{}^{17}1/T_1T`$ at the planar O site, due to the different dependence of their hyperfine form factor in k-space, one can probe whether the in-plane spin correlations are ferromagnetic or antiferromagnetic. It turns out that both $`{}_{}{}^{101}1/T_1T`$ and $`{}_{}{}^{17}1/T_1T`$ increase monotonically down to $`T_{FL}`$, and almost saturate in a Korringa-like behaviour. Moreover, inelastic neutron scattering measurements in the normal state reveal the existence of incommensurate magnetic spin fluctuations located at $`𝐪_\mathrm{𝟎}=(\pm 0.6\pi /a,\pm 0.6\pi /a,0)`$ due to the pronounced nesting properties of the almost one-dimensional $`d_{xz,yz}`$ bands. In fact, the 1D sheets can be schematically described by parallel planes separated by $`\stackrel{ˇ}{q}=\pm 2\pi /3a`$, running both in the $`x`$ and in the $`y`$ directions which give rise to dynamical nesting effects at the wave vectors $`𝐤=(\stackrel{ˇ}{q},k_y)`$, $`𝐤=(k_x,\stackrel{ˇ}{q})`$, and in particular at $`\stackrel{ˇ}{𝐪}=(\stackrel{ˇ}{q},\stackrel{ˇ}{q})`$. Though not as evident as in the $`\chi _{xy}`$ and the NMR measurements, due to the presence of few experimental points, the temperature dependence of the imaginary part of the susceptibility $`\chi ^{\prime \prime }`$, at low energy and at wave vector $`𝐪_\mathrm{𝟎}`$, exhibits a sharp decrease upon temperature increase above a temperature of $``$ T<sub>FL</sub>. Combining the results of NMR, nuclear spin-lattice relaxation rate, and inelastic neutron scattering measurements, it is possible to draw the following physical picture for the Sr<sub>2</sub>RuO<sub>4</sub>. There is a strong enhancement of spin fluctuations above T<sub>FL</sub>, mainly due to ferromagnetic correlations between the electrons in the $`d_{xy}`$ band, as revealed from Knight shifts experiments and nuclear spin relaxation rate, and due to incommensurate contributions from the nesting properties of the almost 1D $`d_{xz,yz}`$ bands. Hence, the main features of the magnetic response turn out to be decoupled for the electrons in the $`d_{xy}`$ and in the $`d_{xz,yz}`$ bands respectively. We consider that the same decoupling manifests in the transport properties. One important experimental observation is that there is a close relation between the change in the magnetic response and the resistivity measurements. Below $``$T<sub>FL</sub> the in-plane resistivity have a quadratic temperature dependence, while above $``$ T<sub>FL</sub> a superlinear term adds to the T<sup>2</sup> contribution . We assume that above $``$ T<sub>FL</sub> the linear contribution in the in-plane resistivity is mainly determined by small momentum scattering in the d<sub>xy</sub> band derived by the ferromagnetic spin fluctuations, and we consider also the slow temperature variation of the relaxation rates ($`T\mathrm{ln}T`$) which comes from the scattering at large momentum transfer due to the incommensurate spin fluctuations for the d<sub>xz,yz</sub> electrons. It is worthwhile pointing out that according to our calculations, the latter does not give any substantial qualitative and quantitative change in the transport properties, mainly due to the limited phase space allowed for the scattering processes. Specifying these considerations to the multi-band system in question, we assume only for the $`\gamma `$-band, a scattering rate proportional to T<sup>2</sup> in calculating the Hall coefficient and the magnetoresistance, with an additional term reproducing the effects of scattering by ferromagnetic spin fluctuations, to determine the zero-field resistivity. For a 2D system, the temperature dependence of the scattering rate due to this mechanism is given by $`\tau _{sf}^1{\displaystyle q^2𝑑qIm\chi (q,\omega )\frac{}{T}n(\omega )},`$ where $`n(\omega )`$ is the Bose distribution and $`\chi (q,\omega )`$ is the dynamical susceptibility. Within the self-consistent spin fluctuation theory, in a paramagnet close to a ferromagnetic instability, one gets $`\tau _{sf}^1T`$. For the other two bands, a quadratic dependence of the scattering rates on the temperature for all the calculated quantities is assumed. To find the transport coefficients, we must calculate the current defined as: $`𝐉={\displaystyle e𝐯g(𝐯)𝑑𝐤}.`$ (1) where $`g(𝐯)`$ is the local distribution of electrons. Confining ourselves to an expansion of order $`B^2`$, one can easily obtain the following general formula $`J_\alpha =\sigma _{\alpha \beta }^{(0)}E_\beta +\sigma _{\alpha \beta \gamma }^{(1)}E_\beta B_\gamma +\sigma _{\alpha \beta \gamma \delta }^{(2)}E_\beta B_\gamma B_\delta ,`$ (2) where the summation convention over the indices of the cartesian components is assumed and $`E_i`$ and $`B_i`$ are the components of the external electric and magnetic field, respectively. To calculate the normal resistivity, the Hall coefficient and the magnetoresistance for a multi-band case, we write the total current as the sum of the contributions coming from the three bands and then we invert the matrix connecting $`𝐉`$ and $`𝐄`$. Neglecting powers above $`B^2`$, we have: $`\rho _0={\displaystyle \frac{1}{\sigma _0^{Tot}}},`$ (3) $`\rho _H=\rho _0^2\sigma _H^{Tot},`$ (4) $`\rho _{MR}=\rho _0^2\left[\sigma _{MR}^{Tot}+\rho _0(\sigma _H^{Tot})^2\right];`$ (5) where $`\sigma _0^{Tot}`$, $`\sigma _H^{Tot}`$ and $`\sigma _{MR}^{Tot}`$ denote the total conductivity, the total Hall conductivity and the second order total conductivity, respectively. The explicit computation of the magnetotransport quantities requires the knowledge of the band spectra as well as the relaxation times for describing the collision of the electrons in the bands produced by the d<sub>xy</sub>, d<sub>xz</sub> and d<sub>yz</sub> Ru orbitals. Referring to the energy spectra, we use the electronic energy band structure of Sr<sub>2</sub>RuO<sub>4</sub> recently calculated by using a simple method combining the extended Hückel theory and the tight-binding approximation. Concerning the relaxation times, in the case of the Hall- and magneto- resistance the following expressions have been considered: $`(\tau _i)^1`$ $`=`$ $`\eta _i+\alpha _iT^2`$ where $`i`$=$`(xz,yz,xy)`$ indicates the band and $`\eta =(2.75,2.75,3.25)`$, and $`\alpha =(0.035,0.04,0.06)`$. The values of $`\eta _i`$ have been chosen in a way to get the experimental observed resistivity at T=4 K of $`0.7\mu \mathrm{\Omega }cm^1`$. The constraint on the values of $`\alpha _i`$ is given by the complicated temperature dependence of $`R_H`$ together with the behaviour of the transverse in-plane magnetoresistance. We notice that the behaviour of the $`R_H`$ can be reproduced only if $`\alpha _{xz}\alpha _{yz}`$ but smaller than $`\alpha _{xy}`$. The calculation is very sensitive to the changes in the time collisions of the two hybridized $`z`$ bands. Indeed, for small relative variation of $`\alpha _{xz}`$ with respect to $`\alpha _{yz}`$ the Hall-coefficient does not show sign changing, being always negative. It is worth pointing out that the explicit expression for the scattering rate is related to several physical mechanisms that give rise to different temperature dependencies of $`\tau _i`$. In particular, while $`\eta _i`$ could be considered as responsible for the residual resistivity, $`\alpha _i`$ could be related to umklapp electron-electron scattering, or inelastic scattering of the electrons by impurities, as well as by phonon mediated interaction between electrons. In the case of electron-electron scattering, an estimation of $`\alpha _i`$ can be obtained by means of the relation $`(\tau _i^{ee})^1(k_BT^2/m_ik_F^i)`$. Using the experimental values for the effective masses as deduced from de Haas-van Alphen experiment, we find that $`\alpha _{xy}`$ is greater than $`\alpha _{xz}`$ and $`\alpha _{yz}`$ and this agrees with our assumption. Concerning the other mechanisms, though the microscopic expression in the case of the phononic and impurity scattering requires a more accurate analysis, we expect that the considerations above are still valid. The fit to the experimental data is reported in Fig.1, for the Hall coefficient and in Fig.2 for the magnetoresistance. The experimental data are taken from for $`R_H`$ and from for the magnetoresistance. In both cases, we find a good agreement between the experimental results and the theoretical prediction indicating that the main contribution to scattering rate follows a $`T^2`$ power law. We notice that paramagnetic and ferromagnetic materials can have a large contribution to the Hall effect mainly due to skew scattering. Indeed, in this case moving charge carriers experience a force due to the magnetic field produced by localized magnetic moment and are scattered asymmetrically. Nevertheless, there is no sign of the saturation of the Hall resistivity one expects when there is a large magnetic contribution to the Hall effect. Therefore, we argue that the experimental data for Sr<sub>2</sub>RuO<sub>4</sub> are dominated by the standard orbital Hall resistivity. We point out that there are other results in literature dealing with the theoretical study of the Hall coefficient in the Sr<sub>2</sub>RuO<sub>4</sub>. In Ref. is presented a theoretical fit of the Hall coefficient based on the assumption of two-carrier system and $`R_H`$ is calculated within the Drude classical model. While, the authors neglect the contribution of one of the electron pockets and the effective electronic structure of the bands, they reproduce the sign change from negative at low temperature to positive at high temperature but fails to give a quantitative agreement with the experimental data. In Ref. , using methods developed by Ong, an expression for the Hall coefficient in multi-band system is derived. We notice that the main assumption of this derivation is that the mean free path is the same for all the Fermi sheets. This hypothesis of isotropic mean-free path is valid at small temperature where the authors obtain a value for $`R_H`$ that compares well with the measured value of the same quantity. Nevertheless, we want to stress that the value of $`R_H`$ is extremely sensitive to details of the k-dependent scattering and/or on the energy spectra, and strongly depends on the temperature. Finally, in Ref. is shown a fit to $`R_H`$ similar to the one here presented. However, in Ref. , the sign change is reproduced with an accuracy weaker than the one presented in the present paper and no mention to other relevant physical quantities is made. With respect to the previous approaches, therefore, we believe the results presented here yield a good qualitative and quantitative agreement with the transport experiments. Let us now discuss the zero-field in plane resistivity. As mentioned in the introduction, $`\rho _{ab}`$ exhibits a Fermi-liquid behaviour up to T<sub>FL</sub>; above this temperature $`\rho _{ab}`$ rises monotonically with a curvature which is weaker than T<sup>2</sup>. We notice that this effect may be probably induced by a spin scattering mechanism which does not contribute to the determination of the Hall coefficient and the magnetoresistance. Therefore, we add to the scattering rate for the $`\gamma `$ band electrons above, a term proportional to the temperature, while we keep unchanged the other two scattering rates. Its temperature dependence is given by: $`(\tau _{xy}^{sf})^1`$ $`=`$ $`\beta T`$ where $`\beta =0.6`$. The resulting fit is reported in Fig.3 where the experimental data are taken from . The quite good agreement between the experimental results and the theoretical prediction gives confidence that spin fluctuations affects only the zero field relaxation time for the xy band. As a final remark, we notice that throughout this paper we have assumed for the relaxation rates an isotropic k-independent form. We have also evaluated the above mentioned quantities assuming for $`\tau _i`$ the suitable form in the case of a tetragonal environment. The results are only slightly modified, so that we have confined ourselves to temperature dependent but k-independent $`\tau _i`$. In summary, we have studied the in-plane normal state magnetotransport quantities of the layered compound Sr<sub>2</sub>RuO<sub>4</sub>. Using the calculated electron energy band structure, we have computed the temperature dependence of the Hall coefficient, the magnetoresistance and the in-plane resistivity by solving the Boltzmann equation for a multi-orbital system. The reasonably good fit of these physical quantities suggests that the assumption of two contributions in the relaxation rate for the xy-electrons used to quantify the galvanomagnetic transport is essentially correct.
warning/0007/math0007059.html	ar5iv	text	# From Flows and Metrics to Dynamics ## 1 Geometric Dynamics Let $`M`$ be an $`n`$-dimensional differentiable manifold. A $`C^{\mathrm{}}`$ vector field $`X`$ on $`M`$ defines the flow $$\frac{dx}{dt}=X(x).$$ $`(1)`$ A semi-Riemann metric $`g`$ on the manifold $`M`$ is a $`C^{\mathrm{}}`$ symmetric tensor field of type (0,2) which assigns to each point $`xM`$ a nondegenerate inner product $`g(x)`$ on the tangent space $`T_xM`$ of signature $`(r,s)`$. The pair $`(M,g)`$ is called a semi-Riemann manifold. The vector field $`X`$ and the semi-Riemann metric $`g`$ determine the energy $$f:MR,f=\frac{1}{2}g(X,X).$$ The vector field (flow) $`X`$ on $`(M,g)`$ is called: 1) timelike, if $`f<0`$; 2) nonspacelike or causal, if $`f0`$; 3) null or lightlike, if $`f=0`$; 4) spacelike, if $`f>0`$. Let $`X`$ be a nonwhere vanishing vector field, of everywhere constant energy. Upon rescaling, it may be supposed that $`f\{1,0,1\}`$. Generally, $``$ is the set of zeros of the vector field $`X`$. If $`\mathrm{}`$, then this rescaling is possible only on $`M`$. Let $``$ be the Levi-Civita connection of $`(M,g)`$. Using the operator $`{\displaystyle \frac{}{dt}}`$ (covariant differentiation along a solution) we obtain the prolongation $$\frac{}{dt}\frac{dx}{dt}=_{\frac{dx}{dt}}X$$ $`(2)`$ of the differential system (1) or of any perturbation of the system (1) obtained adding to the second member $`X`$ a parallel vector field $`Y`$ with respect to the covariant derivative $``$. The prolongation by derivation represents the general dynamics of the flow. The vector field $`Y`$ can be used to illustrate a progression from stable to unstable flows, or converse. The vector field $`X`$, the metric $`g`$, and the connection $``$ determine the external (1,1)-tensor field $$F=Xg^1g(X),$$ $$F_j{}_{}{}^{i}=_jX^ig^{ih}g_{kj}_hX^k,i,j,h,k=1,\mathrm{},n,$$ which characterizes the helicity of vector field (flow) $`X`$. First we write the differential system (2) in the equivalent form $$\frac{}{dt}\frac{dx}{dt}=g^1g(X)\left(\frac{dx}{dt}\right)+F\left(\frac{dx}{dt}\right).$$ $`(2^{})`$ Successively we modify the differential system (2) as follows: $$\frac{}{dt}\frac{dx}{dt}=g^1g(X)(X)+F\left(\frac{dx}{dt}\right),$$ $`(3)`$ $$\frac{}{dt}\frac{dx}{dt}=g^1g(X)\left(\frac{dx}{dt}\right)+F(X),$$ $`(4)`$ $$\frac{}{dt}\frac{dx}{dt}=g^1g(X)(X)+F(X).$$ $`(5)`$ Obviously, the second order systems (3), (4), (5) are prolongations of the first order system (1). Each of them is connected either to the dynamics of the field $`X`$ or to the dynamics of a particle which is sensitive to the vector field $`X`$. Since $$g^1g(X)(X)=gradf,$$ we shall show that the prolongation (3) describes a conservative dynamics of the vector field $`X`$ or of a particle which is sensitive to the vector field $`X`$. The physical phenomenon produced by (4) or (5) was not yet studied. Theorem. 1) If $`F=0`$, then the kinematic system (1) prolonges to a potential dynamical system with $`n`$ degrees of freedom, namely $$\frac{}{dt}\frac{dx}{dt}=gradf.$$ $`(3^{})`$ 2) If $`F0`$, then the kinematic system (1) prolonges to a non-potential dynamical system with $`n`$ degrees of freedom, namely $$\frac{}{dt}\frac{dx}{dt}=gradf+F\left(\frac{dx}{dt}\right).$$ $`(3^{\prime \prime })`$ Corollary. If the metric $`g`$ is chosen such that $`f\{1,0,1\}`$ on $`M`$, then the flow generated by $`X`$ is global and the dynamical systems (3$`^{}`$), (3$`^{\prime \prime }`$) are reduced to $$\frac{^2x}{dt^2}=0,\frac{^2x}{dt^2}=F\left(\frac{dx}{dt}\right).$$ Let us show that the dynamical systems (3$`^{}`$) and (3$`^{\prime \prime }`$) are conservative. To simplify the exposition we identity the tangent bundle $`TM`$ with the cotangent bundle $`T^{}M`$ using the semi-Riemann metric $`g`$. Theorem. 1) The trajectories of the dynamical system (3$`^{}`$) are the extremals of the Lagrangian $$L=\frac{1}{2}g(\frac{dx}{dt},\frac{dx}{dt})+f(x).$$ 2) The trajectories of the dynamical system (3$`^{\prime \prime }`$) are the extremals of the Lagrangian $$L=\frac{1}{2}g(\frac{dx}{dt}X,\frac{dx}{dt}X)=\frac{1}{2}g(\frac{dx}{dt},\frac{dx}{dt})g(X,\frac{dx}{dt})+f(x).$$ 3) The dynamical systems (3$`^{}`$) and (3$`^{\prime \prime }`$) are conservative, the Hamiltonian being the same for both cases, namely $$H=\frac{1}{2}g(\frac{dx}{dt},\frac{dx}{dt})f(x).$$ The restriction of the Hamiltonian $`H`$ to the flow of the vector field $`X`$ is zero. Obviously the values of the Hamiltonian $`H`$ can be positive, negative or zero, even if the metric $`g`$ is a Riemannian metric; therefore, just in this case, there exist boundary-value problems associated to the differential system (3), having three solutions (for example, the first corresponding to constant total energy $`H<0`$, the second for $`H=0`$, and the third for $`H>0`$). For the next theorem we recall that a pregeodesic is a smooth curve which may be reparametrized to be a geodesic. Theorem (Lorentz-Udrişte World-Force Law). 1) Every non-cons-tant trajectory of the dynamical system (3$`^{}`$), which corresponds to a constant value $`H_0`$ of the Hamiltonian, is a pregeodesic of the semi-Riemann-Jacobi manifold $$(M,\overline{}=(_\mathcal{0}+𝒻)).$$ 2) Let $`g_{ij}`$ be the local components of the metric $`g`$ and let $`\mathrm{\Gamma }_{jk}^i`$, $`i,j,k=1,\mathrm{},n`$ be the local components of the connection $``$. Every non-constant trajectory of the dynamical system (3$`^{\prime \prime }`$), which corresponds to a constant value $`H_0`$ of the Hamiltonian, is a horizontal pregeodesic of the semi-Riemann-Jacobi-Lagrange manifold $$(M,\overline{}=(_\mathcal{0}+𝒻),𝒩_𝒿{}_{}{}^{𝒾}=\mathcal{\Gamma }_{𝒿𝓀}^𝒾𝓎^𝓀_𝒿{}_{}{}^{𝒾},𝒾,𝒿,𝓀=\mathcal{1},\mathrm{},𝓃).$$ Corollary. If the metric g is chosen such that $`f\{1,0,1\}`$ on $`M`$ and if we denote $$\alpha ^2=g(\frac{dx}{dt},\frac{dx}{dt}),\beta =k^{1/2}g(X,\frac{dx}{dt}),$$ then every trajectory of the dynamical system (3$`^{\prime \prime }`$), with $`\alpha ^2k^{1/2}+\beta =0`$, is a pregeodesic of the semi-Finsler-Jacobi manifold $$(M,=_𝓌^\mathcal{2}=𝓌\alpha \beta ,𝓌=\text{constant}).$$ ## 2 Hamiltonian Structures on the Tangent Bundle Let $`N`$ be a $`2n`$-dimensional manifold. A nondegenerate and closed 2-form $``$ on $`N`$ is called symplectic form. A manifold $`N`$ with a given symplectic form is called a phase space. Let $`(N,)`$ be a phase space, and $`H:NR`$ be a $`C^{\mathrm{}}`$ real function. We define the Hamilton gradient $`X_H`$ as being the vector field which satisfies $$_p(X_H(x),v)=dH(x)(v),vT_xN,$$ and the Hamilton equations as $$\frac{dx}{dt}=X_H(x).$$ Let $`(M,g)`$ be a semi-Riemann manifold with $`n`$ dimensions. Let $`X`$ be a $`C^{\mathrm{}}`$ vector field on $`M`$, and $`\omega =gF`$ the 2-form associated to the tensor field $`F=Xg^1g(X)`$ via the metric $`g`$. The tangent bundle is usually endowed with the Sasaki metric $`G`$ created by $`g`$. If $`(x^i,y^i)`$ are the coordinates of the point $`(x,y)TM`$ and $`\mathrm{\Gamma }_{jk}^i`$ are the components of the connection induced by $`g_{ij}`$, then $$(\frac{}{x^i}\mathrm{\Gamma }_{ij}^hy^j\frac{}{y^h},\frac{}{y^i}),(dx^j,\delta y^j=dy^j+\mathrm{\Gamma }_{hk}^jy^kdx^h)$$ are dual frames. Also the metric of Sasaki transcribes $$G=g_{ij}dx^idx^j+g_{ij}\delta y^i\delta y^j.$$ Theorem. The dynamical system (3$`^{}`$) lifts to $`TM`$ as a Hamilton dynamical system with respect to the Hamiltonian $$H=\frac{1}{2}g(\frac{dx}{dt},\frac{dx}{dt})f(x)$$ and the symplectic 2-form $$\mathrm{\Omega }_1=g_{ij}dx^i\delta y^j.$$ Hint. $`\eta _1=g_{ij}y^idx^j`$, and $`d\eta _1=\mathrm{\Omega }_1`$. Theorem. The dynamical system (3$`^{\prime \prime }`$) lifts to $`TM`$ as a Hamilton dynamical system with respect to the Hamiltonian $$H=\frac{1}{2}g(\frac{dx}{dt},\frac{dx}{dt})f(x)$$ and the symplectic 2-form $$\mathrm{\Omega }_2=\frac{1}{2}\omega _{ij}dx^idx^j+g_{ij}dx^i\delta y^j.$$ Hint. $`\eta _2=g_{ij}X^idx^j+g_{ij}y^idx^j`$, and $`d\eta _2=\mathrm{\Omega }_2`$. Pendulum Geometric Dynamics. We use the Riemannian manifold $`(R^2,\delta _{ij})`$. The small oscillations of a plane pendulum are described as solutions of the differential system (plane pendulum flow) $$\frac{dx_1}{dt}=x_2,\frac{dx_2}{dt}=x_1.$$ $`(6)`$ In this case $`x_1(t)=0`$, $`x_2(t)=0`$, $`tR`$ is the equilibrium point and $`x_1(t)=c_1\mathrm{cos}t+c_2\mathrm{sin}t`$, $`x_2(t)=c_1\mathrm{sin}tc_2\mathrm{cos}t`$, $`tR`$ is the general solution (family of circles with same centre). Let $$X=(X_1,X_2),X_1(x_1,x_2)=x_2,X_2(x_1,x_2)=x_1,$$ $$f(x_1,x_2)=\frac{1}{2}(x_1^2+x_2^2),\text{rot}X=(0,0,2),divX=0.$$ The pendulum flow conserves the areas. The prolongation by derivation of the kinematic system (6) is $$\frac{d^2x_i}{dt^2}=\underset{j}{}\frac{X_i}{x_j}\frac{dx_j}{dt},i,j=1,2$$ or $$\frac{d^2x_1}{dt^2}=\frac{dx_2}{dt},\frac{d^2x_2}{dt^2}=\frac{dx_1}{dt}.$$ $`(7)`$ This prolongation admits the general solution $$\begin{array}{c}x_1(t)=a_1\mathrm{cos}t+a_2\mathrm{sin}t+h\hfill \\ \\ x_2(t)=a_1\mathrm{sin}ta_2\mathrm{cos}t+k,tR.\hfill \end{array}$$ (family of circles). The pendulum geometric dynamics is described by $$\frac{d^2x_i}{dt^2}=\frac{f}{x_i}+\underset{j}{}\left(\frac{X_i}{x_j}\frac{X_j}{x_i}\right)\frac{dx_j}{dt},i,j=1,2$$ or $$\frac{d^2x_1}{dt^2}=x_12\frac{dx_2}{dt},\frac{d^2x_2}{dt^2}=x_2+2\frac{dx_1}{dt},$$ $`(8)`$ with the general solution $$x_1(t)=b_1\mathrm{cos}t+b_2\mathrm{sin}t+b_3t\mathrm{cos}t+b_4t\mathrm{sin}t$$ $$x_2(t)=b_1\mathrm{sin}tb_2\mathrm{cos}t+b_3t\mathrm{sin}tb_4t\mathrm{cos}t,tR$$ (family of spirals). Using $$\begin{array}{c}L=\frac{1}{2}\left[\left(\frac{dx_1}{dt}\right)^2+\left(\frac{dx_2}{dt}\right)^2\right]+x_2\frac{dx_1}{dt}x_1\frac{dx_2}{dt}+f\hfill \\ \\ H=\frac{1}{2}\left[\left(\frac{dx_1}{dt}\right)^2+\left(\frac{dx_2}{dt}\right)^2\right]f\hfill \\ \\ g_{ij}=(H+f)\delta _{ij},\hfill \\ \\ N_j{}_{}{}^{i}=F_j{}_{}{}^{i}=\delta ^{ih}F_{jh},F_{ij}=\frac{X_j}{x_i}\frac{X_i}{x_j},i,j,h=1,2,\hfill \end{array}$$ the solutions of the differential system (8) are horizontal pregeodesics of the Riemann-Jacobi-Lagrange manifold $$(R^2\{0\},g_{ij},N_j{}_{}{}^{i}).$$ Lorenz Geometric Dynamics. We use the Riemannian manifold $`(R^3,\delta _{ij})`$. The Lorenz flow is a first dissipative model with chaotic behaviour discovered in numerical experiment. Its state equations are $$\begin{array}{c}\frac{dx_1}{dt}=\sigma x_1+\sigma x_2\hfill \\ \\ \frac{dx_2}{dt}=x_1x_3+rx_1x_2\hfill \\ \\ \frac{dx_3}{dt}=x_1x_2bx_3,\hfill \end{array}$$ $`(9)`$ where $`\sigma ,r,b`$ are real parameters. Usually $`\sigma ,b`$ are kept fixed whereas $`r`$ is varied. At $$r>r_0=\frac{\sigma (\sigma +b+3)}{\sigma b1}$$ chaotic behaviour is observed. With $`\sigma =10,b={\displaystyle \frac{8}{3}}`$, the preceding inequality yields $`r_0=24,7368`$. If $`\sigma 0`$ and $`b(r1)>0`$, then the equilibrium points of the Lorenz flow are $$x=0,y=0,z=0;$$ $$x=\pm \sqrt{b(r1)},y=\pm \sqrt{b(r1)},z=r1.$$ Let $$X=(X_1,X_2,X_3),X_1(x_1,x_2,x_3)=\sigma x_1+\sigma x_2,$$ $$X_2(x_1,x_2,x_3)=x_1x_3+rx_1x_2,X_3(x_1,x_2,x_3)=x_1x_2bx_3,$$ $$f=\frac{1}{2}[(\sigma x_1+\sigma x_2)^2+(x_1x_3+rx_1x_2)^2+(x_1x_2bx_3)^2],$$ $$rotX=(2x_1,x_2,rx_3\sigma ).$$ The Lorenz geometric dynamics is described by $$\frac{d^2x_i}{dt^2}=\frac{f}{x_i}+\underset{j}{}\left(\frac{X_i}{x_j}\frac{X_j}{x_i}\right)\frac{dx_j}{dt},i,j=1,2,3$$ or $$\begin{array}{c}\frac{d^2x_1}{dt^2}=\frac{f}{x_1}+(\sigma +x_3r)\frac{dx_2}{dt}x_2\frac{dx_3}{dt}\hfill \\ \\ \frac{d^2x_2}{dt^2}=\frac{f}{x_2}+(rx_3\sigma )\frac{dx_1}{dt}2x_1\frac{dx_3}{dt}\hfill \\ \\ \frac{d^2x_3}{dt^2}=\frac{f}{x_3}+x_2\frac{dx_1}{dt}+2x_1\frac{dx_2}{dt}.\hfill \end{array}$$ $`(10)`$ Using $$\begin{array}{c}L=\frac{1}{2}\underset{i=1}{\overset{3}{}}\left(\frac{dx_i}{dt}\right)^2\underset{i=1}{\overset{3}{}}X_i\frac{dx_i}{dt}+f\hfill \\ \\ H=\frac{1}{2}\underset{i=1}{\overset{3}{}}\left(\frac{dx_i}{dt}\right)^2f\hfill \\ \\ g_{ij}=(H+f)\delta _{ij}\hfill \\ \\ N_j{}_{}{}^{i}=F_j{}_{}{}^{i}=\delta ^{ih}F_{jh},F_{ij}=\frac{X_j}{x_i}\frac{X_i}{x_j},i,j,h=1,2,3,\hfill \end{array}$$ the solutions of the differential system (10) are horizontal pregeodesics of the Riemann-Jacobi-Lagrange manifold $$(R^3,_{𝒾𝒿},𝒩_𝒿{}_{}{}^{𝒾}),$$ where $``$ is the set of equilibrium points. ABC Geometric Dynamics. We use the Riemannian manifold $`(R^3,\delta _{ij})`$. One examples of a fluid velocity that contains exponential stretching and hence instability is the ABC flow, $$\{\begin{array}{c}\frac{dx_1}{dt}=A\mathrm{sin}x_3+C\mathrm{cos}x_2\hfill \\ \\ \frac{dx_2}{dt}=B\mathrm{sin}x_1+A\mathrm{cos}x_3\hfill \\ \\ \frac{dx_3}{dt}=C\mathrm{sin}x_2+B\mathrm{cos}x_1.\hfill \end{array}$$ $`(11)`$ This flow is named after the three mathematicians Arnold, Beltrami and Childress, who have contributed much to our understanding and appreciation of classes of ”chaotic” flows of which the present one is an example. For nonzero values of the constants $`A,B,C`$ the preceding system is not globally integrable. The topology of the flow lines is very complicated and can only be investigated numerically to reveal regions of chaotic behaviour. The $`ABC`$ flow conserves the volumes since the $`ABC`$ field is solenoidal. The $`ABC`$ geometric dynamics is described by $$\frac{d^2x^i}{dt^2}=\frac{f}{x^i}+\underset{j}{}\left(\frac{X_i}{x_j}\frac{X_j}{x_i}\right)\frac{dx_j}{dt},i,j=1,2,3.$$ Since $$f=\frac{1}{2}(A+B+C+2AC\mathrm{sin}x_3\mathrm{cos}x_2+2BA\mathrm{sin}x_1\mathrm{cos}x_3+2CB\mathrm{sin}x_2\mathrm{cos}x_1)$$ $$rotX=X,$$ the $`ABC`$ geometric dynamics is given by the differential system (12): $$\begin{array}{ccc}\frac{d^2x_1}{dt^2}\hfill & =& AB\mathrm{cos}x_1\mathrm{cos}x_3BC\mathrm{sin}x_1\mathrm{sin}x_2(B\mathrm{cos}x_1+C\mathrm{sin}x_2)\frac{dx_2}{dt}+\hfill \\ \multicolumn{3}{c}{}\\ & +& (B\mathrm{sin}x_1+A\mathrm{cos}x_3)\frac{dx_3}{dt}\hfill \end{array}$$ $$\begin{array}{ccc}\frac{d^2x_2}{dt^2}\hfill & =& AC\mathrm{sin}x_2\mathrm{sin}x_3+BC\mathrm{cos}x_1\mathrm{cos}x_2+(B\mathrm{cos}x_1+C\mathrm{sin}x_2)\frac{dx_1}{dt}\hfill \\ \multicolumn{3}{c}{}\\ & & (A\mathrm{sin}x_3+C\mathrm{cos}x_2)\frac{dx_3}{dt}\hfill \end{array}$$ $$\begin{array}{ccc}\frac{d^2x_3}{dt^2}\hfill & =& AC\mathrm{cos}x_3\mathrm{cos}x_2BA\mathrm{sin}x_1\mathrm{sin}x_3(B\mathrm{sin}x_1+A\mathrm{cos}x_3)\frac{dx_1}{dt}+\hfill \\ \multicolumn{3}{c}{}\\ & +& (C\mathrm{cos}x_2+A\mathrm{sin}x_3)\frac{dx_2}{dt}.\hfill \end{array}$$ Using $$\begin{array}{c}L=\frac{1}{2}\underset{i=1}{\overset{3}{}}\left(\frac{dx_i}{dt}\right)^2\underset{i=1}{\overset{3}{}}X_i\frac{dx_i}{dt}+f\hfill \\ \\ H=\frac{1}{2}\underset{i=1}{\overset{3}{}}\left(\frac{dx_i}{dt}\right)^2f\hfill \\ \\ g_{ij}=(H+f)\delta _{ij}\hfill \\ \\ N_j{}_{}{}^{i}=F_j{}_{}{}^{i}=\delta ^{ih}F_{jh},F_{ij}=\frac{X_j}{x_i}\frac{X_i}{x_j},i,j,h=1,2,3\hfill \end{array}$$ the solutions of the differential system (12) are horizontal pregeodesics of the Riemann-Jacobi-Lagrange manifold $$(R^3,_{𝒾𝒿},𝒩_𝒿{}_{}{}^{𝒾}),$$ where $``$ is the set of equilibrium points which is included in the surface of equation $$\mathrm{sin}x_1\mathrm{sin}x_2\mathrm{sin}x_3+\mathrm{cos}x_1\mathrm{cos}x_2\mathrm{cos}x_3=0.$$
warning/0007/hep-th0007026.html	ar5iv	text	# T-dual R-R zero-norm states, D-branes and S-duality of type II string theory ## I Introduction It has been pointed out for a long time that the complete space-time symmetry of string theory is related to the zero-norm state (a physical state that is orthogonal to all physical states including itself) in the old covariant quantization of the string spectrum. This observation had made it possible to explicitly construct many stringy ($`\alpha ^{}\mathrm{}`$) massive symmetry of the theory. This includes the $`w_{\mathrm{}}`$ symmetry of the toy 2D string and the discrete massless and massive T-duality symmetry of closed bosonic string. The authors of show that, in string theory, some target space mirror symmetry of N=2 backgrounds on group manifolds is a Kac-Moody gaugy symmetry. Thus, like T-duality, it should be related to the zero-norm states. On the other hand, the massless and massive SUSY, and some new enlarged spacetime boson-fermion symmetries induced by zero-norm states were also discussed in . It is thus of interest to study the R-R zero-norm state and its relation to D-brane which was recently shown by Polchinski to be the symmetry charge carrier of the propagating R-R forms. Presumably, there should be no R-R zero-norm state in the type II string spectrum since the fundamental string does not interact with the R-R forms. However, to our surprise, it was discovered that there do exist both massless and massive R-R zero-norm states in the type II string spectrum. It was then realized that the degree of freedom of massless R-R zero-norm states does not fit into that of the symmetry parameters of the propagating R-R forms and thus resolved the seeming inconsistency. This observation gives us another justification of the well-known wisdom that perturbative string does not carry the massless R-R charges, although the existence of these R-R zero-norm states remain mysterious. In this paper, we will show that the T-dual R-R zero-norm states serve as the right symmetry parameters of the gauge transformations of the R-R propagating forms. Also one is forced to introduce Type I open string and D-branes into the type II string theory to incorporate these T-dual R-R zero-norm states. Our spacetime zero-norm states argument here is in complementary with the worldsheet string vertex operator argument first given by Binnchi, Pradisi and Sagnotti. They considered R-R one-point function in the (-$`\frac{1}{2}`$,-$`\frac{3}{2}`$) ghost picture on the disk and resulted in a conclusion which was consistent with D-brane as R-R charge carrier. As an important application, we demonstrate that the constant T-dual R-R 0-form zero-norm state, together with the NS-NS singlet zero-norm state which was always neglected in the previous discussions, are responsible for the discrete SL(2,Z) S-duality symmetry of the type II B string theory. This discovery suggests that not only stringy ($`\alpha ^{}\mathrm{}`$) symmetry but also strong-weak ($`g_s\mathrm{}`$) duality symmetry are related to the existence of zero-norm states of the spectrum. ## II T-dual R-R zero-norm states The massless physical NS state of the open superstring is (we use the notation in Ref ) $$\epsilon _\mu b_{\frac{1}{2}}^\mu \|0,k;\text{ }k\epsilon =0,\text{ }k^2=0$$ (1) In addition, there is a singlet zero-norm state $$k_\mu b_{\frac{1}{2}}^\mu \|0,k;\text{ }k^2=0\text{ .}$$ (2) The NS-NS symmetries of graviton and antisymmetry tensor of type II string were derived through the following two zero-norm states $`\epsilon _\mu b_{\frac{1}{2}}^\mu \|0,kk_\mu \underset{\frac{1}{2}}{\overset{\mu }{\stackrel{}{b}}}\|0,k,`$ (3) $`k_\mu b_{\frac{1}{2}}^\mu \|0,k\epsilon _\mu \underset{\frac{1}{2}}{\overset{\mu }{\stackrel{}{b}}}\|0,k`$ (4) in the first order weak field approximation (WFA). The remaining interesting singlet zero-norm state $$k_\mu b_{\frac{1}{2}}^\mu \|0,kk_\mu \underset{\frac{1}{2}}{\overset{\mu }{\stackrel{}{b}}}\|0,k$$ (5) will be discussed in the next section. We now discuss the massless R state. The only propagating spinor is $$\|\stackrel{}{S},ku_\stackrel{}{s};\text{ }F_0\|\stackrel{}{S},ku_\stackrel{}{s}=0\text{ .}$$ (6) The GSO operator in the massless limit reduces to the chirality operator, and only one of the chiral spinor $`8_s(`$or $`8_c)`$ will be projected out. In addition, there is a massless fermionic zero-norm state $$k_\mu \mathrm{\Gamma }_{\stackrel{}{s},\stackrel{}{s}}^\mu \|\stackrel{}{S},k\theta _\stackrel{}{s}\text{ .}$$ (7) Eq.(6) is the only massless solution of the following $$F_0\|\psi ,\text{where }F_1\|\psi =L_0\|\psi =0\text{.}$$ (8) The state in equation (6) is crucial in the discussion of this paper. Note that $`k\mathrm{\Gamma }\|\stackrel{}{S},k\theta _\stackrel{}{s}`$ is left-handed if $`\|\stackrel{}{S},k\theta _\stackrel{}{s}`$ is right-handed and both spinors have exactly the same degree of freedom. The massless propagating R-R states of type II string consist of tensor forms $$G_{\alpha \beta }=\underset{k=0}{\overset{10}{}}\frac{i^k}{k!}G_{\mu _1\mu _2\mathrm{}\mu _k}(\mathrm{\Gamma }^{\mu _1\mu _2\mathrm{}\mu _k})_{\alpha \beta }\text{,}$$ (9) where $`\mathrm{\Gamma }^{\mu _1\mu _2\mathrm{}\mu _k}`$ are the antisymmetric products of gamma-matrix, and $`\alpha ,\beta `$ are spinor indices. There is a duality relation which reduces the number of independent tensor components to up to k=5 form. The on-shell conditions, or two massless Dirac equations, imply G is indeed a field strength and can be written as $$G_{(k)}=dA_{(k1)}$$ (10) which means perturbative string states do not carry the massless R-R symmetry charges. We are now in a position to discuss the R-related symmetry charges. Let’s first introduce the NS-R (R-NS) SUSY zero-norm states $$\text{ }kb_{\frac{1}{2}}\|0,k\|\stackrel{}{S},k\overline{u}_\stackrel{}{s}\text{and }\|\stackrel{}{S},ku_\stackrel{}{s}k\stackrel{~}{b}_{\frac{1}{2}}\|0,k$$ (11) for the II A theory and a trivial modification for the II B theory. The corresponding worldsheet vertex operator in the ($`0,\frac{1}{2}`$) picture for say the first state in equation (10) is $`k_\mu (x^\mu (z)+ik\psi \psi ^\mu )e^{ikx(z)}u_\alpha \stackrel{\alpha }{\stackrel{}{S}}(\overline{z})e^{\frac{1}{2}\stackrel{}{\varphi }}e^{ikx(\overline{z})}`$ (12) $`=`$ $`e^{ikx(z)}u_\alpha \stackrel{\alpha }{\stackrel{}{S}}e^{\frac{1}{2}\stackrel{}{\varphi }}e^{ikx(\overline{z})},`$ (13) which is a worldsheet total derivative and, as in the case of bosonic sector, one can introduce a worldsheet generator and deduce the SUSY current to be $$\underset{\alpha ,\frac{1}{2}}{\overset{}{Q}}=\stackrel{\alpha }{\stackrel{}{S}}e^{\frac{1}{2}\stackrel{}{\varphi }},$$ (14) where $`\stackrel{\alpha }{\stackrel{}{S}}`$and $`\stackrel{}{\varphi }`$ are the right-moving spin field and the bosonized superconformal ghost respectively. This zero-norm state derivation is consistent with the original approach. The advantage of our approach is that one can generalize to derive the enlarged stringy boson-fermion symmetry by using the massive fermion zero-norm state of the spectrum. We give one example here. There exists a m=2 NS-R zero-norm state $$\left[2\theta _{\mu \nu }\alpha _1^\mu b_{\frac{1}{2}}^\nu +k_{[\lambda }\theta _{\mu \nu ]}b_{\frac{1}{2}}^\lambda b_{\frac{1}{2}}^\mu b_{\frac{1}{2}}^\nu \right]\|0,k\stackrel{~}{\alpha }_1^\lambda \|\stackrel{}{S},ku_{\lambda ,\stackrel{}{s}}$$ (15) with $`\theta _{\mu \nu }=\theta _{\nu \mu },`$ $`k^\mu \theta _{\mu \nu }=0`$ and $`[(kd_0)\alpha _1^\mu +d_1^\mu ]u_{\mu ,\stackrel{}{s}}`$ $`=`$ $`0,`$ (16) $`d_0^\mu u_{\mu ,\stackrel{}{s}}`$ $`=`$ $`0.`$ (17) The corresponding vertex operator is calculated to be $`[2\theta _{[\mu \nu ],\lambda \alpha }(x^\mu x^\nu \psi ^\mu \psi ^\nu +ik\psi \psi ^\mu x^\nu )+k_{[\delta }\theta _{\mu \nu ],\lambda \alpha }(3x^\mu +ik\psi \psi ^\mu )`$ (18) $`\psi ^\nu \psi ^\delta ]\overline{}x^\lambda k\overline{\psi }e^{\frac{1}{2}\stackrel{}{\varphi }}\stackrel{}{S}^\alpha e^{ikx(z,\overline{z})}`$ (19) where $`\theta _{\mu \nu ,\lambda \alpha }\theta _{\mu \nu }u_{\lambda \alpha }.`$ It is straight-forward to construct the corresponding ward identity although the symmetry transformation law of the background fields is not easy to write down at this point. We now turn to discuss the R-R zero-norm states. For the massless level, we have the following zero-norm states $$k_\mu \mathrm{\Gamma }_{\stackrel{}{\stackrel{}{s}}\stackrel{}{s}}^\mu \|\stackrel{}{S},k\theta _\stackrel{}{s}\|\stackrel{}{S},ku_\stackrel{}{s}\text{ (II A)}$$ (20) and $$k_\mu \mathrm{\Gamma }_{\stackrel{}{\stackrel{}{s}}\stackrel{}{s}}^\mu \|\stackrel{}{S},k\overline{\theta _\stackrel{}{s}}\|\stackrel{}{S},ku_\stackrel{}{s}\text{ (II B).}$$ (21) These are tensor forms as in equation (8). The on-shell condition on the right mover together with the trivial identity ($`k\mathrm{\Gamma }`$)$`{}_{}{}^{2}\|\stackrel{}{S},k\theta _\stackrel{}{s}=0`$ on the left mover imply, as in equation (9), that $$F_{(k)}=d\omega _{(k1)}.$$ (22) Note that, for the II A (II B) theory, $`\omega _{(p)}`$ in eq(18) does not fit into the gauge symmetry parameters of $`A_{(p)}`$ forms of II A (II B) theory in eq(9) since they share the same tensor index structures. In fact, for a $`p+1`$ form $`A_{(p+1)}`$, one needs a $`p`$ form $`\stackrel{~}{\omega _{(p)}}`$ symmetry parameters, as can be seen from its spacetime coupling to D-brane $$_{\text{world vol of D-brane}}A_{(p+1)}d^{p+1}\xi A_{\mu _1\mu _2\mathrm{}\mu _{p+1}}(x)_1x^{\mu _1}\mathrm{}_{p+1}x^{\mu _{p+1}},$$ (23) which implies a space-time gauge symmetry $$A_{(p+1)}A_{(p+1)}+d\stackrel{~}{\omega }_{(p)}.$$ (24) This justifies that no perturbative type II string state carries the R-R charge. On the other hand, it is well-known that each time we T-dualize in an additional direction the dimension of the D-branes goes down by one and the R-R forms lose an index. To include the right closed string zero-norm state $`\stackrel{~}{\omega }_{(p)}`$, one is thus forced to introduce the type I unoriented open string and T-dualizes k = odd (even) numbers of space-time coordinates and then takes the noncompact limit $`R0`$ for each compatified radius. For k = odd (even), one has type II A (II B) string states in the bulk far away from D-branes. The right $`\stackrel{~}{\omega }_{(p)}\omega _{(p+1)}^{(T)}`$ state, the T-dual R-R zero-norm state, is thus attached to the D p-brane for p= even (odd) in II A (II B) theory. Note that, near the D-branes, the orientation projection of the Type I theory leaves only one linear combination of two SUSY charges(SUSY zero-norm states in eq.(10)) of the Type II theory in the bulk. It is $`Q_\alpha ^{^{}}+(\mathrm{\Pi }_m^k\beta ^m\stackrel{}{Q^{}})_\alpha `$ with $`\beta ^m\mathrm{\Gamma }^m\mathrm{\Gamma }.`$ The T-dual R-R zero-norm states attached in the boundary of the open string 1-loop diagram with D-branes are the T-dual version of the R-R zero-norm states in the bulk of the closed string tree diagram. Our argument resolves the puzzle of seeming unwanted R-R zero-norm states in perturbative type II string spectrum and simultaneously motivates the introduction of D-branes into the theory which is complementary to the argument in Ref. The space-time T-dual R-R zero-norm state has an interesting analogy from worldsheet vertex operator point of view. The authors of considered one point function of R-R vertex operator on the disk. Since the total right + left ghost charge number must add up to $`2`$, one is forced to change the vertex operator in the conventional $`(\frac{1}{2},\frac{1}{2})`$ picture to either $`(\frac{1}{2},\frac{3}{2})`$ or $`(\frac{3}{2},\frac{1}{2})`$ picture. This inverse picture changing involves, among other unrelated things, a factor of $`k\mathrm{\Gamma }`$ , which shifts the field strength to the potential and gives a strong hint that D-brane carries the R-R charge. On the other hand, our space-time T-dual R-R zero-norm states do contain this important $`k\mathrm{\Gamma }`$ factor as can be seen from eqs (16) and (17). This again gives a strong support of our space-time T-dual R-R zero-norm state approach. In the next section, we will see an even more interesting application of these states. ## III Dilaton-Axion Symmetries and SL(2,Z) S-duality According to section II, for II A theory, we have $`A_{(1)},`$ $`A_{(3)}`$ potentials with $`d\omega _{(0)}^{(T)}`$, $`d\omega _{(2)}^{(T)}`$ T-dual zero-norm states and, for II B theory, we have $`A_{(2)}`$, $`A_{(4)}`$ potential with $`d\omega _{(1)}^{(T)}`$, $`d\omega _{(3)}^{(T)}`$ T-dual zero-norm states for their symmetry charge parameters. For completeness we have, in addition, an axion $`A_{(0)}\chi `$ in the II B theory. The corresponding T-dual zero-norm state is naturally identified to be the constant 0-form $`F_{(0)}^{(T)}`$, which is Poincare dual to the constant 10-form $`F_{(10)}^{(T)}d\omega _{(9)}^{(T)}`$. So we have the ”symmetry” $$\chi \chi +F_{(0)}^{(T)}.$$ (25) Note that, in eq. (8), there is a constant 10-form field strength $`G_{(10)}=dA_{(9)}`$ which is Poincare dual to the constant 0-form field strength in II A theory as well. This non-propagating degree of freedom can be included in the massive type II A supergravity and was conjectured to be related to the cosmological constant. See the interesting discussion of this 9-form potential $`A_{(9)}`$ by Polchinski in Ref.. Equation (21) is consistent with the fact that the axion $`\chi `$ is defined up to a constant. The interesting new result here is that we naturally identify this constant to be $`F_{(0)}^{(T)}`$. We now turn to the discussion of NS-NS dilaton $`\varphi `$. Remember we have a Remaining NS-NS singlet zero-norm state in equation (4). The physical meaning of this state will be discussed in the following. In reference each space-time symmetry of the bosonic background field in the first order WFA can be constructed through a superconformal deformation $$(T^{(1)}=\overline{T}^{(1)},T_F^{(1)},\overline{T}_F^{(1)})$$ (26) corresponding to a spacetime zero-norm state. In equation (22), $`T^{(1)}`$($`T_F^{(1)}`$) is the upper component (lower component) of deformation of the superstress tensor in the first order WFA when the background field is turned on. $`\overline{T}^{(1)},(\overline{T}_F^{(1)})`$ is its anti-holomorphic counterpart. It was shown that superconformal deformations constructed from zero-norm states in eq.(3) give the symmetries of graviton and antisymmetric tensor. The superconformal deformation constructed from the zero-norm state in eq.(4), which was neglected in the previous discussion., is calculated to be $`T^{(1)}`$ $`=`$ $`\overline{T}^{(1)}=_\mu _\nu \theta (x^\mu +\stackrel{}{_\lambda }\psi ^\lambda \psi ^\mu )(\overline{}x^\nu +\stackrel{}{_r}\overline{\psi }^r\overline{\psi }^\nu )`$ (27) $`=`$ $`_\mu _\nu \theta x^\mu \overline{}x^\nu ,`$ () $$T_F^{(1)}=\frac{1}{2}_\mu _\nu \theta \overline{}x^\nu \psi ^\mu ,$$ (28) $$\overline{T}_F^{(1)}=\frac{1}{2}_\mu _\nu \theta x^\nu \overline{\psi }^\mu $$ (29) with condition $`\mathrm{}\theta =constant`$, $`\mathrm{}^\mu _\mu .`$ $`\theta (x)`$ in eq(23) is the background field corresponding to the singlet zero-norm state of eq(4). The induced ”symmetry” is calculated to be $$\varphi \varphi +\mathrm{}\theta $$ (30) and $$h_{\mu \nu }h_{\mu \nu }+_\mu _\nu \theta .$$ (31) Equation (25) is merely a change of gauge in the linearized graviton and can be absorbed to the symmetry of the linearized graviton. Equation (24) is the ”symmetry” of the dilaton. The result that $`\mathrm{}\theta `$ is a constant is consistent with the fact that $`\varphi `$ appears in the effecive equation of motion, constructed from vanishing $`\sigma `$model $`\beta `$function, in an overall factor $`e^{2\varphi }`$ other than differentiated. The interesting result here is that we identify the constant $`\mathrm{}\theta `$ to be the zero-norm state in equation (4). This completes the physical effects of all massless zero-norm states in the type II string spectrum. The ”symmetries” presented in equations (21) and (24) were derived in the first order WFA. They can be broken in the higher order correction. However, if one tries to generalize the superconformal deformation to second order in the WFA, one immediately meets the difficulty of nonperturbative nonnormalizibility of the 2d $`\sigma mod`$el, and is forced to introduce counterterms which consist of an infinite number of massive tensor fields. This higher order effect is related to the stringy physics ($`\alpha ^{}\mathrm{}`$) of string theory instead of point particle field theory. In fact, it was known that there exist important stringy bound states called (p,q) string which consists of p F-strings and q D-strings in the II B theory. The coupling of axion $`\chi `$ to (p,q) string, which is an higher order effect and so can not be seen in our first order WFA, breaks the symmetry in equation (21) down to integer shifts. On the other hand, it was known that the symmetry in eq.(24) was broken down to the discrete $`\mathrm{}\theta 2\varphi `$ from the type II B supergravity. If we define $$\rho =\chi +ie^\varphi ,\rho \frac{\theta }{2\pi }+\frac{i}{g_s}\tau ,$$ (32) these two discrete symmetries combine to form the well-known SL(2,Z) S-duality symmetry of II B string. Note that the nonlinearity of SL(2,Z) does not appear in our linear WFA. This is a generic feature of WFA in contract to the usual $`\sigma `$-model loop($`\alpha ^{}`$) expansion. The former contains stringy phenomena (e. g. high energy symmetries) which can not be derived in the loop expansion scheme, while the latter is convenient to obtain the low energy effective field theory of the superstring. This will become clear when one considers the massive states of the string, which are crucial to make string theory different from the usual quantum field theory. An immediate application of this Type II B S-duality is the N=4, d=4 SUSY Yang-Mills S-duality, where dyon with the electric charge p and magnetic charge q can be interpreted as the end points of the (p,q) string on the D3-brane. The $`\tau `$ parameter in this SUSY gauge theory is defined to be $$\tau =\frac{\theta _{YM}}{2\pi }+\frac{i}{g_{YM}^2},$$ (33) and is interpreted to be the constant $`\rho `$ field of II B string in equation (26) associated to a stack of D3-branes. ## IV Conclusion T-dual R-R zero-norm states motivate the introduction of D-branes into Type II string theory. They serve as symmetry charge parameters of R-R tensor forms. The study in this paper reveals again that all space-time symmetries, including the discrete T-duality and S-duality, are related to the zero-norm states in the spectrum. The unified description of S and T dualities makes one to speculate that they are all geometric symmetries (due to the redefinition of string backgrounds) and to conjecture the existence of a bigger discrete U-duality symmetry, and its relation to the zero-norm state. In fact the SL(2,Z) S-duality of II B string led Vafa to propose a 12d F-theory, where $`\tau `$ is the geometric complex structure modulus of torus T<sup>2</sup>. One can even generalize this zero-norm state idea to construct new stringy massive symmetries of string theory. In particular, the existence of some massive R-R zero-norm states and other evidences make us speculate that string may carry some massive R-R charges. Another interesting issue is the identification of D-brane charges with elements of K-theory groups. How T-dual R-R zero-norm states relate to K-theory groups is an interesting question to study. ## V Acknowledgments I would like to thank Pei-Ming Ho and Miao Li for comments. This research is supported by National Science Council of Taiwan, under grant number NSC 89-2112-M-009-006.
warning/0007/cond-mat0007065.html	ar5iv	text	# Mean Field Theory of Spherical Gravitating Systems ## 1 Introduction The observation that a number of different types of astronomical objects appear to be in thermodynamically relaxed states has motivated theorists to understand the thermodynamics and statistical physics of self-gravitating systems. Of particular note are the globular clusters, consisting of about a million stars. Besides having relaxed cores, these structures appear to be organized in two distinct classes characterized by radically different density profiles, x-ray production, and other features. This suggests that globular clusters may exist in different thermodynamic phases. However, in contrast with normal ”chemical” systems, which have been successfully described by thermodynamics at the macroscopic level, both the infinite range and short distance singularity of the Newtonian gravitational potential introduce problems in the statistical theory of phase transitions which make their analysis a challenging task. The description of the system can be simplified by going to the Vlasov limit, i.e. by letting the number of particles become large while controlling the total mass, M, and energy, E. In this limit the system is described by the single particle density $`f(x,v,t)`$ in the $`\mu `$ (position, velocity) space, which is employed by most of the standard treatments, including the present work. We refer to this reduced description as mean field theory (MFT). While MFT avoids the problem of dealing with an N body formulation, the difficulties introduced by the singularity and long range of the potential persist. In the early sixties Antonov investigated the equilibrium behavior of isolated gravitational systems in MFT. To circumvent the problem of escape, he confined the mass to a finite region by introducing a rigid wall. By fixing the total mass and energy, he showed that maximum entropy solutions for $`f`$ are spherically symmetric in position and have the expected Maxwellian velocity dependence. However, he proved that there is no global maximum to the entropy: while extremal solutions can exist, these are at best local maxima. Antonov , as well as Lynden-Bell and Wood , closely investigated the spherically symmetric system confined in a sphere and showed that when the mass and energy are fixed there are no entropy extrema above a critical radius ($`R0.335GM^2/E`$) . When the radius is less than this value, the stability of the extremum solutions were studied by several authors (Lynden-Bell and Katz , Katz , Padmanabhan , and Bavaud ). They found that, in general, above a critical density contrast ($`\rho \left(0\right)/\rho \left(R\right)709`$) all extrema are unstable, i.e. they are not local maxima. Lynden-Bell and Wood coined this phenomena the gravothermal catastrophe and it is also referred to as the Antonov instability. In such a system, there is no upper bound on the entropy and a state of arbitrarily large entropy can be constructed from a centrally concentrated density profile by shifting more of the mass towards the center (core-halo structures have higher entropy). More recently, Kiessling has investigated the thermodynamic stability of the full N body point mass system confined in a spherical box using the canonical ensemble. To avoid the short range singularity he regularized the Newtonian interaction by softening it (letting the potential approach a finite value as the origin is closely approached). He showed that in the limit that the softening vanishes (Newtonian interaction), the canonical equilibrium measure is the superposition of Dirac measures at any temperature, meaning that the system is in a collapsed point mass state. We have to emphasize that this is the equilibrium solution when the system is in thermal equilibrium with a heat bath. Also, based on the results of the finite N-particle system, with proper scaling of the particle mass as we take the mean-field limit, he showed that the single-particle density function is proportional to the Dirac distribution. Therefore the system’s equilibrium state is the collapsed state in the canonical ensemble. Kiessling’s conclusions do not contradict the earlier work applying MFT to the microcanonical ensemble (fixed mass and energy) described above, since no global entropy maximum was found in that case either. In fact, the pure Newtonian potential is never a correct picture in general because of the finite size of stars and atoms. In a nearly hidden appendix of their paper describing the gravothermal catastrophe, it was first pointed out by Lynden-Bell and Wood that if we modify the singular $`1/r`$ Newtonian gravitational potential at the center by introducing a small minimal distance between the particles (the so called hard-sphere model) that complete collapse would be avoided and a global entropy maximum should exist. They further conjectured that in this situation a first order phase transition to a centrally concentrated core-halo configuration would occur as the system energy is reduced. Several authors demonstrated the existence of a phase-transition in special mean-field models with a modified gravitational potential. Hertel and Thirring were the first to show analytically that a gravitating system can undergo a first-order phase transition. Although their system is not purely gravitational and has no singularity, the pair-interaction potential is purely attractive and has a fixed value when the pair of particles are in a given sub-domain. Also Lynden-Bell and Lynden-Bell showed the occurrence of a first-order phase transition in their special gravitational system, consisting of point particles distributed on a shell which cannot shrink into a point mass (inner boundary) or expand to infinity (outer boundary). Kiessling et. al. also applied the hard-sphere model to study planet formation: They were able to explain the existence of observable planets by showing that the mass belonging to the condensed phase is well below the Jeans mass. For finite N-body systems, a gravitational first-order phase transition was first observed dynamically by Youngkins and Miller . They investigated a model consisting of irrotational, concentric, spherical mass shells confined between two rigid spherical boundaries. The system was studied both theoretically in the mean-field limit and numerically by N-body simulations in the microcanonical, canonical, and grand canonical ensembles. The analysis for this one dimensional system showed that the system undergoes a first-order phase transition instead of a gravothermal catastrophe. As expected in gravitational systems, there were some discrepancies in the results for different ensembles; however the numerical N-shell simulations were always in good agreement with the corresponding mean-field predictions. Much earlier, Eddington determined the form of the general stationary solution of the Vlasov equation for a spherical system with an anisotropic velocity distribution which obeys the Schwarzschild law . This model explicitly depends on the square of the angular momenta $`l^2`$, but it also includes the isothermal part: $`f(𝐫,𝐯)e^{\beta ϵ}e^{\gamma l^2}`$. The model was presented much earlier, but it may have been forgotten. Later, several phenomenological models were improved by including Eddington’s anisotropic term (King-Michie models and others ), giving a better fit to the observed density profiles of globular clusters. However, in some cases, a good fit was not obtained for globular clusters with core-halo structures. In other words, a fair number of globular clusters which have a small and very dense core surrounded by a thin halo structure, do not obey these empirical density-fit models. On the other hand, in addition to the system energy, the sum of the angular momenta squares $`L_2=l_i^2`$ is also an integral of the motion for any isolated spherically symmetric system in the mean-field limit , and should not be ignored. The purpose of our work is to present a more general approach for investigating the equilibrium properties of confined, spherical, systems then those of Antonov, Lynden-Bell and Wood, etc. mentioned above which takes into account both integrals of a spherical gravitating system, and to introduce idealized dynamical ”shell” models which also satisfy these constraints. Model gravitating systems consisting of a collection of concentric, infinitesimally thin, spherical shells were first introduced by Henon . They are useful for investigating the initial stages of evolution of a spherically symmetric self gravitating system, before the onset of binary formation arising from three body effects. They have the further advantages of ease and accuracy of algorithm construction, since it is possible to analytically solve for the motion of each shell between encounters, eliminating the need for the tedious and slow, step-wise integration of coupled, nonlinear, differential equations. . In the present work we consider the mean field theory of a system of gravitating point particles moving in three dimensional space, as well as that of thin, rotating, spherical mass shells with angular momentum vectors restricted to manifolds of one, two and three dimensions. We first determine conditions for the equilibrium one particle probability density function $`f(x,p)`$of a unit mass particle (shell) by finding the entropy extrema with respect to the constraints of (1) the normalization, (2) system energy $`E`$, and (3) the sum of the squares of the angular momentum, $`L_2`$. We then show that the introduction of the new integral $`L_2`$ suggests a new type of canonical ensemble ($`T\gamma )`$, in addition to the extension of the microcanonical ensemble ($`EL_2)`$ . A nonlinear differential equation governing the radial density valid for each ensemble is derived for the case of the three dimensional point mass system, and each shell system. We then prove that the radial density of the shell system with angular momentum confined to the Euclidean plane satisfies the identical differential equation as the three dimensional point mass system, and we carefully study the equilibrium solutions for this case numerically. The stability of the extremum solutions of each model is investigated in both microcanonical ($`EL_2)`$ and canonical ($`T\gamma )`$ ensembles. At first glance it is natural to anticipate that the centrifugal barrier associated with the additional constraint will eliminate any tendency for complete core collapse without introducing an inner boundary in the system or changing the gravitational potential by other means. We conclude by investigating the possible presence of a phase-transition which would remove the gravothermal catastrophe . ## 2 The Entropy Extrema ### 2.1 Spherically symmetric point mass systems The primary goal is to evaluate the equilibrium single particle probability density function $`f(x,p)`$ which maximizes the entropy in the mean-field limit. Consider the spherically symmetric isolated point mass system in three space dimensions with total mass $`M`$ confined in a sphere of radius $`b`$. We choose units where $`G=M=b=1`$ and introduce spherical coordinates $`x=(r,\phi ,\vartheta )`$ . The Lagrangian per unit mass of a single particle moving in the mean field potential $`\mathrm{\Phi }(r)`$ (2) is $$L=\frac{1}{2}\dot{r}^2+\frac{1}{2}r^2\left(\dot{\vartheta }^2+\dot{\phi }^2\mathrm{sin}^2\vartheta \right)\mathrm{\Phi }\left(r\right),$$ (1) where, for a Newtonian pair-wise interaction, $`\mathrm{\Phi }\left(r\right)`$ is given by : $$\mathrm{\Phi }\left(r\right)=\left[G(r,r^/)+G(r^/,r)\right]f(x^/,p^/)d^3x^/d^3p^/,$$ (2) where $$G(r,r^/)=\frac{\mathrm{\Theta }\left(rr^/\right)}{r}$$ and, as usual, $`\mathrm{\Theta }\left(r\right)`$ denotes the Heavyside step function. The Hamiltonian is then $$H=\frac{1}{2}v^2+\frac{l_\vartheta ^2}{2r^2}+\frac{l_\phi ^2}{2r^2\mathrm{sin}^2\vartheta }+\mathrm{\Phi }\left(r\right)$$ (3) where $`p=(v,l_\phi ,l_\vartheta )`$ are the corresponding canonical momenta. Our plan is first to construct the entropy extrema, and then verify whether or not the solutions are local maxima. The entropy of the system is $$S=f\mathrm{ln}fd^3xd^3p,$$ (4) and the constraints for which we need to find the extremum are: (1) normalization of $`f`$, (2) energy conservation in the complete system, and (3) conservation of $`L_2`$ : $$1=fd^3xd^3p$$ (5) $$L_2=l^2fd^3xd^3p$$ (6) $$E=f\left(\frac{1}{2}v^2+\frac{1}{2r^2}l^2+\frac{1}{2}\mathrm{\Phi }\left(r\right)\right)d^3xd^3p$$ (7) where $$l^2=\left(l_\vartheta ^2+\frac{l_\phi ^2}{\mathrm{sin}^2\vartheta }\right).$$ Introducing Lagrange multipliers $`\alpha ,`$ $`\beta ,`$ $`\gamma `$ we have an extremum for the functional $`S`$ when $$0=\delta \left(S\alpha \beta E\gamma L_2\right).$$ (8) Taking the first variation of each term in (8), and asserting (5), (7), (6), and (2), we obtain: $$0=\left\{\left(\mathrm{ln}f+1\right)+\alpha +\beta \left[\frac{1}{2}v^2+\frac{1}{2r^2}l^2+\mathrm{\Phi }\left(r\right)\right]+\gamma l^2\right\}\delta fd^3xd^3p$$ from which finally $$0=\left(\mathrm{ln}f+1\right)\alpha \beta \left(\frac{1}{2}v^2+\mathrm{\Phi }\left(r\right)\right)\left(\frac{\beta }{2r^2}+\gamma \right)l^2.$$ (9) Thus the one-particle probability density function (pdf) is $$f=\mathrm{exp}[\left(\alpha +1\right)\beta \left(\frac{1}{2}v^2+\mathrm{\Phi }\right)\left(\frac{\beta }{2r^2}+\gamma \right)l^2]..$$ (10) In order to obtain the radial density $`\rho \left(r\right)`$, we have to integrate $`f`$ over the other variables. To ensure that the integrals over $`v`$ , $`l_\vartheta `$, and $`l_\phi `$ converge, the following conditions are necessary:$`\beta >0`$ and $`\beta /2r^2+\gamma >0`$ at any $`r`$. Therefore the second condition is $`2\gamma /\beta >1/b^2=1`$. Using $`K=\mathrm{exp}\left[\left(\alpha +1\right)\right]`$ , we get $$\rho \left(r\right)=fd^3p𝑑\phi 𝑑\vartheta =_0^{2\pi }_0^\pi K\sqrt{\frac{2\pi }{\beta }}\frac{\pi \mathrm{sin}\vartheta }{\frac{\beta }{2r^2}+\gamma }\mathrm{exp}\left(\beta \mathrm{\Phi }\right)𝑑\phi 𝑑\vartheta $$ $$=\frac{K\left(2\pi \right)^{\frac{5}{2}}}{\sqrt{\beta }}\left(\frac{\beta }{2r^2}+\gamma \right)^1e^{\beta \mathrm{\Phi }\left(r\right)}$$ (11) and, with the Poisson equation in a spherical coordinate system, $$\mathrm{\Delta }\mathrm{\Phi }=\frac{1}{r^2}\frac{d}{dr}\left(r^2\frac{d\mathrm{\Phi }}{dr}\right)=4\pi \rho _V$$ where $`\rho __V`$ is the volumetric mass density. In some cases, it’s more convenient to use the linear (radial) density instead of the volume density: $$\frac{d}{dr}\left(r^2\frac{d\mathrm{\Phi }}{dr}\right)=\rho \left(r\right).$$ (12) Note that because $`M=1`$, the radial probability density function and the linear mass density function are the same. Introducing a new function $`\mathrm{\Psi }=\beta \mathrm{\Phi }`$ and employing 11 we can rewrite (12) as $$\frac{d}{dr}\left(r^2\frac{d\mathrm{\Psi }}{dr}\right)=\rho \left(r\right)=K\left(2\pi \right)^{\frac{5}{2}}\sqrt{\beta }\left(\frac{\beta }{2r^2}+\gamma \right)^1e^{\mathrm{\Psi }\left(r\right)}$$ obtaining a closed equation for $`\mathrm{\Psi }.`$ This in turn can be simplified by introducing constants $`C`$ and $`\mathrm{\Gamma },`$ $$\frac{d}{dr}\left(r^2\frac{d\mathrm{\Psi }}{dr}\right)=C\left(\frac{1}{r^2}+\mathrm{\Gamma }\right)^1e^{\mathrm{\Psi }\left(r\right)},$$ (13) $$C=K\left(2\pi \right)^{\frac{5}{2}}\frac{2}{\sqrt{\beta }}>0,$$ (14) $$\mathrm{\Gamma }=\frac{2\gamma }{\beta }>1.$$ This is the final form of the differential equation for the scaled potential which we will solve using numerical methods. Finally, we can evaluate $`E`$, $`L_2`$, and $`S`$ in terms of the Lagrange multipliers, the density $`\rho \left(r\right)`$, and the potential $`\mathrm{\Phi }\left(r\right).`$ After integrating (7), (6), and (4), we obtain: $$S=\alpha +\frac{5}{2}+\beta _0^b\rho \mathrm{\Phi }𝑑r$$ (15) $$L_2=_0^b\rho \left(\frac{\beta }{2r^2}+\gamma \right)^1𝑑r$$ (16) $$E=\frac{1}{2\beta }+_0^b\rho \left(\frac{1}{\beta +2\gamma r^2}+\frac{\mathrm{\Phi }}{2}\right)𝑑r$$ (17) ### 2.2 Shell systems In this section we consider the system of spherically symmetric, infinitesimally thin, mass shells confined in a sphere with radius $`b.`$ We can define three basic types of the model with dimensions $`d=2,3,4`$. The one-dimensional $`d=1`$ non-rotating shell system is discussed in . In the $`d=2`$ case, every shell rotates about a fixed axis. When $`d=3,`$ the rotational axis of every shell is in a fixed plane, while $`d=4`$ means that every shell can rotate about any arbitrary axis. Therefore we use $`d1`$ angles as coordinates. Let us consider these systems in the mean-field limit again using units where $`M=G=b=1`$. With the potential discussed above (2), the Lagrangian and the Hamiltonian of a unit mass shell in a spherical coordinate system are, respectively,: $$L=\frac{1}{2}\stackrel{.}{r}^2+\frac{1}{3}r^2\underset{k=1}{\overset{n}{}}\dot{\phi }_k^2\mathrm{\Phi }\left(r\right)$$ (18) $$H=\frac{1}{2}v^2+\frac{3}{4}\frac{_{k=1}^nl_k^2}{r^2}+\mathrm{\Phi }\left(r\right)$$ (19) where $`\phi _k`$ $`\left(k=1,2,3\right)`$ are the angles around $`x,`$ $`y`$ and $`z,`$ the $`l_k`$ are the $`x,y,z`$ components of the angular momentum per unit mass, and we have used the fact that the moment of inertia of a shell with unit mass is $`2/3r^2.`$ In the equations above, we use $`n=d1`$ which is the number of degrees of freedom coming only from rotation. This is a generalization of the previous model, and so is the method to find the equilibrium pdf, $`f(x,p)`$. As in the previous section, in order to get the equilibrium solutions first we have to find the extremum of the entropy: $$S=f\mathrm{ln}fd^{n+1}xd^{n+1}p$$ (20) with respect to the three constraints of normalization, $`E,`$ and $`L_2:`$ $$1=fd^{n+1}xd^{n+1}p,$$ $$L_2=f\left(\underset{k=1}{\overset{n}{}}l_k^2\right)d^{n+1}xd^{n+1}p,$$ (21) $$E=f\left(\frac{1}{2}v^2+\frac{3}{4}\frac{_{k=1}^nl_k^2}{r^2}+\frac{1}{2}\mathrm{\Phi }\left(r\right)\right)d^{n+1}xd^{n+1}p.$$ (22) The solution for the variational problem can be obtained easily, and the one-particle density function is now $$f=\mathrm{exp}\left[\left(\alpha +1\right)\right]\mathrm{exp}\left[\beta \left(\frac{1}{2}v^2+\mathrm{\Phi }\right)\right]\mathrm{exp}\left[\left(\frac{3\beta }{4r^2}+\gamma \right)\underset{k=1}{\overset{n}{}}l_k^2\right].$$ (23) Therefore the radial mass density function is $$\rho \left(r\right)=fd^{n+1}pd^n\phi =K\left(2\pi \right)^{\frac{2n+1}{2}}\frac{\pi ^{\frac{n}{2}}}{\sqrt{\beta }}\left(\frac{3\beta }{4r^2}+\gamma \right)^{\frac{n}{2}}\mathrm{exp}\left(\beta \mathrm{\Phi }\right).$$ (24) From the Poisson equation, again using using $`\mathrm{\Psi }=\beta \mathrm{\Phi },`$ we find, $$\frac{d}{dr}\left(r^2\frac{d\mathrm{\Psi }}{dr}\right)=K\left(2\pi \right)^{\frac{2n+1}{2}}\pi ^{\frac{n}{2}}\sqrt{\beta }\left(\frac{3\beta }{4r^2}+\gamma \right)^{\frac{n}{2}}\mathrm{exp}\left(\mathrm{\Psi }\right)$$ (25) which can be simplified to $$\frac{d}{dr}\left(r^2\frac{d\mathrm{\Psi }}{dr}\right)=C\left(\frac{1}{r^2}+\mathrm{\Gamma }\right)^{\frac{n}{2}}e^{\mathrm{\Psi }\left(r\right)},$$ (26) $$C=\frac{2^{\frac{4n+1}{2}}}{(3\beta )^{\frac{n}{2}}}\pi ^{\frac{3n+1}{2}}K\sqrt{\beta }>0,$$ (27) $$\mathrm{\Gamma }=\frac{4\gamma }{3\beta }>1.$$ (28) Comparing the results (26) for the 3 dimensional shell system ($`n=2`$) to that of a 3 dimensional point mass system (13), we can see that these two systems are equivalent in so far as (13 ) and (26) have the same form. It is useful to evaluate $`S,L_2,`$ and $`E`$ in terms of the Lagrange multipliers, $`\rho \left(r\right)`$, and $`\mathrm{\Phi }\left(r\right).`$ The integration of (20), (21), and (22) yields $$S=\alpha +\frac{n+3}{2}+\beta _0^b\rho \mathrm{\Phi }𝑑r$$ (29) $$L_2=\frac{n}{2}_0^b\rho \left(\frac{3\beta }{4r^2}+\gamma \right)^1𝑑r$$ (30) $$E=\frac{1}{2\beta }+_0^b\rho \left(\frac{3n}{8}\frac{1}{\frac{3}{4}\beta +\gamma r^2}+\frac{\mathrm{\Phi }}{2}\right)𝑑r.$$ (31) ## 3 T-$`\gamma `$ ensemble In the previous sections we derived expressions for the entropy extremum solutions for our model systems in terms of the local radial density. A third Lagrange multiplier,$`\gamma `$, was introduced to satisfy the new constraint on $`L_2`$ . Thus specifying $`E`$ and $`L_2`$ (and, of course, $`M`$) defines the analogue to the microcanonical ensemble for these systems. Alternatively, by fixing $`\beta `$ and $`\gamma `$, we can define the analogue of the canonical ensemble, in which $`E`$ and $`L_2`$ are not fixed, but their average is determined by $`\beta `$ and $`\gamma `$ , where $`\gamma =\frac{S}{L_2}`$. We call this the $`T\gamma `$ ensemble. It can be modeled by imagining that the system is in contact with a heat bath with constant $`\beta `$ and also an $`l^2`$ bath at constant $`\gamma `$. The new $`l^2`$ bath corresponds to the fact that we allow some $`l^2`$ exchange between the system and the bath. Since the system is spherically symmetric in position, while the system can also exchange angular momentum with the bath, its vector average will vanish. As an example, we can imagine a globular cluster which is embedded in some large spherically symmetric stellar neighborhood with an isotropic velocity distribution. The mean angular momentum of both system and bath is zero, and only $`l^2`$ and energy exchange can occur (for the moment, we do not take into account the possibility that particles can escape from the cluster). In performing calculations it is more convenient to use the $`T\gamma `$ ensemble than the $`EL_2`$ microcanonical ensemble because in the latter we have to find the Lagrange multipliers from the given $`E`$ and $`L_2`$. As can be seen from 15, 16, and 17 this is a nontrivial and laborious task. In this ensemble the relevant thermodynamic potential is an extension of the Helmholz free energy: $$F=E\frac{1}{\beta }S+\frac{\gamma }{\beta }L_2$$ (32) and equilibrium states, if they exist, minimize $`F.`$ ## 4 Stability From the variational problem, we only know the extremum solutions. In order to separate the unstable solutions from the locally stable, we use the modified method of Poincare’s linear series of equilibria. Following Katz, we can generalize the method to the case of a functional. From now on, we discuss the generalized version of the method which has to be applied to our models to determine whether an extremum solution is stable or unstable. We outline the method below. Let’s assume that we want to find the maximum of the functional $`F^{}(f,s)`$ which depends parametrically on $`s.`$ The function $`f:\mathrm{\Omega }R`$, where $`\mathrm{\Omega }`$ is a compact domain $`\mathrm{\Omega }R^3`$, and the local maximum of the functional is the stable solution. Suppose we partition $`\mathrm{\Omega }`$ and consider the vector $`xR^n`$ with elements $`x_i=f\left(y_i\right)`$ ($`y_i\mathrm{\Omega }`$ is not a coordinate but an indexed element in $`\mathrm{\Omega }`$). Instead of dealing with the functional $`F^{}`$, we can construct a function $`F:R^nR`$ such that $`F(x,s)F^{}(f,s)`$. The accuracy can be controlled by refining the partition (increasing n). The problem is then shifted to finding the extrema of $`F`$: $$_iF(x,s)=0$$ (33) Denote the extremum solutions of the problem by $`x=\left\{X_a\left(s\right)\right\}`$ where $`a,b=1\mathrm{}N`$, labels different extremal solutions. Assume that the first and second derivatives of $`F`$ are continuous in $`x`$ and $`s,`$ and $`\dot{X}_a`$ is continuous as well. Assume also that the matrix $`\left(_i_jF\right)_a`$ has a non-degenerate eigenvalue spectrum, which we may consider to be ordered: $`k_{1a}\left(s\right)<k_{2a}\left(s\right)<`$ . . .$`<k_{na}\left(s\right)<`$ …, and further that $`\left(_i_jF\right)_a`$ is diagonal. (We can always transform it into that form.). Let’s evaluate $`F`$ at the extremum points $`x=X_a\left(s\right)`$ where $$_iF(X_a,s)=0.$$ (34) Then, as the parameter $`s`$ is varied, on the extremum labeled $`a`$, $$0=\frac{d}{ds}_iF(X_a,s)=\left(_s_iF\right)_a+\underset{j=1}{\overset{\mathrm{}}{}}\left(_j_iF\right)_a\dot{X}_a^j=\left(_s_iF\right)_ak_{ia}\left(s\right)\dot{X}_a^i,$$ (35) $$\dot{X}_a^i=\frac{\left(_s_iF\right)_a}{k_{ia}\left(s\right)}.$$ (36) The stability of the extremum is determined by the second derivative of $`F_a\left(s\right)=F(X_a,s)`$, $$\ddot{F}_a=\left(_s^2F\right)_a+\underset{i=1}{\overset{\mathrm{}}{}}\left(_i_sF\right)_a\dot{X}_a^i=\left(_s^2F\right)_a+\underset{i=1}{\overset{\mathrm{}}{}}\frac{\left(_s_iF\right)_a^2}{k_{ia}}.$$ (37) The stability will change only when one of the $`k_{ia}`$ changes sign, and a change in stability occurs only at bifurcation or limit points . Therefore we have to investigate the dependence of $`\ddot{F}_a`$ on $`s`$ in order to decide how the stability changes at an $`ab`$ bifurcation or limit point . If we look at (37), one can see that when a particular $`k_{ia}\left(s_0\right)`$ changes sign from positive to negative, at that point $`\underset{ss_0}{lim}\ddot{F}_a=+\mathrm{}`$ and similarly, from the other $`b`$ branch, $`\underset{ss_0}{lim}\ddot{F}_b=\mathrm{}`$. Going back to our original problem, we have to apply the method to our model systems both in the $`EL_2`$ and $`T\gamma `$ ensembles. The only difficulty is that the method discussed above only allows us to include systems with one control parameter, $`s`$. In our case, we are free to fix either $`E`$ or $`L_2`$ in the $`EL_2`$ ensemble and $`\beta `$ or $`\gamma `$ in the $`T\gamma `$ ensemble and regard the other as the control parameter. Afterwards, we can change the previously fixed parameter and apply the method again for a wide range of parameter sets. In the $`EL_2`$ ensemble a natural choice of parameter is $`E`$ with one fixed value of $`L_2`$, as the entropy has to be a local maxima if the system is in a locally stable state. If we are at the extremum solution points $$\dot{S}=\frac{dS}{dE}=\frac{S}{E}=\beta .$$ (38) For a different value of $`L_2`$, we can clearly see which branch of the extremum solutions are unstable. To obtain a complete description we can use the same method if $`E`$ is fixed: $$\dot{S}=\frac{dS}{dL_2}=\frac{S}{L_2}=\gamma .$$ (39) In the $`T\gamma `$ ensemble, first consider the case where we fix $`\gamma ,`$ and we are looking for the maximum of the functional $`\beta F`$ . On a branch of the extremum solutions $$\frac{dF}{d\beta }=\frac{E}{\beta }+\frac{S}{\beta ^2}\frac{1}{\beta }\frac{S}{E}\frac{E}{\beta }\frac{1}{\beta }\frac{S}{L_2}\frac{L_2}{\beta }\frac{\gamma }{\beta ^2}L_2+\frac{\gamma }{\beta }\frac{L_2}{\beta }.$$ (40) Using $`\frac{S}{E}=\beta `$ and $`\frac{S}{L_2}=\gamma `$, we easily find $$\frac{dF}{d\beta }=\frac{S}{\beta ^2}\frac{\gamma }{\beta ^2}L_2,$$ $$\frac{d\left(\beta F\right)}{d\beta }=\left(F+\beta \frac{dF}{d\beta }\right)=E.$$ (41) If we construct the stable branches for several values of $`\gamma `$, we can build up a general picture of the stability of the extremum solutions. We can apply the method for the case of a fixed $`\beta `$ as well. The extremum solutions are stable when $`\beta F`$ is maximum: $$\frac{dF}{d\gamma }=\frac{E}{\gamma }\frac{1}{\beta }\frac{S}{E}\frac{E}{\gamma }\frac{1}{\beta }\frac{S}{L_2}\frac{L_2}{\gamma }+\frac{1}{\beta }L_2+\frac{\gamma }{\beta }\frac{L_2}{\gamma },$$ (42) $$\frac{dF}{d\gamma }=\frac{1}{\beta }L_2,$$ $$\frac{d\left(\beta F\right)}{d\gamma }=L_2.$$ (43) From the results derived above we have an easy to apply tool to separate the unstable solutions. In the microcanonical ($`EL_2`$) description, we have to inspect the extremum solutions in the $`\beta E`$ plane for several fixed values of $`L_2`$ or in the $`\gamma L_2`$ plane for several fixed values of $`E`$. However, in the $`T\gamma `$ ensemble, we have to inspect the extremum solutions in the $`\left(E\right)\beta `$ plane for fixed $`\gamma `$ values, or we have to consider the extremum solutions in the $`\left(L_2\right)\gamma `$ plane for the case of fixed $`\beta `$. Either way, as we will show below, we can generalize the approach to two parameters. ## 5 Numerical method In order to find the entropy extremum solutions for the systems above, we have to solve (13) or (26), for the 3D point mass system, or the three types of shell system. Generally, all of the differential equations can be written in the form of $`{\displaystyle \frac{dy}{dr}}`$ $`=`$ $`C\left({\displaystyle \frac{1}{r^2}}+\mathrm{\Gamma }\right)^{\frac{n}{2}}e^{\mathrm{\Psi }\left(r\right)}`$ (44) $`{\displaystyle \frac{d\mathrm{\Psi }}{dr}}`$ $`=`$ $`{\displaystyle \frac{y}{r^2}}`$ where $`\mathrm{\Psi }=\beta \mathrm{\Phi }`$ and $`n=d1`$ (as above, d is the dimension of the system). Of course, the parameters $`\mathrm{\Gamma }>1`$ and $`C>0`$ are different case by case. But it’s quite interesting to mention that these systems behave similarly, regardless of whether we deal with a point mass system or a shell system, as long as the dimensions are the same. This is no longer a surprise if we look back at the Hamiltonians. But the similarity does also mean that, for example, a three dimensional point mass system is equivalent to a three dimensional shell system as far as the one particle density function and the stability of solutions are concerned. Also another advantage of the analogy is that we can dynamically model a three dimensional point mass system with a 3 dimensional shell system. Both systems should show the same equilibrium properties in the mean-field limit. Of course, we have to reassign the momentum of inertia of a shell to a different value in order to get exactly the same Hamiltonian in both cases. Equation (44) should be solved with the following boundary conditions: $`y\left(0\right)`$ $`=`$ $`0`$ (45) $`\mathrm{\Psi }\left(1\right)`$ $`=`$ $`\beta `$ Unfortunately, there are two things which make finding the solutions more difficult. First of all, (44) has a singularity at $`r=0`$ and, secondly, the existence and uniqueness of the solution are questionable for any given $`C`$, $`\beta `$, and $`\mathrm{\Gamma }.`$ We cannot simply set $`C`$ and $`\mathrm{\Gamma }`$ to satisfy our constraints of specified $`E`$ and $`L_2`$ because we don’t have explicit forms of $`E`$ and $`L_2`$ in terms of $`C`$ and $`\mathrm{\Gamma }`$ : the constraints are functionals of $`\rho `$ and $`\mathrm{\Phi }`$. To eliminate this problem we can change our boundary condition problem to an initial value problem because, in the latter case, we can ensure the existence and uniqueness of our solution for any $`C>0`$ and $`\mathrm{\Gamma }>1`$. Therefore we choose $`y\left(0\right)`$ $`=`$ $`0`$ (46) $`\mathrm{\Psi }\left(0\right)`$ $`=`$ $`\mathrm{\Psi }_0`$ where $`\mathrm{\Psi }_0(\mathrm{},+\mathrm{})`$, and we are able to construct the solution around $`r=0`$ in the form of a power series of $`\mathrm{\Psi }`$ and $`y`$. From the numerical point of view, we should take the power series of $`\mathrm{\Psi }`$ and $`y`$ in $`r_1`$ (a sufficiently small radius around $`r=0`$), and then continue the integration of (44) numerically from that point with the new initial values $`y\left(r_1\right)`$ $`=`$ $`y_{01}`$ (47) $`\mathrm{\Psi }\left(r_1\right)`$ $`=`$ $`\mathrm{\Psi }_{01}`$ Although the solution we get from (46) is unique, it does not satisfy both (45) and normalization. However, we can prove that we can’t find any $`\chi _1,\chi _2V_0`$ such that $`\mathrm{\Psi }_1=\mathrm{\Psi }_2+c`$, where $`V_0`$ consists of those solutions $`\chi =(y,\mathrm{\Psi })`$ which satisfy the same initial value problem of (44) and (46) with $`C=C_0>0`$ and any $`\mathrm{\Gamma }>1`$. Therefore all of the solutions are unique with fixed $`C`$ and $`\mathrm{\Gamma }`$ or, in other words, all $`\chi V_0`$ are unique. Now let’s consider $`V_1`$, which can be defined the same way as we defined $`V_0`$, but with $`C=C_1C_0`$. We can also prove that for any $`\chi _2V_1`$, there is only one $`\chi _1V_0:\chi _2=\chi _1(0,c)`$. Thus we can find all of the physically unique solutions by solving the initial value problem with a fixed $`C_0.`$ In practice, we used the Bulirsch-Stoer method to integrate the coupled, nonlinear, differential equations from $`y_1`$ and Bode’s five point method to evaluate integrals. Relative errors were controlled to within 10<sup>-12</sup>. ## 6 Numerical results The numerical solutions of (44) are presented in Figure 1 for the case of the spherically symmetric point mass system. The numerical results for the other models are qualitatively very similar, so we will use the point mass system to demonstrate their general features. As we can see, the solutions are bounded and, as a comparison, the results for the familiar isothermal sphere model ($`\mathrm{\Gamma }=0)`$ are presented as well. Examination of Figure 1 clearly demonstrates the existence of an upper bound, $`\beta _c\left(\mathrm{\Gamma }\right),`$ which means that below a critical temperature there’s no extremum solution at fixed $`\mathrm{\Gamma }`$. In fact, this can be proved rigorously from the differential equations. We can also see from the graph that $`\dot{\beta }_c=\frac{d\beta _c}{d\mathrm{\Gamma }}<0`$ and that, when $`\beta \beta _0`$ $`=\underset{\mathrm{\Psi }_0\mathrm{}}{lim}\beta \left(\mathrm{\Psi }_0\right)`$ , the number of solutions goes to infinity. In Figure 2, we show plots of three solutions for a given temperature (we don’t know so far which of them, if any, are stable), and in Figure 3 we plot the volume density profiles of these solutions which have been normalized to the central density. As we can see, solutions represented by smaller $`\mathrm{\Psi }_0`$ are more and more concentrated at the center. As a comparison, we also give the density profiles of the well known isothermal sphere ($`\mathrm{\Gamma }=0`$ ). There’s some difference in the shape of the density profiles in the case where $`\mathrm{\Gamma }0`$ but, in general, the volume density profile is singular as $`\mathrm{\Psi }_0\mathrm{},`$ and has the asymptotic solution $`\rho _V=\frac{1}{4\pi r^2}`$ . Note that the linear density is $`\rho =1`$ in this asymptotic case. In order to see the difference between density profiles of the isothermal sphere ( $`\mathrm{\Gamma }=0`$) and the others, we plot the high temperature solutions (Figure 4). These are the $`\mathrm{\Psi }_0\mathrm{}`$ asymptotic solutions where $`\beta 0`$. If $`\mathrm{\Gamma }>0`$ the relative volume density profiles curve down, but when $`\mathrm{\Gamma }<0`$ the profiles curve up, indicating that if we have an $`l^2`$ reservoir, at higher radius the density profile should change from the homogeneous density profile. In order to show how the $`L_2`$ constraint affects the shape of the density profiles, in Figure. 4 we plot the relative volume density profiles for a high temperature. As we can see in the figure, when $`\gamma 0`$ the density profiles are no longer uniform and, depending on the sign of $`\gamma `$, the density is either increasing $`\left(\mathrm{\Gamma }<0\right)`$ or decreasing $`\left(\mathrm{\Gamma }>0\right)`$ while in the standard case $`\left(\mathrm{\Gamma }=0\right)`$ in the limit of high temperature, the density is uniform. We can understand this behaviour if we recognize that in the limit $`\beta 0`$ gravity can be neglected. Consider 10 when $`\gamma 0`$. Since the kinetic energy contribution is still Maxwellian, the probalility of finding a particle with large $`l`$ is smaller when $`\mathrm{\Gamma }>0`$ $`\left(\gamma >0\right),`$ which means fewer particles will occupy larger radii. From the physical point of view, for a relatively large $`L_2,`$ more particles should concentrate at larger radii in order to maintain the large value, while for relatively small $`L_2,`$ fewer particles should settle at large radii in order to balance the centrifugal forces. For the case where $`\mathrm{\Gamma }<0`$ the situation is the opposite, and the density profile should increase with increasing radius. Of course, while we cannot use this argument for finite temperatures, the origin of the difference in the density profiles when $`\mathrm{\Gamma }0`$ is due to this effect . At finite temperature the tendency persists but the behavior is not guaranteed. In Figures 5 and 6, the stability properties are presented according to the previously discussed Poincare’s linear series of equilibria with fixed $`L_2`$ and $`E`$ in the $`EL_2`$ ensemble. As we can see, the gravothermal catastrophe holds for the $`EL_2`$ ensemble as well. In Figure 5, below a critical energy there’s no extremum solution (here the radius is fixed, not the energy as in). Also, at the same point, $`\dot{\beta }_a\left(E_c\right)=\mathrm{}`$ jumps to $`\dot{\beta }_b\left(E_c\right)=\mathrm{}.`$ ( We follow the spiral in the counter-clockwise direction.) Thus only the branch labelled $`a`$ is stable, and the other branches are increasingly unstable because additional eigenvalues become negative. In the opposite case, when we fix $`E`$ (Figure 6), we can see that above a critical value of $`L_2`$ there’s no extremum solution. The corresponding functions, $`S\left(E\right)`$ and $`S\left(L_2\right),`$ are plotted in Figure 7. We see that the entropy is monotonic up to a critical energy in the stable region of the extremum solutions. In the $`T\gamma `$ ensemble, the results are presented in Figures 8 and 9. First of all, in the case of fixed $`\gamma `$, we have to take a close look at the solutions in the $`\left(E\right)\beta `$ graph. Only the first branch of the solutions are locally stable up to a critical value of $`\beta `$ , say $`\beta _{c\text{,}}`$ because $`\left(\dot{E}\right)_a\left(\beta _c\right)=\mathrm{}`$, and there’s no extremum for $`\beta >\beta _c`$. As a comparison, we selected a value of $`\mathrm{\Gamma }`$ in Figure 2 corresponding to Figure 8; as we wind along the curve to a particular $`\beta _0`$ value, the extremum solutions are represented by a more concentrated set of density profiles. Another result of the stability investigations is that in Figure 1, only those solutions are locally stable which are to the right of the largest maximum of the $`\beta \left(\mathrm{\Psi }_0\right)`$ curves. The corresponding -$`F`$ can be seen in Figure 10. In the stable region, it’s a monotonic function, and therefore there’s no sign of a phase transition. The same holds in the $`EL_2`$ ensemble as can be seen by inspecting the entropy curves in Figure 7. The results for fixed $`\beta `$ are shown in Figure 9. The stable region of extremum solutions becomes unstable at $`\gamma _c`$ and at $`\gamma >\gamma _c`$ there are no extrema. The free energy behaves similarly to the case of fixed $`\gamma `$ : it’s monotonic, and there’s no phase transition. ## 7 Conclusion The main purpose of this work is to study the equilibrium thermodynamics of spherically symmetric self-gravitating systems in the mean field limit. We investigated both the spherically symmetric point mass system and shell systems of differing dimension confined in a sphere. These systems are related to each other both in the form of their Hamiltonian and their equilibrium states. Furthermore, the three dimensional shell system has the interesting and potentially useful property that, with regard to the equilibrium density profile, it is equivalent to the spherically symmetric point mass system. Our description is more general than the standard treatment of the isothermal sphere since, in addition to the energy, we take into account $`L_2`$, the sum of the squares of the individual angular momentum, which is conserved in the mean-field limit for spherically symmetric systems. In this new type of microcanonical description ($`EL_2`$ ensemble) we evaluated the ”equilibrium” one-particle probability density function for each type of system by finding the extrema of the entropy. The resulting pdf’s turned out to be similar to Eddington’s anisotropic density function. Therefore the density profiles obtained here also differ from the isothermal sphere which, in our formulation, is the special case $`\gamma =0.`$ Near the system center, the density profiles are similar to those of the isothermal sphere. However, as the outer boundary is approached, depending on the value of $`\mathrm{\Gamma }`$, deviations can become large, increasing or decreasing depending on the sign of $`\mathrm{\Gamma }`$. The physics behind this behavior is simple: if the system is spun up corresponding to negative $`\mathrm{\Gamma }`$, the outer density increases. If, on the other hand, the radial kinetic energy dominates the rotational energy, $`\mathrm{\Gamma }0`$ and the outer density decreases. In addition to the microcanonical ensemble, the Lagrange multiplier $`\gamma ,`$ which arose from the constraint on $`L_2,`$ yielded a new type of canonical ensemble ($`T\gamma )`$ which corresponds to the system being embedded in a heat bath at temperature $`T`$ and an $`l^2`$ reservoir at $`\gamma `$. But this analogy is only correct if the system is, at the least, in a local equilibrium state, and this issue demonstrates the importance of checking the stability properties of the extremum solutions. The method known as Poincare’s linear series of equilibria proved adequate to analyze the stability of the extremum solutions in both ensembles with a little extra effort. Using it, we showed that only certain types of solutions are locally stable, and others are saddle points. In other words, the gravothermal catastrophe is also present both in the $`EL_2`$ and $`T\gamma `$ ensembles, which means that there is no phase transition in our spherical mean-field models, although one was expected to be present because of the $`L_2`$ constraint. From their description of the gravothermal catastrophe it is easy to imagine that Lynden-Bell and Wood had in mind a dynamical process in which mass was transferred from the halo to a concentrated central core. However, their approach was confined to a comparison of stationary states. Hints of collapse have been seen in some N particle simulations. In future work we plan to study the complete dynamics of the collapse, both analytically and using dynamical simulation of N particle and N shell systems. The shell systems should prove especially useful since they avoid the complication arising from the formation of tight binaries. ## 8 Acknowledgments The authors benefitted from conversations with Igor Prokhorenkov and Michael Kiessling. They also are grateful for the support of the Research Foundation and Department of Information Services of Texas Christian University.
warning/0007/hep-th0007177.html	ar5iv	text	# General brane cosmologies and their global spacetime structure ## I Introduction The idea that our four-dimensional universe might in fact, upon closer inspection, be higher dimensional has always been compelling, although it has been conventionally supposed that the extra dimensions are small and compact. An almost diametrically opposed view, wherein the extra dimensions are large and noncompact has gained ground recently, in which our universe is a ‘defect’ (see for early work) in an anti-de Sitter bulk – the negative cosmological constant providing an effective or exotic compactification of the spacetime. In most models, the universe is taken to be a domain wall, (i.e. of codimension one – although models of higher codimension have been considered ) since this is most naturally supported by heterotic string theory . The original models of Randall and Sundrum focussed on the main features of having either one or two domain wall universes at orbifold fixed points of a single extra dimension; in the case of the two-wall model the second wall had negative energy and tension. The compelling feature of these brane-world models is that gravity has a four-dimensional character on the universe-domain-wall with short-range five-dimensional corrections coming from Kaluza-Klein (KK) modes. Indeed, it was demonstrated explicitly by Garriga and Tanaka that linearized gravity on the brane was indeed Einstein gravity for a single wall universe (see also ). The key feature of the Randall Sundrum (RS) (and indeed many other) models is that the induced metric on the wall (or walls) is flat Minkowski spacetime. In general, a domain wall spacetime is not static, but has a de-Sitter expansion along its spatial directions, however in the presence of a bulk cosmological constant, the effect of the wall energy and tension is neutralized, and a static solution is allowed. However, our universe is not Minkowski spacetime, nor is it a localized linearized perturbation thereof. Therefore, if brane universe models are to be relevant, cosmological brane solutions are the obvious next step. A general homogeneous and isotropic brane cosmology is simply a wall with energy, E, and tension, T, which are no longer fixed by the special domain wall relation $`\text{E}=\text{T}`$. For example, we might wish to set $$\text{E}=\text{E}_0+\rho ;\text{T}=\text{T}_0p$$ (1) and then take $`\rho `$ and $`p`$ to be our ‘cosmological’ energy density and pressure respectively. This is what is conventionally done, with $`\rho `$ and $`p`$ obeying an equation of state (rather than E and T), although it should be noted that when $`\rho `$ and E are of a similar order of magnitude, there is no reason to suppose that the equation of state will continue to hold with this choice of $`\rho `$ and $`p`$, which become rather arbitrary at that point. Much of the work on brane-cosmological models to date has taken a ‘brane-based’ approach, notably the work of Binetruy et. al. (see also ), in which the RS spacetime is generalized to allow for time-dependent cosmological expansion. However, Binetruy et. al. were concerned primarily with deriving the four-dimensional brane-cosmological equations, and for simplicity when finding explicit solutions, chose a Gaussian Normal gauge in which $`g_{nn}`$, the component of the metric normal to the wall, was set to unity. Therefore, while the brane cosmology was readily apparent in their approach, the bulk spacetime structure was less transparent, with coordinate singularities indicating the breakdown of the rather restrictive GN gauge. A more general brane-based approach is epitomized by the work of Shiromizu et. al. (see also Maartens, , for an excellent brane-based description of bulk effects) in which the Gauss-Codazzi formalism is used to obtain a pure-brane ‘Einstein’ equation, in which the bulk geometry is encoded in a single tensor $`_{\mu \nu }`$. While such a brane-based approach has the attraction of being a simple generalization of the standard four-dimensional FRW equations, the effect of the bulk on the brane is somewhat less transparent, with the tensor $`_{\mu \nu }`$ hiding a multitude of sins. An alternate approach for deriving cosmological solutions was taken by Ida , who instead considered a ‘bulk-based’ point of view, in which the most general static AdS solution with the appropriate symmetries: $$ds^2=h(r)dt^2h^1(r)dr^2r^2\left[\frac{d\chi ^2}{1\kappa \chi ^2}+\chi ^2d\mathrm{\Omega }_{II}^2\right]$$ (2) (where $`h(r)=\kappa \frac{\mu }{r^2}+k^2r^2`$) is taken, and the brane becomes an arbitrary boundary of two versions of this spacetime, with possibly differing masses on each side. The Israel-Gauss-Codazzi equations then give the energy and tension of the boundary as a function of its trajectory. The beauty of this approach is that it is quite general and does not rely on any $`Z_2`$ symmetry around the wall itself; it is also straightforward to derive brane-cosmological evolution equations – the cosmological solution then becomes a wall moving in AdS spacetime; the only disadvantage is that it is perhaps a little more abstract than the brane-based approach. A coordinate transformation relating the solutions of to Ida’s solutions was found by by Mukohyama et. al. . Ida’s work in fact generalizes the work of Kraus (see also for related string theoretic work, and , for an AdS/CFT perspective on the issue), who derived the most general $`Z_2`$ symmetric domain wall solutions by taking slicings of (2), and so can apply to models with different cosmological constants such as the Lykken-Randall (LR) model, , for example. What we aim to do in this paper is to bridge the gap between the ‘brane-based’ approach, where the brane represents a fixed boundary in some time-evolving spacetime, and the ‘bulk-based’ approach, where the wall (or spacetime boundary) follows some timelike trajectory in a static bulk spacetime. We start off by considering the most general brane-based formalism, following the approach of Ipser and Sikivie , finding the general wall solutions for a fixed wall of constant spatial curvature embedded in a bulk of constant curvature. We then demonstrate how a fixed wall embedded in a non-static spacetime is strictly equivalent to a moving wall embedded in a static spacetime thereby establishing in full generality the equivalence between the two approaches. This is true for a wall of completely arbitrary equation of state separating spacetimes which may even have a different cosmological constant. Having shown this equivalence we then proceed to find and study the most general wall trajectories, i.e. cosmological evolution equations, in a static AdS or even flat background. The reason for studying such a general set-up is that some more recent generalizations of the RS scenario have involved not only non $`Z_2`$ symmetric walls, such as the ‘Millenium Model’ of Kogan et. al. consisting of two positive tension walls at orbifold fixed points with an additional negative tension wall freely moving inbetween, but also patching together of spacetimes with different cosmological constant, such as the GRS model , in which the central $`Z_2`$ symmetric wall is flanked by two negative tension branes with flat space in the exterior. The curious feature of these apparently contrived models is that gravity not only changes nature at short but also at ultra-large scales, becoming weaker for the Millenium model, and five-dimensional for the GRS model. (In fact gravity for the GRS model has several peculiarities, , which may be ameliorated by a hybrid double-wall variant of the Millenium model, , in which the negative tension wall of the Millenium model is replaced by two negative tension walls with a slice of flat space inbetween.) While we might expect that cosmology in such models is potentially delicate, the beauty of the geometrical approach is that the search for a cosmology becomes a local question of finding a suitable trajectory for the spacetime boundaries, thus the presence of extra walls is irrelevant - unless these walls collide. We therefore do not consider issues such as radion stabilization , or the like, simply examining the possible trajectories (and hence cosmological solutions) of the brane universes. In the next section we set up our formalism and then find the general solution to the Einstein equations for a constant spatial curvature wall in a constant curvature spacetime, generalizing Taub’s solutions to allow for the presence of a cosmological constant in the bulk. We then show how from this brane-based approach we can cross over to the bulk-based approach for a wall of arbitrary trajectory evolving in a static spacetime. In section III we establish the most general cosmological evolution equations for the different classes of spacetime solutions. In the following section we analyze some specific cases of interest in cosmology finding the wall trajectories analytically. We make some concluding remarks in Section V. ## II General wall spacetimes In this section we derive the general spacetime of a brane universe, and the equations of motion it must obey. As per usual, we shall be modelling our four-dimensional Universe, $`𝒰`$, to be an infinitesimally thin wall type defect of constant spatial curvature embedded in a five-dimensional spacetime. In other words, we are looking for a spacetime with planar (or spherical/hyperboloidal) symmetry in three of its spatial directions, which has one (or more) hypersurfaces on which the spacetime curvature has a distributional singularity corresponding to a $`\delta `$-function source for the wall energy-momentum. The most general metric admitting this symmetry can be written in the form $$ds^2=e^{2\nu (\text{t,z})}(B(\text{t,z}))^{2/3}(d\text{t}^2d\text{z}^2)B^{2/3}\left[\frac{d\chi ^2}{1\kappa \chi ^2}+\chi ^2d\mathrm{\Omega }_{II}^2\right]$$ (3) where $`\kappa =0,\pm 1`$ represents the spatial curvature of the 3-spatial sections, and $`\nu ,B`$ will satisfy the bulk Einstein equations (with or without a cosmological constant), as well as appropriate jump conditions at the wall which we will choose to set at $`\text{z}=0`$. We start by summarizing the jump conditions for later use. As is conventional, we denote the normal to $`𝒰`$ by $`n_a`$, in terms of which the first fundamental form, $`h_{ab}`$, of the wall is given by $$h_{ab}=g_{ab}+n_an_b$$ (4) which is simply the projection of the bulk metric on $`𝒰`$. The second fundamental form, or extrinsic curvature of $`𝒰`$ is $$K_{ab}=h_{(a}^ch_{b)}^d_cn_d$$ (5) The intrinsic and extrinsic geometry of $`𝒰`$ are related to the bulk curvature via the Gauss Codazzi equations $`R_{acbd}n^cn^d=K_{ac}K_b^c_nK_{ab}`$ $`=`$ $`R_{ab}^{(4)}R_{cd}h_a^ch_b^dK_{ac}K_b^c+KK_{ab}`$ (7) $`h^{bc}_bK_{ca}h_a^b_bK`$ $`=`$ $`R_{bc}^{(4)}h_a^bn^c`$ (8) $`R^{(4)}K_{ab}K^{ab}+K^2`$ $`=`$ $`2G_{ab}n^an^b`$ (9) where $`R_{ab}^{(4)}`$ is the intrinsic Ricci tensor obtained from $`h_{ab}`$, and $`_n`$ is the Lie derivative with respect to $`n_a`$. Since the wall’s energy-momentum tensor $`T_{ab}^w`$ is a distributional source, we can now integrate (7) across the wall, and defining $$S_{ab}=T_{ab}^w𝑑l$$ (10) we immediately arrive at the Israel junction conditions $$\mathrm{\Delta }K_{ab}=[S_{ab}\frac{1}{3}Sh_{ab}],$$ (11) where $`\mathrm{\Delta }K_{ab}=K_{ab}^+K_{ab}^{}`$ is the jump in the extrinsic curvature, and we have set $`8\pi G_5=1`$. Furthermore taking the sum and difference of (9) and (8) respectively and using (11) we have, $`R^{(4)}\overline{K}_{ab}\overline{K}^{ab}+\overline{K}^2`$ $`=`$ $`{\displaystyle \frac{1}{4}}(S_{ab}S^{ab}{\displaystyle \frac{1}{3}}S^2)+2\overline{G}_{nn}`$ (13) $`\overline{K}_{ab}S^{ab}`$ $`=`$ $`2\mathrm{\Delta }G_{nn}`$ (14) $`_b^{(4)}S^{bc}`$ $`=`$ $`0`$ (15) $`_b^{(4)}\overline{K}_a^b_a^{(4)}\overline{K}`$ $`=`$ $`0`$ (16) where $`\overline{Q}=(Q^++Q^{})/2`$ stands for the mean of a quantity across the wall. We should stress that (II) are in fact integrability conditions and hence will result from the Einstein and junction conditions (11). Since we are looking for cosmological solutions, we shall assume that our brane Universe is made of homogeneous and isotropic matter, hence $$S_{ab}=\text{E}u^au^b+\text{T}(h^{ab}u^au^b)$$ (17) where E is the surface energy density and T the tension of the brane $`𝒰`$. The timelike vector $`u^a`$ is the 5-velocity of an observer comoving with the brane Universe. Obviously if $`\text{E}=\text{T}`$ we have a domain wall, and if $`\text{T}=0`$, a dust wall. We begin by finding the most general bulk solution before examining these boundary junction conditions. The bulk Einstein equations in this case are simply, $$R_{ab}=\frac{2}{3}\mathrm{\Lambda }g_{ab}$$ (18) where $`\mathrm{\Lambda }`$ is the cosmological constant. For the metric (3) these give the following system of partial differential equations, $`B_{,\text{tt}}B_{,\text{zz}}`$ $`=`$ $`\left(2\mathrm{\Lambda }B^{1/3}6\kappa B^{1/3}\right)e^{2\nu }`$ (20) $`\nu _{,\text{tt}}\nu _{,\text{zz}}`$ $`=`$ $`\left({\displaystyle \frac{\mathrm{\Lambda }}{3}}B^{2/3}+\kappa B^{4/3}\right)e^{2\nu }`$ (21) $`\nu _{,\text{z}}B_{,\text{t}}+\nu _{,\text{t}}B_{,\text{z}}`$ $`=`$ $`B_{,\text{tz}}`$ (22) $`2\nu _{,\text{z}}B_{,\text{z}}+2\nu _{,\text{t}}B_{,\text{t}}`$ $`=`$ $`B_{,\text{tt}}+B_{,\text{zz}}`$ (23) The particular way of writing the metric (3) shows the Liouville like character of the system of equations, where (20), (21) are in the presence of $`\mathrm{\Lambda }`$, non-homogeneous coupled wave equations and (22) and (23) will turn out to be integrability conditions. This is not too surprising since we are effectively studying a 1+1 gravity problem by considering the Kaluza-Klein reduction of the constant curvature spacelike dimensions. It proves easiest to rewrite (18) in light-cone coordinates, $$u=\frac{\text{t}\text{z}}{2},v=\frac{\text{t}+\text{z}}{2}$$ (24) in which $`B_{,uv}`$ $`=`$ $`\left(2\mathrm{\Lambda }B^{1/3}6\kappa B^{1/3}\right)e^{2\nu }`$ (26) $`\nu _{,uv}`$ $`=`$ $`\left({\displaystyle \frac{\mathrm{\Lambda }}{3}}B^{2/3}+\kappa B^{4/3}\right)e^{2\nu }`$ (27) $`B_{,u}\left[ln(B_{,u})\right]_{,u}`$ $`=`$ $`2\nu _{,u}B_{,u}`$ (28) $`B_{,v}\left[ln(B_{,v})\right]_{,v}`$ $`=`$ $`2\nu _{,v}B_{,v}`$ (29) First note that (28,29) can be directly integrated to give $$e^{2\nu }=V^{}(v)B_{,u}=U^{}(u)B_{,v}$$ (30) where $`V^{}(v)`$, $`U^{}(u)`$ are arbitrary nonzero functions of $`u,v`$ for generic $`\mathrm{\Lambda },\kappa `$. (If $`\mathrm{\Lambda }=\kappa =0`$, then it is possible for one of $`U^{}(u)`$ or $`V^{}(v)`$ to vanish – see below.) It is straightforward to see then that $`B`$ and $`\nu `$ have the form (where we remind the reader that primes denote ordinary differentiation with respect to the unique variable of the function) $$B=B(U(u)+V(v)),e^{2\nu }=B^{}U^{}V^{}$$ (31) which reduces the remaining PDE (26) to the ODE $`B^{\prime \prime }\left(2\mathrm{\Lambda }B^{1/3}6\kappa B^{1/3}\right)B^{}`$ $`=`$ $`0`$ (32) $`B^{}{\displaystyle \frac{3}{2}}\mathrm{\Lambda }B^{4/3}+9\kappa B^{2/3}`$ $`=`$ $`9\mu `$ (33) where $`\mu `$ is an integration constant. This last relation can be integrated to give $$\frac{1}{\sqrt{9\kappa ^23\mu \mathrm{\Lambda }}}\left[\|r_{}\|\mathrm{tan}^1\frac{r}{\|r_{}\|}r_+\mathrm{coth}^1\frac{r}{r_+}\right]=\left(U+V\right)$$ (34) where $`r=B^{1/3}`$, $`c`$ is a constant of integration, and $$r_\pm ^2=\frac{3\kappa }{\mathrm{\Lambda }}\pm \sqrt{\frac{9\kappa ^2}{\mathrm{\Lambda }^2}\frac{6\mu }{\mathrm{\Lambda }}}.$$ (35) Hence the general solution to (18) is given by, $$ds^2=B^{}U^{}V^{}B^{2/3}(d\text{t}^2d\text{z}^2)B^{2/3}dx_{III}^2$$ (36) where $`B=B(U+V)`$ satisfies (32) or (34), $`U(u)`$ and $`V(v)`$ are arbitrary functions and $`dx_{III}^2`$ stands for the three dimensional constant curvature metric written out explicitly in (3). When $`\mathrm{\Lambda }=\kappa =0`$, (32) gives (wlog) $`B=U+V`$, and there are now in fact two allowed classes of solution, in the terminology of Taub ; the above, (36), being a class II solution. The class I solutions are distinguished by having $`U^{}=0`$ or $`V^{}=0`$, but not both, in which case the general metrics are $`ds^2={\displaystyle \frac{U^{}H(v)}{U^{2/3}}}(d\text{t}^2d\text{z}^2)U^{2/3}d𝐱^2`$ (38) $`ds^2={\displaystyle \frac{V^{}K(u)}{V^{2/3}}}(d\text{t}^2d\text{z}^2)V^{2/3}d𝐱^2`$ (39) for $`U^{}0`$ and $`V^{}0`$ respectively. In fact these class I solutions are flat, as we will see in the next section. In summary Einstein’s equations (18) admit two distinct classes of solutions I, (38), and II, (36), which depend on two arbitrary functions. Having derived the general bulk solutions, now let us examine the constraints, or boundary conditions at $`\text{z}=0`$ imposed by the wall. For what follows we shall assume $`Z_2`$ symmetry in the z-direction to make notation easier and results more transparent. We shall however be dropping this assumption when actually looking for general cosmological solutions in the next section. Reflection symmetry about $`\text{z}=0`$ permits us to consider only $`\text{z}>0`$ in (36) since the metric is an even function with respect to z. Extrinsic curvature components are thus odd functions in z and hence, $$\mathrm{\Delta }K_{ab}=2K_{ab}^+,\overline{K_{ab}}=0$$ In the coordinate system adapted to the wall the intrinsic and extrinsic geometric quantities are greatly simplified, and take the form E $`=`$ $`2e^\nu B^{2/3}_zB=e^\nu {\displaystyle \frac{B^{}(V^{}U^{})}{B^{2/3}}}`$ (41) $`({\displaystyle \frac{2}{3}}\text{E}\text{T})`$ $`=`$ $`2_z\left(B^{1/3}e^\nu \right)=B^{1/3}e^\nu \left[{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{V^{\prime \prime }}{V^{}}}{\displaystyle \frac{U^{\prime \prime }}{U^{}}}\right)+{\displaystyle \frac{1}{2B}}(V^{}U^{})(\mathrm{\Lambda }B^{4/3}6\mu )\right]`$ (42) Hence (II) tell us that given two arbitrary functions $`U`$ and $`V`$ we can determine E and T. We now show that only one of these functions is physical. As a general rule any coordinate transformation in the bulk will shift our boundary making our wall evolving in time. However we can make a class of coordinate transformations in the bulk which leave our wall boundary conditions fixed, namely $$uf(u),vf(v)$$ (43) i.e. the rescaling of the light-cone coordinates by some arbitrary function. This boundary invariance is a result of the two-dimensional conformal symmetry on the t-z plane. Now we choose a particular coordinate transformation, $`fV`$, to fix this non-physical gauge freedom, thereby effectively setting $`Vv`$. Hence there is only one remaining arbitrary function or physical degree of freedom, $`U(u)`$, which completely determines E and T for the wall. There are now two equivalent ways, as it turns out, to look for cosmological solutions. We either, in the spirit of , consider a fixed boundary in the general spacetime (36), using the junction conditions (II), or we make coordinate transformations simplifying the bulk and study the moving wall trajectories in that case. In fact, it turns out that the general class II solution (36) is a static spacetime. Setting $$r=B^{1/3};t=3(vU)$$ (44) (36) becomes, $$ds^2=h(r)dt^2\frac{dr^2}{h(r)}r^2dx_{III}^2$$ (45) where $`h(r)`$ has the familiar form $$h(r)=\frac{B^{}B^{2/3}}{9}=\kappa \frac{\mathrm{\Lambda }}{6}r^2\frac{\mu }{r^2}$$ (46) recognised as the generalised “Schwarzschild-AdS” solutions. Note however that in making our background static our wall is no longer fixed, since the trajectory $`\text{z}=0`$, or $`u=v`$, gives a nontrivial relationship between $`r`$ and $`t`$: $`r=R(t)`$ from (44). The reader should note here, that time dependence of the metric (36) in the bulk, portrayed by arbitrary function $`U(u)`$, and static boundary (II) has been transformed into time dependence for the boundary portrayed by the arbitrary trajectory $`r=R(t)`$, in a static bulk spacetime (45). Hence all essential information of the bulk metric is transferred onto the boundary conditions and vice-versa. As a final remark, we note that the generalization to domain walls in arbitrary dimension is straightforward, simply altering the $`\mu /r^2`$ to $`\mu /r^{(n3)}`$ in (46) (see for a different approach). ## III Cosmological evolution equations In this section we derive the equations of motion for a general (i.e. not necessarily symmetric) wall bounding two regions of spacetime with possibly differing cosmological constants. The reason for setting up the problem in its full generality is that we will be able to use the same set of equations to deal not only with non $`Z_2`$ symmetric walls, such as in the Millenium model, but also with walls which separate negative and zero cosmological constant spacetimes. This latter case will of course also require consideration of both class I and class II solutions for the $`\mathrm{\Lambda }=\kappa =0`$ spacetime. We begin with a wall bounding two class II spacetimes. The starting point is the general solution (45). To get a wall solution, we simply take a general boundary $$X^a=(t(\tau ),R(\tau ),\chi ,\theta ,\varphi )$$ (47) and compute the extrinsic curvatures on each side. Note that $`\tau `$ is the proper time with respect to an observer comoving with the wall, i.e. $$\dot{t}^2h(r)\dot{R}^2h^1(r)=1$$ (48) and is therefore well defined and intrinsic to the wall, whereas $`t`$ is the coordinate time in the bulk, and need not agree with the coordinate time on the other side of the wall. The only other subtlety arises in the choice of sign for the outward normal $$n_a=\pm (\dot{R}(\tau ),\dot{t}(\tau ),\mathrm{𝟎})$$ (49) This corresponds to the fact that we have a choice of which part of the spacetime we wish to keep (i.e. $`r<R`$ or $`r>R`$). This depends on the physical solution required, i.e. whether we have a positive or negative energy wall, and what sort of asymptotic solution we require. In general, positive energy walls match across two interior spacetimes and correspond to the lower ($``$) choice of sign in the outward normal and vice versa. It is also possible to have a wall matching an interior and exterior spacetime (for example the spherical collapsing wall, , used by Gogberashvili in an early version of the RS scenario), however, for clarity in what follows unless explicitly stated, the formulae will be valid for +($``$) energy walls matching two interior (exterior) patches of the Schwarzschild-ADS spacetime. Computing the extrinsic curvature components on each side of the wall gives E $`=`$ $`{\displaystyle \frac{3}{R}}\left(\dot{t}_+h_++\dot{t}_{}h_{}\right)`$ (51) $`{\displaystyle \frac{2}{3}}\text{E}\text{T}`$ $`=`$ $`\pm \left({\displaystyle \frac{\ddot{R}+\frac{1}{2}h_+^{}}{\dot{t}_+h_+}}+{\displaystyle \frac{\ddot{R}+\frac{1}{2}h_{}^{}}{\dot{t}_{}h_{}}}\right)`$ (52) where $`t_\pm `$ represent the coordinate times on each side of the wall, $`\dot{t}_\pm h_\pm `$ being given by (48) in each case. Note that we can also express (52) in a different fashion as an energy conservation equation, $$\frac{D\text{E}}{d\tau }+3\frac{\dot{R}}{R}(\text{E}\text{T})=0$$ (53) These represent the general equations of motion for an arbitrary wall separating two regions of spacetime with possibly different cosmological constants. Note that at this point E and T are unconstrained physically, and are simply given by the above expressions in terms of $`R(\tau )`$, which, other than representing a non-spacelike hypersurface, is completely arbitrary. This is worth stressing – that we can take any such $`R(\tau )`$, and we will get a wall solution with energy and tension determined by the embedding. Only when we impose physical constraints on our energy and tension do we get constraints on our possible hypersurface $`R(\tau )`$. Before exploring such constraints however, we simply note that (III) can also be expressed in terms of brane cosmological evolution equations (cf. Binetruy et. al. ) which are $`\left({\displaystyle \frac{\dot{R}}{R}}\right)^2={\displaystyle \frac{\text{E}^2}{36}}+{\displaystyle \frac{9(\mathrm{\Delta }k^2)^2}{4\text{E}^2}}\overline{k^2}{\displaystyle \frac{\kappa }{R^2}}+{\displaystyle \frac{\overline{\mu }}{R^4}}{\displaystyle \frac{9\mathrm{\Delta }k^2\mathrm{\Delta }\mu }{2R^4\text{E}^2}}+{\displaystyle \frac{9(\mathrm{\Delta }\mu )^2}{4\text{E}^2R^8}}`$ (55) $`\ddot{R}=(2\text{E}3\text{T})\left({\displaystyle \frac{R\text{E}}{36}}+{\displaystyle \frac{9\mathrm{\Delta }k^2\mathrm{\Delta }\mu }{2\text{E}^3R^3}}\right)+R\left({\displaystyle \frac{9(\mathrm{\Delta }k^2)^2(4\text{E}3\text{T})}{4\text{E}^3}}\overline{k^2}\right){\displaystyle \frac{\overline{\mu }}{R^3}}{\displaystyle \frac{27\text{T}(\mathrm{\Delta }\mu )^2}{4\text{E}^3R^7}}`$ (56) where we have now put $`6k_\pm ^2=\mathrm{\Lambda }_\pm `$. These are the completely general brane cosmological evolution equations for a brane separating two ‘ADS’ regions, and were first derived by Stoica et. al. . Since we are particularly interested in the cosmological solutions for branes in models with regions of vacuum spacetime, we must also consider the case of a planar wall bordering on a class I spacetime which we may write as $$ds^2=H(v)drdvr^2d𝐱^2$$ (57) This spacetime is actually flat , as can be seen by the coordinate transformation $`𝐱^{}`$ $`=`$ $`r𝐱`$ (59) $`t^{}z^{}`$ $`=`$ $`r`$ (60) $`t^{}+z^{}`$ $`=`$ $`{\displaystyle H𝑑v}+r𝐱^2`$ (61) which gives the Minkowski metric in the starred coordinate system. Clearly to match with an ADS spacetime we need to set $`r=R(\tau )`$, and compute the relevant extrinsic curvature quantities. For a wall bounding a class II and class I spacetime, we find that the junction conditions are now E $`=`$ $`{\displaystyle \frac{3}{R}}\left(\dot{R}+\dot{t}_{}h_{}\right)`$ (63) $`{\displaystyle \frac{2}{3}}\text{E}\text{T}`$ $`=`$ $`\pm \left({\displaystyle \frac{\ddot{R}}{\dot{R}}}+{\displaystyle \frac{\ddot{R}+\frac{1}{2}h_{}^{}}{\dot{t}_{}h_{}}}\right)`$ (64) Therefore we now find a somewhat different set of cosmological equations for the evolution of the wall: $`{\displaystyle \frac{\dot{R}}{R}}`$ $`=`$ $`\pm \left({\displaystyle \frac{3k^2}{2\text{E}}}{\displaystyle \frac{\text{E}}{6}}{\displaystyle \frac{3\mu _{}}{2\text{E}R^4}}\right)`$ (66) $`\ddot{R}`$ $`=`$ $`{\displaystyle \frac{k^2R}{2\text{E}^2}}(9k^2\text{E}^2)+{\displaystyle \frac{(2\text{E}3\text{T})R}{36\text{E}^3}}(81k^4\text{E}^4)+{\displaystyle \frac{\mu _{}}{2\text{E}^3R^3}}\left[9k^2(3\text{T}2\text{E})\text{E}^3\right]{\displaystyle \frac{27\text{T}\mu _{}^2}{4\text{E}^3R^7}}`$ (67) where we see that the ‘Friedman’ equation is now linear in $`\dot{R}`$. Finally, if we are matching two class I spacetimes, the junction conditions give the particularly simple relations E $`=`$ $`{\displaystyle \frac{6\dot{R}}{R}}`$ (69) T $`=`$ $`\left({\displaystyle \frac{4\dot{R}}{R}}+{\displaystyle \frac{2\ddot{R}}{\dot{R}}}\right)`$ (70) In the next section we will compute cosmological trajectories for general walls, however, before doing so it is useful to verify that this approach is indeed valid by cross-checking it against a few simple known solutions. For example, a planar symmetric domain wall in pure Einstein gravity is known to have a particularly simple interpretation in terms of the matched interiors of two hyperboloids in Minkowski spacetime . In the context of these equations, a $`\mathrm{\Lambda }=\kappa =0`$ spacetime can match class I or class II spacetimes. For the matching of two class I spacetimes we must use (III) with $`\text{E}=\text{T}`$. Clearly this has the solution $`R=e^{\text{E}\tau /6}`$, giving a de-Sitter like induced metric in agreement with the standard domain wall metric in wall-based coordinates (see e.g. ). Inverting this relation to find $`R(v)`$ in terms of the bulk coordinates gives $`R=18/\text{E}^2v`$. Finally, inverting (57) gives the particularly simple bulk equation of motion for the domain wall trajectory: $`t^2z^2𝐱^2=36/\text{E}^2`$, i.e. a hyperboloid in Minkowski spacetime. The ‘planar’ domain wall is therefore a boundary between the interior of two Minkowski hyperboloids as required. ## IV Cosmological wall solutions In this section we apply the general equations worked out in the previous section to find a variety of cosmological solutions. Let us begin by considering the particularly simple equation of state of a domain wall, $`\text{E}=\text{T}`$, and find the most general trajectories. First of all, note that (53) implies that E is a constant, therefore the brane-Friedmann equation, (55), can be written as $$\frac{R^4}{4}\left(\frac{dR^2}{d\tau }\right)^2=aR^8\kappa R^6+bR^4+c$$ (71) where $`a`$ $`=`$ $`{\displaystyle \frac{\text{E}^2}{36}}+{\displaystyle \frac{9(\mathrm{\Delta }k^2)^2}{4\text{E}^2}}\overline{k^2}`$ (73) $`b`$ $`=`$ $`\overline{\mu }{\displaystyle \frac{9\mathrm{\Delta }k^2\mathrm{\Delta }\mu }{2\text{E}^2}}`$ (74) $`c`$ $`=`$ $`{\displaystyle \frac{9(\mathrm{\Delta }\mu )^2}{4\text{E}^2}}`$ (75) which can be integrated in general, giving $`R^2`$ in terms of elliptic functions, although the expression is not particularly illuminating. If one of $`\kappa ,c`$ is zero, the solutions are very simple to write down, but for general $`a,b,c,\kappa `$, the qualitative behaviour of the solutions can be deduced from (71). The constant ‘$`a`$’ can be seen to measure the departure from criticality of the domain wall – i.e. that value which allows a static planar domain wall solution in the absence of an ADM mass term. We will therefore denote $`a=0`$ walls as critical, and $`a>(<)\mathrm{\hspace{0.17em}0}`$ walls as super-(sub-)critical. Clearly, if there is a nonzero ADM mass in the bulk there is always a solution which expands outward from an initial singularity, $`R=0`$, which is where the wall trajectory touches the central singularity of the black hole. Whether or not this solution expands indefinitely, or there is a final singularity depends on the roots of (71). For $`a0`$, there are zero or two (possibly repeated) roots for $`R^2>0`$. If $`b>\frac{9}{32a}`$ for $`\kappa =1`$, or $`b>0`$ otherwise, then there are no positive roots and the cosmology expands indefinitely. Otherwise there can be two separate roots, in which case there is a cosmology with initial and final singularities, and a nonsingular cosmology which contracts in to a minimum radius and re-expands. If there is a repeated root, $`R_c^2`$, then the cosmology asymptotes $`R_c`$ exponentially at late times either as a contracting branch, or an expanding branch with an initial singularity. If $`a<0`$ there is always one root and our universe starts and ends its life on the black hole singularity. Finally, there is always a static solution at any positive root of the quartic. For planar domain walls, i.e. $`\kappa =0`$, we have the exact solution $$R^4=\{\begin{array}{cc}\frac{1}{4a}\left(e^{4\sqrt{a}(\tau \tau _0)}2b+(b^24ac)e^{4\sqrt{a}(\tau \tau _0)}\right)\hfill & a>0\hfill \\ \frac{\sqrt{b^24ac}}{2\|a\|}\mathrm{sin}\left\{4\sqrt{\|a\|}(\tau \tau _0)\right\}\frac{b}{2\|a\|}\hfill & a<0\hfill \\ 4b(\tau \tau _0)^2\frac{c}{b}\hfill & a=0\hfill \end{array}$$ (76) which clearly demonstrates the expansion outwards from an initial singularity, and for $`a>0`$ shows the late time inflationary nature of the cosmology. If $`c=0`$, i.e. the mass terms on either side of the wall are the same, the exact solutions are $$R^2=\{\begin{array}{cc}\frac{1}{4a}\left(e^{2\sqrt{a}(\tau \tau _0)}+2\kappa +(\kappa ^24ab)e^{2\sqrt{a}(\tau \tau _0)}\right)\hfill & a>0\hfill \\ \frac{\sqrt{\kappa ^24ab}}{2a}\mathrm{sin}\left\{2\sqrt{\|a\|}(\tau \tau _0)\right\}+\frac{\kappa }{2a}\hfill & a<0\hfill \\ \kappa (\tau \tau _0)^2+\frac{b}{\kappa }\hfill & a=0\hfill \end{array}$$ (77) where we have taken $`\kappa 0`$: if $`\kappa =0`$ use (76). A nice analysis of supercritical domain wall solutions and their bulk embeddings was given by Khoury et. al. in . For the case of a general wall we need some sort of equation of state $`\text{T}(\text{E})`$, for which we will follow convention in setting $$\text{E}=\text{E}_0+\rho ;\text{T}=\text{E}_0p$$ (78) where $`p=(\gamma 1)\rho `$ would be a typical equation of state, although we re-emphasize that for $`\rho \text{E}`$ there is no reason to suppose that such an equation of state will be accurate, or even possible. The conservation of energy momentum (53) implies that $`\rho =\rho _0R^{3\gamma }`$ and hence E is no longer a constant. In order to deal with this additional difficulty (in the non-$`Z_2`$-symmetric case) we keep only terms linear in $`\rho `$ which boils down to looking at wall trajectories for late times. For example in the case of radiation cosmology, $`\gamma =4/3`$, the Friedmann equation (55) reduces to, $$\left(\frac{\dot{R}}{R}\right)^2=\frac{\text{E}_0^2}{36}+\frac{9(\mathrm{\Delta }k^2)^2}{4\text{E}_0^2}\overline{k^2}\frac{\kappa }{R^2}+\frac{1}{R^4}\left(\overline{\mu }\frac{9\mathrm{\Delta }k^2\mathrm{\Delta }\mu }{2\text{E}_0^2}+\frac{\text{E}_0\rho _0}{18}\frac{9\rho _0(\mathrm{\Delta }k^2)^2}{2\text{E}_0^3}\right)+O(R^8)$$ (79) The above equation has exactly the same solution as (77) for $`\kappa 0`$ where however the constants $`a`$ and $`b`$ are given by, $$a=\frac{\text{E}_0^2}{36}+\frac{9(\mathrm{\Delta }k^2)^2}{4\text{E}_0^2}\overline{k^2};b=\overline{\mu }\frac{9\mathrm{\Delta }k^2\mathrm{\Delta }\mu }{2\text{E}_0^2}+\frac{\text{E}_0\rho _0}{18}\frac{9\rho _0(\mathrm{\Delta }k^2)^2}{2\text{E}_0^3};$$ (80) Another important class of solutions are those with $`Z_2`$ symmetry, for which the cosmological Friedman equation has the form of (71) with the differences ‘$`\mathrm{\Delta }`$’ of all quantities being zero. The Friedmann equation (55) becomes $$\left(\frac{\dot{R}}{R}\right)^2=\frac{\text{E}_0^2}{36}k^2\frac{\kappa }{R^2}+\frac{\text{E}_0\rho _0}{18R^{3\gamma }}+\frac{\mu }{R^4}+\frac{\rho _0^2}{36R^{6\gamma }}$$ (81) where the above is now an exact expression. For example, a radiation cosmology would have $`\gamma =4/3`$, and the equation of motion would be as (71), but with the constants $`a_r,b_r,c_r,`$ now defined as $$a_r=\frac{\text{E}_0^2}{36}k^2;b_r=\mu +\frac{\text{E}_0\rho _0}{18};c_r=\frac{\rho _0^2}{36}$$ (82) and so a generic planar solution would be given by (76) with these new values for the constants $`a,b`$, and $`c`$. We see now the interplay between the contribution of the cosmological matter energy density $`\rho _0`$, which contributes to both $`b`$ and $`c`$, and the “CFT” contribution, $`\mu `$, which appears only in $`b`$. This changes the way in which the initial singularity, $`R=0`$, is approached, for example, Gubser’s solution has $`a=0=\rho _0`$, and $`R=\mu ^{\frac{1}{4}}\sqrt{2(\tau \tau _0)}`$, as opposed to $`R=b^{\frac{1}{4}}[4(\tau \tau _0)^2c/b^2]^{\frac{1}{4}}`$. A description of $`Z_2`$ symmetric cosmological solutions and their embeddings was given by Mukohyama et. al. . Let us now turn to the less well-explored question of cosmological solutions in models with slices of vacuum and ADS spacetimes, such as the noncompact GRS model, or compact models such as those in . Without loss of generality, we will take $`k_+=0`$, $`k_{}=k`$, for which the general equations of motion are given by (III) or (III). For simplicity, since we are unlikely to be interested in cosmologies on negative tension walls, let us consider the case where the boundary has the equation of state of a domain wall, i.e. $`\text{E}=\text{T}`$, and let $`\epsilon =(\text{E}+3k)`$ be the departure from ‘criticality’ of this domain wall ($`\epsilon >0`$ being supercritical). We start by considering the matching of two class II spacetimes. Note that for $`\kappa =0,1`$, in order to have the correct spacetime signature, we require $`\mu _+<0`$, and we have a naked singularity in the full bulk solution. For $`\kappa =1`$, $`\mu _+>0`$ gives the five-dimensional Schwarzschild solution. Substituting the values of $`k_\pm `$ gives for the constants in (71) $$a=\frac{\epsilon ^2}{36\text{E}^2}\left(\epsilon +2\text{E}\right)^2;b=\mu _++\frac{\epsilon (\epsilon +2\text{E})\mathrm{\Delta }\mu }{2\text{E}^2};c=\frac{9(\mathrm{\Delta }\mu )^2}{4\text{E}^2}$$ (83) Note that unlike the generic domain wall case, $`a>0`$ for both sub and super-critical walls. For $`\kappa =1`$, $`\mu _+>0`$ and we again have solutions which expand out (or contract into) the singularity at $`r=0`$. Depending on the roots of the quartic, there are solutions which expand indefinitely, recontract, or indeed asymptote a constant radius. The critical domain wall always recollapses. For $`\kappa =1`$, $`\mu _+<0`$, and the qualitative behaviour is similar, however the critical domain wall now expands indefinitely if $`27c>4\mu _+`$. For the planar ($`\kappa =0`$) walls, the exact solutions are $$R^4=\{\begin{array}{cc}\frac{1}{4a}\left(e^{\pm 4\sqrt{a}(\tau \tau _0)}2b+(b^24ac)e^{4\sqrt{a}(\tau \tau _0)}\right)\hfill & a>0\hfill \\ 4b(\tau \tau _0)^2\frac{c}{b}\hfill & a=0\hfill \end{array}$$ (84) For the critical domain wall, $`b=\mu _+<0`$, and we see that the domain wall universe has an initial and final singularity, for $`\tau _0\frac{\mathrm{\Delta }\mu }{4k\mu _+}<\tau <\tau _0+\frac{\mathrm{\Delta }\mu }{4k\mu _+}`$. Inverting (48) to find the trajectory in terms of coordinate time in the vacuum spacetime to the right of the wall shows that these singularities occur at finite coordinate time, and the wall actually expands out of, and re-collapses into, the central naked singularity in this spacetime. This corresponds to figure 1, and would appear to be a rather undesirable spacetime from the five-dimensional point of view. For the subcritical domain wall however, $`b<0`$ with $`b^24ac>0`$, hence the cosmology has no initial singularity, but merely follows a modestly corrected $`\mathrm{sinh}`$ trajectory. For the supercritical domain wall, there are two possibilities, depending on the relative magnitudes of $`\mu _\pm `$. If $`9k^2\|\mu _+\|\mu _{}(\epsilon ^2+6k\epsilon )`$, then the cosmology is completely nonsingular as for the subcritical case, however, if $`9k^2\|\mu _+\|<\mu _{}(\epsilon ^2+6k\epsilon )`$, then there is an initial (or final) singularity with the universe inflating away (deflating towards) it. Finally, for a planar asymptotically vacuum spacetime, the other possibility is matching to a class I spacetime for which we need (III). Setting $`\text{E}=\text{T}=\epsilon 3k<0`$ again we have $$\frac{\dot{R}}{R}=\frac{\epsilon (\epsilon +2\text{E})}{6\text{E}}\frac{3\mu _{}}{2\text{E}R^4}$$ (85) which has the general solution $$R^4=\{\begin{array}{cc}\mathrm{exp}\left\{\frac{2\epsilon (\epsilon +2\text{E})}{3\text{E}}(\tau \tau _0)\right\}+\frac{9\mu _{}}{\epsilon (\epsilon +2\text{E})}\hfill & \epsilon 0\hfill \\ \frac{2\mu _{}}{k}(\tau \tau _0)\hfill & \epsilon =0\hfill \end{array}$$ (86) A supercritical wall therefore has $`\dot{R}>0`$ and expands outward from an initial singularity (where it touches the Schwarzschild singularity on the ($``$) side of the wall). A subcritical wall however has $`\dot{R}_<^>0`$ for $`R_{}^{4}{}_{>}{}^{<}\frac{9\mu _{}}{\epsilon (\epsilon +2\text{E})}`$. This shows that there are two solutions, one of which has an inital singularity and expands outward to $`R_c^4=\frac{9\mu _{}}{\epsilon (\epsilon +2\text{E})}`$, and another which contracts to $`R_c`$. A critical domain wall has a power law expansion outward from an initial singularity. It is interesting to transform into the starred coordinate system to follow these wall trajectories in Minkowski spacetime: $`𝐱^{}`$ $`=`$ $`R(\tau )𝐱`$ (88) $`t^{}z^{}`$ $`=`$ $`R(\tau )`$ (89) $`t^{}+z^{}`$ $`=`$ $`{\displaystyle \frac{d\tau }{\dot{R}(\tau )}}+R(\tau )𝐱^2=F(R(\tau ))+R(\tau )𝐱^2`$ (90) For a critical domain wall $`F(R)=4k^2R^7/7\mu _{}^2`$, and for a noncritical domain wall bordering pure ADS, $`F(R)=1/4R`$. We see therefore that this latter case is the familiar hyperboloid in Minkowski spacetime, although it is now the interior which is excised for this negative tension wall. The critical domain wall satisfies $$𝐱^2=t^2z^2\frac{4k^2}{7\mu _{}^2}(t^{}z^{})^8$$ (91) which is a deformed hyperboloid. At fixed time, the 3-sphere is now squashed into an egg-shape in the $`\|𝐱^{}\|,z^{}`$ directions. Again, it is the interior of this squashed S<sup>3</sup> which is excised. On the ADS side, the wall expands outward from the initial black hole singularity and at late times expands as $`Rt_{}^{1/3}`$ (c.f. $`Rt`$ for Gubser’s critical $`Z_2`$ symmetric solution). The appropriate matching is shown in figure 2. Finally it is worth stressing that these trajectories are independent of whether or not there is an additional positive tension wall present – the existence of the wall trajectory is a local question. We can therefore add in a positive domain wall in the usual way by simply adding another boundary to the Schwarzschild ADS spacetime as shown by the dotted line in figure 2. ## V Discussion Starting from the most general five-dimensional metric with homogeneous and isotropic spatial 3-sections, we have shown that the most general cosmological brane solution matches two Schwarzschild-ADS spacetimes, or a Schwarzschild-ADS spacetime with a class I or II vacuum spacetime. It is worth emphasizing that the three dimensional homogeneity and isotropy ensures that the bulk metric solution is invariant under two dimensional conformal symmetries. By inputting a boundary or wall we in principle break these conformal symmetries, however, it turns out that the junction conditions are such that only half are broken. This in turn ensures that the only remaining degree of freedom is the wall’s trajectory itself. Hence all the essential information of the wall’s dynamics is contained on the wall trajectory itself (the boundary) the bulk being a specific static spacetime playing a ‘background’ role. We then derived the generalised FRW equations for a wall with general energy and tension, then applied these to a range of cases, focusing on the previously unexplored case of a wall bordering a vacuum region on one side. The interesting feature of these latter solutions is that if two class II spacetimes are matched, the bulk can contain timelike naked singularities on the vacuum side of the wall. For a more satisfactory nonsingular solution, one must match the planar wall to a class I spacetime. Here the wall trajectories are reminiscent of the vacuum domain wall of Vilenkin, Ipser and Sikivie, in that they are deformed hyperboloids, although it is the interior of the hyperboloid which is excised for this negative energy wall. If we now wish to construct a ‘cosmological’ GRS solution, i.e. one which has a $`Z_2`$ symmetric positive energy central wall with matter residing on it, and a negative energy outer wall, then we must combine one of the $`Z_2`$ solutions of (81) with a class II/II or class II/I solution for the outer wall. In each case, we must be careful to keep the ($``$) wall to the right of the ($`+`$) wall, $`R_{}<R_+`$. For example, if we take the planar domain wall, the central domain wall trajectory, $`R_+`$, is given by (76) with the constants $`a,b,c`$ in (82). Using the class II/I solutions of (84) for the ($``$) wall shows that we cannot match arbitrary walls, since for example a subcritical central wall would collide with a critical or supercritical outer wall bringing our universe to an abrupt and catastrophic end. Two generic critical walls are compatible however, with $$\{\begin{array}{cc}R_+^4=4\left(\mu +\frac{k\rho _0}{3}\right)(\tau _+\tau _{{}_{0}{}^{}+})^2\frac{\rho _0^2}{36\mu +12k\rho _0}\hfill & \\ R_{}^4=\frac{2\mu }{k}(\tau _{}\tau _{}_{0}{}^{})\hfill & \end{array}$$ (92) as is a critical central wall with a subcritical outer one. A supercritical central wall can exist with a critical or subcritical outer wall, or a supercritical outer wall with $`2\sqrt{\text{E}_+^236k^2}>(9k^2\text{E}_{}^2)/\|\text{E}_{}\|`$; and a subcritical central wall can exist with a subcritical outer wall provided $`\frac{\sqrt{b_r^24a_rc_r}b_r}{2\|a_r\|}>\frac{9\mu }{(9k^2\text{E}_{}^2)}`$. For example we can match the critical Gubser solution, $`R_+^4=4\mu \tau _+^2`$, with the critical class II/I solution $`R_{}^4=2\mu \tau _{}/k`$. However, a pure critical planar radiation cosmology, $`R_+^4=\rho _0(16k^2\tau _+^21)/12k`$ cannot be matched to any class I asymptotically flat spacetime across a ($``$) wall, as the only possible solution is an exponentially expanding or contracting one for a super/sub critical wall respectively. Therefore, while there are some constraints on the central wall cosmologies, these are not overwhelmingly restrictive. This leads us to the issue of the radion in these cosmological models. There has been a great deal of debate about the radion in the context of cosmological models (see e.g. ) and in the case of quasi-localised gravity (see e.g. ). We have not supposed any stabilization mechanism, , for the extra dimension, but simply looked for free cosmological solutions. Since the radion leads to anti-gravity in the original GRS model, it is interesting to see how it manifests itself here. The wavefunction of the radion is straightforward to find in linearized gravity, , and it behaves as a scalar field living on the brane worldvolume which satisfies a massless equation of motion, $`^2f=0`$. In the context of a homogeneous and isotropic cosmology, this has the simple solution $`f=f_0t`$, where $`t`$ is the coordinate time on the brane, and we have chosen to have zero displacement for $`t=0`$. Now consider a supercritical $`Z_2`$-symmetric positive energy planar domain wall in pure ADS spacetime; this follows a trajectory $$R=\mathrm{exp}\{\sqrt{\text{E}^236k^2}(\tau \tau _0)/6\}=\left[1\frac{k^2\sqrt{\text{E}^236k^2}}{\text{E}}(tt_0)\right]^1$$ (93) in terms of either the brane proper time, $`\tau `$, or the local bulk coordinate time, $`t`$. We see therefore that for small $`t`$ and $`\sqrt{\text{E}^236k^2}`$, $$R1+\frac{k\sqrt{3k}\sqrt{\text{E}6k}}{3}(tt_0)$$ (94) we can therefore identify $`f_0=k^{3/2}\sqrt{\text{E}6k}\sqrt{3}`$, i.e. $`\delta \text{E}=\text{E}6k=3f_0^2k^{3/2}`$, which has the correct dependence on $`f_0`$ for a standard energy momentum tensor for the radion. (Note we cannot have a subcritical planar domain wall.) This means that a negative energy wall requires $`\text{E}<6k`$, and nominally a negative radion energy on that wall. Obviously these solutions do not see any evidence of the fifth dimension “opening up” , or any instability due to the radion, however, this is because of the high degree of symmetry of the set-up. If we were to break isotropy or homogeneity, the integrable nature of the system would likely be destroyed, and the simplicity of the description of the wall as a trajectory in bulk Schwarzschild-ADS would disappear. Destruction of these symmetries would lead to a brane-bulk interaction (as evidenced by the extreme case of the as yet undiscovered black hole on the brane solution – see for discussions on this topic) through which the more familiar effects of the radion would be recovered. Of course it is tempting to enquire whether the so-called “missing mode” of free motion of a wall can be identified in this set-up. For a non-gravitating defect, the effective equation of motion is given by the relation $`\overline{K}=0`$, i.e. the wall (in this case) is a minimal hypersurface. An equivalent perturbative description would state that the displacement of the wall satisfies a massless wave equation on the wall. There have been claims in the literature that when gravity is included this ‘free motion’ disappears, as would appear to agree with the rejection of the free ‘radion’ solution $$\delta g_{\mu \nu }=\frac{(r^2r^2)}{4k}\overline{K}_{\mu \nu }$$ (95) in the usual Randall-Sundrum perturbation analysis. However, in this simple picture, each wall has its own ‘radion’ corresponding to motion through the ambient spacetime. Indeed, a perturbative analysis of the thick four-dimensional domain wall indicates no such disappearance of this motion with the coupling to gravity. What therefore can be going on? A clue perhaps lies in the effective equation of state for a freely moving wall. As Carter and Vilenkin , argued for the cosmic string, a freely moving defect will have an altered effective equation of state. The effective energy per unit area will increase, and the tension will decrease. Generalizing the Carter-Vilenkin formula for the domain wall gives the effective equation of state $$\text{E}\text{T}^3=\text{E}_0^4$$ (96) for a 3-brane domain wall. Clearly for E close to $`\text{E}_0`$, this gives an equation of state of a radiation cosmology on the background domain wall. Now we see how a freely moving wall might no longer be a small perturbation of a ‘straight’ one: the averaged equation of state of the freely moving wall causes the wall to follow a centre of mass trajectory associated with a radiation universe, which is not, at late times (when the radiation universe approximation is particularly good) a small perturbation of the unmoving wall. Of course, a freely moving wall will not be isotropic and homogeneous except at the very large scale, and it is also possible that there are brane-bulk interactions which further complicate the issue. Finally, the problem of brane cosmological perturbation theory is very important to understand if we wish to do real cosmology. The setting up of the formalism for perturbation theory on the brane is already underway ; by considering the cosmological domain wall in terms of its global spacetime structure, the role of the bulk, its interactions with matter on the brane, and the interpretation of some of the gauge invariant variables in terms of bulk physical quantities will be more clearly elucidated. ## Acknowledgements It is a great pleasure to thank Roberto Emparan, David Fairlie, Jihad Mourad, Simon Ross, Valery Rubakov, and Douglas Smith for helpful and enlightening discussions. C.C. was supported by PPARC, and R.G. by the Royal Society.
warning/0007/quant-ph0007124.html	ar5iv	text	# Generalization of Grover’s Algorithm to Multiobject Search in Quantum Computing, Part II: General Unitary Transformations ## 1 Introduction This paper is a continuation from \[References\] on quantum computing algorithms for multiobject search. L.K. Grover’s first papers \[References, References\] on “quantum search for a needle in a haystack” have stimulated broad interest in the theoretical development of quantum computing algorithms. Let an unsorted database consist of $`N`$ objects $`\{w_j1jN\}`$; each object $`w_j`$ is stored in a quantum computer (QC) memory as an eigenstate $`\|w_j`$, $`j=1,2,\mathrm{},N`$, with $`\{\|w_j1jN\}`$ forming an orthonormal basis of a Hilbert space $``$. Let $`\|w`$ be an element of $``$ which is the (single) object to be searched. Grover’s algorithm in \[References, References\] is to utilize a unitary operator $$UI_sI_w$$ (1.1) where $`I_w`$ $`𝑰2\|ww\|,(𝑰\text{ the identity operator on })`$ (1.2) $`I_s`$ $`𝑰2\|ss\|,\|s{\displaystyle \frac{1}{\sqrt{N}}}{\displaystyle \underset{i=1}{\overset{N}{}}}\|w_i,`$ (1.3) to perform the iterations $`U^m\|s`$, which will lead to the target state $`\|w`$ with probability close to 1 after approximately $`\frac{\pi }{4}\sqrt{N}`$ number of iterations. The algorithm is of optimal order. In a more recent paper \[References\], Grover showed that the state $`\|s`$ in (1.3) can be replaced by any quantum state $`\|\gamma `$ with nonvanishing amplitude for each object $`w_j`$ and, correspondingly, the Walsh-Hadamard operator previously used by him to construct the operator $`I_s`$ can be replaced by a sufficiently general nontrivial unitary operator. Grover’s new “search engine” in \[References\] is a unitary operator taking the form $$U=I_\gamma V^1I_wV:$$ (1.4) where $`V`$ is an arbitrary unitary operator. The object $`w`$ will be attained (with probability close to 1) by iterating $`U^m\|\gamma `$. This seems to give the algorithm/software designer large flexibility in conducting quantum computer search and code development. It increases the variety of quantum computational operations that can feasibly be performed by practical software. In particular, it opens the possibility of working with an initial state $`\|\gamma `$ (in place of $`\|s`$) that is other than a superposition of exactly $`N=2^n`$ ($`n=`$ number of qubits) alternatives. This suggests a new paradigm in which the whole dataset (not just the key) is encoded in the quantum apparatus. This new point of view may also overcome some of the practical difficulties noted by Zalka \[References\] in searching a physical database by Grover’s method. In the next section, we study the generalization of (1.4) to multiobject search. ## 2 Multiobject Search Algorithm Using a General Unitary Transformation Let $`\{\|w_i1iN\}`$ be the basis of orthonormal eigenstates representing an unsorted database $`w_i`$, $`1iN`$, as noted in §I. We inherit much of the notation in \[References\]: let $`f`$ be an oracle function such that $$f(w_i)=\{\begin{array}{cc}1,\hfill & 1i\mathrm{},\hfill \\ 0,\hfill & \mathrm{}+1iN,\hfill \end{array}$$ where $`w_i`$, $`i=1,2,\mathrm{},\mathrm{}`$, represent the multiobjects under search. We wish to find at least one $`w_i`$, for $`i=1,2,\mathrm{},\mathrm{}`$. Let $`\|\gamma `$ be any unit vector in $``$, and let $`L\text{span}\{\|w_i1i\mathrm{}\}`$. Define $$I_\gamma =𝑰2\|\gamma \gamma \|:,$$ and $$I_L\|w_j=(1)^{f(w_j)}\|w_j,j=1,2,\mathrm{},N,$$ and $`I_L`$ is then uniquely extended linearly to all $``$ with the representation $$I_L=𝑰2\underset{i=1}{\overset{\mathrm{}}{}}\|w_iw_i\|.$$ Both $`I_\gamma `$ and $`I_L`$ are unitary operators. Let $`V`$ be any unitary operator on $``$. Now, define $$U=I_\gamma V^1I_LV.$$ (2.1) Then $`U`$ is a unitary operator; it degenerates into Grover’s operator $`U`$ in (1.4) when $`\mathrm{}=1`$ and further into the old Grover’s operator $`U`$ in (1.1) if $`V𝑰`$. The unit vector $`\|\gamma `$ is arbitrary except that we require $`V\|\gamma L`$. (Obviously, any $`\|\gamma `$ such that $`w_i\|\gamma 0`$ for all $`i=1,2,\mathrm{},N`$, will work, including $`\|\gamma \|s`$ in (1.3).) If $`V\|\gamma L`$, then $$V\|\gamma =\underset{j=1}{\overset{\mathrm{}}{}}g_i\|w_i,g_i,\underset{j=1}{\overset{\mathrm{}}{}}\|g_i\|^2=1.$$ A measurement of the state $`V\|\gamma `$ will yield an eigenstate $`\|w_j`$, for some $`j:1j\mathrm{}`$, with probability $`\|g_j\|^2`$. Thus the search task would have been completed. Thus, let us consider the nontrivial case $`V\|\gamma L`$. This implies $`\|\gamma V^1(L)`$ and, hence, $$\stackrel{~}{L}\text{span}(\{\|\gamma \}V^1(L))$$ (2.2) is an $`(\mathrm{}+1)`$-dimensional subspace of $``$. It effects a reduction to a lower dimensional invariant subspace for the operator $`U`$, according to the following. ###### Lemma 2.1. Assume that $`\gamma \|\gamma =1`$ and $`V\|\gamma L`$. Then $`U(\stackrel{~}{L})=\stackrel{~}{L}`$. ###### Proof. For any $`j:1j\mathrm{}`$, denote $$\mu _{\gamma ,j}=w_j\|V\|\gamma .$$ * We have, for $`j:1j\mathrm{}`$, $`U(V^1\|w_j)`$ $`=I_\gamma V^1\left(I2{\displaystyle \underset{i=1}{\overset{\mathrm{}}{}}}\|w_iw_i\|\right)\|w_j`$ $`=I_\gamma V^1(\|w_j)`$ $`=I_\gamma V^1\|w_j`$ $`=(I2\|\gamma \gamma \|)V^1\|w_j`$ $`=V^1\|w_j2(\gamma \|V^1\|w_j)\|\gamma `$ $`=V^1\|w_j2\overline{\mu }_{\gamma ,j}\gamma \stackrel{~}{L};`$ (2.3) * $`U\|\gamma `$ $`=I_\gamma V^1\left(I2{\displaystyle \underset{i=1}{\overset{\mathrm{}}{}}}\|w_iw_i\|\right)(V\|\gamma )`$ $`=(I2\|\gamma \gamma \|)\left[\|\gamma 2{\displaystyle \underset{i=1}{\overset{\mathrm{}}{}}}(w_i\|V\|\gamma )V^1\|w_i\right]`$ $`=\|\gamma +2{\displaystyle \underset{i=1}{\overset{\mathrm{}}{}}}\mu _{\gamma ,i}V^1\|w_i4{\displaystyle \underset{i=1}{\overset{\mathrm{}}{}}}\mu _{\gamma ,i}\overline{\mu }_{\gamma ,i}\|\gamma `$ $`=\left(14{\displaystyle \underset{i=1}{\overset{\mathrm{}}{}}}\|\mu _{\gamma ,i}\|^2\right)\|\gamma +2{\displaystyle \underset{i=1}{\overset{\mathrm{}}{}}}\mu _{\gamma ,i}V^1\|w_i\stackrel{~}{L}.\mathit{}`$ (2.4) By Lemma 2.1, the Hilbert space $``$ admits an orthogonal direct sum decomposition $$=\stackrel{~}{L}\stackrel{~}{L}^{}$$ such that $`\stackrel{~}{L}^{}`$ is also an invariant subspace of $`U`$. In our subsequent iterations, the actions of $`U`$ will be restricted to $`\stackrel{~}{L}`$, as the following Lemma 2.2 has shown. Therefore we can ignore the complementary summand space $`\stackrel{~}{L}^{}`$. ###### Lemma 2.2. Under the same assumptions as Lemma 2.1, we have $`U^m\|\gamma \stackrel{~}{L}`$ for $`m^+\{0,1,2,\mathrm{}\}`$. ###### Proof. It follows obviously from by (2.2) and Lemma 2.1. ∎ Consider the action of $`U`$ on $`\stackrel{~}{L}`$. Even though $`\|\gamma ,V^1\|w_i`$, $`i=1,\mathrm{},\mathrm{}`$, form a basis of $`\stackrel{~}{L}`$, these vectors are not mutually orthogonal. We have $`U\left[\begin{array}{c}\|\gamma \\ V^1\|w_1\\ V^1\|w_2\\ \mathrm{}\\ V^1\|v_{\mathrm{}}\end{array}\right]`$ $`=\left[\begin{array}{ccccc}14\underset{i=1}{\overset{\mathrm{}}{}}\|\mu _{\gamma ,j}\|^2& 2\mu _{\gamma ,1}& 2\mu _{\gamma ,2}& \mathrm{}& 2\mu _{\gamma ,\mathrm{}}\\ 2\overline{\mu }_{\gamma ,1}& 1& 0& \mathrm{}& 0\\ 2\overline{\mu }_{\gamma ,2}& 0& 1& & 0\\ \mathrm{}& \mathrm{}& & \mathrm{}& \mathrm{}\\ 2\overline{\mu }_{\gamma ,\mathrm{}}& 0& 0& \mathrm{}& 1\end{array}\right]\left[\begin{array}{c}\|\gamma \\ V^1\|w_1\\ V^1\|w_2\\ \mathrm{}\\ V^1\|w_{\mathrm{}}\end{array}\right],`$ (2.5) $`\left[\begin{array}{c}\|\gamma \\ V^1\|w_1\\ V^1\|w_2\\ \mathrm{}\\ V^1\|w_{\mathrm{}}\end{array}\right],`$ according to (2.3) and (2.4). Therefore, with respect to the basis $`\{\|\gamma ,V^1\|w_ii=1,\mathrm{},\mathrm{}\}`$, the matrix representation of $`U`$ on $`\stackrel{~}{L}`$ is $`^T`$, the transpose of $``$. These two $`(\mathrm{}+1)\times (\mathrm{}+1)`$ matrices $``$ and $`^T`$ are nonunitary, however, because the basis $`\{\|\gamma ,V^1\|w_i\|`$, $`i=1,2,\mathrm{},\mathrm{}\}`$ is not orthogonal. This fact is relatively harmless here, as we can further effect a reduction of dimensionality by doing the following. Define a unit vector $$\|\mu =2\underset{j=1}{\overset{\mathrm{}}{}}\mu _{\gamma ,j}V^1\|w_j/a,a\left(4\underset{j=1}{\overset{\mathrm{}}{}}\|\mu _{\gamma ,j}\|^2\right)^{1/2}>0.$$ (2.6) ###### Theorem 2.3. Let $`𝒱\text{span}\{\|\gamma ,\|\mu \}`$. Then $`𝒱`$ is a two-dimensional invariant subspace of $`U`$. We have $$U\left[\begin{array}{c}\|\gamma \\ \|\mu \end{array}\right]=M\left[\begin{array}{c}\|\gamma \\ \|\mu \end{array}\right],M\left[\begin{array}{cc}1a^2& a\\ a& 1\end{array}\right].$$ (2.7) Consequently, with respect to the basis $`\{\|\gamma ,\|\mu \}`$ in $`𝒱`$, the matrix representation of $`U`$ is $`M^T`$. ###### Proof. Using (2.3), we have $`U\|\mu `$ $`=2{\displaystyle \underset{j=1}{\overset{\mathrm{}}{}}}\mu _{\gamma ,j}V^1\|w_j{\displaystyle \frac{1}{a}}2{\displaystyle \underset{j=1}{\overset{\mathrm{}}{}}}\|\mu _{\gamma ,j}\|^2{\displaystyle \frac{1}{a}}\|\gamma `$ $`=\|\mu a\|\gamma .`$ Again, from the definition of $`\|\mu `$ in (2.6), we see that (2.4) gives $$U\|\gamma =(1a^2)\|\gamma +a\|\mu .$$ Therefore (2.7) follows. ∎ Theorem 2.3 gives a dramatic reduction of dimensionality to 2, i.e., the dimension of the invariant subspace $`𝒱`$. Again, we note that the matrices $`M`$ and $`M^T`$ in (2.7) are not unitary. Any vector $`\|v𝒱`$ can be represented as $$\|v=c_1\|\gamma +c_2\|\mu ,$$ and so $`U\|v`$ $`=U(c_1\|\gamma +c_2\|\mu `$ $`=c_1[(1a^2)\|\gamma +a\|\mu ]+c_2[a\|\gamma +\|\mu ],`$ and thus $$U\|v=M^T\left[\begin{array}{c}c_1\\ c_2\end{array}\right]=\left[\begin{array}{cc}1a^2& a\\ a& 1\end{array}\right]\left[\begin{array}{c}c_1\\ c_2\end{array}\right],$$ (2.8) where the first component of the vector on the right hand side of (2.8) corresponds to the coefficient of $`\|\gamma `$ while the second component corresponds to the coefficient of $`\|\mu `$. Therefore $$U^m\|\gamma =\left[\begin{array}{cc}1a^2& a\\ a& 1\end{array}\right]^m\left[\begin{array}{c}1\\ 0\end{array}\right].$$ (2.9) The above can be viewed geometrically (\[References\]) as follows: $`M^T\left[\begin{array}{c}1\\ 0\end{array}\right]=\left[\begin{array}{c}1a^2\\ a\end{array}\right]`$, for $`a>0`$ very small, $`a\mathrm{sin}a`$, and therefore $`\left[\begin{array}{c}1a^2\\ a\end{array}\right]`$ is a vector obtained from the unit vector $`\left[\begin{array}{c}1\\ 0\end{array}\right]`$ by rotating it counterclockwise with angle $`a`$. It takes approximately $$m\frac{\pi /2}{a}=\frac{\pi }{2a}=\pi /4\left[\underset{j=1}{\overset{\mathrm{}}{}}\|\mu _{\gamma ,j}\|^2\right]^{1/2}$$ rotations to closely align the vector $`U^m\|\gamma `$ with $`\|\mu V^1L^1`$. Thus $`V(U^m\|\gamma )`$ deviates little from the subspace $`L=\text{span}\{\|w_ii=1,2,\mathrm{},\mathrm{}\}`$. A measurement of $`VU^m\|\gamma `$ gives one of the eigenstates $`\|w_j`$, for some $`j:1j\mathrm{}`$, with probability nearly equal to 1, and the task of multiobject search is completed with this large probability.
warning/0007/math0007043.html	ar5iv	text	# On the Motive of the Hilbert scheme of points on a surface ## 1. Introduction Let $`S`$ be a smooth quasiprojective variety over an algebraically closed field $`k`$ of characterstic $`0`$. In this note we determine the class $`[S^{[n]}]`$ of the Hilbert scheme $`S^{[n]}`$ of subschemes of length $`n`$ on $`S`$ in the Grothendieck ring $`K_0(V_k)`$ of $`k`$-varieties. The result expresses $`[S^{[n]}]`$ in terms of the classes of the symmetric powers $`S^{(l)}=S^l/𝔖_l`$. Here $`𝔖_l`$ is the symmetric group acting by permutation of the factors of $`S^l`$. ###### Theorem 1.1. $$[S^{[n]}]=\underset{\alpha P(n)}{}[S^{(\alpha )}\times 𝔸^{n\|\alpha \|}].$$ Here $`P(n)`$ is the set of partitions of $`n`$. We write a partition $`\alpha P(n)`$ as $`(1^{a_1},2^{a_2},\mathrm{},n^{a_n})`$, where $`a_i`$ is the number of occurences of $`i`$ in $`\alpha `$. Then the length $`\|\alpha \|`$ of $`\alpha `$ is the sum of the $`a_i`$ and $`S^{(\alpha )}=S^{(a_1)}\times \mathrm{}\times S^{(a_n)}.`$ In the case $`k=`$ the cohomology of $`S^{[n]}`$ has been studied by a number of authors \[E-S1\],\[Gö1\],\[G-S\],\[Ch1\],\[N\],\[Gr\],\[L\],\[dC-M1\]. In particular the Betti numbers and the Hodge structure have been determined. The class $`[X]`$ of a smooth projective variety over $``$ (or more generally of a projective variety with finite quotient singularities) determines its Hodge structure; so Theorem 1.1 gives a new and elementary proof of the corresponding formulas. By a result of \[Gi-Sou\] and \[Gu-Na\] our result implies that the same formula holds in the Grothendieck ring of effective Chow motives. Similar arguments apply to the incidence variety $`S^{[n,n+1]}S^{[n]}\times S^{[n+1]}`$. At the end we give some applications to moduli spaces of rank two sheaves on surfaces. While I was finishing this paper, the preprint \[dC-M2\] appeared, in which the Chow groups and the Chow motive of $`S^{[n]}`$ are determined over any field using different methods. Our approach is mostly motivated by \[Ch1\] and by \[dB\]. Lemma 4.4 plays an important rôle in this paper. I am very thankful to B. Totaro who proved it for me. B. Totaro also pointed out a mistake in an earlier version of the paper and explained to me how to deduce Conjecture 2.5 from a conjecture of Beilinson and Murre. I thank M.S. Narasimhan and K. Paranjape for very useful discussions. ## 2. Grothendieck rings of varieties and motives and Chow groups In this paper let $`k`$ be an algebraically closed field of characteristic $`0`$. Let $`K_0(V_k)`$ be the Grothendieck ring of $`k`$-varieties. This is the abelian group generated by the isomorphism classes of $`k`$-varieties with the relation that $`[XY]=[X][Y]`$, when $`Y`$ is a closed subvariety of $`X`$. The addition and multiplication in this ring are given by the disjoint union and the product of varieties. Let $`M_k`$ be the category of effective Chow motives over $`k`$. For the precise definitions and some results about motives see e.g. \[Ma\],\[Kl\],\[Sch\]. Let $`A^l(X)`$ be the $`l`$-th Chow group of the variety $`X`$ with $``$-coefficients. Let $`X`$, $`Y`$ be smooth projective varieties. Assume $`X`$ has dimension $`d`$. The group $`Hom_C(X,Y):=A^d(X\times Y)`$ is the group of correspondences from $`X`$ to $`Y`$ of degree $`0`$. An object in $`M_k`$ is a pair $`(X,p)`$ where $`X`$ is a smooth projective variety over $`k`$ and $`pHom_C(X,X)`$ with $`p^2=p`$. The morphisms are $`Hom((X,p),(Y,q))=qHom_C(X,Y)p`$. There is a contravariant functor $`h`$ from the category of smooth projective varieties to $`M_k`$ by sending $`X`$ to $`(X,[\mathrm{\Delta }_X])`$ (where $`\mathrm{\Delta }_XX\times X`$ is the diagonal) and $`f:XY`$ to the class of the transpose of its graph $`[\mathrm{\Gamma }_f^{}]A^{}(Y\times X)`$. We define $`(X,p)(Y,q):=(XY,pq)`$ $`(X,p)(Y,q):=(X\times Y,p\times q)`$. The Grothendieck ring $`K_0(M_k)`$ is the quotient of the free abelian group on the isomorphism classes $`[M]`$ of effective Chow motives by the subgroup generated by elements $`[M][M^{}][M^{\prime \prime }]`$ whenever $`MM^{}M^{\prime \prime }`$. We denote by $`[N]`$ the class of a motive in $`K_0(M_k)`$. ###### Theorem 2.1. \[Gi-Sou\],\[Gu-Na\]. Let $`k`$ be a field of characteristic zero. There exists a unique ring homomorphism $`\overline{h}:K_0(V_k)K_0(M_k)`$ with $`\overline{h}([X])=[h(X)]`$ for $`X`$ smooth projective. The Lefschetz motive $`L`$ is defined by $`h(^1)=1L`$, where $`1:=h(pt)`$ for $`pt=\text{Spec }k`$. By $`[^1]=[pt]+[𝔸^1]`$ we see that $`[L]=\overline{h}([𝔸^1])`$. So Theorem 1.1 immediately implies the identity $$[h(S^{[n]})]=\underset{\alpha P(n)}{}\overline{h}([S^{(\alpha )}])[L^{(n\|\alpha \|)}].$$ in $`K_0(M_k)`$. If a finite group $`G`$ acts on a $`k`$-variety $`X`$, the motive $`h(X/G)`$ of the quotient can be defined as $`(X,_{gG}[g])`$ where $`[g]`$ is the graph of the action by $`g`$. Therefore we can associate two a priory different elements of $`K_0(M_k)`$ to the quotient $`X/G`$, namely $`[h(X/G)]`$ and $`\overline{h}([X/G])`$. ###### Theorem 2.2. \[dB\], Chapter 2, \[dB-Na\]. $`[h(X/G)]=\overline{h}([X/G])`$. Therefore we obtain ###### Corollary 2.3. $$[h(S^{[n]})]=\underset{\alpha P(n)}{}[h(S^{(\alpha )})L^{(n\|\alpha \|)}].$$ The Chow groups of a Chow motive $`N`$ are defined by $`A^l(N)=Hom(L^l,N)`$ and for a smooth projective variety $`X`$ we have $`A^l(h(X))=A^l(X)`$. ###### Remark 2.4. If $`N`$ and $`M`$ are motives with $`[N]=[M]`$, then there exists a motive $`P`$ with $`NP=MP`$. By the definition of $``$ and of Chow groups of motives is is evident that as graded vector spaces $`A^{}(NP)=A^{}(N)A^{}(P)`$. Therefore $`A^{}(N)A^{}(P)=A^{}(M)A^{}(P)`$. If the Chow groups of $`P`$ are finite dimensional, it follows that $`A^{}(N)=A^{}(M)`$. We expect that this result holds without the restriction of finite dimensionality. ###### Conjecture 2.5. If $`N`$ and $`M`$ are effective Chow motives with $`[N]=[M]`$ in $`K_0(M_k)`$, then $`M`$ and $`N`$ are isomorphic. In particular they have the same Chow groups with rational coefficients. In a previous version of this paper the result of Remark 2.4 was claimed without the restriction of finite dimensionality. The mistake was pointed out to me by B. Totaro. He also explained to me the following argument how Conjecture 2.5 follows (over any field $`k`$) from the following conjecture of Beilinson and Murre. ###### Conjecture 2.6. (see \[Ja\] Conj. 2.1.). Let $`H^{}`$ be a Weil cohomology theory. For each smooth projective variety $`X`$ over $`k`$ and all $`j0`$, there exists a descending filtration $`F^{}`$ on $`A^j(X)`$ such that 1. $`F^0A^j(X)=A^j(X)`$ and $`F^1A^j(X)`$ is the kernel of the cyle map $`A^j(X)H^{2j}(X)`$, 2. $`F^rA^i(X)F^sA^j(X)F^{r+s}A^{i+j}(X)`$ for the intersection product, 3. $`F^{}`$ is respected by $`f^{}`$, $`f_{}`$ for morphisms $`f:XY`$, 4. $`F^lA^j(X)=0`$ for $`l0`$. Now we assume Conjecture 2.6 and show Conjecture 2.5. Let $`M=(X,p)`$ be an effective Chow motive. Let $`R:=End(M)A^{}(X\times X)`$. The cycle map induces a homomorphism $`REnd(H^{}(X))`$. Let $`I`$ be the kernel. By the definition of the composition of correspondences and parts 2. and 3. of Conjecture 2.6, we see that for $`fI`$, $`f^nF^nA^{}(X\times X)`$. So, by part 4., $`I`$ is nilpotent. Our aim is to show that this implies that effective Chow motives satisfy the Krull-Schmidt Theorem: Every effective Chow motive is the direct sum of finitely many indecomposable motives, whose isomorphism classes are uniquely defined. This immediately implies Conjecture 2.5: If $`M,NM_k`$ with $`[M]=[N]`$, then $`MPNP`$ for $`PM_k`$. By the Krull-Schmidt Theorem it follows that $`MN`$. In the theorem in Section 3.3 in \[Ga-Ro\] it is shown that an additive category $`𝒞`$ whose isomorphism classes form a set satisfies the Krull Schmidt Theorem if the following holds: Every idempotent in $`𝒞`$ splits and for each object $`A`$ in $`𝒞`$, if we write $`R:=End(A)`$, then $`R/rad(R)`$ is semisimple and all idempotents in $`R/rad(R)`$ are the images of idempotents in $`R`$. Here $`rad(R)`$ is the Jacobson radical of $`R`$. We check these conditions for the category $`M_k`$ of effective Chow motives. Idempotents split because $`M_k`$ is pseudoabelian. For a motive $`M=(X,p)`$ let as above $`R:=End(M)`$ and $`I=ker(REnd(H^{}(X))`$. Then $`I`$ is nilpotent and $`R/I`$ is a finite dimensional $``$-algebra. Since $`I`$ is nilpotent, $`Irad(R)`$. So $`R/rad(R)`$ is a finite dimensional $``$-algebra with radical $`0`$. So it is semisimple. Furthermore we see that $`rad(R)`$ is nilpotent: $`rad(R/I)`$ is nilpotent because $`R/I`$ is finite dimensional. The result follows for $`rad(R)`$ because $`I`$ is nilpotent. Then Theorem 1.7.3 in \[Be\] implies that all idempotents of $`R/rad(R)`$ lift to idempotents of $`R`$. Corollary 2.3 and Conjecture 2.5 imply the formulas $`h(S^{[n]})`$ $`={\displaystyle \underset{\alpha P(n)}{}}h(S^{(\alpha )})L^{(n\|\alpha \|)},`$ $`A^i(S^{[n]})`$ $`={\displaystyle \underset{\alpha P(n)}{}}A^{i+\|\alpha \|n}(S^{(\alpha )}).`$ These formulas have been shown in \[dC-M2\] over an arbitrary field. ## 3. The stratification of $`S^{(n)}`$ and $`S^{[n]}`$ Let $`\omega _n:S^{[n]}S^{(n)}`$ be the Hilbert-Chow morphism, which associates to each subscheme $`Z`$ its support with multiplicities. $`S^{[n]}`$ and $`S^{(n)}`$ have a natural stratification parameterized by the partitions of $`n`$, which has been used before \[Gö1\],\[G-S\],\[Ch1\],\[dC-M1\] to compute the cohomology of $`S^{[n]}`$. Let $`P(n)`$ be the set of all partitions of $`n`$. A partition $`\alpha =(n_1,\mathrm{},n_r)`$ is also written as $`\alpha =(1^{a_1},\mathrm{},n^{a_n})`$, where $`a_i`$ is the number of occurences of $`a_i`$ in $`\alpha `$. We write $`\|\alpha \|=r=_ia_i`$. The corresponding locally closed stratum $`S_\alpha ^{(n)}`$ is the set of all zero-cycles $`\xi =n_1x_1+\mathrm{}n_rx_r`$ with $`x_1,\mathrm{},x_r`$ distinct points of $`S`$. We put $`S_\alpha ^{[n]}:=\omega _n^1(S_\alpha ^{[n]})_{red}.`$ The strata $`S_\alpha ^{(n)}`$ are smooth, but the $`S_\alpha ^{[n]}`$ and the closures $`\overline{S_\alpha ^{(n)}}`$ usually are singular. There is a natural map $`h_\alpha :S^{(\alpha )}S^{(n)};(\xi _1,\mathrm{},\xi _n)_{i=1}^ni\xi _i`$ whose image is the closure $`\overline{S_\alpha ^{(n)}}`$; in fact it is easy to see that it is the normalization of $`\overline{S_\alpha ^{(n)}}`$. Let $`g_\alpha :S^{(\alpha )}\times 𝔸^{n\|\alpha \|}S^{(n)}`$ be the composition of the projection to $`S^{(\alpha )}`$ with $`h_\alpha `$, and let $`g:_{\alpha P(n)}S^{(\alpha )}\times A^{n\|\alpha \|}S^{(n)}`$ be the map induced by the $`g_\alpha `$. A more precise version of Theorem 1.1 is the following. ###### Proposition 3.1. $`[g^1(S_\beta ^{(n)})]=[S_\beta ^{[n]}]`$ in $`K_0(V_k)`$ for all $`\beta P(n)`$. Theorem 1.1 follows from Proposition 3.1 by summing over all $`\beta P(n)`$. ## 4. Proof of the main result We will determine both sides of the equality in Proposition 3.1. We need some preliminaries. ###### Remark 4.1. In the Grothendieck ring of $`k`$-varieties we have: 1. If $`f:XY`$ is a Zariski locally trivial fibre bundle with fibre $`F`$, then $`[X]=[Y][F]`$ (stratify $`Y`$ such that $`f`$ is trivial over the strata). 2. If $`f:XY`$ is a bijective morphism, then $`[X]=[Y]`$. (There is a dense open subset $`UX`$, on which $`f`$ is an isomorphism (here we use $`char(k)=0`$), replacing $`X`$ by $`XU`$ and $`Y`$ by $`Yf(U)`$, we can argue by induction over the dimension.) ###### Remark 4.2. If a quasiprojective variety $`X`$ has a decomposition $`X=X_1\mathrm{}X_l`$ into locally closed subvarieties, then it is immediate to see that $`_{n_1+\mathrm{}+n_l=n}_{i=1}^lX_i^{(n_i)}`$ is a decomposition of $`X^{(n)}`$ into locally closed subvarieties. Therefore we can define $`[X]^{(n)}=[X^{(n)}]`$ for $`X`$ a variety, and $`([X_1]+\mathrm{}+[X_l])^{(n)}=_{n_1+\mathrm{}+n_l=n}_{i=1}^l[X_i]^{(n_i)}.`$ ###### Notation 4.3. Let $`X`$ be a $`k`$-variety. If $`Y`$ is a variety with a natural morphism to $`X^{(m)}`$ for some $`m>0`$, we write $`Y_{}`$ for the preimage of the open subvariety of all zero cycles $`\xi X^{(m)}`$ whose support consists of $`m`$ distinct points. The following lemma was proved for me by Burt Totaro. ###### Lemma 4.4. Let $`X`$ be a variety over $`k`$. Let $`p:(X\times 𝔸^l)^{(n)}X^{(n)}`$ be the obvious projection. Then $`[(X\times 𝔸^l)^{(n)}]=[X^{(n)}\times 𝔸^{nl}]`$, and $`[p^1(X_\alpha ^{(n)})]=[X_\alpha ^{(n)}\times 𝔸^{nl}]`$ for all $`\alpha P(n)`$. ###### Proof. The first statement follows from the second. It is enough to treat the case $`l=1`$; the general case follows by trivial induction. There is a cartesian diagram $$\begin{array}{ccc}X_{}^{\|\alpha \|}\times _{i=1}^n((𝔸^1)^{(i)})^{a_i}& \stackrel{\overline{p}}{}& X_{}^{\|\alpha \|}\\ & & \\ p^1(X_\alpha ^{(n)})& \stackrel{p}{}& X_\alpha ^{(n)}\end{array}$$ By the fundamental theorem of symmetric functions the fibre product is $`X^{\|\alpha \|}\times _i(𝔸^i)^{a_i}=X_{}^{\|\alpha \|}\times 𝔸^n.`$ So we get an étale trivialization of $`p`$. Any two trivializations are related by the action of the group $`𝔖_{a_1}\times \mathrm{}\times 𝔖_{a_n}`$ by reordering the factors in the $`X^{a_i}`$. This acts on the fibres of $`\overline{p}`$ by reordering the factors $`(𝔸^1)^{(i)}`$ in the $`((𝔸^1)(i))^{a_i}`$. Choosing an origin in $`𝔸^1`$ determines an origin in $`𝔸^n`$, and the action becomes linear on on $`k^n`$, so $`p^1(X_\alpha ^{(n)})`$ is an étale locally trivial vector bundle over $`X_\alpha ^{(n)}`$. Therefore, by Hilbert Theorem 90 \[Se\] p. 1.24, it is locally trivial in the Zariski topology. Thus we get $`[p^1(X_\alpha ^{(n)}]=[X_\alpha ^{(n)}\times 𝔸^n]`$. ∎ Now we determine the right hand side in Proposition 3.1. Let $`R_n=Hilb^n(𝔸^2,0)`$ be the punctual Hilbert scheme of subschemes of length $`n`$ of $`𝔸^2`$ concentrated in $`0`$. Then by \[E-S1\] $`R_n`$ has a cell decomposition and $$[R_n]=\underset{\beta P(n)}{}[𝔸^{n\|\beta \|}].$$ ###### Lemma 4.5. $`[S_\alpha ^{[n]}]=[\left(_{i=1}^n(S\times R_i)^{(a_i)}\right)_{}].`$ ###### Proof. By Lemma 2.1.4 of \[Gö2\] $`S_{(l)}^{[l]}`$ is a locally trivial fibre bundle over $`S`$ with fibre $`R_l`$, thus $`[S_{(l)}^{[l]}]=[S\times R_l]`$. There is a natural morphism $`f:\left(_{i=1}^n(S_{(i)}^{[i]})^{(a_i)}\right)_{}S_\alpha ^{[n]}`$ defined on $`T`$-valued points by sending $`(Z_1,\mathrm{},Z_n)`$ to $`_{i=1}^nZ_i`$. $`f`$ is obviously invariant under the action of $`𝔖_{a_1}\times \mathrm{}\times 𝔖_{a_n}`$ by permuting the factors in the $`(S_{(i)}^{[i]})^{(a_i)}`$, and the induced morphism from the quotient to $`S_\alpha ^{[n]}`$ induces a bijection on $`k`$-valued points. This implies $`[S_\alpha ^{[n]}]=\left[\left(_{i=1}^n(S_{(i)}^{[i]})^{(a_i)}\right)_{}\right]`$. ∎ ###### Notation 4.6. 1. For any $`xS`$ and any $`\xi S^{(n)}`$ we call $`m_x(\xi )`$ the multiplicity with which $`x`$ occurs in $`\xi `$. 2. We denote $`P:=_{n>0}P(n)`$. For $`\alpha =(1^{a_1},2^{a_2},\mathrm{})P(n)`$, $`\beta :=(1^{b_1},2^{b_2},\mathrm{})P(m)`$ and $`l_0`$ we denote $`l\alpha :=(1^{la_1},2^{la_2}\mathrm{})P(nl)`$, $`\alpha +\beta :=(1^{a_1+b_1},2^{a_2+b_2}\mathrm{})P(n+m)`$. ###### Lemma 4.7. $`[S_\alpha ^{[n]}]=\left[_f\left(_{\beta P}S^{(f(\beta ))}\right)_{}\times 𝔸^{n_{\beta P}f(\beta )\|\beta \|}\right].`$ Here $`f`$ runs through the $`f:P_0`$ with $`_{\beta P(i)}f(\beta )=a_i`$ for all $`i`$. ###### Proof. By Lemma 4.5 we get $`[S_\alpha ^{[n]}]=\left[\right(_{i=1}^n(_{\beta _iP(i)}S\times 𝔸^{i\|\beta _i\|})^{(a_i)})_{}].`$ By Remark 4.2 this implies $`[S_\alpha ^{[n]}]=\left[\right(_{i=1}^n_{f_i}_{\beta _iP(i)}(S\times 𝔸^{i\|\beta _i\|})^{(f_i(\beta _i)})_{}],`$ where the $`f_i`$ rund through the $`f_i:P(i)_0`$ with $`_{\beta P(i)}f_i(\beta )=a_i`$. The result follows by Lemma 4.4. ∎ Now we determine the left hand side in Proposition 3.1. Let $`\beta =(1^{b_1},2^{b_2},\mathrm{})=(n_1,\mathrm{},n_r)P(n)`$. ###### Lemma 4.8. $`h_\alpha ^1S_\beta ^{(n)}_f\left(_{\gamma P}S^{(f(\gamma ))}\right)_{}.`$ Here the sum is over all functions $`f:P_0`$ with $`_{\gamma P(i)}f(\gamma )=b_i`$ and $`_{\gamma P}f(\gamma )\gamma =\alpha `$. Using that $`g_\alpha ^1(S_\beta ^{(n)})=h_\alpha ^1(S_\beta ^{(n)})\times 𝔸^{n\|\alpha \|}`$, Proposition 3.1 follows immediately from Lemma 4.7 and Lemma 4.8 by summing over all $`\beta `$. ###### Proof. Any $`\xi =(\xi _1,\mathrm{},\xi _n)h_\alpha ^1S_\beta ^{(n)}`$ induces a map $`f_\xi :P_0`$ as follows. Let $`h_\alpha (\xi )=_{i=1}^rn_ix_i`$. For all $`xS`$ let $`\gamma _x(\xi ):=_{j=1}^n(j^{m_x(\xi _j)})`$. For each $`i=1,\mathrm{}r`$ we have $`_{j=1}^njm_{x_i}(\xi _j)=n_i`$, so $`\gamma _{x_i}(\xi )P(n_i)`$; furthermore $`_i\gamma _{x_i}(\xi )=\alpha `$. We define $`f_\xi :P_0`$ by $`f_\xi (\gamma ):=\mathrm{\#}\left\{xS\right\|\gamma _x(\xi )=\gamma \}.`$ Then $`_{\gamma P(j)}f_\xi (\gamma )=b_j`$ and $`_{\gamma P}f_\xi (\gamma )\gamma =\gamma _{x_1}(\xi )+\mathrm{}+\gamma _{x_r}(\xi )=\alpha .`$ Now fix $`f:P_0`$ with the above properties. Let $`S_f^{(\alpha )}:=\left\{\xi S^{(\alpha )}\right\|f_\xi =f\}`$. We claim that $`S_f^{(\alpha )}\left(_{\gamma P}S^{(f(\gamma ))}\right)_{}.`$ For $`\xi S_\mathrm{\Gamma }^{(\alpha )}`$ define $`\varphi (\xi )=(\varphi (\xi )_\gamma )_{\gamma P}`$ by letting $`\varphi (\xi )_\gamma S^{(f(\gamma ))}`$ be the sum over all $`xS`$ with $`\gamma _x(\xi )=\gamma `$. For $`\zeta =(\zeta _\gamma )_{\gamma P}`$ with $`\zeta _\gamma S^{(f(\gamma ))}`$ let $`\psi (\zeta ):=(\xi _1,\mathrm{},\xi _n)`$ with $`\xi _i=_{\gamma P}c_i\zeta _\gamma S^{(a_i)},`$ where we write $`\gamma =(1^{c_1},2^{c_2},\mathrm{})`$. It is straightforward from the definitions that $`\varphi `$ and $`\psi `$ are inverse to each other. ∎ ###### Example 4.9. 1. Let $`S`$ be a projective rational surface. Then $`[S]=[𝔸^0]+b[𝔸^1]+[𝔸^2]`$ for suitable $`b>0`$, and $$\underset{n0}{}[S^{[n]}]t^n=\underset{l>0}{}\frac{1}{(1[𝔸^{l1}]t^l)(1[𝔸^l]t^l)^b(1[𝔸^{l+1}]t^l)}.$$ 2. Let $`S`$ be a ruled surface over a curve $`C`$. Then $$\underset{n0}{}[S^{[n]}]t^n=\underset{l>0}{}\left(\underset{m0}{}[C^{(m)}\times 𝔸^{m(l1)}]t^{ml}\right)\left(\underset{m0}{}[C^{(m)}\times 𝔸^{ml}]t^{ml}\right).$$ 3. If $`\widehat{S}`$ is the blowup of $`S`$ in a point then $$\underset{n0}{}[\widehat{S}^{[n]}]t^n=\frac{_{n0}[S^{[n]}]t^n}{_{l>0}(1[𝔸^l]t^l)}.$$ This follows from Theorem 1.1 and Lemma 4.4. ## 5. The incidence variety Similar but simpler arguments to those for $`S^{[n]}`$ can be used for the incidence variety $`S^{[n,n+1]}:=\left\{(Z,W)S^{[n]}\times S^{[n+1]}\right\|ZW\}`$, which plays a rôle in inductive arguments for $`S^{[n]}`$ \[E-S2\],\[E-G-L\]. The Hodge numbers of $`S^{[n,n+1]}`$ were computed in \[Ch1\]. ###### Theorem 5.1. $`[S^{[n,n+1]}]={\displaystyle \underset{l=0}{\overset{n}{}}}[S\times S^{[l]}\times 𝔸^{nl}].`$ By Theorem 2.1 this immediately implies ###### Corollary 5.2. $`[h(S^{[n,n+1]})]={\displaystyle \underset{l=0}{\overset{n}{}}}[h(S\times S^{[l]})L^{(nl)}].`$ For the proof we introduce a stratification of $`S^{[n,n+1]}`$. Let $$\overline{\omega }:S^{[n,n+1]}S\times S^{(n)},(Z,W)(\omega _{n+1}(W)\omega _n(Z),\omega _n(Z)).$$ For $`0mn`$ let $`(S\times S^{(n)})_m:=\left\{(x,\xi )S\times S^{(n)}\right\|m_x(\xi )=m\}`$, and let $`S_m^{[n,n+1]}:=\overline{\omega }^1((S\times S^{(n)})_m)_{red}`$. The $`(S\times S^{(n)})_m`$ and the $`S_m^{[n,n+1]}`$ form stratifications of $`S\times S^{(n)}`$ and $`S^{[n,n+1]}`$ respectively into locally closed subvarieties. Let $$\overline{g}:\underset{m=0}{\overset{n}{}}S\times S^{[m]}\times 𝔸^{nm}S\times S^{(n)};(x,Z,a)(x,(nm)x+\omega _m(Z)).$$ Then Theorem 5.1 follows from the following. ###### Proposition 5.3. $`[S_m^{[n,n+1]}]=[\overline{g}^1((S\times S^{(n)})_m).]`$ ###### Proof. If $`X`$ is a variety with a natural map to $`S\times S^{(m)}`$ for some $`m0`$, we will write $`X_0`$ for the preimage of the locus of $`(x,\xi )`$ with $`xsupp(\xi )`$. Let $`(𝔸^2,0)^{[n,n+1]}:=\left\{(Z,W)(𝔸^2)^{[n,n+1]}\right\|ZW,supp(W)=\{0\}\}`$ with the reduced structure. In \[Ch2\] it is shown that $`(𝔸^2,0)^{[n,n+1]}`$ has a cell decomposition. Her formula for the numbers of cells of different dimensions implies that $`[(𝔸^2,0)^{[n,n+1]}]=_{l=0}^n[R_l\times 𝔸^{nl}].`$ We now determine $`[S_m^{[n,n+1]}]`$. First it is easy to see analoguously to the case of $`S_{(n)}^{[n]}`$ in \[Gö2\] that $`S_n^{[n,n+1]}`$ is a locally trivial fibre bundle over $`S`$ with fibre $`(𝔸^2,0)^{[n,n+1]}`$. Therefore $$[S_n^{[n,n+1]}]=\underset{l=0}{\overset{n}{}}[S\times R_l\times 𝔸^{nl}]=\underset{l=0}{\overset{n}{}}[S_{(l)}^{[l]}\times 𝔸^{nl}].$$ There is a natural morphism $`\sigma :(S_m^{[m,m+1]}\times S^{[nm]})_0S_m^{[n,n+1]}`$ given on $`T`$-valued points by $`((Z,W),X)(ZX,WX)`$. $`\sigma `$ is obviously a bijection on $`k`$ valued points. Thus we get $$[S_m^{[n,n+1]}]=[(S_m^{[m,m+1]}\times S^{[mn]})_0]=\underset{l=0}{\overset{m}{}}[(S_{(l)}^{[l]}\times S^{[nm]})_0\times 𝔸^{ml}].$$ Now we determine $`\overline{g}^1(S\times S^{(n)})_m)`$. Let $`(S\times S^{[m]})_l=\left\{(x,Z)S\times S^{[m]}\right\|len_x(Z)=l\}.`$ Then $`\overline{g}^1((S\times S^{(n)})_m)=_{l=0}^m(S\times S^{[nm+l]})_l\times 𝔸^{ml}`$. Furthermore we have a morphism $`\varphi :(S_{(l)}^{[l]}\times S^{[ml]})_0(S\times S^{[m]})_l`$ sending $`(Y,Z)`$ to $`(supp(Y),YZ)`$, which is bijective on $`k`$-valued points. Thus $`[g^1((S\times S^{(n)})_m)]=_{l=0}^m[(S_{(l)}^{[l]}\times S^{[nm]})_0\times 𝔸^{ml}].`$ ## 6. Moduli of stable sheaves Let $`S`$ be a projective surface over $`k`$. Fix $`CPic(S)`$ and let $`H`$ be an ample line bundle on $`S`$. We denote by $`NS(S)`$ the Picard group of $`S`$ modulo numerical equivalence. The moduli space $`M_S^H(C,d)`$ of $`H`$-semistable rank $`2`$ torsion-free sheaves $`E`$ with $`det(E)=C`$ and $`c_2C^2/4=d`$ depends on $`H`$ via a system of walls and chambers. This dependence has been studied and used by various authors (e.g. \[Q1\],\[F-Q\],\[E-G\],\[Gö3\]). A class $`\xi NS(S)+C/2`$ is of type $`(C,d)`$ if $`0<\xi ^2d`$ and there exists an ample divisor $`H`$ with $`\xi H=0`$. In this case we say that $`H`$ lies on the corresponding wall. Ample divisors $`H,LPic(S)`$ are separated by $`\xi `$ if $`(\xi H)(\xi L)<0`$. Assume that $`H`$ and $`L`$ do not lie on a wall of type $`(C,d)`$. If they are not separated by a class of type $`(C,d)`$ we say that they lie in the same chamber of type $`(C,d)`$. In this case $`M_S^L(C,d)=M_S^H(C,d)`$. More generally let $`FPic(S)`$ be nef with $`F^20`$ and $`FC`$ odd. Then $`F+nH`$ is ample for any ample divisor $`H`$, and for $`n`$ sufficiently large the chamber of $`F+nH`$ does not depend on $`H`$. We will write $`M_S^F(C,d):=M_S^{F+nH}(C,d)`$ for $`n0`$. Let $`L`$, $`H`$ be ample divisors not on a wall of type $`(C,d)`$. If $`2\xi +K_S`$ is not effective for all $`\xi `$ separating $`L`$ and $`H`$ (we say that $`L`$ and $`H`$ are separated only by good walls), then $`M_S^L(C,d)`$ is obtained from $`M_S^H(C,d)`$ by successively blowing up along projective space bundles over products $`S^{[l]}\times S^{[n]}`$ of Hilbert schemes of points followed by blowdowns of the exceptional divisor to another projective space bundle over $`S^{[l]}\times S^{[n]}`$. Therefore the proof of Theorem 3.4 in \[Gö3\] shows: ###### Proposition 6.1. Let $`L`$, $`H`$ be ample divisors not on a wall of type $`(C,d)`$ separated only by good walls. Then $`[M_S^L`$ $`(C,d)][M_S^H(C,d)]=[Pic^0(S)]`$ $`{\displaystyle \underset{\xi }{}}[(SS)^{[d+\xi ^2]}]\left([^{d\xi ^2+\xi K_S\chi (𝒪_S)1}][^{d\xi ^2\xi K_S\chi (𝒪_S)1}]\right).`$ The sum is over all $`\xi `$ of type $`(C,d)`$ with $`\xi L>0>\xi H`$, and we use the convention $`^1=\mathrm{}`$. ###### Corollary 6.2. Under the assumptions of Proposition 6.1, if $`S`$ is a rational surface, then $`[M_S^L(C,d)][M_S^H(C,d)]`$ is a $``$-linear combination of the $`[𝔸^l]`$ with $`l4d3`$. This follows from Proposition 6.1, and Example 4.9. ###### Corollary 6.3. If $`K_S`$ is numerically trivial, then $`[M_S^H(C,d)]`$ does not depend on $`H`$ as long as $`H`$ does not lie on a wall. In \[G-Z\] and \[Gö4\] Theta functions for indefinite lattices were introduced to study the wallcrossing. Let $`\mathrm{\Gamma }`$ be $`NS(S)`$ with the negative of the intersection form as quadratic form, which we denote by $`,`$. Then for $`F,GPic(S)`$ with $`F^20`$, $`G^20`$, $`FG>0`$, we define $$\mathrm{\Theta }_{\mathrm{\Gamma },C}^{F,G}(\tau ,x):=\underset{\xi \mathrm{\Gamma }+C/2}{}(\mu (\xi ,F\mu (\xi ,G)q^{\xi ,\xi /2}e^{2\pi i\xi ,x}.$$ Here $`\mu (t)=1`$ of $`t0`$, and $`\mu (t)=0`$ otherwise, and $`q=e^{2\pi i\tau }`$ for $`\tau `$ in the complex upper half plane $``$ and $`x\mathrm{\Gamma }_{}`$. This function is defined on a suitable open subset of $`\times \mathrm{\Gamma }_{}`$ and has a meromorphic extension to the whole of $`\mathrm{\Gamma }_{}`$. For a meromorphic function $`f:\times \mathrm{\Gamma }_{}`$ and $`v\mathrm{\Gamma }_Z`$ we write $$f\|_v(\tau ,x):=q^{v,v/2}e^{2\pi iv,x}f(\tau ,x+v\tau ).$$ Then $`\mathrm{\Theta }_{\mathrm{\Gamma },C}^{F,G}(\tau ,x)=\mathrm{\Theta }_\mathrm{\Gamma }^{F,G}\|_{C/2}(\tau ,x)`$, where we have written $`\mathrm{\Theta }_\mathrm{\Gamma }^{F,G}:=\mathrm{\Theta }_{\mathrm{\Gamma },0}^{F,G}.`$ The reason for introducing these theta functions was that they can be expressed in terms of standard theta functions in case $`F^2=G^2=0`$. In the rest of this section we write $`y:=e^{2\pi iz}`$ for $`z`$ a complex variable. Recall the standard theta functions $$\mathrm{\Theta }_{\mu ,\nu }(\tau ,z):=\underset{nZ}{}(1)^{n\nu }q^{(n+\mu /2)^2/2}y^{n+\mu /2}(\mu ,\nu \{0,1\}).$$ If $`F^2>0`$ and $`G^2>0`$ or $`FC`$ and $`GC`$ are odd, then for every $`L\mathrm{\Gamma }`$, $`\mathrm{\Theta }_{\mathrm{\Gamma },C}^{F,G}(\tau ,Lz)`$ is a power series in $`q^{1/8}`$ with coefficients Laurent polynomials in $`y^{1/2}`$. We will write $`\mathrm{\Theta }_{\mathrm{\Gamma },C}^{F,G}(2\tau ,K_Sz)^{}`$ for the power series in $`t^{1/4}`$ with coefficients Laurent polynomials in $`[𝔸^1]`$, which we obtain by replacing $`y`$ by $`[𝔸^1]`$ and $`q`$ by $`[𝔸^2]t`$ in $`\mathrm{\Theta }_{\mathrm{\Gamma },C}^{F,G}(2\tau ,K_Sz)`$. There are no half integer powers of $`[𝔸^1]`$, because $`K_S`$ is characteristic. The following follows from Theorem 3.4 in \[Gö3\] in the same way as Theorem 4.1 in \[Gö4\]. ###### Corollary 6.4. Assume $`H,L`$ do not lie on a wall of type $`(C,d)`$ for any $`d`$, and are separated only by good walls. Then, in $`K_0(V_k)`$ $`{\displaystyle \underset{d0}{}}`$ $`([M_S^H(C,d)][M_S^L(C,d)])t^d`$ $`=[Pic^0(S)]\left({\displaystyle \underset{n0}{}}[S^{[n]}][𝔸^n]t^n\right)^2{\displaystyle \frac{\mathrm{\Theta }_{\mathrm{\Gamma },C}^{L,H}(2\tau ,K_Sz)^{}}{[𝔸^1]^{\chi (𝒪_S)}([𝔸^1]1)}}.`$ On the r.h.s. we mean that, after multiplying out, all negative powers of $`[𝔸^1]`$ vanish. In \[L-Q1\], \[L-Q2\] a blowup formula was proven for the Euler numbers and the virtual Hodge numbers of moduli spaces of stable rank $`2`$ sheaves on surfaces. In \[Ka\] a blowup formula is proven for principal bundles. We can show that the formula of \[L-Q1\], \[L-Q2\] holds in $`K_0(V_k)`$. Let $`H`$ be ample on $`S`$ and assume that $`CH`$ is odd. Let $`\widehat{S}`$ be the blowup of $`S`$ in a point and denote by $`E`$ the exceptional divisor, we denote by $`H`$ also the pullback of $`H`$ to $`\widehat{S}`$. ###### Theorem 6.5. Assume $`k=`$. Let $`a\{0,1\}`$ then $`{\displaystyle \frac{_{d0}[M_{\widehat{S}}^H(C+aE,d)]t^d}{_{d0}[M_S^H(C,d)]t^d}}`$ $`={\displaystyle \frac{_n[𝔸^{\left(\genfrac{}{}{0pt}{}{2n+a+1}{2}\right)}]t^{(n+\frac{a}{2})^2}}{_{l>0}(1[𝔸^{2l}]t^l)}}.`$ ###### Proof. In \[L-Q1\] the authors use virtual Hodge polynomials $`e(X:x,y)`$ in order to show that there exists a universal power series $`Z_a(x,y,t)`$ such that $$\underset{d0}{}e(M_{\widehat{S}}^H(C+aE,d),:x,y)t^d=Z_a(x,y,z)\left(\underset{d0}{}e(M_S^H(C,d):x,y)t^d\right).$$ To do this they only use the basic property of virtual Hodge polynomials that $`e(XY:x,y)=e(X:x,y)e(Y:x,y)`$ for $`Y`$ a closed subvariety of $`X`$. So their proof shows that there is a universal power series $`Y_a(t)`$ in $`K_0(V_k)[[t]]`$, such that $$\underset{d0}{}[M_{\widehat{S}}^H(C+aE,d)]t^d=Y_a(t)\left(\underset{d0}{}[M_S^H(C,d)]t^d\right).$$ In the paper \[L-Q2\] they compare the wallcrossing on a rational ruled surface and its blowup in a point to determine $`Z_a(x,y,t)`$. We can translate their argument into our language, where it proves the theorem: Let $`H`$, $`L`$ be ample on $`S`$ with $`CH`$ and $`CL`$ odd. Write $`M_d:=[M_S^H(C,d)][M_S^L(C,d)]`$ and $`M_{a,d}:=[M_{\widehat{S}}^H(C+aE,d)][M_{\widehat{S}}^L(C+aE,d)]`$. Write $`\mathrm{\Gamma }=H^2(S,)`$. Then, as $`Y_a(t)`$ is universal, we get, using Corollary 6.4, $`Y_a(t)`$ $`={\displaystyle \frac{_{d0}M_{a,d}t^d}{_{d0}M_dt^d}}={\displaystyle \frac{\mathrm{\Theta }_{\mathrm{\Gamma }E,C+aE}^{L,H}(2\tau ,K_{\widehat{S}}z)^{}\left(_{n0}[\widehat{S}^{[n]}]t^n\right)^2}{\mathrm{\Theta }_{\mathrm{\Gamma },C}^{L,H}(2\tau ,K_Sz)^{}\left(_{n0}[S^{[n]}]t^n\right)^2}}`$ By definition it is obvious that $`\mathrm{\Theta }_{\mathrm{\Gamma }E,C+aE}^{L,H}(\tau ,K_{\widehat{S}}z)=\theta _{a,0}(\tau ,z)\mathrm{\Theta }_{\mathrm{\Gamma },C}^{L,H}(\tau ,K_Sz)`$. The result follows by Example 4.9.(3) and the identity $$\theta _{a,0}(2\tau ,z)^{}=\underset{n}{}[𝔸^{\left(\genfrac{}{}{0pt}{}{2n+a+1}{2}\right)}]t^{(n+\frac{a}{2})^2}.$$ Corollary 6.4 and Theorem 6.5 imply in particular that the computations of \[Gö4\] hold in the Grothendieck group of varieties (replacing $`y`$ in the formulas there by $`[𝔸^1]`$). As there, we use the following elementary fact. ###### Remark 6.6. Let $`S`$ be the blowup of a ruled surface in finitely many points and $`F`$ the pullback of a fibre of the ruling. Assume $`FC`$ is odd. Then $`M_S^F(C,d)=\mathrm{}`$ for all $`d`$. In particular we get the following results. ###### Corollary 6.7. Let $`S`$ be a rational surface. 1. Let $`CPic(S)\{0\}`$. Let $`H`$ be an ample divisor not on a wall of type $`(C,d)`$ and assume that $`K_SH0`$, then $`[M_S^H(C,d)]`$ is a $``$-linear combination of $`[𝔸^l]`$ with $`l4d3`$. (\[Gö4\], Prop.4.9), 2. Let $`S`$ be the blowup of $`^2`$ in 9 points, with exceptional divisors $`E_1,\mathrm{},E_9`$ and let $`H`$ be the pullback of the hyperplane class. Let $`F:=3HE_1\mathrm{}E_9`$, and let $`CH^2(S,)`$ with $`C^2`$ odd. Then $`[M_S^F(C,d)]=[S^{[2d3/2]}]`$. (\[Gö4\], Thm.7.3). ###### Remark 6.8. For $`S`$ a rational surface it follows from \[dC-M2\] that the rational Chow groups of the $`S^{[n]}`$ are finite dimensional. By Remark 2.4 and the discussion before Proposition 6.1 it follows that under the conditions of Corollary 6.7.1. the calculations in \[Gö4\] hold also in the Chow ring of $`M_S^H(C,d)`$. In particular in 2. we get that $`M_S^F(C,d)`$ and $`S^{[2d3/2]}`$ have the same Chow groups with rational coefficients. As a final example we determine the classes of some moduli spaces over a ruled surface $`S`$ over an elliptic curve $`E`$ with a section $`\sigma `$ with minimal self-intersection $`\sigma ^2=1`$. ###### Proposition 6.9. Let $`F`$ be the class of a fibre of the ruling. We write $`G:=2\sigma F`$. Let $`CNS(S)`$ with $`C^2`$ odd. Then $`[M_S^G(C,d)]`$ is the coefficient of $`t^{2d1/2}`$ in $`[E]`$ $`{\displaystyle \underset{m>0}{}}{\displaystyle \underset{i=0,1}{}}(1+[E]\left({\displaystyle \underset{l1}{}}[^{l1}][𝔸^{m(2li)}]t^{2lm}\right))^2`$ $`{\displaystyle \underset{n>0}{}}((1[𝔸^{n1}]t^n)(1[𝔸^n]t^n)^2(1[𝔸^{n+1}]t^n))^{(1)^n}`$ ###### Proof. $`C^2`$ odd implies $`CF`$ odd and $`CG`$ odd, therefore Remark 6.6 implies that $`M_S^F(C,d)=\mathrm{}`$. By Proposition 6.1 $`[M_S^G(C,d)]`$ depends only on the numerical equivalence class of $`C`$. Therefore we can assume that $`C=\sigma =\frac{G+F}{2}`$ of $`C=\sigma F=\frac{GF}{2}`$. We write $`\mathrm{\Gamma }=NS(S)`$ with the negative of the intersection form. We note that $`K_S=G`$. By Corollary 6.4 $`[M_S^G(C,d)]`$ is the coefficient of $`t^d`$ in $$[E]\left(\underset{n0}{}[S^{[n]}][𝔸^n]t^n\right)^2\frac{\mathrm{\Theta }_{\mathrm{\Gamma },\frac{F+G}{2}}^{F,G}(2\tau ,Gx)^{}+\mathrm{\Theta }_{\mathrm{\Gamma },\frac{FG}{2}}^{F,G}(2\tau ,Gx)^{}}{[𝔸^1]1}.$$ F or $`n1`$ $`E^{(n)}`$ is a locally trivial bundle over $`E`$ with fibre $`^{n1}`$, thus $`[E^{(n)}]=[E\times ^{n1}]`$. Using Example 4.9, we only need to show that $`\mathrm{\Theta }_{\mathrm{\Gamma },\frac{F+G}{2}}^{F,G}`$ $`(2\tau ,Gz)+\mathrm{\Theta }_{\mathrm{\Gamma },\frac{FG}{2}}^{F,G}(2\tau ,Gz)`$ $`=(y^{1/2}y^{1/2})q^{1/4}{\displaystyle \underset{n>0}{}}\left((1q^{n/2}y^1)(1q^{n/2})^2(1q^{n/2}y)\right)^{(1)^n}.`$ Let $`L`$ be the lattice generated by $`F/2,G/2`$. $`0`$ and $`G/2`$ are a basis of $`L`$ modulo $`\mathrm{\Gamma }`$. Therefore $`\mathrm{\Theta }_\mathrm{\Gamma }^{F,G}+\mathrm{\Theta }_{\mathrm{\Gamma },G}^{F,G}=\mathrm{\Theta }_L^{F,G}`$ and $`\mathrm{\Theta }_{\mathrm{\Gamma },\frac{F+G}{2}}^{F,G}+\mathrm{\Theta }_{\mathrm{\Gamma },\frac{FG}{2}}^{F,G}=\mathrm{\Theta }_{L,\frac{F+G}{2}}^{F,G}`$. By $`F^2=G^2=0`$ and $`FG=2`$, formula (2.14) in \[Gö4\] gives that $$\mathrm{\Theta }_L^{F,G}(2\tau ,x)=\frac{\eta (\tau )^3\theta _{1,1}(\tau ,(F+G)/2,x)}{\theta _{1,1}(\tau ,F/2,x)\theta _{1,1}(\tau ,G/2,x)}.$$ Easy computations give $$\mathrm{\Theta }_{L,\frac{F+G}{2}}^{F,G}(2\tau ,x)=\mathrm{\Theta }_L^{F,G}\|_{\frac{F+G}{2}}(2\tau ,x)=\frac{\eta (\tau )^3\theta _{1,1}(\tau ,(F+G)/2,x)}{\theta _{0,1}(\tau ,F/2,x)\theta _{0,1}(\tau ,G/2,x)},$$ $$\mathrm{\Theta }_{L,\frac{F+G}{2}}^{F,G}(2\tau ,Gz)=\frac{\eta (\tau )^3\theta _{1,1}(\tau ,z)}{\theta _{0,1}(\tau ,z)\theta _{0,1}(\tau ,0)}.$$ The result now follows from the product formulas $`\theta _{1,1}(\tau ,z)`$ $`=q^{\frac{1}{8}}(y^{\frac{1}{2}}y^{\frac{1}{2}}){\displaystyle \underset{n>0}{}}(1q^n)(1q^ny)((1q^ny^1),`$ $`\theta _{0,1}(\tau ,z)`$ $`={\displaystyle \underset{n>0}{}}(1q^n)(1q^{n\frac{1}{2}}y)((1q^{n\frac{1}{2}}y^1).`$ ## 7. Remarks and Speculations Recently there has been a lot of interest in motivic integration \[Ko\], \[D-L\], \[Lo\]. This is a method to determine the class of a variety in a ring $`\widehat{M}_k`$ which is obtained from $`K_0(V_k)`$ via localization at the class of $`[𝔸^1]`$ and a suitable completion. In many cases the result of motivic integration is an identity between the classes of smooth projective varieties or projective varieties with finite quotient singularities in $`\widehat{M}_k`$, and one would expect that the identity does indeed hold in $`K_0(V_k)`$. Using Conjecture 2.5 this conjecturally gives an isomorphism of motives and of the Chow groups with rational coefficients. Two important instances of this are the following: (1) Let $`X,Y`$ be a smooth projective birational $`k`$-varieties with $`K_X=K_Y=0`$. Then motivic integration is used to show that $`X`$ and $`Y`$ have the same Hodge numbers (\[Ko\],\[D-L\]) (that they have the same Betti numbers was shown before in \[Ba\] via $`p`$-adic integration). In fact they have the same class in $`\widehat{M}_k`$. It should be true that $`[X]=[Y]`$ in $`K_0(V_k)`$. This happens e.g. in some cases for moduli spaces of K3-surfaces: Let $`S`$ be a K3 surface with $`Pic(S)=L`$ for $`L`$ an ample divisor. Then $`[M_S^L(L,L^2/2+3)]=[S^{[L^2/2+3]}]`$. This follows from the proof of Proposition 1.9 in \[G-H\] We write $`M:=M_S^L(L,L^2/2+3)`$, $`X:=S^{[L^2/2+3]}`$. There is a diagram of birational maps $`M\stackrel{\varphi }{}N\stackrel{\psi }{}X`$ together with stratifications $`M=M_l`$, $`N=N_l`$, $`X=X_l`$ with $`N_l=\varphi ^1M_k=\psi ^1X_l`$, such that $`N_lM_l`$ and $`N_lX_l`$ are $`^{l1}`$-bundles. (2) In a similar way the results of motivic integration on the McKay correspondence \[D-L\],\[R\] make it seem likely that the following holds in the Grothendieck ring of varieties. Let $`X`$ be a smooth projective variety acted upon by a finite group $`GSl(n,)`$ such that the action preserves the canonical divisor. Let $`Y`$ be a crepant resolution of $`X/G`$. Choosing an eigenbasis, we can write $`g=diag(ϵ^{a_1},\mathrm{},ϵ^{a_n})`$, where $`ϵ`$ is a primitive $`r^{th}`$ root of unity for $`r`$ the order of $`g`$. Write $`a(g)=\frac{1}{r}a_i`$ and let $`C(g)`$ be the centralizer of $`gG`$. Then one should have in the Grothendieck ring of varieties $`[Y]=_{[g]}[X^g/C(g)][𝔸^{a(g)}]`$. Here $`[g]`$ runs through the conjugacy classes elements of $`G`$. In particular we should have $`A^i(Y)=_{[g]}A^{ia(g)}(X^g)^{C(g)}`$. Using \[Gö5\], Theorem 1.1 and the main result of \[dC-M2\] say that this is true for the resolution $`\omega _n:S^{[n]}S^{(n)}`$ of the symmetric power $`S^{(n)}=S^n/𝔖_n`$.
warning/0007/hep-th0007017.html	ar5iv	text	# Untitled Document Göteborg-ITP-98-09 hep-th/0007017 January 1998 A superspace approach to branes and supergravity\* Bengt E.W. Nilsson Institute for Theoretical Physics Chalmers University of Technology and Göteborg University S-412 96 Göteborg, Sweden Abstract: Recent developments in string and M theory rely heavily on supersymmetry suggesting that a revival of superspace techniques in ten and eleven dimensions may be advantageous. Here we discuss three topics of current interest where superspace is already playing an important role and where an improved understanding of superspace might provide additional insight into the issues involved. \*In Proceedings of the 31st International Symposium Ahrenshoop, September 2-6, 1997, Buckow, Germany. Eds H. Dorn, D. Lüst and G. Weigt (Wiley 1998). 1. Introduction It is well-known that besides the string (the $`p=1`$ brane) itself also branes of other dimensionality ($`d=p+1`$) play a central role in the non-perturbative structure of the theory. Using Dp-branes and T5-branes (having respectively vector and second rank antisymmetric tensor potentials on the world surface) non-perturbative relations between all string theories can be established as well as a connection with 11d supergravity. This hints at the existence of an underlying more profound formulation of the whole theory. Some aspects of this so called M theory are captured by M(atrix) theory which has its origin in D-particle (D0-brane) physics, but have many features in common with the 11d membrane (see e.g. the talk by B. de Wit). The three topics discussed below are related to the role superspace is playing in this context. We start by discussing one of the direct implications of formulating the fundamental theory in terms of M(atrix) theory, namely the presence of higher order corrections (e.g. $`R^n`$) in 11d supergravity (for a recent review see ). Here we will have reason to recall certain facts in 10d supergravity established in the 1980’s . Then we review some recent results concerning the realization of $`\kappa `$-symmetry for various branes focusing on the 11d membrane ($`p=2`$), the D3-brane in 10d type IIB supergravity, and the T5-brane in 11d superspace. In the latter case we present a new action which avoids some of the problematic features of the lagrangian constructed previously in . As a final topic we consider the generalization of this setup, based on bosonic world sheets embedded into target superspaces, to a situation where both target and the embedded world surface are superspaces . In the relation between non-linearities of the tensor dynamics on the world sheet and non-linearly realized (super)symmetries is analyzed. This paper also contains a detailed account of general embeddings plus some new results for T5-branes in 7d. 2. Superspace and higher order corrections in 11d One property of M theory that has been highlighted by the recent developments in M(atrix) theory is the higher order corrections in terms of e.g. the curvature tensor that must occur in the low energy supergravity lagrangian in 11d. Since 11d supersymmetry is tremendously restrictive (there is only one multiplet whose spin content does not exceed spin two, a cosmological term is not possible , etc) but also very messy to deal with in terms of component fields, it might be worth the effort to develop further the superspace techniques that were introduced in the 1970’s. The heart of the matter is the question of how to realize supersymmetry off-shell. This is fairly well understood in the case of $`N=1`$ supersymmetry in 10d , a subject that we will therefore have reason to come back to. The field content of 11d supergravity is $`e_{m}^{}{}_{}{}^{a},\psi _m^\alpha ,c_{mnp}`$ where $`m,n,..`$ are 11d vector indices and $`\alpha =1,..,32`$ is a Majorana spinor index. These fields describe 128 bosonic and 128 fermionic degrees of freedom. The field equations are obtained from the corresponding superfields and their super-Bianchi identities. To this end we combine $`e_{m}^{}{}_{}{}^{a}`$ and $`\psi _m^\alpha `$ into the supervielbein $`E_{M}^{}{}_{}{}^{A}`$ where the superworld index $`M=(m,\mu )`$ and the supertangent space index $`A=(a,\alpha )`$. Normally one introduces also a superfield $`C_{MNP}`$ with $`c_{mnp}`$ as its first component. However, it was recently clarified that this is not necessary, and we will refrain from doing so. Then the field strengths are just the supercurvature $`R_{A}^{}{}_{}{}^{B}`$ and the supertorsion $`T^A=DE^A=dE^AE^B\omega _{B}^{}{}_{}{}^{A}`$. Their Bianchi identities read $`DT^A=E^BR_{B}^{}{}_{}{}^{A}`$ and $`DR_{A}^{}{}_{}{}^{B}=0`$ but in fact it is only necessary to consider the first identity since the second one is automatically satisfied. The meaning of “solving the Bianchi identities” is as follows. By subjecting the torsion components ($`T_{\alpha \beta }^{}{}_{}{}^{c}`$, etc) to certain constraints the Bianchi identities cease to be identities and become equivalent to the field equations. One should remember however that a given set of gauge fields can often be given a variety of different kinds of dynamics, related either to different sets of constraints or to differences in the Bianchi identities. As shown recently by Howe (on-shell) 11d supergravity is obtained from the $`single`$ constraint ($`\mathrm{\Gamma }^c`$ is an 11d Dirac matrix) $$T_{\alpha \beta }^{}{}_{}{}^{c}=2i(\mathrm{\Gamma }^c)_{\alpha \beta }$$ $`(2.1)`$ which is invariant under super-Weyl rescalings. Furthermore, no off-shell formulation of this theory is known. This situation should be compared to what is known for $`N=1`$ supergravity in 10d. When the Bianchi identities are solved on-shell one finds that all physical fields appear at different $`\theta `$ levels in a scalar superfield denoted $`\mathrm{\Phi }(Z)`$, where $`Z=(x,\theta )`$. Note that 10d supergravity contains a scalar and a spinor, apart from $`e_{m}^{}{}_{}{}^{a},\psi _m^\alpha ,B_{mn}`$. In contrast to 11d, one must in this case constrain several torsion components. This theory can be coupled to superYang-Mills using the same torsion constraints and the Bianchi identity $`DH=trF^2`$ for the three-form $`H=dB`$. In fact, as proven in even the $`R^2`$ term needed for anomaly cancellation can be dealt with without altering the constraints. However, also $`R^4`$ and higher terms are present in the field theory of the low energy 10d superstring. In general such terms cannot be incorporated into the superspace Bianchi identities if the on-shell constraints are used. Fortunately, in this case it is known how to proceed since the off-shell field content is known. It consists of a superconformal gravity multiplet and an unconstrained scalar auxiliary superfield $`w`$ . To account for this new superfield the constraints must be modified to $$T_{\alpha \beta }^{}{}_{}{}^{c}=2i(\mathrm{\Gamma }^c)_{\alpha \beta }+2i(\mathrm{\Gamma }^{c_1\mathrm{}c_5})_{\alpha \beta }X_{}^{c}{}_{c_1\mathrm{}c_5}{}^{}$$ $`(2.2)`$ where $`X`$ is in the representation $`1050^+`$ of $`so(1,9)`$ appearing at level $`\theta ^4`$ in $`w`$. All higher order corrections, like $`R^4`$, that can occur must be compatible with supersymmetry and fit somewhere in the solution of the Bianchi identities that follow from the off-shell constraints above. E.g. the $`R^4`$ term related by supersymmetry (see ) to the anomaly term $`BX_8`$ can be added as follows : $$S^{D=10,N=1}=\frac{1}{(\kappa _{10})^2}d^{10}xd^{16}\theta E\mathrm{\Phi }(w+c)$$ $`(2.3)`$ where $`c`$ is a constant proportional to $`\zeta (3)`$, and where the $`w`$ term is the kinetic one and the $`c`$ term is the supersymmetrization of $`R^4`$. $`E`$ is the superspace measure. Turning to 11d the situation changes dramatically since it is not known how to solve the Bianchi identities off-shell or how to write down an off-shell action in components. This makes it very hard to address questions in 11d supergravity concerning the higher order corrections. We will here introduce the equivalent of $`w`$ in 10d into the 11d supertorsion by means of the relaxed constraint $$T_{\alpha \beta }^{}{}_{}{}^{c}=2i\mathrm{\Gamma }_{}^{c}{}_{\alpha \beta }{}^{}+2i\mathrm{\Gamma }_{}^{d_1d_2}{}_{\alpha \beta }{}^{}X_{}^{c}{}_{d_1d_2}{}^{}+2i\mathrm{\Gamma }_{}^{d_1\mathrm{}d_5}{}_{\alpha \beta }{}^{}X_{}^{c}{}_{d_1\mathrm{}d_5}{}^{}$$ $`(2.4)`$ where the tensors in the last two terms are in the representations 429 and 4290 which appear at level $`\theta ^4`$ in an unconstrained 11d scalar superfield. A preliminary analyzes of the Bianchi identities indicates an “off-shell situation” where new terms appear in the torsion which could account for the higher order corrections. In particular terms generated by anomalies and the presence of branes should be investigated. Note that the T5-brane produces in the supersymmetry algebra a five-form central charge, a fact that should be compared to the extra terms in the torsion. Further studies will hopefully tell if these techniques can be utilized in the endeavour to extract an 11d supergravity theory from M(atrix) theory or perhaps directly from the 11d branes. 3. $`\kappa `$ symmetric branes as bosonic surfaces The relevant branes in 11d are the membrane and the T5-brane, while in 10d type II theories there are also Dp-branes intermediate between p-branes and T-branes. As we will see below, for branes with vector or tensor fields propagating on the world surface $`\kappa `$-symmetry is technically more complicated than for ordinary p-branes. Let us as an example of an ordinary p-brane consider the 11d membrane . After the elimination (see ) of the independent world sheet metric by means of its algebraic field equation the action reads: $$S_3=d^3\xi [\sqrt{detg}\epsilon ^{ijk}B_{ijk}]$$ $`(3.1)`$ where $`g_{ij}`$ is the pull-back of the target space metric, i.e. $`g_{ij}=\mathrm{\Pi }_i^a\mathrm{\Pi }_j^b\eta _{ab}`$, and the $`\xi ^i`$’s are three bosonic coordinates on the world sheet. The background superfields $`E_{M}^{}{}_{}{}^{A}`$ and $`B_{MNP}`$ of 11d supergravity enter via the pull-backs $`\mathrm{\Pi }_i^A=_iZ^ME_{M}^{}{}_{}{}^{A}`$ and $`B_{ijk}=\mathrm{\Pi }_i^A\mathrm{\Pi }_j^B\mathrm{\Pi }_k^CB_{CBA}`$. In order for this action to be supersymmetric the number of bosonic (here 11-3=8) and fermionic ($`32\times \frac{1}{2}`$) world sheet on-shell degrees of freedom must match. This requires the presence of the local fermionic $`\kappa `$-symmetry giving another factor of $`\frac{1}{2}`$ in the fermionic count. Its existence relies on the possibility to construct a projection operator $`\frac{1}{2}(1+\mathrm{\Gamma })`$ with $`\mathrm{\Gamma }^2=1`$. Here $`\mathrm{\Gamma }=\frac{1}{6\sqrt{g}}ϵ^{ijk}\mathrm{\Gamma }_{ijk}`$ where $`\mathrm{\Gamma }_{ijk}`$ is the pull-back of $`\mathrm{\Gamma }_{abc}=\mathrm{\Gamma }_{[a}\mathrm{\Gamma }_b\mathrm{\Gamma }_{c]}`$. This structure can be found also in the case of Dp-branes , but is now more involved due to the presence of the field strength $`F_{ij}`$. E.g. the D3-brane in the 10d type IIB theory has an action that reads $$S_{D3}=d^4\xi \sqrt{det(g+e^{\frac{\varphi }{2}}F)}+e^FC$$ $`(3.2)`$ In this case there are, apart from the dilaton $`\varphi `$, two kinds of background potentials $`B`$ and $`C`$, coming from the NS-NS and the R-R sector, respectively. In $`S_{D3}`$, $`F_{ij}=F_{ij}B_{ij}`$ where $`B_{ij}`$ is the pullback of $`B_{MN}`$ in the 10d IIB target space theory. The last term in the action is constructed as a formal sum of forms of different rank and the integral is supposed to pick up only the four-form in this case. That the action is $`\kappa `$-symmetric can then be shown using the $`\mathrm{\Gamma }`$ matrix $`(\mathrm{\Gamma }^2=1)`$ $$\mathrm{\Gamma }=\frac{ϵ^{ijkl}}{\sqrt{det(g+F)}}(\frac{1}{24}\mathrm{\Gamma }_{ijkl}I\frac{1}{4}F_{ij}\mathrm{\Gamma }_{kl}J+\frac{1}{8}F_{ij}F_{kl})$$ $`(3.3)`$ The $`SL(2;Z)`$ symmetry of the IIB theory mixes the $`B`$ and $`C`$ potentials and indeed a more symmetric version of the action exists . In that version all background potentials have associated world sheet field strengths to which they couple as $`FB`$. The third case to be discussed here is the T5-brane in 11d. This brane has an additional complication in that the three-form field strength is self-dual. Finding a covariant action for such a field has been a long-standing problem. A first attempt at a solution was given recently by Bandos et al in ) and involves an auxiliary scalar field $`a`$ entering the lagrangian through a factor $`\frac{1}{(a)^2}`$. Objections against using such an action in quantum calculations have been formulated (see also ). Another action that does not make use of such a scalar was subsequently presented in . Although the problems associated with the scalar are gone, the other objections in probably remain. Nevertheless, since this action exhibits some new features it might be of some interest. The action is ($`F`$ is an independent six-form field strength) $$S_{T5}=d^6\xi \sqrt{g}\lambda (1+\frac{1}{12}F_{ijk}F^{ijk}\frac{1}{24}k_{ij}k^{ij}+\frac{1}{72}(trk)^2(F)^2)$$ $`(3.4)`$ where $`k_{ij}=\frac{1}{2}F_{i}^{}{}_{}{}^{jk}F_{jkl}`$. Note that, as explained in , the self-duality relation $$(F)F_{ijk}=F_{ijk}\frac{1}{2}k_{[i}^{}{}_{}{}^{l}F_{jk]l}+\frac{1}{6}(trk)F_{ijk}$$ $`(3.5)`$ does not arise as a field equation but as a result of demanding $`\kappa `$-symmetry, and must not be inserted into the action. 4. Superworld sheets in target superspace It is possible to reformulate the brane dynamics, following from actions of the kinds described in the previous section, in terms of world sheet superfields. This can be done by embedding superworld sheets into target superspace. Besides making the world sheet supersymmetry manifest this has the further advantages of explaining the origin of $`\kappa `$-symmetry and the projection matrix as well as of providing the connection between superembeddings on one hand, and Goldstone fermions and non-linearly realized supersymmetries on the supersheets on the other hand . In Bagger and Galperin showed how the 4d Born-Infeld action for an abelian gauge field can be obtained by embedding a (4d,N=1) superspace into a (4d,N=2) one. The broken fermionic translations turn into non-linear supersymmetry transformations on the Goldstone fermions arising from half of the fermionic coordinates that are turned into dependent variables. The non-linearities of the Born-Infeld action are then seen to be a consequence of demanding consistency with the extra non-linear supersymmetries. In this programme was taken over to the T5-brane in 7d superspace. This led to an equation for one of the supertorsion components whose rather complicated solution indicates that this way of analyzing this system is not the most efficient one. Fortunately once it is realized that this form of the torsion equation can also be obtained in the much more general formalism known as the embedding formalism developed in these problems can be circumvented. The central equation is the torsion pullback equation (8,9,10) $$D_AE_{B}^{}{}_{}{}^{\underset{¯}{C}}(1)^{AB}D_BE_{A}^{}{}_{}{}^{\underset{¯}{C}}+T_{AB}^{}{}_{}{}^{C}E_{C}^{}{}_{}{}^{\underset{¯}{C}}=(1)^{A(A+\underset{¯}{B})}E_{B}^{}{}_{}{}^{\underset{¯}{B}}E_{A}^{}{}_{}{}^{\underset{¯}{A}}T_{\underset{¯}{A}\underset{¯}{B}}^{}{}_{}{}^{\underset{¯}{C}}$$ $`(4.1)`$ where underlined indices refer to target superspace while the other indices are connected with the superworld sheet. As shown in inserting constraints on the torsion components turns this equation into the equations of motion for the world sheet fields. In particular, the highly non-linear dynamics of the T5-branes in 11d and in 7d can be obtained this way. Acknowledgement: I wish to thank my coauthors on refs. for very nice and fruitful collaborations. References M.B. Green, hep-th/9712195. P. Howe, H. Nicolai and A. Van Proeyen, Phys. Lett. 112B (1982) 446. B.E.W. Nilsson, Phys. Lett. 175B (1986) 319. P.S. Howe and A. 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warning/0007/math0007136.html	ar5iv	text	# 1. Introduction ## 1. Introduction A tiling of a checkerboard with dominoes is a way of putting dominoes on the board so that no square of the board is uncovered and no two dominoes overlap. Given a local pattern (see figure 1 for examples), a location in the board, and the shape and size of the board, how many tilings of the board have the given pattern at the given location? (Alternatively, we can substitute “bond” for “domino” and “particle” for “square”, and ask for the probability of local patterns in a system of particles each of which bonds with exactly one of its neighbors.) Suppose that the squares of the board are very small compared to the board itself. For some board shapes, the probability of finding a pattern at a given location will be the same for almost all locations. This is the case for the square board. (See figures 4 to 5, where tiles are colored according to their direction and parity for the sake of clarity; see figure 6 for the coloring scheme.) There are some boards, however, for which the probability does depend on the location. Consider, for example, the Aztec diamond, that is, the board whose boundary is a square tilted $`45`$ degrees (figures 8 and 8). In random tilings of the Aztec diamond, we usually find brick-wall patterns outside the inscribed circles, and more complicated behavior inside the circle. (See figures 10 to 13.) The probabilities of local patterns in a rectangular board were computed recently . Until now, there was no other board for which the probabilities of all local patterns were known. Many experiments and some important partial results had shown that, as already stated, the probabilities of patterns in the Aztec diamond depend on location. This qualitative difference between the Aztec diamond and the rectangular board made the former as worthy of analysis as the latter. The main result of this work is an expression for the probability of any local pattern in a random tiling of the Aztec diamond. The expression is a determinant of size proportional to the number of squares in the pattern, just like Kenyon’s expression for the probabilities in the rectangural board, ###### Main Result 1. The probability of a pattern covering white squares $`v_1,v_2,\mathrm{}v_k`$ and black squares $`w_1,w_2,\mathrm{}w_k`$ of an Aztec diamond of order $`n`$ is equal to the absolute value of $$\left\|c(v_i,w_j)\right\|_{i,j=1,2,\mathrm{}k}.$$ The coupling function $`c(v,w)`$ at white square $`v`$ and black square $`w`$ is $$2^n\underset{j=0}{\overset{x_i1}{}}\mathrm{kr}(j,n,y_i1)\mathrm{kr}(y_i1,n1,n(j+x_ix_i))$$ for $`x_i>x_i`$ and $$2^n\underset{j=x_i}{\overset{n}{}}\mathrm{kr}(j,n,y_i1)\mathrm{kr}(y_i1,n1,n(j+x_ix_i))$$ for $`x_ix_i`$, where $`(x_i,y_i)`$ and $`(x_i,y_i)`$ are the coordinates of $`v`$ and $`w`$, respectively, in the coordinate system in figure 17, and the Krawtchouk polynomial $`\mathrm{kr}(a,b,c)`$ is the coefficient of $`x^a`$ in $`(1x)^c(1+x)^{bc}`$. Our line of attack is as follows. 1. Reduce the problem of finding probabilities of patterns to an enumerative problem; 2. Reduce the enumerative problem to a simpler one involving Aztec diamonds with two holes rather than arbitrary even-area holes; 3. Compute the weighted number of tilings of an Aztec diamond with two holes. The first two steps involve known techniques, and were already considered to be a plausible strategy by other researchers. The third step is new. Henry Cohn is currently analyzing the case of the board with infinitely small squares by approximating the sum of Krawtchouk polynomials in our main result as an integral for $`n\mathrm{}`$. His results will be presented in a later, joint version of this paper. ## 2. The Kasteleyn Matrix What is the probability of finding a pattern in a tiling of a given board? It is equal, by definition, to the number of tilings of the board with the given pattern, divided by the total number of tilings of the board. Clearly, the number of tilings of the board with the given pattern depends only on the squares occupied by the pattern, and not by how they are tiled by the pattern: there is a one-to-one correspondence between tilings with the given pattern, and tilings of the board with the squares covered by the pattern removed. (See figure 2.) Thus, what we want to know is the number of tilings of the board with the squares covered by a pattern removed, divided by the number of tilings of the board. Kasteleyn showed that the number of tilings of a board can be expressed as the absolute value of a determinant with half as many rows as there are squares in the board. Thus, in a sense, our problem is already solved: since determinants can be computed in time polynomial on the number of rows, we can compute the probability of any pattern in a board in time polynomial on the size of the board. Unfortunately, there are two disadvantages to this approach. The first one is that, if the Aztec diamond in question has side $`n`$, we will have to compute a determinant of side $`\frac{n^2}{2}`$, and that demands a considerable amount of time and space. Even more seriously, a sequence of determinants of varying size is very hard to analyze asymptotically. We would like to analyze such a sequence in order to determine the probability of a pattern near a point for Aztec diamonds of very high order, that is, of very fine “grain” (see figure 13). Thus, a Kasteleyn determinant is not good enough. A determinant whose size depended only on the size of the pattern, and not on the size of the board, would be much easier to manipulate. In this section, we will prove that the number of tilings of a board equals the absolute value of the determinant of its Kasteleyn matrix, and then show how this implies that the probability of a pattern in a random tiling of a board is equal to a minor of the inverse of the Kasteleyn matrix of the board. This minor has side proportional to the number of squares in the pattern. Finally, we will show that the problem of finding an entry in the inverse of a Kasteleyn matrix can be reduced to an enumerative problem concerning an Aztec diamond with one black hole and one white hole. We will henceforth refer, not to boards, but to their dual graphs, which can be seen as having a vertex at the center of every square and an edge perpendicular to every edge between two squares. For example, the Aztec diamond will mean for us the object in figure 8 and not the object in figure 8. This convention will simplify graph-theoretical arguments considerably. The results in this section are in part a modern formulation of Kasteleyn’s work, and in part a codification of local folklore, as crafted and passed down by R. Kenyon, J. Propp, D. Wilson and others. ### 2.1. The Kasteleyn-Wilson matrix Consider a finite subgraph G of the infinite square lattice with as many white as black vertices, where the infinite square lattice is colored as a checkerboard. Let $`(v_1,\mathrm{},v_n)`$ be its white vertices and $`(w_1,\mathrm{},w_n)`$ its black vertices. (Any ordering from $`1`$ to $`n`$ can be chosen.) The Kasteleyn-Wilson matrix $$K((v_1,\mathrm{},v_n),(w_1,\mathrm{},w_n))$$ (or $`K(G)`$, by abuse of language) is defined to be $`\left\|a_{i,j}\right\|_1^n`$, where $`a_{i,j}`$ is * $`0`$, if $`\{v_i,w_j\}`$ is not an edge of G; * $`1`$, if $`\{v_i,w_j\}`$ is a horizontal edge of G; * $`(1)^k`$, if $`\{v_i,w_j\}`$ is a vertical edge going from row $`l`$ to row $`l+1`$ and there are $`k`$ vertices in row $`l`$ to the left of the edge. We will show that the number of perfect matchings of $`G`$ is equal to the absolute value of the determinant of $`K(G)`$. As could be exprected, the Kasteleyn-Wilson matrix is only one of several Kasteleyn matrices $`K(G)`$, that is, determinants whose absolute values are equal to the perfect matchings of $`G`$. We will prove the enumerative property of the Kasteleyn-Wilson matrix because it is true for any subgraph $`G`$ of the infinite square lattice, and not only for the Aztec diamond. Fortunately, the proof can be applied to other Kasteleyn matrices with minimal cases. In fact, in the last subsection of this section, and in following sections, we will use the following convention, which is more convenient for our purposes than Wilson’s: * $`0`$, if $`\{v_i,w_j\}`$ is not an edge of G; * $`1`$, if $`\{v_i,w_j\}`$ is a horizontal edge of G; * $`(1)^k`$, if $`\{v_i,w_j\}`$ is a vertical edge going from row $`l`$ to row $`l+1`$ and there are $`k`$ vertical edges from row $`l`$ to row $`l+1`$ to the left of the edge. This convention is valid for the Aztec diamond and for any other subgraph $`G`$ of the infinite square lattice such that any lattice vertex inside any loop in $`G`$ is also a vertex of $`G`$. Thus, only the material in the following two subsections is valid for any subgraph of the infinite square lattice. The rest of this work is specific to the Aztec diamond, although the techniques used are applicable to other boards. ### 2.2. Why does $`K`$ give us the number of perfect matchings? We want to show that the number of perfect matchings of a subgraph $`G`$ of the infinite square lattice is equal to the absolute value of the determinant of $`K(G)`$. We can express $`det(K)`$ as $$\underset{\pi P(\{1,2,\mathrm{},n\})}{}\mathrm{sgn}(\pi )\underset{i=1}{\overset{n}{}}a_{i,\pi (i)},$$ (1) where $`P(\{1,2,\mathrm{},n\})`$ is the set of all permutations of $`\{1,2,\mathrm{}n\}`$. Define a map $`f`$ from the set of all perfect matchings of $`G`$ to $`P(\{1,2,\mathrm{},n\})`$ as follows. Any perfect matching can be expressed in the form $$\{\{v_1,w_{k(1)}\},\{v_2,w_{k(2)}\},\mathrm{},\{v_n,w_{k(n)}\}\},$$ (2) where $`k`$ is a map from $`\{1,2,\mathrm{}n\}`$ to itself. The map $`f`$ takes $`\{\{v_1,w_{k(1)}\},\mathrm{},\{v_n,w_{k(n)}\}\}`$ to $`k`$. It is clear that $`k`$ is a permutation; otherwise there would be unpaired vertices, as well as vertices belonging to more than one pair. Thus, $`f`$ is well defined. Moreover, 1. $`f`$ is injective: If two matchings had the same map $`k`$, they would be the same matching. 2. every element $`k`$ of the image of $`f`$ satisfies $`(_{i=1}^na_{i,k(i)})0`$: if $`v_i,w_{k(i)}`$ is an edge, then $`a_{i,k(i)}`$ must be non-zero. 3. if a permutation $`\pi `$ of $`1,2,\mathrm{},n`$ satisfies $`_{i=1}^na_{i,\pi (i)}0`$, then, for every $`1in`$, $`(v_i,w_{\pi (i)})`$ is a valid edge. Moreover, for $`i_1i_2`$, $`\pi (i_1)\pi (i_2)`$, and thus $`(v_{i_1},w_{\pi (i_1)})`$ and $`(v_{i_2},w_{\pi (i_2)})`$ do not have any vertices in common. Therefore $$\{\{v_1,w_{k(1)}\},\{v_2,w_{k(2)}\},\mathrm{},\{v_n,w_{k(n)}\}\}$$ (3) is a perfect matching. Hence $`f`$ is a one-to-one and onto map from the set of all perfect matchings of $`G`$ to the set of all permutations $`k`$ of $`1,2,\mathrm{},n`$ satisfying $`_{i=1}^na_{i,k(i)}0`$. Therefore there are as many perfect matchings as there are non-zero terms in $$det(K)=\underset{\pi P(\{1,2,\mathrm{},n\})}{}\mathrm{sgn}(\pi )\underset{i=1}{\overset{n}{}}a_{i,\pi (i)}.$$ (4) Every non-zero term is equal to either $`1`$ or $`1`$. To prove that the absolute value of $`det(K)`$ equals the number of perfect matchings, we have to show that all non-zero terms have the same sign. Let $`M=\{\{v_i,w_{\pi (i)}\}\}_{i=1}^n`$ and $`M=\{\{v_i,w_{\pi (i)}\}\}_{i=1}^n`$ be two perfect matchings of $`G`$. By the definition of perfect matching, every vertex of $`G`$ is in one edge of $`M`$ and in one edge of $`M`$. It follows that every vertex of $`G`$ is either in two edges or in no edges of $`(MM)(MM)`$. Therefore $`(MM)(MM)`$ consists entirely of loops, that is, it is a collection of disjoint sets of the form $$\{\{v_{i_1},w_{j_1}\},\{w_{j_1},v_{i_2}\},\{v_{i_2},w_{j_2}\},\mathrm{}\{v_{i_m},w_{j_m}\},\{w_{j_m},v_{i_1}\}\}.$$ (5) (See figure 14.) If a vertex is in two edges of $`(MM)(MM)`$, one of these two edges must be in $`M`$, and the other one in $`M`$. We can assume without loss of generality that $`\{v_{i_1},w_{j_1}\}`$ is in $`M`$, and hence $`\{w_{j_1},v_{i_2}\}`$ is in $`M`$, $`\{v_{i_2},w_{j_2}\}`$ is in $`M`$, and so on. Then, on one hand, $`j_l=\pi (i_l)`$ for $`1lm`$, and, on the other hand, $`j_l=\pi (i_{l+1})`$ for $`1lm1`$, $`j_m=\pi (i_1)`$. Hence $`i_2=((\pi )^1\pi )(i_1)`$, $`i_3=((\pi )^1\pi )(i_2)`$,…$`i_1=((\pi )^1\pi )(i_m))`$. Thus every loop in $`(MM)(MM)`$ induces a cycle in $`(\pi )^1\pi `$. It is easy to see that, conversely, for every cycle in $`(\pi )^1\pi `$ there is a loop in $`(MM)(MM)`$ that induces it. Because the sign of a permutation is equal to the product over all its cycles of $`(1)`$ to the power of the length of the cycle minus one, and because a cycle has length equal to half the number of edges of the loop inducing it, we have $$\mathrm{sgn}((\pi )^1\pi )=\underset{\mathrm{}L}{}(1)^{\mathrm{len}(\mathrm{})/21},$$ (6) where $`L`$ is the set of all loops of $`(MM)(MM)`$ and $`\mathrm{len}(\mathrm{})`$ is the number of edges in loop $`\mathrm{}`$. From this equation, from $$\mathrm{sgn}((\pi )^1\pi )=\frac{\mathrm{sgn}(\pi )}{\mathrm{sgn}(\pi )},$$ (7) and from the fact that all $`a_{i,j}`$ are $`1`$ or $`(1)`$, it follows that the result we want to prove in this section, namely, $$\mathrm{sgn}(\pi )\underset{i=1}{\overset{n}{}}a_{i,\pi (i)}=\mathrm{sgn}(\pi )\underset{i=1}{\overset{n}{}}a_{i,\pi (i)},$$ (8) is equivalent to $$\underset{\mathrm{}L}{}(1)^{\mathrm{len}(\mathrm{})/21}=\underset{i=1}{\overset{n}{}}a_{i,\pi (i)}\underset{i=1}{\overset{n}{}}a_{i,\pi (i)}.$$ (9) Now, $$\underset{i=1}{\overset{n}{}}a_{i,\pi (i)}\underset{i=1}{\overset{n}{}}a_{i,\pi (i)}=\underset{\mathrm{}L}{}(\underset{iI(\mathrm{})}{}a_{i,\pi (i)}\underset{iI(\mathrm{})}{}a_{i,\pi (i)})$$ (10) , where $`I(\mathrm{})`$ is the set of indices $`i`$ of all white vertices in loop $`\mathrm{}`$. Therefore it is enough for us to prove that $$(1)^{\mathrm{len}(\mathrm{})/21}=\underset{iI(\mathrm{})}{}a_{i,\pi (i)}\underset{iI(\mathrm{})}{}a_{i,\pi (i)}$$ (11) for every loop $`\mathrm{}`$ in $`(MM)(MM)`$. $`_{iI(\mathrm{})}a_{i,\pi (i)}_{iI(\mathrm{})}a_{i,\pi (i)}`$ is the product $`a_{i,j}`$ over all $`1i,jn`$ such that $`\{v_i,w_j\}`$ is an edge of loop $`\mathrm{}`$. Since $`a_{i,j}=1`$ for $`\{v_i,w_j\}`$ horizontal, we can restrict the product to $`1i,jn`$ such that $`\{v_i,w_j\}`$ is a vertical edge of loop $`\mathrm{}`$. What is, specifically, the product of $`a_{i,j}`$ over all $`1i,jn`$ such that $`\{v_i,w_j\}`$ is a vertical edge of loop $`\mathrm{}`$ having a vertex on row $`y_0`$ and another in row $`y_0+1`$, where $`y_0`$ is a given? Let $`\{(x_1,y_0),(x_1,y_0+1)\},\{(x_2,y_0),(x_2,y_0+1)\},\mathrm{}\{(x_m,y_0),(x_m,y_0+1)\}`$ be all such edges. (We are referring vertices by their Cartesian coordinates.) Since the horizontal line $`\{(t,y+\frac{1}{2}):t(\mathrm{},\mathrm{})\}`$ crosses the loop an even number of times, $`m`$ must be even. Let $`m=2m_0`$. Since all $`a_{i,j}`$ are $`1`$ or $`(1)`$, we have $$\underset{j=1}{\overset{m}{}}a_{J_0(x_j,y_0),J_1(x_j,y_0+1)}=\underset{j=1}{\overset{m_0}{}}\frac{a_{J_0(x_{2j},y_0),J_1(x_{2j},y_0+1)}}{a_{J_0(x_{2j1},y_0),J_1(x_{2j1},y_0+1)}},$$ (12) where $`J_0(z,w)`$ is the index of the white vertex of coordinates $`(z,w)`$, and $`J_1(z,w)`$ is the index of the black vertex of coordinates $`(z,w)`$. Since $`a_{J_0(x_i,y_0),J_1(x_i,y_0+1)}`$ is equal to the number of vertices of $`G`$ on row $`y_0`$ and to the left of $`x_i`$, and $`a_{J_0(x_{i+1},y_0),J_1(x_{i+1},y_0+1)}`$ is equal to the number of vertices of $`G`$ on row $`y_0`$ and to the left of $`x_{i+1}`$, $$\frac{a_{J_0(x_{i+1},y_0),J_1(x_{i+1},y_0+1)}}{a_{J_0(x_i,y_0),J_1(x_i,y_0+1)}}$$ (13) is equal to the number of vertices of $`G`$ on row $`y_0`$ and to the left of $`x_{i+1}`$ but not of $`x_i`$. This is the same as the number of vertices of the infinite square grid on row $`y_0`$ and to the left of $`x_{i+1}`$ but not of $`x_i`$, minus the number of vertices in the grid but not in $`G`$, on row $`y_0`$ and to the left of $`x_{i+1}`$ but not of $`x_i`$. The number of vertices of the grid to the left of $`x_{i+1}`$ but not of $`x_i`$ is equal to the number of squares of the grid contained between the edges $`((x_i,y),(x_i,y+1))`$ and $`((x_{i+1},y),(x_{i+1},y+1))`$. It is clear that, when $`i`$ is even, the squares between the edges $`((x_i,y),(x_i,y+1))`$ and $`((x_{i+1},y),(x_{i+1},y+1))`$ are in the interior of the loop, as are the vertices on row $`y`$ of the grid which do not belong to $`G`$ and which are the left of $`x_{i+1}`$ but not of $`x_i`$. (Since any vertex on the loop, and, specifically, $`x_i`$ and $`x_{i+1}`$, must be in $`G`$, we can hencefort refer to these vertices as “the vertices on row $`y`$ of the grid which do not belong to $`G`$ and which lie between $`x_i`$ and $`x_{i+1}`$”.) Conversely, every square between rows $`y`$ and $`y+1`$ and in the interior of the loop lies between edges $`((x_i,y),(x_i,y+1))`$ and $`((x_{i+1},y),(x_{i+1},y+1))`$ for some odd $`i`$, and, moreover, every vertex on row $`y`$ of the grid, not in $`G`$, and in the interior of the loop lies between $`x_i`$ and $`x_{i+1}`$ for some odd $`i`$. Hence the number of squares between edges $`((x_i,y),(x_i,y+1))`$ and $`((x_{i+1},y),(x_{i+1},y+1))`$ for $`i`$ odd equals the number of squares which are both between rows $`y`$ and $`y+1`$ and in the interior of the loop, and, furthermore, the number of vertices not in $`G`$ lying on row $`y`$ between $`x_i`$ and $`x_{i+1}`$ for $`i`$ odd equals the number of vertices not in $`G`$ lying on row $`y`$ and in the interior of the loop. Thus $$\underset{j=1}{\overset{m}{}}a_{J_0(x_j,y_0),J_1(x_j,y_0+1)}=\underset{j=1}{\overset{m_0}{}}\frac{a_{J_0(x_{2j},y_0),J_1(x_{2j},y_0+1)}}{a_{J_0(x_{2j1},y_0),J_1(x_{2j1},y_0+1)}}$$ (14) equals $`(1)`$ to the power of the number of squares which are both in the interior of the loop and between rows $`y`$ and $`y+1`$ minus the number of such vertices in the interior on row $`y`$. Hence, the product of this expression over all rows, that is, $$\frac{_{iI}a_{i,\pi (i)}}{_{iI}a_{i,\pi (i)}}$$ (15) equals $`(1)`$ to the power of the number of squares in the interior of the loop minus the number of vertices lying on row $`y`$ and in the interior, but not in $`G`$. Therefore $$(1)^{\frac{\mathrm{len}(\mathrm{})}{2}1}\frac{_{iI(\mathrm{})}a_{i,\pi (i)}}{_{iI(\mathrm{})}a_{i,\pi (i)}}$$ (16) is equal to $`(1)`$ to the power of the length of the loop, divided by 2, minus $`1`$, plus the number of squares inside the loop, minus the number of vertices in the grid and inside the loop, but not in $`G`$. Pick’s theorem states that, for polygons whose vertices belong to a square grid, $`A=I+B/21,`$ where $`A`$ is the area enclosed by the polygon, $`I`$ is the number of grid points inside the polygon, and $`B`$ is the number of grid points on the boundary of the polygon. Hence our expression is equal to $`(1)`$ to the power of the number of vertices inside the loop, minus the number of vertices in the grid and inside the loop, but not in $`G`$. This is the same as the number of vertices of $`G`$ inside the loop. Now, the vertices inside the loop are matched only among themselves, not with the vertices outside the loop, in either $`M`$ or $`M`$. For a set of vertices to be matched only among themselves, there must be an even number of vertices in the set. Therefore the number of vertices of $`G`$ inside the loop must be even, and $`(1)`$ raised to the power of this number must be $`1`$. Hence $$(1)^{\frac{\mathrm{len}(\mathrm{})}{2}1}\frac{_{iI(\mathrm{})}a_{i,\pi (i)}}{_{iI(\mathrm{})}a_{i,\pi (i)}}=1$$ (17) as we desired to prove. We conclude, by taking the product over all loops $`\mathrm{}L`$, that $$\frac{\mathrm{sgn}(\pi )_{i=1}^na_{i,\pi (i)}}{\mathrm{sgn}(\pi )_{i=1}^na_{i,\pi (i)}}=1.$$ (18) Therefore the terms of $$det(K)=\underset{\pi P(\{1,2,\mathrm{},n\})}{}\mathrm{sgn}(\pi )\underset{i=1}{\overset{n}{}}a_{i,\pi (i)}$$ (19) all have the same sign. It follows that the number of perfect matchings, that is, of permutations $`\pi `$ for which $`_{i=1}^na_{i,\pi (i)}`$ is non-zero, is equal to the absolute value of $`det(K)`$. ### 2.3. Why does $`K^1`$ give us the probabilities of patterns? Let us have a graph $`G`$ and a subgraph $`H`$ whose every vertex has degree $`1`$. There is a one-to-one correspondence between perfect matchings of $`G`$ having $`H`$ as a subgraph and perfect matchings of $`GH`$. (By $`GH`$ we mean the graph $`(V,E)`$, where $`V`$ is the set of all vertices of $`G`$ not in $`H`$ and $`E`$ is the set of all edges of $`G`$ between two vertices of $`G`$ not in $`H`$.) Therefore the probability that a random perfect matching of $`G`$ have the set of edges of $`H`$ as a subset is equal to the number of perfect matchings of $`GH`$ divided by the number of perfect matchings of $`G`$. What is, then, the number of perfect matchings of $`GH`$? One way to compute it is to construct a Kasteleyn-Wilson matrix for $`GH`$. Another way, which will soon prove its virtues, is to take the minor of the Kasteleyn-Wilson matrix for $`G`$ resulting from the deletion of the rows and columns corresponding to $`H`$. Certainly, this minor is not the same as the Kasteleyn-Wilson matrix for $`GH`$; the signs of the matrix entries are different. Nevertheless, the absolute value of the determinant of the minor is equal to the number of perfect matchings of $`GH`$. To prove this, we need to do the same as in the previous section, namely, show that all non-zero terms of the expression of the determinant as a sum over permutations have the same sign. If we proceed in the same way as before, we arrive at a point where the only difference is that we have $`(1)`$ to the power of the number of vertices of $`G`$ inside a loop instead of the number of vertices of $`GH`$ inside the same loop. This difference is no difference if there is an even number of vertices of $`H`$ inside the loop. Since the vertices of $`H`$ inside the loop cannot be connected with the vertices of $`H`$ outside the loop, the vertices of $`H`$ inside the loop are paired among themselves. Thus we have, as we wanted, that there is an even number of vertices of $`H`$ inside the loop. This is enough for us to show that the terms in the expression of the minor of $`K(G)`$ as a sum over permutations do not cancel. Therefore the number of perfect matchings of $`GH`$ is equal to the absolute value of the determinant of the minor of $`K(G)`$ lacking the rows and columns corresponding to the vertices of $`H`$. If we assign the same probability to every perfect matching of $`G`$, the probability that a random perfect matching will have $`H`$ as a subgraph will be equal to the number of perfect matchings of $`GH`$ divided by the number of perfect matchings of $`G`$, that is, to the absolute value of the determinant of the aforementioned minor of $`K(G)`$ divided by the absolute value of the determinant of $`K(G)`$. By Jacobi’s rule, a corollary of Cramer’s rule, this is equal to the absolute value of the determinant of the minor of $`(K(G)^1)^T`$ consisting of those rows and columns omitted from the minor in the numerator, that is, of rows $`1b_1<b_2<\mathrm{}b_mn`$ and columns $`1c_1<c_2<\mathrm{}c_mn`$, where $`v_{b_1},v_{b_2},\mathrm{}v_{b_m}`$ and $`w_{c_1},w_{c_2},\mathrm{}w_{c_m}`$ are the vertices of $`H`$. This is a clear improvement over the expression $`\frac{K(GH)}{K(G)}`$. Instead of dealing with determinants of the size of $`G`$, we deal with a determinant of the size of $`H`$. As we explained in the introduction, we are interested in finding the probability of small subgraphs, or “local patterns”, in a large graph $`G`$. We want to find what happens when we have an infinite sequence of $`G`$’s whose number of vertices goes to infinity. Now it is enough for us to examine a fixed number of entries in each $`(K(G_i)^1)^T`$, and determine their asymptotic behavior as $`i\mathrm{}.`$ ### 2.4. How can we find the entries of $`K^1`$ for an Aztec diamond by counting perfect matchings? Suppose that we have to compute the determinant of the minor of $`K(G)^1`$ consisting of rows $`1b_1<b_2<\mathrm{}<b_mn`$ and columns $`1c_1<c_2<\mathrm{}<c_mn`$. To do so, we have to compute the entry $`((K(G)^1)^T)_{b_i,c_j}`$ for $`1i,jm`$. In other words, in order to compute the probability of any local pattern, it suffices to be able to compute an arbitrary entry of the inverse Kasteleyn matrix. If we have a sequence of graphs $`\{G_k\}_{k=1}^{\mathrm{}}`$ (such as, for example, Aztec diamonds of higher and higher order) we will know the asymptotic behavior of the probabilities of local patterns if we know the asymptotic behavior of the entries in the sequence of matrices $`\{K(G_k)^1\}_{k=1}^{\mathrm{}}`$. We can, of course, compute a first minor<sup>1</sup><sup>1</sup>1That is, a minor that has all columns of the matrix of which it is a minor, but one, and all rows but one. of $`K(G)`$, and then apply Cramer’s rule, whenever we want an entry of $`(K(G)^1)^T`$. Unfortunately, it seems very hard to obtain asymptotic expressions directly from the minors. We will reduce the problem of computing the minor of $`K(G)`$ resulting from the deletion of row $`i`$ and column $`j`$ to an enumerative problem whose solution we will be able to represent in a form other than a determinant. It would seem, at first sight, that we can use the same kind of argument we used to answer the previous question, and show that such a minor is equal (up to sign) to the number of all perfect matchings of $`G`$ with white vertex $`i`$ and black vertex $`j`$ deleted. Unfortunately, this is not the case. A first minor without row $`i`$ and column $`j`$ is equal to a sum whose number of terms is equal to number of perfect matchings of $`G`$ with white vertex $`i`$ and black vertex $`j`$ deleted. The problem is that, while every term has absolute value $`1`$, not every term has the same sign. It is easy to show, by the same kind of reasoning we employed in our answer to the first question, that the matter of whether or not two terms have the same sign can be determined by examining the loops in the superimposition of the two matchings corresponding to the two terms. The terms have the same sign if and only if there is an even number of loops having one of the two deleted vertices, but not the other, in their interiors. ###### Definition 1. The Aztec diamond of order $`n`$ is a planar graph consisting of vertices $$\{(2r+1,2s):0r<n,0sn\}\{(2r,2s+1):0rn,0s<n\}$$ (20) and of edges $`\{`$ $`\{\{(2r,2s+1),(2r+1,2s+2)\},\{(2r+1,2s+2),(2r+2,2s+1)\},`$ (21) $`\{(2r+2,2s+1),(2r+1,2s)\},\{(2r+1,2s),(2r,2s+1)\}\}:0r<n,0s<n\}`$ in Cartesian coordinates. It will soon become apparent that, for, our purposes, the system of coordinates in Figure 17 is more convenient than Cartesian coordinates. It will also become clear why we draw the Aztec diamond as if on an infinite square grid tilted $`45`$ degrees from the “natural” direction. For now, let us notice that, if we color the Aztec diamond as a checkerboard, all vertices on a column have the same color, as do all vertices on a row. Let us also use the system of coordinates in Figure 17 instead of the Cartesian system, and refer to the edge consisting of white vertex $`(x_0,y_0)`$ and black vertex $`(x_1,y_1)`$ as $`((x_0,y_0),(x_1,y_1))`$. We can reformulate the rule for comparing terms’ signs so that it does not mention loops. ###### Lemma 1. Let $`A`$ and $`B`$ be two perfect matchings of the Aztec diamond of order $`n`$ with the white vertex at $`(w_0,w_1+d_1)`$ and the black vertex at $`(w_0+d_0,w_1)`$ deleted<sup>2</sup><sup>2</sup>2We shall henceforth use the system of coordinates in Figure 17 instead of the Cartesian system.. The following two conditions are equivalent: 1. $`ABAB`$ has an even number of loops containing exactly one of the two deleted vertices. 2. $`w(A)w(B)mod2`$, where $`w(T)`$ is the number of edges of the form $`((i1,w_1+1),(i,w_1))`$, $`1<i<w_0+d_0`$; $`((i,w_1+1),(i,w_1))`$, $`1i<w_0+d_0`$; $`((i,w_1+d_1),(i,w_1+d_11)`$, $`1i<w_0`$, and $`((i,w_1+d_1),(i+1,w_1+d_11))`$, $`1i<w_0`$ in a perfect matching $`T`$. ###### Proof. The sum $`w(A)+w(B)`$ is congruent, modulo $`2`$, to the total number of edges in $`ABAB`$ consisting of a black vertex $`(x,y)D`$ and white vertex $`(x1,y+1)`$ or $`(x,y+1)`$. Given a loop $`\mathrm{}`$ in $`ABAB`$, the number of edges in it consisting of a black vertex $`(x,y)E`$ and white vertex $`(x1,y+1)`$ or $`(x,y+1)`$ is equal to the number of times the loop crosses a ray with its end slightly below the deleted black vertex and with diagonal direction with respect to the square grid. (See figure 15.) This number is even if the deleted black vertex is in the exterior of the loop, and odd if it is in the interior. The same holds for $`F`$ and the deleted white vertex. Hence a loop has an even number of edges consisting of a black vertex $`(x,y)D`$ and white vertex $`(x1,y+1)`$ or $`(x,y+1)`$ if and only if it has exactly one of the two deleted vertices in its interior. We conclude the proof by summing over all loops. ∎ It follows that $`\|(1)^{w(T)}\|=\|_\pi \mathrm{sgn}(\pi )_{i=1}^na_{i,\pi (i)}\|`$, where the sum on the left is over all perfect matchings $`T`$ of the Aztec diamond with white vertex at $`(w_0,w_1+d_1)`$ and black vertex at $`(w_0+d_0,w_1)`$, and $`\left\|a_{i,j}\right\|_1^n`$ is the first minor of the Kasteleyn matrix $`K(G)`$ of the (intact) Aztec diamond $`G`$ where the row corresponding to the white vertex at $`(w_0,w_1+d_1)`$ and the column corresponding to the black vertex at $`(w_0+d_0,w_1)`$. Thus, if we compute $`\|(1)^{w(T)}\|`$ (and this is essentially an enumerative problem, which we shall solve enumeratively), we will know the absolute value of $`_\pi \mathrm{sgn}(\pi )_{i=1}^na_{i,\pi (i)}`$, which will give us the absolute value of the entry of $`(K(G)^1)^T`$ to be computed, but not its sign. We will find the sign now. It is better, for this particular goal, to express the entry of the inverse Kasteleyn matrix of the Aztec diamond, not as the ratio of the determinant of the minor of the Kasteleyn matrix resulting from the deletion of column $`i_0`$ and row $`j_0`$, divided by the determinant of the Kasteleyn matrix, multiplied by $`(1)^{i_0+j_0}`$, but rather as the ratio of the determinant $`\left\|b_{i,j}\right\|_1^{n1}`$ to the determinant of the Kasteleyn matrix $`K(G)=\left\|a_{i,j}\right\|_1^n`$, where $`b_{i_0,j}=0`$ for $`jj_0`$, $`b_{i,j_0}=0`$ for $`ii_0`$, $`b_{i_0,j_0}=1`$, $`b_{i,j}=a_{i,j}`$ for $`ii_0`$, $`jj_0`$. We can then ask whether, for a permutation $`\pi `$ of $`\{1,2,\mathrm{},n1\}`$, $`\mathrm{sgn}(\pi )_ib_{i,\pi (i)}`$ has the same sign as $`(1)^{w(T)}`$, where $`T`$ is the perfect matching corresponding to $`\pi `$. (The answer will be the same for all permutations, so we have to ask it for only one permutation (or matching) we choose.) If the signs are the same, then $`(1)^{w(T)}=(((K(G))^1)^T)_{i_0,j_0}\left\|K(G)\right\|_1^n`$; if the sign are different, then $`(1)^{w(T)}=(((K(G))^1)^T)_{i_0,j_0}\left\|K(G)\right\|_1^n`$. For our search for the sign, we need to fix an ordering of the black and white vertices of the Aztec diamond. The ordering in Figure 17 will prove itself convenient. We want to prove that, for any $`(w_0,w_1,d_0,d_1)`$, $$(1)^{w(T)}=(1)^{d_0+d_1+1}\underset{\pi }{}\mathrm{sgn}(\pi )\underset{i=1}{\overset{n}{}}a_{i,\pi (i)}.$$ The most straightforward method of proof consists of comparing the signs of $`(1)^{w(T)}`$ and $`\mathrm{sgn}(\pi )_{i=1}^na_{i,\pi (i)}`$ for some tiling $`T`$ and its corresponding permutation $`\pi `$. Unfortunately, this is also quite cumbersome, as which tilings are possible depends on the order of $`w_0`$, $`w_0+d_0`$, $`w_1`$ and $`w_1+d_1`$. In order not to bore the reader with twenty-four different cases, we will proceed explicitly only for the six cases corresponding to $`d_0,d_1>0`$, which can be treated as four different cases. #### 2.4.1. Case 1: $`w_1w_0<w_1+d_1`$ We choose to examine the matching having $`/`$-edges at $`((i,w_1),(i,w_1))`$, $`w_1i<w_0+d_0`$, $`((w_1,j),(w_1,j))`$, $`w_1<j<w_1+d_1`$, and $`((i,w_1+d_1),(i+1,w_1+d_11))`$, $`w_1i<w_0`$, where $`((x_0,y_0),(x_1,y_1))`$ denotes the edge consiting of a white vertex at $`(x_0,y_0)`$ and a black vertex at $`(x_1,y_1)`$. All other vertices are covered by $``$-edges: $`((x,y),(x+1,y))`$ for $`xy`$, given that at most one of the two conditions $`w_1x<w_0+d_0`$, $`y=w_1`$ is fulfilled, and $`((x,y),(x,y1))`$ for $`x<y`$, given that at most one of $`x=w_1`$, $`w_1<y<w_1+d_1`$ holds, as well as at most one of $`w_1x<w_0`$, $`y=w_1+d_1`$. (Note that we denote by $`/`$-edges and $``$-edges what we called “vertical edges” and “horizontal edges” before. This is just a change in notation with the purpose that the fact that we now draw the Aztec diamond tilted $`45`$ degrees will not confuse the reader.) It is easy to verify that this is a perfect matching. For this matching, $`(1)^{w(T)}`$ is $`(1)^{w_0+w_1}`$. Each matching corresponds to a permutation of $`\{1,2,\mathrm{},n\}`$. (We have permutations of $`\{1,2,\mathrm{},n\}`$, and not of $`\{1,2,\mathrm{},n1\}`$, because we are working with $`\left\|b_{i,j}\right\|_1^n`$, and not with a first minor of $`K(G)=\left\|a_{i,j}\right\|_1^n`$. The convenience of this choice, which we made without justification, will be clear by the end of this paragraph.) The permutation $`\pi `$ corresponding to the matching we are now examining consists of one cycle. The length of this cycle is equal to the number of vertical lozenges, plus one. (This is not true in general. However, it is true in cases similar to the one we are currently examining, in which white vertices $`\{f_1,f_2,\mathrm{}f_m\}`$ and black vertices $`\{g_1,g_2,\mathrm{}g_m\}`$, which would be paired in the form $`(f_i,g_i)`$, $`1im`$ in the all-$``$ matching, are paired in the form $`(f_i,g_{i+1})`$, $`1i<m`$, and vertices $`f_m`$ and $`g_1`$ are missing.) Hence $`\mathrm{sgn}(\pi )`$ is equal to $`(1)`$ to the power of the number of vertical lozenges, that is, $`2(w_0w_1)+d_0+d_11`$ lozenges. We now need only examine $`_{i=1}^mb_{i,\pi (i)}`$. Among the edges in our perfect matching, only $`((i,w_1+d_1),(i+1,w_1+d_11))`$, $`w_1i<w_0`$, correspond to entries equal to $`(1)`$ in the determinant $`\left\|b_{i,j}\right\|`$. Therefore $`_{i=1}^mb_{i,\pi (i)}=(1)^{w_0w_1}`$. We conclude that our matching $`T`$, for which $`(1)^{w(T)}=(1)^{w_0+w_11}`$, contributes $`(1)^{2(w_0w_1)+d_0+d_11}(1)^{w_0w_1}`$, to the determinant. Therefore an arbitrary matching $`T`$ contributes $$(1)^{w(T)(w_0+w_11)}(1)^{2(w_0w_1)+d_0+d_11}(1)^{w_0w_1},$$ (22) that is, $$(1)^{w(T)}(1)^{d_0+d_1+1},$$ (23) to the determinant. Therefore $$(1)^{w(T)}=(1)^{d_0+d_1+1}\underset{\pi }{}\mathrm{sgn}(\pi )\underset{i=1}{\overset{n}{}}a_{i,\pi (i)}$$ for case 1. #### 2.4.2. Case 2: $`w_0w_1+d_1`$ We choose to examine the matching having $`//`$-edges at $`((i,w_1),(i,w_1))`$, $`w_1<i<w_0+d_0`$, at $`((w_0+1,j),(w_0,j+1))`$, $`w_1+d_1jw_0`$, at $`((w_1,j),(w_1,j))`$, $`w_1jw_0`$, and at $`((i+1,w_0),(i,w_0+1))`$, $`w_1i<w_0`$. All other vertices are covered by $``$-edges. The $`(1)^{w(T)}`$ of this matching is $`(1)^{d_1+1}`$. The sign of the permutation is $$(1)^{(w_0+d_0w_11)+(w_0(w_1+d_1)+1)+(w_0w_1+1)+(w_0w_1)},$$ (24) that is, $`(1)^{d_0+d_1+1}`$. The product $`_{i=1}^mb_{i,\pi (i)}`$ is $`(1)^{(w_0(w_1+d_1)+1)+(w_0w_1)}`$, that is, $`(1)^{d_0+d_1}`$. Hence each matching $`T`$ contributes $$(1)^{w(T)}(1)^{d_0+d_1+1}$$ (25) to the determinant. Therefore $$(1)^{w(T)}=(1)^{d_0+d_1+1}\underset{\pi }{}\mathrm{sgn}(\pi )\underset{i=1}{\overset{n}{}}a_{i,\pi (i)}$$ for case 2. #### 2.4.3. Case 3: $`w_0<w_1<w_0+d_0`$ We choose to examine the matching having $`//`$-edges at $`((w_0,j),(w_0,j))`$, $`w_0j<w_1+d_1`$, at $`((i,w_0),(i,w_0))`$, $`w_0<iw_0+d_0`$, and at $`((w_0+d_0,j+1),(w_0+d_0+1,j))`$, $`w_0j<w_1`$. All other vertices are covered by $``$-edges. The $`(1)^{w(T)}`$ of this matching is $`(1)^{w_0+w_1+1}`$. The sign of the permutation is $`(1)^{d_0+d_1}`$. The product $`_{i=1}^mb_{i,\pi (i)}`$ is $`(1)^{w_1w_0}`$. Therefore each matching $`T`$ contributes $$(1)^{w(T)}(1)^{d_0+d_1+1}$$ (26) to the determinant. Therefore $$(1)^{w(T)}=(1)^{d_0+d_1+1}\underset{\pi }{}\mathrm{sgn}(\pi )\underset{i=1}{\overset{n}{}}a_{i,\pi (i)}$$ for case 3. #### 2.4.4. Case 4: $`w_1w_0+d_0`$ We choose to examine the matching having $`//`$-edges at $`((w_0,j),(w_0,j))`$, $`w_0j<w_1+d_1`$, at $`((i,w_0),(i,w_0))`$, $`w_0<iw_1`$, at $`((w_1,j+1),(w_1+1,j))`$, $`w_0jw_1`$, and at $`((i,w_1+1),(i+1,w_1))`$, $`w_0+d_0i<w_1`$. All other vertices are covered by $``$-edges. The $`(1)^{w(T)}`$ of this matching is $`(1)^{d_0+1}`$. The sign of the permutation is $`(1)^{d_0+d_1+1}`$. The product $`_{i=1}^mb_{i,\pi (i)}`$ is $`(1)^{d_0+1}`$. Therefore each matching $`T`$ contributes $$(1)^{w(T)}(1)^{d_0+d_1+1}$$ (27) to the determinant. Therefore $$(1)^{w(T)}=(1)^{d_0+d_1+1}\underset{\pi }{}\mathrm{sgn}(\pi )\underset{i=1}{\overset{n}{}}a_{i,\pi (i)}$$ for case 4. ## 3. An enumerative problem Consider an Aztec diamond of side, or order, $`n`$, with black vertex $`(w_0+d_0,w_1)`$ and white vertex $`(w_0,w_1+d_1)`$ missing. Our task in this section is to compute $$\underset{T}{}(1)^{w(T)}$$ (28) where $`T`$ ranges over all perfect matchings of this Aztec diamond with two missing vertices. Enumeratively speaking, we must count the number of perfect matchings, counting as “negative” any matching $`T`$ for which $`(1)^{w(T)}=(1)`$. The function $`w(T)`$, as defined in the previous section, gives, for a perfect matching $`T`$, the number of edges consisting of a black vertex $`(x,y)D`$ and white vertex $`(x1,y+1)`$ or $`(x,y+1)`$. $`D`$ is the union of two set of black vertices: $`E=\{(i,w_1):i<w_0+d_0\}`$ and $`F=\{(i,w_1+d_1):iw_0\}`$. Throughout this section, we will assume $`d_0>0`$, $`d_1>0`$. At the end it will become clear that this assumption involves no loss of generality. One of our tools for attacking the special case $`d_0>0`$, $`d_1>0`$ will be the EKLP Lemma, a classical result in which we will soon explain. First of all, we must define two kinds of subsets of the Aztec diamond. They will be our intermediate objects of study. ###### Definition 2. An $`n\times m`$ black-edged Aztec rectangle with dents at $`1x_1<x_2<\mathrm{}<x_mn+1`$ is a graph consisting of white vertices $`(i,j),`$ $`1in,`$ $`1jm`$, black vertices $`(i,j),`$ $`1in+1`$, $`1jm1`$, and black vertices $`(i,m)`$ such that $`1km`$ we have $`ix_k`$. An $`n\times m`$ white-edged Aztec rectangle with teeth at $`1y_1<y_2<\mathrm{}<y_mn`$ is a graph consisting of white vertices $`(i,j),`$ $`1in,`$ $`1jm,`$ black vertices $`(i,j),`$ $`1in+1,`$ $`1jm,`$ and white vertices $`(i,m+1)`$ s.t. $`1km`$ with $`i=x_k`$. Let the number of perfect matchings of an $`n\times m`$ black-edged Aztec rectangle with dents at $`1x_1<x_2<\mathrm{}<x_mn+1`$ be called $`D_{n,m}(x_1,x_2,\mathrm{},x_m).`$ What is the number of perfect matchings of an $`n\times m`$ white-edged Aztec rectangle $`Y`$ with teeth at $`1y_1<y_2<\mathrm{}<y_mn`$? Every white vertex (or “tooth”) $`(y_i,m+1)`$ must be covered by a horizontal or vertical edge, which will also cover black vertex $`(y,m)`$ or $`(y+1,m)`$, respectively. Thus, any matching of $`Y`$ is composed of * a matching of a black-edged $`n\times m`$ Aztec rectangle with dents at $`1x_1<x_2<\mathrm{}<x_mn+1,`$ where $`x_iy_i`$ is either $`0`$ or $`1`$; * the unique matching of the remaining region. Therefore the number of matchings of $`Y`$ is the $`m`$-fold sum $$E_{n,m}(y_1,\mathrm{},y_m)=\underset{x_i=y_i\text{ or }y_i+1}{}D_{n,m}(x_1,\mathrm{},x_m)$$ (29) where we assume that $`D_{n,m}(x_1,x_2,\mathrm{},x_m)=0`$ when any two $`x_i`$’s are equal. In the same way that, given $`D_{n,m},`$ we have found $`E_{n,m},`$ given $`E_{n,m},`$ we can find $`D_{n,m+1}`$. We have again an $`m`$-fold sum: $$D_{n,m+1}(x_1,\mathrm{},x_{m+1})=\underset{x_iy_i<x_{i+1}}{}E_{n,m}(y_1,\mathrm{},y_m)$$ (30) where we assume that $`E_{n,m}(x_1,x_2,\mathrm{},x_m)=0`$ when any two $`x_i`$’s are equal. We will now be able to prove the following by induction. The proof in does not use Lemma 3 explicitly. ###### Lemma 2 (EKLP Lemma). The number of matchings of an $`n\times m`$ black-edged Aztec rectangle with dents $`1x_1<x_2<\mathrm{}<x_mn+1`$ is $$D_{n,m}(x_1,x_2,\mathrm{},x_m)=\frac{2^{\frac{m(m1)}{2}}}{(m1)!!}\left\|x_i^{j1}\right\|_1^m,$$ (31) where $`k!!=1!\mathrm{\hspace{0.17em}2}!\mathrm{\hspace{0.17em}3}!\mathrm{}(k1)!k!`$. The number of matchings of an $`n\times m`$ white-edged Aztec rectangle with dents at $`1y_1<y_2<\mathrm{}<y_mn`$ is $$E_{n,m}(y_1,y_2,\mathrm{},y_m)=\frac{2^{\frac{m(m+1)}{2}}}{(m1)!!}\left\|y_i^{j1}\right\|_1^m$$ (32) For the proof of this lemma, we need a special case of another lemma that we will later state in its full generality. ###### Lemma 3 (Lemma A (special case)). Let $`A`$ be an operator carrying polynomials to polynomials of the same or lesser degree: $$\begin{array}{ccccc}x^0& & & & a_{0,0}x^0\\ x^1& & & a_{1,1}x^1+& a_{1,0}x^0\\ x^2& & a_{2,2}x^2+& a_{2,1}x^1+& a_{2,0}x^0\\ \mathrm{}\end{array}$$ (33) Then $$\left\|A(x^{j1})(x_i)\right\|_1^m=a_{0,0}a_{1,1}\mathrm{}a_{m1,m1}\left\|x_i^{j1}\right\|_1^m,$$ (34) where $$A(x^j)(x_i)=a_{j,j}x_i^j+a_{j,j1}x_i^{j1}+\mathrm{}+a_{j,0}x_i^0.$$ (35) ###### Proof. The determinant on the left side of (34) can be obtained from the determinant on the right side by elementary column operations. ∎ ###### Proof of EKLP Lemma. The base case $`(m=1)`$ is trivial. The inductive step $`D_{n,m}E_{n,m}`$ follows directly from (30) and from Lemma A (special case). The inductive step $`E_{n,m}D_{n,m+1}`$ requires some more work. $`D_{n,m+1}(x_1,\mathrm{}x_{m+1})`$ $`={\displaystyle \underset{x_iy_i<x_{i+1}}{}}E_{n,m}(y_1,\mathrm{},y_m)`$ (36) $`={\displaystyle \underset{x_iy_i<x_{i+1}}{}}{\displaystyle \frac{2^{\frac{m(m+1)}{2}}}{(m1)!!}}\left\|y_i^{j1}\right\|_1^m`$ $`={\displaystyle \frac{2^{\frac{m(m+1)}{2}}}{(m1)!!}}\left\|x_i^{j1}+(x_i+1)^{j1}+\mathrm{}(x_{i+1}1)^{j1}\right\|_1^m`$ $`={\displaystyle \frac{2^{\frac{m(m+1)}{2}}}{(m1)!!}}\left\|{\displaystyle \frac{1}{j}}(B_j(x_{i+1})B_j(x_i))\right\|_1^m,`$ where $`B_r(y)`$ is the Bernoulli polynomial of degree $`r`$, which has the property $`_{n=M}^{N1}n^q=(B_{q+1}(N)B_{q+1}(M))/(q+1)`$. For a definition of $`B_r(y)`$, consult . The only property of $`B_r(y)`$ we need to know is that it is a polynomial of degree r whose leading coefficient is one. We will adopt, for the sake of convenience, the convention $`B_0(y)=y^0`$. We can now continue: $`D_{n,m+1}(x_1,\mathrm{}x_{m+1})`$ $`={\displaystyle \frac{2^{\frac{m(m+1)}{2}}}{m!!}}\left\|B_j(x_{i+1})B_j(x_i)\right\|_1^m`$ (37) $`={\displaystyle \frac{2^{\frac{m(m+1)}{2}}}{m!!}}\left\|\begin{array}{ccccc}1& B_1(x_1)& B_2(x^1)& \mathrm{}& B_m(x_1)\\ 1& B_1(x_2)& B_2(x_2)& \mathrm{}& B_m(x_2)\\ \mathrm{}\\ 1& B_1(x_{m+1})& B_2(x_{m+1})& \mathrm{}& B_m(x_{m+1})\end{array}\right\|,`$ (38) where we have pasted a new column onto the left edge of the determinant and a new row onto the top edge, and then added the first row to the second one, the second to the third one, and so on, in succession $`={\displaystyle \frac{2^{\frac{m(m+1)}{2}}}{m!!}}\left\|x_i^{j1}\right\|_1^{m+1},`$ (39) by Lemma A (special case). ∎ Let us now count matchings with weights $`1`$ and $`(1)`$: a matching may count as one matching or as minus one matchings. Consider an $`n\times m`$ black-edged Aztec rectangle with dents at $`1x_1<x_2<\mathrm{}x_{m1}n+1`$, and, in addition, a dent at $`w_0`$. Let us multiply the total number of matchings by $`(1)`$ if there is an odd number of $`x_i`$’s smaller than $`w_0`$. (Here we are counting either all matchings as positive matchings or all as negative matchings.) What is, then, this total, weighted number of matchings? Conveniently, it is $$\frac{2^{\frac{m(m1)}{2}}}{(m1)!!}\left\|\begin{array}{c}w_0^{j1}\\ x_1^{j1}\\ x_2^{j1}\\ \mathrm{}\\ x_{m1}^{j1}\end{array}\right\|$$ (40) We have just taken (31) and shifted the $`w_0^j`$ row corresponding to the dent at $`w_0`$ as many positions upwards as there are $`x_i`$’s smaller than $`w_0`$. What is the weighted number of matchings of a white-edged $`n\times m`$ Aztec rectangle $`Y`$ with teeth at $`y_1,y_2,\mathrm{}y_{m1}`$ and a black hole at $`(w_0,m)`$? (The weight of each matching is $`(1)`$ to the power of the number of its edges consist of a black vertex $`(i,w_1)`$ and white vertex $`(i1,w_1+1)`$ or $`(i,w_1+1)`$, where $`1i<w_0+d_0`$.) A cursory examination makes clear that a matching of $`Y`$ will have weight $`1`$ if and only if the black-edged $`n\times m`$ Aztec rectangle the matching outlines has an even number of dents with indices lower than $`w`$. Thus the weighted number of matchings is $$\underset{x_i=y_iory_i+1}{}\frac{2^{\frac{m(m1)}{2}}}{(m1)!!}\left\|\begin{array}{c}w_0^{j1}\\ x_1^{j1}\\ x_2^{j1}\\ \mathrm{}\\ x_{m1}^{j1}\end{array}\right\|,$$ (41) which is the same as $$\frac{2^{\frac{m(m1)}{2}}}{(m1)!!}\left\|\begin{array}{c}w_0^{j1}\\ y_1^{j1}+(y_1+1)^{j1}\\ y_2^{j1}+(y_2+1)^{j1}\\ \mathrm{}\\ y_{m1}^{j1}+(y_{m1}+1)^{j1}\end{array}\right\|.$$ (42) Here we face a difficulty. We would like to use Lemma 3. Unfortunately, we have $`w_0^{j1}`$ on the top row. What we need is to find something which is transformed into $`w_0^{j1}`$ by the linear operations taking a row of the form $`x^{j1}`$ to $`x^{j1}+(x+1)^{j1}`$. First of all, we need a stronger version of Lemma 3. ###### Lemma 4 (Lemma A, stronger version). Let A be an operator carrying polynomials to polynomials of the same or lesser degree: $$\begin{array}{ccccc}x^0& & & & a_{0,0}x^0\\ x^1& & & a_{1,1}x^1+& a_{1,0}x^0\\ x^2& & a_{2,2}x^2+& a_{2,1}x^1+& a_{2,0}x^0\\ \mathrm{}\end{array}$$ (43) Then $$\left\|\underset{1kl_i}{}b_{k,i}A(x^{j1})(x_{k,i})\right\|_1^m=a_{0,0}\mathrm{}a_{m1,m1}\left\|\underset{1kl_i}{}b_{k,i}x_{k,i}^{j1}\right\|_1^m,$$ (44) for any $`l_i`$, $`b_{k,i},`$ $`1kl_i,`$ $`1im`$, where $$A(x^j)(x_i)=a_{j,j}x_i^j+a_{j,j1}x_i^{j1}+\mathrm{}+a_{j,0}x_i^0.$$ (45) ###### Proof. The special case followed from the fact that the same sequence of column operations transforms the row $`A(x^{j1})(y)`$ into $`y^j`$ and the row $`A(x^j)(z)`$ into $`z^j`$, for any values $`y`$,$`z`$. Therefore the same sequence transforms the row $`aA(x^j)(y)+bA(x^j)(z)`$ into $`ay^j+bz^j`$. ∎ We need $`a_k`$’s such that $`_ka_kA(x^j)(v_k)=w_0^j`$ for $`0j<m`$, where $`A`$ is the operator taking $`x^j`$ to $`x^j+(x+1)^j`$. ###### Definition 3. $`\mathrm{\Delta }`$ is the operator taking $`x^j`$ to $`(x+1)^jx^j`$. Then $`A=(2I+\mathrm{\Delta })`$, where $`I`$ is the identity operator, and we have, formally, $`(2I+\mathrm{\Delta })^1`$ $`={\displaystyle \frac{1}{2}}(I+{\displaystyle \frac{\mathrm{\Delta }}{2}})^1`$ (46) $`={\displaystyle \frac{1}{2}}{\displaystyle \underset{j=0}{\overset{\mathrm{}}{}}}(1)^j{\displaystyle \frac{\mathrm{\Delta }}{2}}^j`$ Since $`\mathrm{\Delta }^j`$ vanishes on polynomials of degree smaller than $`j`$, and we are dealing with polynomials of degree at most $`m1`$, we can use just the first $`m`$ terms of the series: $$(1/2)\underset{j=0}{\overset{m1}{}}(1)^j(\mathrm{\Delta }/2)^j$$ (47) The reader can easily check that, for any $`x^k`$ with $`k<m`$, $$(2I+\mathrm{\Delta })((\frac{1}{2}\underset{j=0}{\overset{m1}{}}(1)^j(\frac{\mathrm{\Delta }}{2})^j)(x^k))$$ (48) gives $`x^k`$ plus a constant times $`\mathrm{\Delta }^mx^k`$, and, since $`m>k`$, $`\mathrm{\Delta }^mx^k`$ is zero. One can also check the same for $$((\frac{1}{2}\underset{j=0}{\overset{m1}{}}(1)^j(\frac{\mathrm{\Delta }}{2})^j)(x^k))(2I+\mathrm{\Delta }).$$ (49) It is quite convenient that $`(2I+\mathrm{\Delta })`$ and $$((\frac{1}{2}\underset{j=0}{\overset{m1}{}}(1)^j(\frac{\mathrm{\Delta }}{2})^j)(x^k))$$ (50) commute, and that their composition is the same as the identity operator for the domain we are interested in. We will use the shorthand $`(2I+\mathrm{\Delta })^1`$ for $$((\frac{1}{2}\underset{j=0}{\overset{m1}{}}(1)^j(\mathrm{\Delta }/2)^j)(x^k))$$ (51) without any compuctions. If we define $`a_0,a_1,\mathrm{},a_{m1}`$ by $$(\frac{1}{2}\underset{j=0}{\overset{m1}{}}(1)^j(\mathrm{\Delta }/2)^j)(f(x))=\underset{i=0}{\overset{m1}{}}a_if(x+i),$$ (52) then $`{\displaystyle \underset{i=0}{\overset{m1}{}}}a_i((2I+\mathrm{\Delta })(x^j)(w_0+i))`$ $`=((({\displaystyle \frac{1}{2}}{\displaystyle \underset{j=0}{\overset{m1}{}}}(1)^j(\mathrm{\Delta }/2)^j)(2I+\mathrm{\Delta }))(x^j))(w_0)`$ (53) $`=(x^j)(w_0)`$ $`=w_0^j,`$ for $`0j<m,`$ as we desired. Hence, by Lemma A, $$\left\|\begin{array}{c}w_0^{j1}\\ y_1^{j1}+(y_1+1)^{j1}\\ y_2^{j1}+(y_2+1)^{j1}\\ \mathrm{}\\ y_{m1}^{j1}+(y_{m1}+1)^{j1}\end{array}\right\|$$ (54) is equal to $$2^m\left\|\begin{array}{c}(\frac{1}{2}_{k=0}^{m1}(1)^k(\frac{\mathrm{\Delta }}{2})^k)(x^{j1})(w_0)\\ y_1^{j1}\\ y_2^{j1}\\ \mathrm{}\\ y_{m1}^{j1}\end{array}\right\|,$$ (55) or, in shorthand, $$2^m\left\|\begin{array}{c}((2I+\mathrm{\Delta })^1(x^{j1}))(w_0)\\ y_1^{j1}\\ y_2^{j1}\\ \mathrm{}\\ y_{m1}^{j1}\end{array}\right\|$$ (56) Now, for every polynomial $`p`$, $$p(w_0)=(((I+\mathrm{\Delta })^{w_01})(p))(1).$$ (57) Hence we can write $$2^m\left\|\begin{array}{c}((I+\mathrm{\Delta })^{w_01}((2I+\mathrm{\Delta })^1(x^{j1})))(1)\\ y_1^{j1}\\ y_2^{j1}\\ \mathrm{}\\ y_{m1}^{j1}\end{array}\right\|$$ (58) and eliminate all terms of degree $`m`$ or higher. Hence the weighted number of matchings of a white-edged $`n\times m`$ Aztec rectangle with teeth at $`y_1,y_2,\mathrm{}y_{m1}`$ and a black hole at $`(w_0,m)`$ is $$\frac{2^{\frac{m(m+1)}{2}}}{(m1)!!}\left\|\begin{array}{c}((I+\mathrm{\Delta })^{w_01}((2I+\mathrm{\Delta })^1(x^{j1})))(1)\\ y_1^{j1}\\ y_2^{j1}\\ \mathrm{}\\ y_{m1}^{j1}\end{array}\right\|$$ (59) What would be the weighted number of matchings of a black-edged $`n\times (m+1)`$ Aztec rectangle with dents at $`x_1,x_2,\mathrm{}x_m`$ and a black hole at $`(w_0,m)`$? This is the sum we have to simplify: $$\underset{x_iy_i<x_{i+1}}{}\frac{2^{\frac{m(m+1)}{2}}}{(m1)!!}\left\|\begin{array}{c}((I+\mathrm{\Delta })^{w_01}((2I+\mathrm{\Delta })^1(x^{j1})))(1)\\ y_1^{j1}\\ y_2^{j1}\\ \mathrm{}\\ y_{m1}^{j1}\end{array}\right\|$$ (60) This is equal to $`\frac{2^{\frac{m(m+1)}{2}}}{(m1)!!}`$ times $`{\displaystyle \underset{x_iy_i<x_{i+1}}{}}`$ $`\left\|\begin{array}{c}((I+\mathrm{\Delta })^{w_01}((2I+\mathrm{\Delta })^1(x^{j1})))(1)\\ y_1^{j1}\\ y_2^{j1}\\ \mathrm{}\\ y_{m1}^{j1}\end{array}\right\|`$ (61) $`=\left\|\begin{array}{c}((I+\mathrm{\Delta })^{w_01}((2I+\mathrm{\Delta })^1(x^{j1})))(1)\\ x_1^{j1}+(x_1+1)^{j1}+\mathrm{}+(x_21)^{j1}\\ x_2^{j1}+(x_2+1)^{j1}+\mathrm{}+(x_31)^{j1}\\ \mathrm{}\\ x_{m1}^{j1}+(x_{m1}+1)^{j1}+\mathrm{}+(x_m1)^{j1}\end{array}\right\|`$ (62) $`=\left\|\begin{array}{c}((I+\mathrm{\Delta })^{w_01}((2I+\mathrm{\Delta })^1(x^{j1})))(1)\\ \frac{1}{j}(B_j(x_2)B_j(x_1))\\ \frac{1}{j}(B_j(x_3)B_j(x_2))\\ \mathrm{}\\ \frac{1}{j}(B_j(x_m)B_j(x_{m1}))\end{array}\right\|`$ (63) $`=\left\|\begin{array}{cc}1& \frac{1}{j1}B_{j1}(x_1)\\ 0& ((I+\mathrm{\Delta })^{w_01}((2I+\mathrm{\Delta })^1(x^{j2})))(1)\\ 0& \frac{1}{j1}B_{j1}(x_2)B_{j1}(x_1)\\ \mathrm{}& \mathrm{}\\ 0& \frac{1}{j1}B_{j1}(x_m)B_{j1}(x_{m1})\end{array}\right\|`$ (64) $`=\left\|\begin{array}{cc}0& ((I+\mathrm{\Delta })^{w_01}((2I+\mathrm{\Delta })^1(x^{j2})))(1)\\ 1& \frac{1}{j1}B_{j1}(x_1)\\ 0& \frac{1}{j1}B_{j1}(x_2)B_{j1}(x_1)\\ \mathrm{}& \mathrm{}\\ 0& \frac{1}{j1}B_{j1}(x_m)B_{j1}(x_{m1})\end{array}\right\|`$ (65) $`=\left\|\begin{array}{cc}0& ((I+\mathrm{\Delta })^{w_01}((2I+\mathrm{\Delta })^1(x^{j2})))(1)\\ 1& \frac{1}{j1}B_{j1}(x_1)\\ 1& \frac{1}{j1}B_{j1}(x_2)\\ \mathrm{}& \mathrm{}\\ 1& \frac{1}{j1}B_{j1}(x_m)\end{array}\right\|`$ (66) Since $`\{((I+\mathrm{\Delta })^{w_01}((2I+\mathrm{\Delta })^1(x^{j2})))(1)\}_{j=2,3,\mathrm{}m+1}`$ is a linear combination of rows of the form $`\{k^{j2}\}_{j=2,3,\mathrm{}m+1}`$, $`1km`$, it is enough to show how to simplify $$\left\|\begin{array}{cc}0& k^{j2}\\ 1& \frac{1}{j1}B_{j1}(x_1)\\ 1& \frac{1}{j1}B_{j1}(x_2)\\ \mathrm{}& \mathrm{}\\ 1& \frac{1}{j1}B_{j1}(x_m)\end{array}\right\|$$ (67) This happens to be easier than one would expect. As a particular case of $$\underset{n=M}{\overset{N1}{}}n^q=\frac{B_{q+1}(N)B_{q+1}(M)}{q+1},$$ we have $$k^{j1}=\frac{1}{j}(B_j(k+1)B_j(k)).$$ Therefore $`\left\|\begin{array}{cc}0& k^{j2}\\ 1& \frac{1}{j1}B_{j1}(x_1)\\ 1& \frac{1}{j1}B_{j1}(x_2)\\ \mathrm{}& \mathrm{}\\ 1& \frac{1}{j1}B_{j1}(x_m)\end{array}\right\|`$ $`=\left\|\begin{array}{cc}11& \frac{1}{j1}(B_{j1}(k+1)B_{j1}(k))\\ 1& \frac{1}{j1}B_{j1}(x_1)\\ 1& \frac{1}{j1}B_{j1}(x_2)\\ \mathrm{}& \mathrm{}\\ 1& \frac{1}{j1}B_{j1}(x_m)\end{array}\right\|`$ (68) $`={\displaystyle \frac{1}{m!}}\left\|\begin{array}{c}(k+1)^{j1}k^{j1}\\ x_1^{j1}\\ x_2^{j1}\\ \mathrm{}\\ x_m^{j1}\end{array}\right\|`$ (69) $`={\displaystyle \frac{1}{m!}}\left\|\begin{array}{c}\mathrm{\Delta }(x^{j1})(k)\\ x_1^{j1}\\ x_2^{j1}\\ \mathrm{}\\ x_m^{j1}\end{array}\right\|`$ (70) Thus, we now know that the sequence of elementary column operations we have to apply to the matrix having $`x_i^j`$ on its lower rows in order to make it into a matrix having $`\{1,\frac{1}{1}B_1(x_i),\frac{1}{2}B_2(x_i),\mathrm{}\frac{1}{m}B_m(x_i)\}`$ on its lower rows transforms the row $`\{\mathrm{\Delta }(x^{j1})(k)\}_1^{m+1}`$ into the row $`\{0,k^0,k^1,\mathrm{}k^m\}`$. What row would be transformed into the row $`\{0,`$ $`((I+\mathrm{\Delta })^{w_01}((2I+\mathrm{\Delta })^1(x^0))(1),((I+\mathrm{\Delta })^{w_01}((2I+\mathrm{\Delta })^1(x^1))(1),\mathrm{},`$ $`((I+\mathrm{\Delta })^{w_01}((2I+\mathrm{\Delta })^1(x^{m1}))(1)\}\mathrm{?}`$ Let us express $`(I+\mathrm{\Delta })^{w_01}((2I+\mathrm{\Delta })^1)(f)(x)`$ ($`f`$ any polynomial of degree at most $`m1`$) in the form $`_{k=0}^{m1}a_kf(x+k)`$ for some numbers $`a_1,a_2,\mathrm{}a_m`$. If rows $`y_1,y_2,\mathrm{}y_r`$ are transformed into rows $`z_1,z_2,\mathrm{}z_r`$, respectively, then the row $`_{i=1}^rb_iy_i`$ must be transformed into the row $`_{i=1}^rb_iz_i`$. Therefore the row $`\{_{k=0}^{m1}a_k(\mathrm{\Delta }(x^{j1}))(k+1)\}_1^{m+1}`$ is transformed into $$\{0,\underset{k=0}{\overset{m1}{}}a_k(k+1)^0,\underset{k=0}{\overset{m1}{}}a_k(k+1)^1,\mathrm{},\underset{k=0}{\overset{m1}{}}a_k(k+1)^{m1}\},$$ which is the same as $`\{0,`$ $`(I+\mathrm{\Delta })^{w_01}((2I+\mathrm{\Delta })^1(x^0))(1),(I+\mathrm{\Delta })^{w_01}((2I+\mathrm{\Delta })^1(x^1))(1),\mathrm{},`$ $`(I+\mathrm{\Delta })^{w_01}((2I+\mathrm{\Delta })^1(x^{m1}))(1)\}.`$ Now, what is $`_{k=0}^{m1}a_k(\mathrm{\Delta }(x^j))(k+1)`$ for $`0j<m`$? It is $$((I+\mathrm{\Delta })^{w_01}((2I+\mathrm{\Delta })^1(\mathrm{\Delta }(x^j))))(1).$$ Hence the weighted number of matchings of a black-edged $`n\times (m+1)`$ Aztec rectangle with dents at $`x_1,\mathrm{},x_m`$ and a black hole at $`(w_0,m)`$ is $$\frac{2^{\frac{m(m+1)}{2}}}{m!!}\left\|\begin{array}{c}((I+\mathrm{\Delta })^{w_01}(2I+\mathrm{\Delta })^1\mathrm{\Delta })(x^{j1})(1)\\ x_1^{j1}\\ x_2^{j1}\\ \mathrm{}\\ x_m^{j1}\end{array}\right\|.$$ (71) In the same way we have arrived at this result, and using it as a base case for induction, it is easy to prove the following. ###### Lemma 5. The weighted number of matchings of a white-edged $`n\times (m+d_1)`$ Aztec rectangle with teeth at $`y_1,\mathrm{}y_{m+d_11}`$ and a black hole at $`(w_0,m)`$ is $$(1)^{d_1}\frac{2^{\frac{(m+d_1)(m+d_1+1)}{2}}}{(m+d_11)!!}\left\|\begin{array}{c}((I+\mathrm{\Delta })^{w_01}(2I+\mathrm{\Delta })^{(d_1+1)}\mathrm{\Delta }^{d_1})(x^{j1})(1)\\ y_1^{j1}\\ y_2^{j1}\\ \mathrm{}\\ y_{m+d_11}^{j1}\end{array}\right\|$$ (72) and that the weighted number of matchings of a black-edged $`n\times (m+d_1)`$ Aztec rectangle with dents at $`x_1,x_2,\mathrm{},x_{m+d_11}`$ and a black hole at $`(w_0,m)`$ is $$(1)^{d_1}\frac{2^{\frac{(m+d_1)(m+d_11)}{2}}}{(m+d_11)!!}\left\|\begin{array}{c}((I+\mathrm{\Delta })^{w_01}(2I+\mathrm{\Delta })^{d_1}\mathrm{\Delta }^{d_1})(x^{j1})(1)\\ y_1^{j1}\\ y_2^{j1}\\ \mathrm{}\\ y_{m+d_11}^{j1}\end{array}\right\|.$$ (73) Notice that we have used the fact that $`(2I+\mathrm{\Delta })^1`$ and $`\mathrm{\Delta }`$ commute. They do so because $`(2I+\mathrm{\Delta })^1`$, on the domain of polynomials of degree lower than a given bound, is shorthand for a finite sum of powers of $`\mathrm{\Delta }`$. We can now attack our main objective, namely, the enumeration of matchings of an Aztec diamond with a black hole at $`(w_0+d_0,w_1)`$ and a white hole at $`(w_0,w_1+d_1)`$, where some matchings are counted as negative matchings. Let us first consider the special case $`w_0=1`$. What happens at white and black vertices $`(x,y)`$ with $`x=1`$? The hole at $`(1,w_1+d_1)`$ forces a zig-zag pattern covering all those vertices. (See figure 18.) None of the edges covering these vertices counts towards the weight of the matching. Thus, the weighted number of matchings of an Aztec diamond of order $`n`$ with a black hole at $`(1+d_0,w_1)`$ and a white hole at $`(1,w_1+d_1)`$ is the same as the weighted number of matchings of a $`(n1)\times n`$ white-edged Aztec rectangle with teeth at $`\{1,2,\mathrm{}n\}`$ and a black hole at $`(d_0,w_1)`$. By Lemma 5, this is equal to $$(1)^{nw_1}\frac{2^{\frac{n(n+1)}{2}}}{(n1)!!}\left\|\begin{array}{c}((I+\mathrm{\Delta })^{d_01}(2I+\mathrm{\Delta })^{(nw_1+1)}\mathrm{\Delta }^{nw_1})(x^{j1})(1)\\ 1^{j1}\\ 2^{j1}\\ \mathrm{}\\ (n1)^{j1}\end{array}\right\|.$$ (74) The top row is a linear combination of rows of the form $`\{\mathrm{\Delta }^k(x^{j1})(1)\}`$, $`1jn`$. All such rows with $`k>(n1)`$ vanish. All such rows with $`k<(n1)`$ are linear combinations of the rows $$\{i^{j1}\},1jn,1in1,$$ Since an $`n`$ by $`n`$ determinant whose last $`n1`$ rows are $`\{1^j\},\{2^j\},\mathrm{}\{(n1)^j\}`$ and whose first row is one of 3 vanishes, we can discard all terms of the form $$\{\mathrm{\Delta }^k(x^{j1})(1)\},1jn,k<(n1)$$ from the top row of the determinant in (74). Hence we must consider only the term of the form $$\{C\mathrm{\Delta }^{n1}(x^{j1})(1)\},1jn$$ in the top row of the determinant. Thus (74) is equal to $`(1)^{nw_1}`$ $`{\displaystyle \frac{2^{\frac{n(n+1)}{2}}}{(n1)!!}}(C(1)^{n1}(n1)!!)`$ $`=C(1)^{w_1+1}2^{\frac{n(n+1)}{2}}.`$ times the coefficient of $`\mathrm{\Delta }^{n1}`$ in $$(I+\mathrm{\Delta })^{d_01}(2I+\mathrm{\Delta })^{(n+1w_1)}\mathrm{\Delta }^{nw_1}.$$ Since, in the task of finding this coefficient, $`\mathrm{\Delta }`$ plays a purely symbolic role, we might as well use a symbol which does not denote an operator, such as $`x`$. It will also be convenient to use the following notation, due to Richard Stanley. ###### Definition 4. The coefficient of $`x^j`$ in the formal power series $`f(x)`$ on $`x`$ is denoted by $`[x^j](f(x))`$. Our task is to compute $`[x^{n1}]((1+x)^{d_01}(2+x)^{(n+1w_1)}x^{nw_1}`$. We have $`[x^{n1}]((1+x)^{d_01}(2+x)^{(n+1w_1)}x^{nw_1})`$ $`=[x^{w_11}]((1+x)^{d_01}(2+x)^{(n+1w_1)})`$ $`=[x^{w_11}]((2+x)^{w_11}{\displaystyle \frac{(1+x)^{d_01}}{(2+x)^n}})`$ $`=[(2y)^{w_11}]((2+2y)^{w_11}{\displaystyle \frac{(1+2y)^{d_01}}{(2+2y)^n}})`$ $`=2^{(w_11)}[y^{w_11}]((2+2y)^{w_11}{\displaystyle \frac{(1+2y)^{d_01}}{(2+2y)^n}})`$ $`=2^{(w_11)}2^{w_11}2^n[y^{w_11}]((1+y)^{w_11}{\displaystyle \frac{(1+2y)^{d_01}}{(1+y)^n}})`$ $`=2^n[y^{w_11}]((1+y)^{w_11}{\displaystyle \frac{(1+2y)^{d_01}}{(1+y)^n}})`$ Now, for any formal power series $`f(y)`$ on $`y`$, $`[y^{w_11}]((1+y)^{w_11}f(y))`$ $`={\displaystyle \underset{j=0}{\overset{w_11}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{w_11}{j}}\right)[y^j](f(y))`$ (75) $`={\displaystyle \underset{j=0}{\overset{w_11}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{(w_11j)+j}{j}}\right)[y^j](f(y))`$ $`={\displaystyle \underset{j=0}{\overset{w_11}{}}}(([z^{w_11j}]({\displaystyle \frac{1}{(1z)^{j+1}}}))([y^j](f(y))))`$ $`={\displaystyle \underset{j=0}{\overset{w_11}{}}}(([z^{w_11}]({\displaystyle \frac{z^j}{(1z)^{j+1}}}))([y^j](f(y))))`$ $`={\displaystyle \underset{j=0}{\overset{w_11}{}}}(([z^{w_1}]({\displaystyle \frac{z^{j+1}}{(1z)^{j+1}}}))([y^{j+1}](yf(y))))`$ $`=[z^{w_1}](({\displaystyle \frac{z}{1z}})f({\displaystyle \frac{z}{1z}}))`$ $`=[z^{w_11}](({\displaystyle \frac{1}{1z}})f({\displaystyle \frac{z}{1z}})).`$ Therefore $`((1+x)^{d_01}(2+x)^{(n+1w_1)})`$ $`=2^n[y^{w_11}]((1+y)^{w_11}{\displaystyle \frac{(1+2y)^{d_01}}{(1+y)^n}})\text{(by (}\text{3}\text{))}`$ (76) $`=2^n[z^{w_11}]({\displaystyle \frac{1}{1z}}{\displaystyle \frac{(1+\frac{2z}{1z})^{d_01}}{(1+\frac{z}{1z})^n}})`$ $`=2^n[z^{w_11}]({\displaystyle \frac{1}{1z}}(1+{\displaystyle \frac{2z}{1z}})^{d_01}(1z)^n)`$ $`=2^n[z^{w_11}]({\displaystyle \frac{1}{1z}}({\displaystyle \frac{1+z}{1z}})^{d_01}(1z)^n)`$ $`=2^n[z^{w_11}]((1+z)^{d_01}(1z)^{(n1)(d_01)})`$ An expression of this form is called a Krawtchouk polynomial (, p. 130). We have just proven ###### Lemma 6. The weighted number of matchings of an Aztec diamond of order $`n`$ with a black hole at $`(d_0+1,w_1)`$ and a white hole at $`(1,w_1+d_1)`$, $`d_0,d_11`$, is $$(1)^{w_1+1}[z^{w_11}]((1+z)^{d_01}(1z)^{(n1)(d_01)})2^{\frac{n(n1)}{2}}.$$ We can now work on the general case. Every matching of an Aztec diamond of order $`n`$ with a black hole at $`(w_0+d_0,w_1)`$ and a white hole at $`(w_0,w_1+d_1)`$ can be subdivided into one, and only one, pair of matchings of the following form. 1. The first item of the pair is a matching of a white-edged $`n\times (w_1+d_11)`$ Aztec rectangle with dents at $`y_1,y_2,\mathrm{}y_{w_1+d_12}`$ and with a black hole at $`(w_0+d_0,w_1)`$. 2. The second item of the pair is a matching of a white-edged $`n\times (n(w_1+d_11))`$ Aztec rectangle with dents at $`z_1,z_2,\mathrm{}z_{n(w_1+d_11)}`$. The numbers $`y_1,y_2,\mathrm{}y_{w_1+d_12}`$, $`z_1,z_2,\mathrm{}z_{n(w_1+d_11)}`$ and $`w_0`$ are all distinct and cover all of the interval $`\{1,2,\mathrm{},n\}`$. Thus every matching can be described as a partition of $`\{1,2,\mathrm{},n\}\{w_0\}`$ into two subsets, a matching of an Aztec rectangle with the first subset as its set of dents, and a matching of an Aztec rectangle with the second subset as its set of dents. As we stated in the previous section, we will weigh each matching $`T`$ by a factor of $`(1)^{w(T)}`$, where $`w(T)`$ is the sum of the number of edges of the form $`((i1,w_1+1),(i,w_1))`$ or $`((i,w_1+1),(i,w_1))`$ for $`1i<w_0+d_0`$ and the number of edges of the form $`((i,w_1+d_1),(i,w_1+d_11))`$ or $`((i,w_1+d_1),(i+1,w_1+d_11))`$ for $`1i<w_0`$. Thus the weight $`w(T)`$ of a matching is equal to the weight of the matching of the $`n\times (w_1+d_11)`$ white-edged Aztec rectangle which it induces, plus the number of teeth of this Aztec rectangle whose indices are lower than $`w_0`$. Hence the weighted number of matchings of an Aztec diamond of order $`n`$ with a black hole at $`(w_0+d_0,w_1)`$ and a white hole at $`(w_0,w_1+d_1)`$ is equal to the sum of $`(1)^{d_11}(1)^{t(w_0,y_1,\mathrm{},y_{w_1+d_12})}{\displaystyle \frac{2^{\frac{(w_1+d_11)(w_1+d_1)}{2}}}{(w_1+d_12)!!}}`$ (77) $`\left\|\begin{array}{c}((I+\mathrm{\Delta })^{w_0+d_01}((2I+\mathrm{\Delta })^{d_1}\mathrm{\Delta }^{d_11}(x^{j1})))(1)\\ y_1^{j1}\\ y_2^{j1}\\ \mathrm{}\\ y_{w_1+d_12}^{j1}\end{array}\right\|`$ $`{\displaystyle \frac{2^{\frac{(n(w_1+d_11))((n(w_1+d_11))+1)}{2}}}{(n(w_1+d_11)1)!!}}\left\|z_i^{j1}\right\|_1^{n(w_1+d_11)}`$ over all partitions of $`\{1,2,\mathrm{},n\}\{w_0\}`$ into two sets $$\{y_1,y_2,\mathrm{}y_{w_1+d_12}\},\{z_1,z_2,\mathrm{}z_{n(w_1+d_11)}\},$$ where $`y_1<y_2<\mathrm{}<y_{w_1+d_12}`$ and $`z_1<z_2<\mathrm{}<z_{n(w_1+d_11)}`$, and $`t(w_0,y_1,\mathrm{},y_k)`$ is equal to how many of $`y_1,y_2,\mathrm{}y_k`$ are less than $`w_0`$. We may dispose of the inconvenient $`t(w_0,y_1,\mathrm{},y_{w_1+d_12})`$ by expressing the same result as the sum of $`(1)^{d_1}{\displaystyle \frac{2^{\frac{(w_1+d_11)(w_1+d_1)}{2}}}{(w_1+d_12)!!}}\left\|\begin{array}{cc}0& ((I+\mathrm{\Delta })^{w_0+d_01}((2I+\mathrm{\Delta })^{d_1}\mathrm{\Delta }^{d_11}(x^{j2})))(1)\\ \delta _{v_1,w_0}& v_1^{j2}\\ \delta _{v_2,w_0}& v_2^{j2}\\ \mathrm{}\\ \delta _{v_{w_1+d_11},w_0}& v_{w_1+d_11}^{j2}\end{array}\right\|`$ (78) $`{\displaystyle \frac{2^{\frac{(n(w_1+d_11))((n(w_1+d_11))+1)}{2}}}{(n(w_1+d_11)1)!!}}\left\|z_i^{j1}\right\|_1^{n(w_1+d_11)}`$ over all partitions of $`\{1,2,\mathrm{},n\}`$ into two sets $`\{v_1,\mathrm{},v_{w_1+d_11}\}`$, $`z_1,\mathrm{}z_{n(w_1+d_11)}`$, where $`v_1<v_2<\mathrm{}<v_{w_1+d_11}`$ and $`z_1<z_2<\mathrm{}<z_{n(w_1+d_11)}`$, and where $`\delta _{i,j}`$ is equal to $`1`$ for $`i=j`$, $`0`$ for $`ij`$. The following lemma follows immediately from Laplace’s development of a determinant (, pp. 22-25). ###### Lemma 7. Let us have a determinant $`\left\|b_{i,j}\right\|_1^{m_1+m_2}`$, where $`b_{i,j}=c_{i,j}`$ for $`jm_1`$, $`b_{i,j}=(1)^{i1}d_{i,jm_1}`$ for $`j>m_1`$. Then $$\left\|b_{i,j}\right\|_1^{m_1+m_2}=(1)^{\frac{m_1+m_2}{2}\frac{m_1}{2}}\left\|c_{x_i,j}\right\|_1^{m_1}\left\|d_{y_i,j}\right\|_1^{m_2}$$ (79) where the sum is over all partitions of $`\{1,2,\mathrm{},m_1+m_2\}`$ into two sets $`\{x_1,x_2,\mathrm{},x_{m_1}\}`$, $`\{y_1,y_2,\mathrm{},y_{m_2}\}`$, where $`x_1<x_2<\mathrm{}<x_{m_1}`$ and $`y_1<y_2<\mathrm{}<y_{m_2}`$, and where $`z`$ is the largest integer less than or equal to $`z`$. We can now express our sum over partitions as a determinant. The weighted number of matchings of an Aztec diamond with two holes is $$(1)^{d_1}\frac{2^{\frac{(w_1+d_11)(w_1+d_1)}{2}}}{(w_1+d_12)!!}\frac{2^{\frac{(n(w_1+d_11))((n(w_1+d_11))+1)}{2}}}{(n(w_1+d_11)1)!!}(1)^{\frac{n+1}{2}\frac{w_1+d_1}{2}}\left\|d_{i,j}\right\|_1^{n+1},$$ (80) where * $`d_{i,1}=0`$ for all $`1in+1`$, $`iw_0+1`$, * $`d_{w_0+1,1}=1`$, * $`d_{1,j}=((I+\mathrm{\Delta })^{w_0+d_01}((2I+\mathrm{\Delta })^{d_1}\mathrm{\Delta }^{d_11}(x^{j2})))(1)`$ for all $`1<jw_1+d_1`$, * $`d_{1,j}=0`$ for all $`w_1+d_1<jn+1`$, * $`d_{i,j}=(i1)^{j2}`$ for $`i>1`$, $`1<jw_1+d_1`$, * $`d_{i,j}=(1)^{i1}(i1)^{j(w_1+d_1+1)}`$ for $`i>1`$, $`j>w_1+d_1`$ The task ahead is to compute the determinant $`\left\|d_{i,j}\right\|_1^{n+1}`$. For convenience, we will refer to it as $`D(w_0,d_0,w_1,d_1)`$. From (80) and Lemma 6, it follows that $`D(1,d_0,w_1,d_1)`$ $`=((1)^{d_1}{\displaystyle \frac{2^{\frac{(w_1+d_11)(w_1+d_1)}{2}}}{(w_1+d_12)!!}}{\displaystyle \frac{2^{\frac{(n(w_1+d_11))((n(w_1+d_11))+1)}{2}}}{(n(w_1+d_11)1)!!}}(1)^{\frac{n+1}{2}\frac{w_1+d_1}{2}})^1`$ $`(1)^{w_1+1}[z^{w_11}]((1+z)^{d_01}(1z)^{(n1)(d_01)})2^{\frac{n(n1)}{2}}.`$ We will reduce the general case to the special case $`w_0=0`$ by expressing $`D(w_0,d_0,w_1,d_1)D(w_01,d_0,w_1,d_1)`$, as the product of $`D(1,d_0+w_01,w_1,d_1)`$ times something else. If we take the determinant $`\left\|d_{i,j}\right\|_1^{n+1}`$ for $`D(w_01,d_0,w_1,d_1)`$, and add one to the bases of all powers, we obtain $$D(w_01,d_0,w_1,d_1)=\left\|g_{i,j}\right\|_1^{n+1}$$ (81) where * $`g_{i,1}=0`$ for all $`1in+1`$, $`iw_0`$; * $`g_{w_0,1}=1`$; * $`g_{1,j}=((I+\mathrm{\Delta })^{w_0+d_01}((2I+\mathrm{\Delta })^{d_1}\mathrm{\Delta }^{d_11}(x^{j2})))(1)`$ for all $`1<jw_1+d_1`$ ; notice how we have raised the exponent of $`(I+\mathrm{\Delta })`$ from $`w_0+d_02`$ to $`w_0+d_01`$; * $`g_{1,j}=0`$ for all $`w_1+d_1<jn+1`$; * $`g_{i,j}=i^{j2}`$ for $`i>1`$, $`1<jw_1+d_1`$, * $`g_{i,j}=(1)^{i1}i^{j(w_1+d_1+1)}`$ for $`i>1`$, $`j>w_1+d_1`$, The $`n`$-tuple $`\left\|g_{n+1,j}\right\|_2^{n+1}`$ is a linear combination of the $`n`$-tuples $`\left\|g_{i,j}\right\|_2^{n+1}`$, $`2i<n+1`$ and of the $`n`$-tuple $`\{1\}_2^{n+1}`$, which is what $`g_{1,j}`$ would be if the pattern for $`\{g_{n,j}\}_2^{n+1},`$ $`\{g_{n1,j}\}_2^{n+1},`$$`\{g_{2,j}\}_2^{n+1}`$ were continued. ###### Lemma 8. Let $`a_k`$ be the coefficient of $`x^k`$ in $`((x1)^{w_1+d_11}(x+1)^{n(w_1+d_1)+1})`$. Then $$\underset{k=0}{\overset{n}{}}a_k(k+1)^{j2}=0\text{ for }1<jw_1+d_1,$$ (82) $$\underset{k=0}{\overset{n}{}}a_k(1)^k(k+1)^{j(w_1+d_1+1)}=0\text{ for }w_1+d_1+1jn+1,$$ (83) ###### Proof. $$\mathrm{\Delta }^{w_1+d_11}x^{j2}=0\text{ for }1<jw_1+d_1$$ (84) implies $$\mathrm{\Delta }^{w_1+d_11}(\mathrm{\Delta }+2I)^{n(w_1+d_1)+1}x^{j2}=0\text{ for }1<jw_1+d_1$$ (85) If we take the value at $`x=1`$, we obtain $$\underset{k=0}{\overset{n}{}}a_k(k+1)^{j2}=0\text{ for }1<jw_1+d_1.$$ (86) Similarly, $$\mathrm{\Delta }^{n(w_1+d_1)+1}x^{j(w_1+d_1+1)}=0\text{ for }w_1+d_1+1jn+1$$ (87) implies $$(\mathrm{\Delta }+2I)^{w_1+d_11}(\mathrm{\Delta })^{n(w_1+d_1)+1}x^{j(w_1+d_1+1)}=0\text{ for }w_1+d_1+1jn+1.$$ (88) If we let $`H`$ be the operator taking $`x^j`$ to $`(x+1)^j`$, we can write $$(I+H)^{w_1+d_11}(IH)^{n(w_1+d_1)+1}x^{j(w_1+d_1+1)}=0\text{ for }w_1+d_1+1jn+1$$ (89) Clearly the coefficient of $`H^k`$ in $$(I+H)^{w_1+d_11}(IH)^{n(w_1+d_1)+1}$$ is equal to $`(1)^k`$ times the coefficient of $`H^k`$ in $$(IH)^{w_1+d_11}(I+H)^{n(w_1+d_1)+1},$$ which is equal to $`(1)^{w_1+d_11}`$ times the coefficient of $`H^k`$ in $$(HI)^{w_1+d_11}(H+I)^{n(w_1+d_1)+1},$$ that is, $`a_k`$. Hence $`(1)^{w_1+d_11}{\displaystyle \underset{k=0}{\overset{n}{}}}(1)^ka_k(k+1)^{j(w_1+d_1+1)}`$ $`={\displaystyle \underset{k=0}{\overset{n}{}}}((1)^{w_1+d_11}(1)^ka_k)(k+1)^{j(w_1+d_1+1)}`$ $`=(I+H)^{w_1+d_11}(IH)^{n(w_1+d_1)+1}x^{j(w_1+d_1+1)}(1)`$ $`=(\mathrm{\Delta }+2I)^{w_1+d_11}(\mathrm{\Delta })^{n(w_1+d_1)+1}x^{j(w_1+d_1+1)}`$ $`=(\mathrm{\Delta }+2I)^{w_1+d_11}\mathrm{\Delta }^{n(w_1+d_1)+1}x^{j(w_1+d_1+1)}`$ $`=0\text{ for }w_1+d_1+1jn+1.`$ Therefore $$\underset{k=0}{\overset{n}{}}a_k(1)^k(k+1)^{j(w_1+d_1+1)}=0\text{ for }w_1+d_1+1jn+1.$$ Hence, if we add $`a_n`$ times the $`n`$th row, $`a_{n1}`$ times the $`(n1)`$th row,…$`a_2`$ times the second row to the bottom row of $`\left\|g_{i,j}\right\|_1^{n+1}`$, we obtain the row $`\{h_j\}_1^{n+1}`$, where * $`h_0=[x^{w_01}]((x1)^{w_1+d_11}(x+1)^{n+1(w_1+d_1)})`$, * $`h_j=(1)(1)^{w_1+d_11}\text{ for }1<jw_1+d_1`$, * $`h_j=(1)(1)^{w_1+d_11}\text{ for }j>w_1+d_1`$. After switching signs and shifting the bottom row of $`D(w_01,d_0,w_1,d_1)`$ (now $`\{h_j\}_1^{n+1}`$) to the second-to-topmost place, we obtain $$D(w_01,d_0,w_1,d_1)=(1)(1)^{w_1+d_11}(1)^{n1}\left\|k_{i,j}\right\|_1^{n+1},$$ (90) where * $`k_{i,1}=0`$ for all $`1in+1`$, $`iw_0+1`$, $`i2`$; * $`k_{2,1}=(1)(1)^{w_1+d_11}[x^{w_01}]((x1)^{w_1+d_11}(x+1)^{n+1(w_1+d_1)})`$; * $`k_{w_0+1,1}=1`$; * $`k_{1,j}=((I+\mathrm{\Delta })^{w_0+d_01}((2I+\mathrm{\Delta })^{d_1}\mathrm{\Delta }^{d_11}(x^{j2})))(1)`$ for all $`1<jw_1+d_1`$; * $`k_{1,j}=0`$ for $`w_1+d_1+1jn+1`$; * $`k_{i,j}=(i1)^{j2}`$ for $`i>1`$, $`1<jw_1+d_1`$; * $`k_{i,j}=(1)^i(i1)^{j(w_1+d_1+1)}`$ for $`i>1`$, $`jw_1+d_1+1`$. We multiply the $`n(w_1+d_1)+1`$ rightmost columns by $`(1)`$, obtaining $$D(w_01,d_0,w_1,d_1)=\left\|l_{i,j}\right\|_1^{n+1},$$ (91) where * $`l_{i,1}=0`$ for all $`1in+1`$, $`iw_0+1`$, $`i2`$; * $`l_{2,1}=(1)(1)^{w_1+d_11}[x^{w_01}]((x1)^{w_1+d_11}(x+1)^{n+1(w_1+d_1)})`$; * $`l_{w_0+1,1}=1`$; * $`l_{1,j}=((I+\mathrm{\Delta })^{w_0+d_01}((2I+\mathrm{\Delta })^{d_1}\mathrm{\Delta }^{d_11}(x^{j2})))(1)`$ for all $`1<jw_1+d_1`$; * $`l_{1,j}=0`$ for $`w_1+d_1+1jn+1`$; * $`l_{i,j}=(i1)^{j2}`$ for $`i>1`$, $`1<jw_1+d_1`$; * $`l_{i,j}=(1)^{i1}(i1)^{j(w_1+d_1+1)}`$ for $`i>1`$, $`jw_1+d_1+1`$. Therefore $$D(w_0,d_0,w_1,d_1)D(w_01,d_0,w_1,d_1)=\left\|r_{i,j}\right\|_1^{n+1},$$ (92) where * $`r_{i,1}=0`$ for all $`1in+1`$, $`i2`$; * $`r_{2,1}=(1)^{w_1+d_11}[x^{w_01}]((x1)^{w_1+d_11}(x+1)^{n+1(w_1+d_1)})`$; * $`r_{1,j}=((I+\mathrm{\Delta })^{w_0+d_01}((2I+\mathrm{\Delta })^{d_1}\mathrm{\Delta }^{d_11}(x^{j2})))(1)`$ for all $`1<jw_1+d_1`$; * $`r_{1,j}=0`$ for $`w_1+d_1+1jn+1`$; * $`r_{i,j}=(i1)^{j2}`$ for $`i>1`$, $`1<jw_1+d_1`$; * $`r_{i,j}=(1)^{i1}(i1)^{j(w_1+d_1+1)}`$ for $`i>1`$, $`jw_1+d_1+1`$. This is equal to $`r_{2,1}D(1,w_0+d_01,w_1,d_1)`$ $`=(1)^{w_1+d_11}[x^{w_01}]((x1)^{w_1+d_11}(x+1)^{n+1(w_1+d_1)})`$ (93) $`D(1,w_0+d_01,w_1,d_1)`$ $`=[x^{w_01}]((1x)^{w_1+d_11}(1+x)^{(n(w_1+d_11))})`$ $`D(1,w_0+d_01,w_1,d_1)`$ Therefore $`D(w_0,d_0,w_1,d_1)`$ $`=({\displaystyle \underset{j=1}{\overset{w_01}{}}}D(j+1,d_0,w_1,d_1)D(j,d_0,w_1,d_1))`$ (94) $`+D(1,d_0,w_1,d_1)`$ $`={\displaystyle \underset{j=1}{\overset{w_01}{}}}([x^j](1x)^{w_1+d_11}(1+x)^{n(w_1+d_11)})D(1,j+d_0,w_1,d_1)`$ $`+D(1,d_0,w_1,d_1)`$ $`={\displaystyle \underset{j=0}{\overset{w_01}{}}}([x^j](1x)^{w_1+d_11}(1+x)^{n(w_1+d_11)})D(1,j+d_0,w_1,d_1)`$ By (80), it follows that the weighted number of matchings $`(1)^{w(T)}`$ of an Aztec diamond of order n with a black hole at $`(w_0+d_0,w_1)`$ and a white hole at $`(w_0,w_1+d_1)`$ is equal to $`(1)^{d_1}`$ $`{\displaystyle \frac{2^{\frac{(w_1+d_11)(w_1+d_1)}{2}}}{(w_1+d_12)!!}}{\displaystyle \frac{2^{\frac{(n(w_1+d_11))((n(w_1+d_11))+1)}{2}}}{(n(w_1+d_11)1)!!}}(1)^{\frac{n+1}{2}\frac{w_1+d_1}{2}}`$ $`{\displaystyle \underset{j=0}{\overset{w_01}{}}}([x^j](1x)^{w_1+d_11}(1+x)^{n(w_1+d_11)})D(1,j+d_0,w_1,d_1)`$ $`={\displaystyle \underset{j=0}{\overset{w_01}{}}}(([x^j](1x)^{w_1+d_11}(1+x)^{n(w_1+d_11)})`$ $`((1)^{d_1}{\displaystyle \frac{2^{\frac{(w_1+d_11)(w_1+d_1)}{2}}}{(w_1+d_12)!!}}{\displaystyle \frac{2^{\frac{(n(w_1+d_11))((n(w_1+d_11))+1)}{2}}}{(n(w_1+d_11)1)!!}}`$ $`(1)^{\frac{n+1}{2}\frac{w_1+d_1}{2}}D(1,j+d_0,w_1,d_1))`$ The term within parentheses including $`D(1,j+d_0,w_1,d_1)`$ is equal to the number of matchings of an Aztec diamond of order n with a black hole at $`(j+d_0+1,w_1)`$ and a white hole at $`(1,w_1+d_1)`$. By Lemma 6, this number is equal to $`(1)^{w_1+1}2^{\frac{n(n1)}{2}}[z^{w_11}]((1+z)^{j+d_01}(1z)^{(n1)(j+d_01)})`$ Hence $`(1)^{w(T)}`$ is equal to $`(1)^{w_1+1}2^n{\displaystyle \underset{j=0}{\overset{w_01}{}}}(`$ $`[x^j]((1x)^{w_1+d_11}(1+x)^{n(w_1+d_11)})`$ $`[z^{w_11}]((1+z)^{j+d_01}(1z)^{(n1)(j+d_01)})),`$ From this and from (2.4) the result we have sought follows immediately. ###### Proposition 9. The entry in the inverse of the Kasteleyn matrix of an Aztec diamond of order $`n`$ corresponding to a black square at $`(w_0+d_0,w_1)`$ and a white square at $`(w_0,w_1+d_1)`$, $`d_0,d_1>0`$, is $`(1)^{d_0+d_1+w_1}2^n{\displaystyle \underset{j=0}{\overset{w_01}{}}}(`$ $`[x^j]((1x)^{w_1+d_11}(1+x)^{n(w_1+d_11)})`$ (95) $`[z^{w_11}]((1+z)^{j+d_01}(1z)^{(n1)(j+d_01)})),`$ where $`[x^j](p(x))`$ is the coefficient of $`x^j`$ in the polynomial $`p(x)`$. (Alternatively, this can be called the value of the coupling function of the Aztec diamond of order $`n`$ at the black hole $`(w_0+d_0,w_1)`$ and the white hole $`(w_0,w_1+d_1)`$.) In order to deal with the cases $`d_00`$ and $`d_10`$, we merely need to flip the Aztec diamond so as to make $`d_0,d_1>0`$ and compute the weighted number of tilings in the manner we have described. Of course, we have to account for the fact that the weighting has to be computed differently. For $`d_00`$, we also have to express $`D(w_0,d_0,w_1,d_1)`$ as $$\underset{j=w_0}{\overset{n}{}}D(j+1,d_0,w_1,d_1)D(j,d_0,w_1,d_1)$$ and not as $$D(1,d_0,w_1,d_1)+\underset{j=1}{\overset{w_01}{}}D(j+1,d_0,w_1,d_1)D(j,d_0,w_1,d_1))$$ as we did in (94). These are the only two details worth mention in the otherwise trivial derivation of the following result from Proposition 9. ###### Corollary 10. The coupling function of the Aztec diamond of order $`n`$ at the black square $`(w_0+d_0,w_1)`$ and the white square $`(w_0,w_1+d_1)`$ is $`(1)^{d_0+d_1+w_1}2^n{\displaystyle \underset{j=0}{\overset{w_01}{}}}(`$ $`[x^j]((1x)^{w_1+d_11}(1+x)^{n(w_1+d_11)})`$ (96) $`[z^{w_11}]((1+z)^{j+d_01}(1z)^{(n1)(j+d_01)}))`$ for $`d_0>0`$, $`(1)^{d_0+d_1+w_1}2^n({\displaystyle \underset{j=w_0}{\overset{n}{}}}(`$ $`[x^j]((1x)^{w_1+d_11}(1+x)^{n(w_1+d_11)})`$ (97) $`[z^{w_11}]((1+z)^{j+d_01}(1z)^{(n1)(j+d_01)})))`$ for $`d_00`$. When we take a minor of the inverse Kasteleyn matrix, the factors $`(1)^{d_0+d_1+w_1}`$, multiplied, give the same product in every term of the expression of the minor in a form such as $`_\pi \mathrm{sgn}(\pi )_{i=1}^kc_{i,\pi (i)}`$. Thus we can leave them out, and our main result follows. ###### Theorem 11. The probability of a pattern covering white squares $`v_1,v_2,\mathrm{}v_k`$ and black squares $`w_1,w_2,\mathrm{}w_k`$ of an Aztec diamond of order $`n`$ is equal to the absolute value of $$\left\|c(v_i,w_j)\right\|_{i,j=1,2,\mathrm{}k}.$$ The coupling function $`c(v,w)`$ at white square $`v`$ and black square $`w`$ is $$2^n\underset{j=0}{\overset{x_i1}{}}\mathrm{kr}(j,n,y_i1)\mathrm{kr}(y_i1,n1,n(j+x_ix_i))$$ for $`x_i>x_i`$ and $$2^n\underset{j=x_i}{\overset{n}{}}\mathrm{kr}(j,n,y_i1)\mathrm{kr}(y_i1,n1,n(j+x_ix_i))$$ for $`x_ix_i`$, where $`(x_i,y_i)`$ and $`(x_i,y_i)`$ are the coordinates of $`v`$ and $`w`$, respectively, in the coordinate system in figure 17, and the Krawtchouk polynomial $`\mathrm{kr}(a,b,c)`$ is the coefficient of $`x^a`$ in $`(1x)^c(1+x)^{bc}`$. ## 4. In perspective Proposition 9 is valuable in itself, in that it gives us an efficient algorithm for computing an arbitrary entry in the inverse Kasteleyn matrix of the Aztec diamond. Figures 19 to 24 show the absolute value of the entry as a function of $`w_0`$ and $`w_1`$, for fixed $`d_0`$ and $`d_1`$. As we showed in section 2, given a pattern consisting of $`k`$ vertices, we can compute the probability of its occurence at any point in a Aztec diamond of given order by computing $`(k/2)^2`$ entries of the inverse Kasteleyn matrix. Thus, for a fixed pattern, the time required for computing its probability is equal to a constant times the time required for computing an entry of the inverse Kasteleyn matrix using Proposition 9. Whether computation time grows quadratically on the order of the Aztec diamond, or somewhat faster, depends on whether multiplying integers is assumed to take constant time. What is clear is that we now have an algorithm that is much more efficient than computing the entries of an inverse Kasteleyn matrix by actually inverting the matrix or computing minors. Proposition 9 gives us an expression that is more closed than an entry in the inverse of a Kasteleyn matrix. What do we mean by this? There are few tools available that would allow us to obtain asymptotic expressions for a sequence of entries in a sequence of inverses of arbitrary matrices. For finding the asymptotics of sums such as (95), however, there are many well-developed analytical techniques. At the time of this writing, Henry Cohn is working on some minor problems involved in applying the saddle-point technique to the asymptotics of (95). Once he superates these difficulties (something that seems to be about to happen), the goals set in the introduction will have been achieved completely. ## 5. Acknowledgements The author would like to thank Matthew Blum for converting the output of the urban-renewal program ren.c to figures 4 to 13. He would also like to thank James Propp and Ira Gessel for their support, and Henry Cohn, in advance, for finishing his work on the asymptotics of the herein described results and for writing the still unwritten continuation to this paper.
warning/0007/cond-mat0007110.html	ar5iv	text	# Depinning Transition of a Two Dimensional Vortex Lattice in a Commensurate Periodic Potential \[ ## Abstract We use Monte Carlo simulations of the 2D one component Coulomb gas on a triangular lattice, to study the depinning transition of a 2D vortex lattice in a commensurate periodic potential. A detailed finite size scaling analysis indicates this transition to be first order. No significant changes in behavior were found as vortex density was varied over a wide range. \] The theory of defect mediated melting of a two dimensional (2D) solid, introduced by Kosterlitz and Thouless (KT) and developed by Nelson and Halperin (NH) and Young in the late 70’s, has remained a topic of active investigation. Numerous numerical studies have reported conflicting results as to whether the 2D melting transition is indeed a 2nd order KT type transition, or whether it is 1st order . The most recent results have supported the KT scenario . In the same paper in which they developed this theory of 2D melting in a continuum, Nelson and Halperin also considered the case of a 2D solid in a commensurate periodic potential. In this case, they argued that the 2D solid would have two distinct phases: a “pinned” solid with long range translational correlations at low temperatures, and a “floating” solid with algebraic translational correlations, similar to that found in the continuum, at intermediate temperatures. By mapping the problem to a vector Coulomb gas, they argued that this depinning transition was a 2nd order KT type transition similar to that of melting, with a universal discontinuous jump in the exponent $`\eta `$ characterizing the algebraic correlations of the floating solid phase. Despite the wide attention given to the continuum 2D melting problem, this depinning transition has been very little studied. Only recently has the floating solid phase been observed in numerical studies of 2D vortex lattices , and related XY models . In this paper we present the first detailed finite size scaling analysis of this depinning transition. We treat the specific case of logarithmically interacting 2D vortices. This is not only a system of considerable recent interest in connection with high temperature superconductors, but also has a unique advantage for numerical simulations: it is an incompressible system. We can simulate at constant density, yet there will still be a sharp transition temperature, rather than the finite temperature interval of coexisting phases that one would have for other interactions (should the transition be 1st order). We can therefore avoid the controversy, that arose in 2D melting simulations, as to whether one should use a constant pressure rather than a constant volume ensemble. For this vortex system, our results are quantitatively consistent with a 1st order depinning transition. Our model is the one component classical Coulomb gas on a triangular lattice, given by the Hamiltonian, $$=\frac{1}{2}\underset{i,j}{}(n_if)G_{ij}(n_jf).$$ (1) Here $`n_i=0`$, $`1`$ is the integer charge at site $`i`$ of an $`L\times L`$ periodic triangular lattice; $`f`$ is a uniform background charge density; charge neutrality fixes the number of integer charges to $`_in_i=fL^2`$. $`G_{ij}`$ is the 2D Coulomb potential for a discrete triangular lattice with periodic boundary conditions , and the sum is over all pairs of sites. $`G_{ij}\mathrm{ln}\|𝐫_i𝐫_j\|`$ for distances large compared to the lattice spacing, but small compared to $`L`$. The charges $`n_i`$ model logarithmically interacting vortices in the phase of a superconducting wavefunction, $`f`$ is the number of applied magnetic flux quanta per unit cell of the triangular lattice, and restricting the vortices to the sites of the triangular lattice models the periodic pinning potential . We study the Hamiltonian (1) using Monte Carlo (MC) simulations. Details of our simulation methods follow those described in Ref.. Our main results are for the fixed vortex density of $`f=1/100`$. To decide the order of the depinning transition at $`T_c`$, we study the finite size scaling of the energy density histogram $`P(E)`$ measured close to $`T_c`$. For a first order transition, one expects to find a bimodal $`P(E)`$, with peaks that sharpen to two separated $`\delta `$functions as $`L\mathrm{}`$; these give the differing energies of the two coexisting phases. In Fig. 1 we plot the normalized $`P(E)`$ for sizes $`L=120`$, $`160`$, and $`200`$, near $`T_c0.00229`$. Analysis of the vortex structure function, $$S(𝐤)=\frac{1}{fL^2}\underset{i,j}{}\mathrm{e}^{i𝐤(𝐫_i𝐫_j)},$$ (2) clearly shows that the lower (upper) energy peak corresponds to the pinned (floating) vortex lattice. For $`L=120`$, $`160`$, and $`200`$ respectively, we have carried out $`0.77`$, $`1.5`$ and $`4.7\times 10^8`$ MC passes though the lattice to compute averages. These resulted in $`2219`$, $`1186`$ and $`841`$ hops, respectively, between the two peaks. Having a large number of hops is important to achieve good equilibration of the relative weights of the two peaks (we have also done simulations for $`L=80`$, $`100`$, $`140`$, and $`180`$). Qualitatively, we see that these peaks sharpen as $`L`$ increases, consistent with a 1st order transition. We now make this observation quantitative using several different criteria. For an infinite system, the bimodal $`P(E)`$ should exist only precisely at $`T_c`$, where the two phases coexist. For finite $`L`$, this coexistence region where $`P(E)`$ is noticeably bimodal persists over a finite temperature interval $`\mathrm{\Delta }T`$. Since the relative weight of each peak is determined by the total free energy difference between the two phases, $`\mathrm{\Delta }FL^d`$, we expect the scaling $`\mathrm{\Delta }T1/\mathrm{\Delta }F=1/L^d=1/L^2`$ in 2D. Using standard methods to extrapolate the histograms of Fig. 1 to nearby temperatures, we define the upper (lower) limit of the coexistence region, $`T_{c\mathrm{max}(\mathrm{min})}`$, as the temperature at which the height of the lower (upper) peak has decreased to $`1\%`$ the height of the upper (lower) peak. In Fig. 2 we plot the results. We see that $`T_{c\mathrm{max}}`$ and $`T_{c\mathrm{min}}`$ converge to a common value as $`L`$ increases, and that $`\mathrm{\Delta }T1/L^2`$ as expected for a 1st order transition. Next, we consider in detail how the two peaks sharpen as $`L`$ increases. For a 1st order transition at $`T_c`$, the width $`\sigma `$ of each peak is determined by ordinary non-critical finite size fluctuations giving, $`\sigma 1/L^{d/2}=1/L`$ in 2D. The total width of the bimodal $`P(E)`$, $`\sigma _{\mathrm{tot}}`$, approaches a constant (this is equivalent to the familiar observation that the specific heat $`cL^d\sigma _{\mathrm{tot}}/TL^d`$ at a 1st order transition). To verify these scaling behaviors, we need to deconvolve the $`P(E)`$ of Fig. 1 into two separate peaks, and to determine $`T_c(L)`$ by the criteria that the two peaks subtend equal areas . For large enough $`L`$, we expect each peak to have a Gaussian shape, and we have found this to give a good fit for the floating solid peak. The pinned solid peak however lies too close to the ground state energy $`E_0`$; even for our biggest size $`L=200`$, the states in this peak correspond to a few discrete excitations above the ground state. We therefore fit the pinned solid peak to the ad hoc form, $`(EE_0)\mathrm{exp}[(EE_1)^2/2\sigma ^2]`$, where $`E_1`$ and $`\sigma `$ are fitting parameters. Using data from only the far side of each peak, we fit $`P(E)`$ to the sum of a Gaussian and the modified Gaussian above, to get the fitted curves shown in Fig. 1. From these fits we determine the probabilities that a state with a given value of $`E`$ belongs to the pinned phase, the floating phase, or neither (the later being the “transition” states, consisting of large domains of one phase in a background of the other, which give rise to the transitions between the two phases). We then go though the ensemble of states that enter into our averages, and probabilistically assign each state to the pinned phase, the floating phase, or neither. With this deconvolution, we then extrapolate to the $`T_c(L)`$ which gives equal weight to the two phases, and we can then compute the widths $`\sigma _{\mathrm{float}}`$, $`\sigma _{\mathrm{pin}}`$, and $`\sigma _{\mathrm{tot}}`$ at this temperature. Our result for $`T_c(L)`$ is shown in Fig. 2$`a`$. Our results for the widths are shown in Fig. 3. We find $`\sigma _{\mathrm{tot}}`$ constant, and $`\sigma _{\mathrm{float}}`$, $`\sigma _{\mathrm{pin}}1/L`$, as expected for a first order transition . We can also use this decomposition to compute the entropy jump per vortex at the transition, $`\mathrm{\Delta }s=(E_{\mathrm{float}}E_{\mathrm{pin}})/(fT_c)`$. Our results, shown in Fig. 4, give a $`\mathrm{\Delta }s(L)`$ that saturates to a constant as $`L`$ increases, again as expected for a 1st order transition. Next we consider the translational correlations. In the floating solid phase correlations are algebraic , $`\mathrm{exp}[i𝐊(𝐫_i𝐫_j)]\|𝐫_i𝐫_j\|^{\eta _K(T)}`$, with $`𝐊`$ a reciprocal lattice vector, and $`\eta _K\eta \|𝐊\|^2/\|𝐊_1\|^2`$, with $`𝐊_1`$ the smallest non-zero $`𝐊`$. Substituting this into the structure function, Eq.(2) results in, $$\frac{S(𝐊)}{L^2}L^{\eta _K}\mathrm{e}^{(\eta \mathrm{ln}L)\|𝐊\|^2/\|𝐊_1\|^2}.$$ (3) According to the NH theory , $`\eta `$ should take a discontinuous jump from zero in the pinned solid to a universal finite value at $`T_c`$ in the floating solid. For an incompressible system such as ours, in which the bulk modulus $`\lambda =\mathrm{}`$, the predicted jump is, $$\eta (T_c^+)=4f,$$ (4) where $`f`$ is the vortex density. From Eq.(3) we see that $`S(𝐊_1)/L^2`$ vs. $`L`$ plotted on a log-log scale should yield a straight line of negative slope $`\eta (T)`$. In Fig. 5a we show such plots for several values of $`T`$ near $`T_c`$, using $`S(𝐊_1)`$ averaged over all configurations encountered in the simulation. The straight lines are a least squares fit to the assumed algebraic form. We see that the fit is not very good, except perhaps at the highest temperatures. If, however, we average $`S(𝐊_1)`$ separately over only the states in the floating solid phase, and over only the states in the pinned solid phase, we see the expected behavior as shown in Fig. 5b. In the pinned phase, the curves saturate to a finite value as $`L`$ increases, reflecting the long range order of this phase. In the floating phase, the curves give good straight line fits, showing algebraically decaying correlations. From these fits, we extract the exponent $`\eta `$ of the floating phase, which we plot vs. $`T`$ in Fig. 6. We see that $`\eta (T_c^+)`$ is close to, but noticeably above, the NH prediction of Eq.(4). Since the NH mapping onto a vector Coulomb gas $`inverts`$ the temperature scale, this is consistent with a 1st order transition pre-empting the NH defect unbinding transition, with Eq.(4) serving as a lower bound on $`\eta (T)`$. The results above are all consistent with a 1st order transition rather than the 2nd order NH prediction. However one can question whether this is an artifact of the particular density $`f=1/100`$ that we have used. Hattel and Wheatley in particular have argued that, as $`f`$ decreases in the vortex system, the core energy of the defects that lead to depinning increases, and hence even if depinning is 1st order at some $`f`$, it must ultimately become 2nd order as $`f`$ decreases. To check this prediction we have carried out simulations for a variety of densities $`f=49`$, $`64`$, $`100`$, $`196`$, $`400`$. Rather than do finite size scaling for each case, we study only systems with a fixed total number of vortices $`N_v`$, corresponding to system sizes $`L=\sqrt{N_v/f}`$. We choose $`N_v=144`$ because our finite size studies of $`f=1/100`$ indicated this to be sufficiently large to be within the asymptotic large $`L`$ limit. Analyzing the data from these simulations in exactly the same manner as described above, we show in Fig. 7 our result for the entropy jump per vortex at depinning, $`\mathrm{\Delta }s`$, as a function of density $`f`$. We see that $`\mathrm{\Delta }s`$ is only weakly dependent on $`f`$, extrapolating to a finite $`\mathrm{\Delta }s0.60`$ as $`f0`$. Finally, we consider the translational correlations for the different $`f`$. Since we have simulated only for fixed values of $`fL^2`$, we cannot extract $`\eta `$ by scaling with $`L`$ as in Fig. 5b. However we can estimate $`\eta `$ as follows. From Eq.(3) we expect a Gaussian envelope for the peaks $`S(𝐊)`$ as a function of $`\|𝐊\|`$; the width gives a measure of $`\eta `$. This envelope is not purely Gaussian; a correction exists due to a $`\|𝐊\|`$ dependence of a prefactor in the $`L^{\eta _K}`$ scaling of Eq.(3). Making a simple approximate correction for this effect, as done in Ref. , we plot our estimate for the correlation exponent, $`\eta ^{}`$, in Fig. 8. We show results using $`S(𝐊)`$ averaged over all states, only the floating states, and only the pinned states, plotting the data in scaled units of $`\eta ^{}/f`$ vs. $`T/T_c`$. For $`f=1/100`$, our estimate $`\eta _{\mathrm{float}}^{}`$ is slightly larger than the more correct determination in Fig. 6. However the main point of Fig. 8 is to observe that the estimates $`\eta ^{}`$ all collapse to a common curve as the density $`f`$ decreases. There is no evidence for a continuing decrease in the value of $`\eta (T_c^+)`$ as $`f`$ decreases, as might be the case if the 1st order transition was weakening and the NH prediction of Eq.(4) was approached more closely. To conclude, a finite size scaling analysis for the density $`f=1/100`$ is strongly consistent with the depinning transition being 1st order. We further find no evidence that this 1st order transition changes in any significant way if the density $`f`$ is decreased. We cannot rule out the possibility of different behavior for softer interactions than the logarithm considered here. We would like to thank D. R. Nelson and M. Franz for very helpful conversations. This work was supported by the Engineering Research Program of the Office of Basic Energy Sciences at the Department of Energy, grant DE-FG02-89ER14017.
warning/0007/math0007111.html	ar5iv	text	# Discreteness of spectrum for the magnetic Schrödinger operators. I. ## 1 Introduction and main results ### 1.1 Notations and preliminaries The main object of this paper is a magnetic Schrödinger operator in $`\text{}^n`$ and its generalizations. In the simplest case it has the form (1.1) $$H_{a,V}=\underset{j=1}{\overset{n}{}}P_j^2+V,$$ where (1.2) $$P_j=\frac{1}{i}\frac{}{x^j}+a_j,$$ and $`a_j=a_j(x)`$, $`V=V(x)`$, $`x=(x^1,\mathrm{},x^n)\text{}^n`$. We assume that $`a_j`$ and $`V`$ are real-valued functions. For simplicity we will assume now that $`a_jC^1(\text{}^n)`$, $`VL_{loc}^{\mathrm{}}(\text{}^n)`$ (the later means that $`V`$ is measurable and locally bounded). Then $`H_{a,V}`$ is well defined on $`C_c^{\mathrm{}}(\text{}^n)`$ (the set of all complex-valued $`C^{\mathrm{}}`$ functions with a compact support in $`\text{}^n`$), and it is an unbounded symmetric operator in $`L^2(\text{}^n)`$. We will always impose (explicitly or implicitly) conditions which insure that the operator $`H_{a,V}`$ is essentially self-adjoint. For instance, the condition $`V0`$ is sufficient (see e.g. H. Leinfelder and C.G. Simader where this is proved under most general local regularity conditions on $`a_j`$ and $`V`$). But some negative potentials (even mildly blowing up to $`\mathrm{}`$ when $`x\mathrm{}`$) will do as well - see e.g. T. Ikebe and T. Kato , A. Iwatsuka , M. Shubin for several versions of this fact. For the sake of convenience of the reader we give in Section 2 a very short proof of the fact which is important for us: the semi-boundedness below for the operator $`H_{a,V}`$ (on $`C_c^{\mathrm{}}(\text{}^n)`$) implies that it is essentially self-adjoint (this is an extension of the Povzner–Wienholtz theorem – see , , , ) and it is proved for the case of operators on any complete Riemannian manifold in . We will also denote by $`H_{a,V}`$ the corresponding self-adjoint operator in $`L^2(\text{}^n)`$. Actually the condition $`VL_{loc}^{\mathrm{}}(\text{}^n)`$ is not necessary for our study. For example, it will be sufficient to have $`VL_{loc}^2(\text{}^n)`$ and locally semi-bounded below. Moreover, it is sufficient to have $`VL_{loc}^1(\text{}^n)`$ and locally semi-bounded below. In this case we have to impose conditions which guarantee that the corresponding quadratic form $`h_{a,V}`$ is semi-bounded below and consider the operator defined by this form. This does not make any difference in our arguments because we can work with the quadratic form only. But we prefer not to get the reader distracted by unimportant details. On the other hand working with capacities usually requires $`V`$ to be locally semi-bounded below. So this condition often can not be removed. We will say that a self-adjoint operator $`H`$ in a Hilbert space $``$ has a discrete spectrum if its spectrum consists of isolated eigenvalues of finite multiplicities. It follows that the only accumulation points of these eigenvalues can be $`\pm \mathrm{}`$. Equivalently we may say that $`H`$ has a compact resolvent. Our goal will be to provide conditions (mainly sufficient but sometimes necessary and sufficient) for the discreteness of the spectrum of $`H_{a,V}`$. We will abbreviate the discreteness of the spectrum of $`H_{a,V}`$ by writing $`\sigma =\sigma _d`$. Actually under the conditions we will impose the operator will be always semi-bounded below, so the only accumulation point of the eigenvalues will be $`+\mathrm{}`$. Let us recall first some facts concerning the Schrödinger operator $`H_{0,V}=\mathrm{\Delta }+V`$ without magnetic field (i.e. the operator (1.1) with $`a=0`$). It is a classical result of K. Friedrichs (see also e.g. , Theorem XIII.67, or , Theorem 3.1) that the condition (1.3) $$V(x)+\mathrm{}\text{as}x\mathrm{}$$ implies $`\sigma =\sigma _d`$ (for $`H_{0,V}`$). Now assume that (1.4) $$V(x)C$$ with a constant $`C`$, i.e. the potential $`V`$ is semi-bounded below. Without loss of generality we can assume then that $`V0`$. Let us formulate a simple necessary condition for the discreteness of the spectrum. Denote by $`B(x,r)`$ the open ball with the radius $`r>0`$ and the center at $`x\text{}^n`$. Then $`\sigma =\sigma _d`$ for $`H_{0,V}`$ implies that for every fixed $`r>0`$ (1.5) $$_{B(x,r)}V(y)𝑑y+\mathrm{}\text{as}x\mathrm{}.$$ This observation was made in a remarkable paper by A. Molchanov who proved that in case $`n=1`$ this condition is in fact necessary and sufficient (assuming the semi-boundedness of the potential $`V`$). More importantly, A. Molchanov found a necessary and sufficient condition for $`\sigma =\sigma _d`$ to hold, again assuming (1.4). This condition is intermediate between (1.3) and (1.5). It is formulated in terms of the Wiener capacity which we will denote $`\mathrm{cap}`$ (see e.g. for necessary properties of the capacity and more details). In case $`n=2`$ the capacity of a set $`FB(x,r)`$ is always taken relative to a ball $`B(x,R)`$ of a fixed radius $`R>r`$. (Expositions of Molchanov’s and more general results can be found in .) A. Molchanov proved that $`H_{0,V}`$ has a discrete spectrum if and only if there exist $`c>0`$ and $`r_0>0`$ such that for any $`r(0,r_0)`$ ($`M_c`$) $$\underset{𝐹}{inf}\left\{_{B(x,r)F}V(y)𝑑y\right\|\mathrm{cap}(F)c\mathrm{cap}(B(x,r))\}+\mathrm{}\text{as}x\mathrm{}.$$ (In this case we will say that the function $`V`$ satisfies $`(M_c)`$, or that the Molchanov condition $`(M_c)`$ holds for $`V`$. Later we will impose this condition on some other functions.) Note that $`(M_c)`$ implies $`(M_c^{})`$ for any $`c^{}<c`$. Hence we can equivalently write that $`(M_c)`$ is satisfied for all $`c(0,c_0)`$ with a positive $`c_0`$. In fact A. Molchanov provides a particular value of $`c`$ (e.g. $`c=2^{2n6}`$ would do – see ), though it is by no means precise. Note also that $`\mathrm{cap}(B(x,r))`$ can be explicitly calculated. It equals $`c_nr^{n2}`$ if $`n3`$. If $`n=2`$ then $`\mathrm{cap}(B(x,r))`$ asymptotically equals $`c_2\left(\mathrm{log}(1/r)\right)^1`$ as $`r0`$. Hence in the formulation of the Molchanov condition $`(M_c)`$ we can replace $`\mathrm{cap}(B(x,r))`$ by $`r^{n2}`$ if $`n3`$ and by $`\left(\mathrm{log}(1/r)\right)^1`$ if $`n=2`$. A simple argument given in (see also Sect.3 of this paper) shows that if $`H_{0,V}`$ has a discrete spectrum, then the same is true for $`H_{a,V}`$ whatever the magnetic potential $`a`$. Therefore the condition ($`M_c`$) together with (1.4) is sufficient for the discreteness of spectrum of $`H_{a,V}`$. This means that a magnetic field can only improve the situation from our point of view. Papers by J. Avron, I. Herbst and B. Simon , A. Dufresnoy and A. Iwatsuka provide some quantitative results which show that even in case $`V=0`$ the magnetic field can make the spectrum discrete. In this paper we will improve the results of the above mentioned papers. In particular we will add the capacity into the picture, so in many cases our conditions become necessary and sufficient in case when there is no magnetic field, i.e. when $`a=0`$. We will also make both electric and magnetic fields work together to achieve the discreteness of spectrum. Unfortunately we can not provide efficient necessary and sufficient conditions of the discreteness of the spectrum when both fields are present (or even if the magnetic field only is present). The conditions which we can give always contain some hypotheses which are hard to check (unless $`a=0`$ when they become trivial). Some of these conditions will be discussed in a future continuation of this paper. It is convenient to consider the magnetic potential as a 1-form $`a`$ with components $`a_j`$: (1.6) $$a=a_jdx^j,$$ where we use the Einstein summation convention (i.e. the summation over all repeated suffices is understood). Now the magnetic field is a 2-form $`B`$ which is defined as (1.7) $$B=da=\frac{a_j}{x^k}dx^kdx^j=\frac{1}{2}B_{jk}dx^jdx^k,$$ where $`B_{kj}=B_{jk}`$. Obviously (1.8) $$B=\underset{j<k}{}B_{jk}dx^jdx^k,$$ and (1.9) $$B_{jk}=\frac{a_k}{x^j}\frac{a_j}{x^k},$$ so in the standard vector analysis notation $`B=\text{curl}a`$. The functions $`B_{jk}`$ will be called components of the magnetic field $`B`$. In case $`n=2`$ the magnetic field has essentially one non-trivial component $`B_{12}=B_{21}`$ and in this case we will denote $`B=B_{12}`$. We will need a norm of $`B`$ which is defined as (1.10) $$\|B\|=\left(\underset{j<k}{}\|B_{jk}\|^2\right)^{1/2}.$$ Note that the components of the magnetic field show up in the commutation relations (1.11) $$[P_j,P_k]=\frac{1}{i}B_{jk},$$ where $`[A,B]=ABBA`$ for operators $`A,B`$ in the same Hilbert space (at the moment we assume that all operations are performed on the same domain, e.g. $`C_c^{\mathrm{}}(\text{}^n)`$ for $`P_j`$ and $`P_k`$). The relation (1.11) allows to apply uncertainty principle type arguments in investigating the spectrum. An important fact is the gauge invariance of the spectrum of $`H_{a,V}`$: this spectrum does not depend of the choice of the magnetic potential $`a`$ provided the magnetic field $`B`$ is fixed. Namely, if $`a,a^{}`$ are two magnetic potentials with $`da=da^{}=B`$, then $`\sigma (H_{a,V})=\sigma (H_{a^{},V})`$ for any $`V`$. To see this note that by the Poincaré Lemma (see e.g. 4.18 in ) we have $`a^{}=a+d\varphi `$, where $`\varphi C^1(\text{}^n)`$ is defined up to an additive constant and can be assumed real-valued. Then the corresponding operators $$P_j^{}=\frac{1}{i}\frac{}{x^j}+a_j^{}$$ are related with $`P_j`$ by the formulae $$P_j^{}=e^{i\varphi }P_je^{i\varphi }.$$ Therefore $$H_{a^{},V}=e^{i\varphi }H_{a,V}e^{i\varphi },$$ and the operators $`H_{a^{},V}`$ and $`H_{a,V}`$ are unitarily equivalent, hence have the same spectra. For the weakest requirements on the magnetic potential $`a`$ the gauge invariance was established by H. Leinfelder . By this reason it is more natural for spectral theory to formulate the conditions on $`H_{a,V}`$ in terms of $`B,V`$ rather than $`a,V`$. Let us assume that we are given a magnetic potential $`a=a_jdx^j`$, $`a_jC^1(\text{}^n)`$. For a function $`uC^1(\text{}^n)`$ (or, more generally, for a locally Lipschitz function) define magnetic differential as (1.12) $$d_au=du+iua\mathrm{\Lambda }^1(\text{}^n).$$ It is also convenient to identify this complex-valued 1-form with the corresponding complex vector field which is called magnetic gradient: (1.13) $$_au=(\frac{u}{x^1}+ia_1u,\mathrm{},\frac{u}{x^n}+ia_nu)=(iP_1u,\mathrm{},iP_nu).$$ We will denote by $`\|\|`$ the usual euclidean norm of vectors or 1-forms. ### 1.2 Localization Necessary and sufficient conditions of discreteness of spectrum for $`H_{a,V}`$ can be formulated in terms of bottoms of Dirichlet or Neumann spectra on balls of a fixed radius or on cubes of a fixed size. We will call these facts localization results. The first result of this kind about usual Schrödinger operators (without magnetic field) is due to A. Molchanov (see also for a more general theorem on manifolds). A. Iwatsuka proved a localization theorem for magnetic Schrödinger operators. The bottoms of Dirichlet and Neumann spectra for the operator $`H_{a,V}`$ in an open set $`\mathrm{\Omega }\text{}^n`$ are defined in terms of its quadratic form which we will denote $`h_{a,V}`$: (1.14) $$h_{a,V}(u,u)=_\mathrm{\Omega }(\|_au\|^2+V\|u\|^2)𝑑x.$$ This form is well defined e.g. for all $`uL^2(\mathrm{\Omega })`$ such that $`P_juL^2(\mathrm{\Omega })`$, $`j=1,\mathrm{},n`$, the derivatives are understood in the sense of distributions, and $`V\|u\|^2L^1(\mathrm{\Omega })`$. In particular $`h_{a,V}(u,u)`$ is well defined for all $`uC_c^{\mathrm{}}(\mathrm{\Omega })`$. Denote by $`(u,v)`$ the usual scalar product of $`u`$ and $`v`$ in $`L^2(\mathrm{\Omega })`$. It is easy to check that the following gauge invariance relation holds: (1.15) $$h_{a+d\varphi ,V}(u,u)=h_{a,V}(e^{i\varphi }u,e^{i\varphi }u),$$ for any $`\varphi C^1(\text{}^n)`$ and $`u`$ as above. Now we can define (1.16) $$\lambda (\mathrm{\Omega };H_{a,V})=\underset{𝑢}{inf}\left\{\frac{h_{a,V}(u,u)}{(u,u)},uC_c^{\mathrm{}}(\mathrm{\Omega })0\right\},$$ (1.17) $$\mu (\mathrm{\Omega };H_{a,V})=\underset{𝑢}{inf}\{\frac{h_{a,V}(u,u)}{(u,u)},u(C^{\mathrm{}}(\mathrm{\Omega })\{0\})L^2(\mathrm{\Omega }),_\mathrm{\Omega }V\|u\|^2dx>\mathrm{}\},$$ i.e. $`\lambda (\mathrm{\Omega };H_{a,V})`$ and $`\mu (\mathrm{\Omega };H_{a,V})`$ are bottoms of the Dirichlet and Neumann spectra (of $`H_{a,V}`$) respectively, in the usual variational understanding (see e.g. , ). The relation (1.15) obviously implies that the numbers $`\lambda (\mathrm{\Omega };H_{a,V})`$ and $`\mu (\mathrm{\Omega };H_{a,V})`$ are gauge invariant, i.e. they do not change if we replace $`a`$ by $`a+d\varphi `$ for any $`\varphi C^1(\text{}^n)`$. The following theorem slightly extends the result of A. Iwatsuka removing the requirement $`V0`$ and allowing non-continuous minorant functions. ###### Theorem 1.1 The following conditions are equivalent: $`(a)`$ $`H_{a,V}`$ is essentially self-adjoint, semi-bounded below and has a discrete spectrum; $`(b)`$ $`\lambda (B(x,r);H_{a,V})+\mathrm{}`$ as $`x\mathrm{}`$ for any fixed $`r>0`$; $`(c)`$ there exists $`r>0`$ such that $`\lambda (B(x,r);H_{a,V})+\mathrm{}`$ as $`x\mathrm{}`$; $`(d)`$ there exists a real valued function $`\mathrm{\Lambda }C(\text{}^n)`$ such that $`\mathrm{\Lambda }(x)+\mathrm{}`$ as $`x\mathrm{}`$ and the operator inequality (1.18) $$H_{a,V}\mathrm{\Lambda }(x)$$ holds in the sense of quadratic forms $`(`$on $`C_c^{\mathrm{}}(\text{}^n)`$$`)`$; $`(e)`$ there exists a measurable function $`\mathrm{\Lambda }:\text{}^n\text{}`$ such that $`\mathrm{\Lambda }(x)+\mathrm{}`$ as $`x\mathrm{}`$ and (1.18) holds. If we additionally assume that the electric potential $`V`$ is semi-bounded below then we can add the bottoms of the Neumann spectrum to the picture as was first done by A. Molchanov , though without magnetic field. This will be important for some arguments in this paper. ###### Theorem 1.2 If $`V`$ is semi-bounded below then the conditions in Theorem 1.1 are also equivalent to the following conditions: $`(f)`$ $`\mu (B(x,r);H_{a,V})+\mathrm{}`$ as $`x\mathrm{}`$ for any fixed $`r>0`$; $`(g)`$ there exists $`r>0`$ such that $`\mu (B(x,r);H_{a,V})+\mathrm{}`$ as $`x\mathrm{}`$. Finally we can weaken the requirement on $`\mathrm{\Lambda }`$ by use of capacity through the Molchanov condition: ###### Theorem 1.3 Let us assume that $`V`$ is semi-bounded below. Then the conditions $`(a)(g)`$ in Theorems 1.1 and 1.2 are equivalent to the existence of a measurable function $`\mathrm{\Lambda }:\text{}^n\text{}`$ such that $`1)`$ $`\mathrm{\Lambda }`$ is semi-bounded below and satisfies $`(M_c)`$ with some $`c>0`$; $`2)`$ the operator inequality $`H_{a,V}\mathrm{\Lambda }(x)`$ holds in the same sense as in Theorem 1.1. For $`\mathrm{\Lambda }=V`$ and $`a=0`$ this gives a sufficiency part of the Molchanov theorem. It also implies the following result from : ###### Corollary 1.4 If $`V`$ is semi-bounded below and $`H_{0,V}`$ has a discrete spectrum $`(`$or, equivalently, $`V`$ satisfies $`(M_c)`$ with some $`c>0`$$`)`$, then $`H_{a,V}`$ also has a discrete spectrum for any magnetic potential $`a`$. Our proof of this statement is in fact purely variational and it is completely different from the proof given in where semigroup methods are used. However the regularity requirements for $`a`$ and $`V`$ are much weaker in . Finally let us formulate a convenient sufficient condition which does not require $`V`$ or $`\mathrm{\Lambda }`$ to be semi-bounded below. ###### Theorem 1.5 Assume that there exists a measurable function $`\mathrm{\Lambda }:\text{}^n\text{}`$ such that the following conditions are satisfied: $`1)`$ the operator inequality (1.19) $$H_{a,0}\mathrm{\Lambda }(x)$$ holds in the same sense as above; $`2)`$ there exists $`\delta [0,1)`$ such that the effective potential (1.20) $$V_{\text{eff}}^{(\delta )}(x)=V(x)+\delta \mathrm{\Lambda }(x)$$ is semi-bounded below and satisfies $`(M_c)`$ with some $`c>0`$. Then $`H_{a,V}`$ is essentially self-adjoint, semi-bounded below and has a discrete spectrum. ### 1.3 Sufficient conditions ($`n=2`$) The case $`n=2`$ is much simpler than the general case because the magnetic field does not change direction (and may only change sign). We will identify the magnetic field with its component $`B_{12}`$. Note that by changing enumeration of coordinates we can change sign of $`B=B_{12}`$. Some uncertainty principle related arguments lead to the following fact established by J. Avron, I. Herbst and B. Simon : ###### Theorem 1.6 If $`n=2`$ and $`\|B(x)\|\mathrm{}`$ as $`x\mathrm{}`$, then $`\sigma =\sigma _d`$ for $`H_{a,0}`$ (hence for $`H_{a,V}`$ with arbitrary $`V0`$). This is the simplest result which shows that magnetic field alone may cause a localization (i.e. a discrete spectrum) for a quantum particle. Classically this can be understood from the fact that a strong magnetic field causes a fast rotation of a charged moving particle without changing its kinetic energy, and in this way the field impedes possible escape of the particle to infinity. Note that $`\|B(x)\|\mathrm{}`$ means in fact that either $`B`$ or $`B`$ tend to $`+\mathrm{}`$ as $`x\mathrm{}`$. Hence the following theorem improves the result above: ###### Theorem 1.7 Assume that the following conditions are satisfied: $`(a)`$ there exists $`C>0`$ such that $`B(x)C`$, $`x\text{}^2`$; $`(b)`$ the Molchanov condition $`(M_c)`$ is satisfied for $`B(x)`$ or, equivalently, for $`\|B(x)\|`$ $`(`$instead of $`V(x)`$$`)`$ with some $`c>0`$. Then $`H_{a,0}`$ has a discrete spectrum. The simplest explicit result which takes into account both electric and magnetic field is given by the following ###### Theorem 1.8 Assume that $`n=2`$ and there exists $`\delta [1,1]`$ such that for the effective potential (1.21) $$V_{\text{eff}}^{(\delta )}(x)=V(x)+\delta B(x)$$ we have $`V_{\text{eff}}^{(\delta )}(x)+\mathrm{}`$ as $`x\mathrm{}`$. Then $`H_{a,V}`$ is essentially self-adjoint, semi-bounded below and has a discrete spectrum. V. Ivrii noticed that the range of possible $`\delta `$ is precise here: the conclusion does not hold if we take any $`\delta [1,1]`$. The following theorem strengthens Theorem 1.7 by taking into account the influence of the electric potential $`V`$: ###### Theorem 1.9 Assume that $`n=2`$ and there exists $`\delta (1,1)`$ such that the effective potential $`V_{eff}`$ given by (1.21) is semi-bounded below and satisfies the Molchanov condition $`(M_c)`$ with some $`c>0`$. Then $`H_{a,V}`$ is essentially self-adjoint, semi-bounded below and has a discrete spectrum. Clearly Theorem 1.9 implies Theorem 1.7 (take $`V=0`$). It also implies the sufficiency of the Molchanov condition in case $`B=0`$. Note however that the conditions of the theorem do not imply any growth of $`V`$ or $`B`$. If $`\|B\|`$ itself tends to $`\mathrm{}`$ as $`x\mathrm{}`$, then $`V`$ is even allowed to go to $`\mathrm{}`$, though in this case $`\|V\|`$ must be dominated by $`\delta \|B\|`$ with some positive $`\delta <1.`$ Theorem 1.9 also implies Theorem 1.8 except for the extreme values $`\delta =\pm 1`$. Unfortunately we are not aware whether any of the conditions in Theorem 1.9 is necessary (though they are necessary if $`B=0`$ and $`V`$ is semi-bounded below due to the Molchanov result quoted above). Even Theorem 1.8, which does not require any use of capacity and can be proved by elementary means (see Sect.5.1), seems to be absent in the literature. A stronger result without capacity can be obtained if we replace capacity by the Lebesgue measure. Namely, let us formulate the corresponding condition for a semi-bounded below function $`V`$: there exists $`r_0>0`$ such that for any $`r(0,r_0)`$ ($`\stackrel{~}{M}_{c,N}`$) $$\underset{𝐹}{inf}\left\{_{B(x,r)F}V(y)𝑑y\right\|\mathrm{mes}(F)cr^N\}+\mathrm{}\text{as}x\mathrm{}.$$ It follows from a well-known estimate of measure by capacity, that for any $`c>0`$ and $`N>0`$ the condition $`(\stackrel{~}{M}_{c,N})`$ implies $`(M_c^{})`$ for some $`c^{}>0`$ (see the proof of part 2) of Theorem 6.1 in ). Therefore Theorem 1.9 implies the following ###### Corollary 1.10 Assume that there exists $`\delta (1,1)`$ such that the effective potential $`V_{\text{eff}}^{(\delta )}`$ given by (1.21) is semi-bounded below and satisfies $`(\stackrel{~}{M}_{c,N})`$ with some $`c>0`$ and $`N>0`$. Then $`H_{a,V}`$ is essentially self-adjoint, semi-bounded below and has a discrete spectrum. Finally note that an elementary argument given in the proof of Corollary 6.2 in gives a sufficient condition which is stronger than ($`\stackrel{~}{M}_{c,N}`$) but very easy to check. ###### Corollary 1.11 Assume that there exists $`\delta (1,1)`$ such that $`V_{eff}^{(\delta )}`$ is semi-bounded below and for any $`A>0`$ and any small $`r>0`$ (1.22) $$\mathrm{mes}\{y\|yB(x,r),V_{eff}^{(\delta )}(y)A\}0\text{as}x\mathrm{}.$$ Then $`H_{a,V}`$ is essentially self-adjoint, semi-bounded below and has a discrete spectrum. ### 1.4 Sufficient conditions ($`n3`$) The behavior of the spectrum of the magnetic Schrödinger operator in dimensions $`n3`$ is much more complicated than in dimension 2 because of possible varying direction of $`B`$. In particular none of the results formulated above for $`n=2`$ holds for $`n3`$. A. Dufresnoy gave the first example of the operator $`H_{a,0}`$ with (1.23) $$\|B(x)\|\mathrm{}\text{as}x\mathrm{},$$ and yet with non-compact resolvent (or, equivalently, with non-discrete spectrum) in an arbitrary dimension $`n3`$. J. Avron, I. Herbst and B. Simon gave sufficient conditions for the discreteness of the spectrum of $`H_{a,0}`$ which in addition to (1.23) require that the direction of $`B`$ varies sufficiently slowly. A more explicit condition of this kind (in terms of estimates for derivatives of the direction of $`B`$) was given in . Later A. Iwatsuka improved these results, giving almost precise estimates of this kind. He also produced a series of spectacular examples. One of them shows that no growth condition for $`\|B(x)\|`$ (i.e. a condition of the form $`\|B(x)\|\rho (x)`$ with a fixed continuous function $`\rho `$) would suffice for the discreteness of the spectrum of $`H_{a,0}`$. This is in drastic contrast with the results for $`n=2`$ formulated above. Another example given in shows that the condition (1.23) is also not necessary for the discreteness of spectrum of $`H_{a,0}`$. (This is of course less surprising because you can expect that integrally small perturbations of $`B`$ should not generally affect the discreteness of spectrum.) However we will show in Sect.5 that some explicit sufficient conditions for the discreteness of spectrum can still be formulated in terms of effective potentials similarly to the case $`n=2`$. The appropriate effective potentials will include both electric and magnetic fields. They will incorporate information about the direction of the magnetic field. The above mentioned results of J. Avron, I. Herbst and B. Simon, A. Dufresnoy and A. Iwatsuka will follow if we impose additional conditions on the direction of the magnetic field. These conditions allow us to simplify the form of the effective potential. Here we will formulate just one result coming from this approach (others can be found in Sect.5). By $`\mathrm{Lip}(X)`$ we will denote the set of all Lipschitz functions on any metric space. We will only use $`X`$ which are locally compact. In this case $`\mathrm{Lip}_{\text{loc}}(X)`$ will denote the space of functions which are locally Lipschitz on $`X`$. Recall that Lipschitz functions on an open subset in $`\text{}^n`$ are exactly the ones which have bounded distributional derivatives (see e.g. ). Let us define a smoothed direction of the magnetic field as follows: (1.24) $$A_{jk}(x)=\chi (\|B(x)\|)\frac{B_{jk}(x)}{\|B(x)\|},$$ where $`\chi \mathrm{Lip}([0,\mathrm{}))`$, $`\chi (r)=0`$ if $`r1/2`$, $`\chi (r)=1`$ if $`r1`$, $`\chi (r)=2r1`$ if $`1/2r1`$, so $`0\chi (r)1`$ and $`\|\chi ^{}(r)\|2`$ for all $`r`$. ###### Theorem 1.12 Let us assume that $`B_{jk}\mathrm{Lip}(\text{}^n)`$ for all $`j,k`$; $`A_{jk}`$ are defined by (1.24), and a positive measurable function $`X(x)`$ in $`\text{}^n`$ satisfies (1.25) $$\underset{k=1}{\overset{n}{}}\left\|\frac{A_{kj}(x)}{x^k}\right\|X(x),x\text{}^n,$$ for all $`j`$. Then for any $`\epsilon >0`$ and $`\delta [0,1)`$ define an effective potential (1.26) $$V_{\text{eff}}^{(\delta ,\epsilon )}(x)=V(x)+\frac{\delta }{n1+\epsilon }\|B(x)\|\frac{n\delta }{4\epsilon (n1+\epsilon )}X^2(x).$$ If there exist $`\epsilon >0`$ and $`\delta [0,1)`$, such that $`V_{\text{eff}}^{(\delta ,\epsilon )}`$ is bounded below and satisfies the Molchanov condition $`(M_c)`$ with some $`c>0`$ (in particular, this holds if $`V_{\text{eff}}^{(\delta ,\epsilon )}(x)+\mathrm{}`$ as $`x\mathrm{}`$), then $`H_{a,V}`$ is essentially self-adjoint, semi-bounded below and has a discrete spectrum. Note that in case $`B=0`$ and $`V`$ semi-bounded below the condition on $`V`$ in this theorem becomes necessary and sufficient due to the Molchanov theorem. In Section 5 we will establish that imposing different regularity conditions on $`B`$, e.g. as in , leads to different types of effective potentials so that results similar to Theorem 1.9 hold. Here we will give the simplest example of this kind. As explained above it is natural to impose some a priori conditions of regularity on $`B`$. We will mainly formulate them in the form of estimates for derivatives of $`B`$. We will always assume that $`B\mathrm{Lip}_{\text{loc}}`$, i.e. $`B_{jk}\mathrm{Lip}_{\text{loc}}`$ for all $`j,k`$. Following A. Iwatsuka we will use the estimates of the form ($`B_\alpha `$) $$\|B(x)\|C(1+\|B(x)\|)^\alpha ,x\text{}^n,$$ where $`\alpha >0`$, $`C>0`$ and $`B`$ means vector whose components are all possible first order derivatives $`B_{jk}/x^l`$. We will write that $`B`$ satisfies ($`B_\alpha `$) if there exists $`C>0`$ such that this estimate is satisfied with the given $`\alpha `$. A little bit stronger condition used in and is ($`B_\alpha ^0`$) $$\|B(x)\|(1+\|B(x)\|)^\alpha 0\text{as}x\mathrm{}.$$ A. Dufresnoy proved that the conditions (1.23) and $`(B_{3/2}^0)`$ imply that the spectrum of $`H_{a,0}`$ is discrete. In fact he proved that instead of $`(B_{3/2}^0)`$ it is sufficient to require a slightly weaker condition $$\left\|\left(\frac{B(x)}{\|B(x)\|}\right)\right\|\|B(x)\|^{1/2}0\text{as}x\mathrm{}.$$ A. Iwatsuka proved that in fact $`(B_2^0)`$ together with (1.23) imply the discreteness of spectrum for $`H_{a,0}`$. He also provided an example which shows that $`(B_2^0)`$ can not be replaced by $`(B_2)`$ (hence it can not be replaced by $`(B_\alpha )`$ or $`(B_\alpha ^0)`$ with any $`\alpha >2`$). The following theorem compared with the above mentioned results takes into account the behavior of $`V`$. ###### Theorem 1.13 Assume that $`(a)`$ $`B`$ satisfies $`(B_{3/2}^0)`$, and $`(b)`$ there exists $`\delta [0,1)`$ such that the effective potential (1.27) $$V_{\text{eff}}^{(\delta )}=V+\frac{\delta }{n1}\|B\|$$ is semi-bounded below and satisfies the Molchanov condition $`(M_c)`$ with a small $`c>0`$ (possibly depending on $`\delta `$). Then $`H_{a,V}`$ is essentially self-adjoint, semi-bounded below and has a discrete spectrum. ###### Corollary 1.14 Assume that $`(a)`$ $`B`$ satisfies $`(B_{3/2}^0)`$ and $`(b)`$ there exists $`\delta [0,1)`$ such that $`V_{\text{eff}}^{(\delta )}(x)+\mathrm{}`$ as $`x\mathrm{}`$ for the effective potential $`V_{\text{eff}}^{(\delta )}`$ defined by (1.27). Then $`H_{a,V}`$ is essentially self-adjoint, semi-bounded below and has a discrete spectrum. Due to the arguments given in Sect.6.1 of it is also possible to replace capacity by the Lebesgue measure. Namely, let us introduce the corresponding condition for a semi-bounded below function $`V`$: ($`\stackrel{~}{M}_c`$) $$\underset{𝐹}{inf}\left\{_{B(x,r)F}V(y)𝑑y\right\|\mathrm{mes}(F)cr^n\}+\mathrm{}\text{as}x\mathrm{}.$$ As shown in Sect.6.1 of , for any $`c>0`$ there exists $`c^{}>0`$ such that $`(\stackrel{~}{M}_c)`$ implies $`(M_c^{})`$. Therefore we have the following ###### Corollary 1.15 Assume that $`B`$ satisfies $`(B_{3/2}^0)`$ and $`(b)`$ there exists $`\delta [0,1)`$ such that the the effective potential (1.27) is semi-bounded below and there exists $`c>0`$ such that the condition ($`\stackrel{\mathit{~}}{M}_c`$) holds for $`V_{\text{eff}}^{(\delta )}`$ (instead of $`V`$); Then $`H_{a,V}`$ is essentially self-adjoint, semi-bounded below and has a discrete spectrum. Finally let us formulate an analogue of Corollary 1.11 for $`n3`$. ###### Corollary 1.16 Assume that $`B`$ satisfies $`(B_{3/2}^0)`$ and $`(b)`$ there exists $`\delta [0,1)`$ such that the effective potential (1.27) is semi-bounded below and for any $`A>0`$ and any small $`r>0`$ (1.28) $$\mathrm{mes}\{y\|yB(x,r),V_{\text{eff}}^{(\delta )}(y)A\}0\text{as}x\mathrm{}.$$ Then $`H_{a,V}`$ is essentially self-adjoint, semi-bounded below and has a discrete spectrum. Note that if we assume that $`V`$ is semi-bounded below, then we can always take $`V+\|B\|`$ as an effective potential in all statements above. ### 1.5 Acknowledgments Most part of this work was done during our stay in the program Research in Pairs of Oberwolfach Forschungsinstitut für Mathematik in May-June 1999. We gratefully acknowledge the generous support of this Institute and its kind and helpful personnel. The second author was also partially supported by NSF grant DMS-9706038. The authors are grateful to M.S. Birman, J.-M. Bismut, V. Ivrii, J. Lott, V.G. Maz’ya and M.Z. Solomyak for useful discussions. ## 2 Essential self-adjointness The goal of this section is to give a simple proof of the following simplest version of the theorem on essential self-adjointness of any semi-bounded magnetic Schrödinger operator (see for other versions). ###### Theorem 2.1 Assume that $`a_j\mathrm{Lip}_{\text{loc}}(\text{}^n)`$, $`VL_{loc}^{\mathrm{}}(\text{}^n)`$ and the corresponding magnetic Schrödinger operator $`H_{a,V}`$ is semi-bounded below on $`C_c^{\mathrm{}}(\text{}^n)`$, i.e. there exists a constant $`C\text{}`$ such that (2.1) $$(H_{a,V}u,u)C(u,u),uC_c^{\mathrm{}}(\text{}^n).$$ Then $`H_{a,V}`$ is essentially self-adjoint. Proof. We will extend the Wienholtz proof of the Povzner theorem as it is explained by I.M. Glazman . Let us recall that if $`g\mathrm{Lip}(\mathrm{\Omega })`$ where $`\mathrm{\Omega }`$ is an open subset in $`\text{}^n`$, then $`g/x^jL^{\mathrm{}}(\mathrm{\Omega })`$, $`j=1,\mathrm{},n`$ (see e.g. , Sect.1.1), where the derivatives are understood in the sense of distributions (but also exist almost everywhere). This implies that the operator $`H_{a,V}`$ is well defined on $`C_c^{\mathrm{}}(\text{}^n)`$ (and maps this space into $`L^2(\text{}^n)`$) as well as on $`L^2(\text{}^n)`$ (which it maps to the space of distributions on $`\text{}^n`$). Note that adding $`(C+1)I`$ to $`H_{a,V}`$ we can assume that $`H_{a,V}I`$ on $`C_c^{\mathrm{}}(\text{}^n)`$, i.e. (2.2) $$h_{a,V}(u,u)(u,u),uC_c^{\mathrm{}}(\text{}^n).$$ If this is true then it is well known (see e.g. ) that the essential self-adjointness of $`H_{a,V}`$ is equivalent to the fact that the equation (2.3) $$H_{a,V}u=0$$ has no non-trivial solutions in $`L^2(\text{}^n)`$ (understood in the sense of distributions). Assume that $`u`$ is such a solution. First note that it is in $`W_{loc}^{2,2}(\text{}^n)`$ (i.e. has distributional derivatives of order $`2`$ in $`L_{loc}^2(\text{}^n)`$) due to a simple elliptic regularity argument (see e.g. Lemma 4.1 in for more details). Let us take a cut-off function $`\varphi _RC_c^{\mathrm{}}(\text{}^n)`$ with the following properties: $`0\varphi _R1;`$ $`\varphi _R=1\text{on}B(0,R)\text{and}0\text{on}\text{}^nB(0,2R);`$ $`\epsilon _R:=\underset{x\text{}^n}{sup}\|\varphi _R(x)\|0\text{as}R\mathrm{}.`$ Then denoting $`u_R=\varphi _Ru`$ we see that $`u_R`$ is in the domain of the minimal operator associated with $`H_{a,V}`$, hence (2.4) $$u_R^2(H_{a,V}u_R,u_R).$$ Let us calculate $`H_{a,V}u_R`$ using the Leibniz type formula for $`P_j`$: (2.5) $$P_j(fg)=(P_jf)g+f(D_jg),$$ where $`D_j=i/x^j`$. Applying this formula twice to calculate $`P_j^2(\varphi _Ru)`$ and summing up we easily obtain due to (2.3): $`H_{a,V}u_R`$ $`=\varphi _RH_{a,V}u2(\varphi _R)(_au)u\mathrm{\Delta }\varphi _R`$ $`=2(\varphi _R)(_au)u\mathrm{\Delta }\varphi _R`$ $`=2\varphi _Ru2i(a\varphi _R)u\mathrm{\Delta }\varphi _R.`$ Therefore due to 2.4 we have (2.7) $$\varphi _Ru^2_\text{}^n(2\varphi _Ru2i(a\varphi _R)uu\mathrm{\Delta }\varphi _R)\varphi _R\overline{u}𝑑x.$$ Since $`(H_{a,V}u_R,u_R)`$ is real, we can replace the right hand side here by the complex conjugate expression. Adding the two estimates obtained in such a way and dividing by 2 we see that the term with the magnetic potential $`a`$ cancels and we get, applying integration by parts: $`\varphi _Ru^2`$ $`{\displaystyle _\text{}^n}\left[\varphi _R(\varphi _R)(\overline{u}u+u\overline{u})\varphi _R(\mathrm{\Delta }\varphi _R)\|u\|^2\right]𝑑x`$ $`={\displaystyle _\text{}^n}\left[\varphi _R(\varphi _R)(\|u\|^2)\varphi _R(\mathrm{\Delta }\varphi _R)\|u\|^2\right]𝑑x`$ $`={\displaystyle _\text{}^n}\left[\varphi _R(\mathrm{\Delta }\varphi _R)\|u\|^2+\|\varphi _R\|^2\|u\|^2\varphi _R(\mathrm{\Delta }\varphi _R)\|u\|^2\right]𝑑x`$ $`={\displaystyle _\text{}^n}\|\varphi _R\|^2\|u\|^2𝑑x.`$ In particular we obtain using the conditions on $`\varphi _R`$ above: $$_{B(0,R)}\|u\|^2𝑑x\epsilon _R^2_{B(0,2R)}\|u\|^2𝑑x.$$ Allowing $`R`$ to go to $`+\mathrm{}`$ we see that $`u^2=0`$, hence $`u0`$. Remark. The local condition $`VL_{loc}^{\mathrm{}}(\text{}^n)`$ can be considerably weakened. For example, it is sufficient to require that $`V=V_++V_{}`$, where $`V_+0`$, $`V_+L_{loc}^2(\text{}^n)`$, $`V_{}0`$, $`V_{}L_{loc}^p(\text{}^n)`$ with $`p=2`$ if $`n3`$, $`p>2`$ if $`n=4`$, and $`p=n/2`$ if $`n5`$. (See e.g. .) ## 3 Localization In this section we will prove different localization theorems which were formulated in Sect.1.2 and provide an important preliminary material related to compactness arguments and estimates of the bottoms of the Dirichlet and Neumann spectra. We will use notations from previous sections. In particular, $`a,V`$ will always denote magnetic and electric potential with the same regularity as in Section 1. We will start with the following elementary and well known diamagnetic inequality (see e.g. ): ###### Lemma 3.1 Let $`a`$ be an arbitrary magnetic potential (with components from $`C^1`$). Let $`u`$ be a complex valued Lipschitz function in an open set $`U\text{}^n`$. Then $`\|u\|`$ is also Lipschitz and (3.1) $$\|\|u\|\|\|_au\|\text{a.e.,}$$ where a.e. means almost everywhere with respect to the Lebesgue measure. Let us assume that $`H_{a,V}`$ is bounded below, hence essentially self-adjoint due to Theorem 2.1. Without loss of generality we can assume hereafter that $`H_{a,V}I`$ (or that the estimate (2.2) is satisfied). The essential self-adjointness of $`H_{a,V}`$ means that $`C_c^{\mathrm{}}(\text{}^n)`$ is its core, i.e. the closure of $`H_{a,V}`$ from the initial domain $`C_c^{\mathrm{}}(\text{}^n)`$ is a self-adjoint operator in $`L^2(\text{}^n)`$. It follows that $`C_c^{\mathrm{}}(\text{}^n)`$ is also a core for the corresponding quadratic form. Denote (3.2) $$=\{uC_c^{\mathrm{}}(\text{}^n)\|h_{a,V}(u,u)1\}.$$ ###### Lemma 3.2 $`\sigma =\sigma _d`$ if and only if $``$ is precompact in $`L^2(\text{}^n)`$. Proof. The proof is the same as the proof of Lemma 2.2 in and it is essentially abstract. Clearly $`\sigma =\sigma _d`$ is equivalent to saying that for $`H=H_{a,V}`$ we have (3.3) $$\{u\|u\mathrm{Dom}(H^{1/2}),H^{1/2}u1\}$$ is precompact in $`L^2(\text{}^n)`$. Here $`\mathrm{Dom}(H^{1/2})`$ also coincide with the domain of the quadratic form which is the closure of $`h_{a,V}`$. Since $`C_c^{\mathrm{}}(\text{}^n)`$ is a core for the quadratic form too, we see that precompactness of the set (3.3) is equivalent to the precompactness of $``$. ###### Lemma 3.3 (Small Tails Lemma) Let us assume as above that $`H_{a,V}`$ is essentially self-adjoint and semi-bounded below so that (2.2) holds. Then $`\sigma =\sigma _d`$ if and only if the following small tails condition is satisfied: for any $`\epsilon >0`$ there exists $`R>0`$ such that (3.4) $$_{\text{}^nB(0,R)}\|u\|^2𝑑x<\epsilon \text{for any}u,$$ or, in other words, $$_{\text{}^nB(0,R)}\|u\|^2𝑑x0\text{as}R\mathrm{},\text{uniformly in}u.$$ Proof. Again the proof is similar to the proof of Lemma 2.3 in , though a small additional argument is needed to avoid using semi-boundedness of $`V`$ which was used in . Clearly $`\sigma =\sigma _d`$ (or precompactness of $``$) implies the small tails condition because any precompact set has a $`\epsilon `$-net for any $`\epsilon >0`$. Vice versa, assume that the small tails condition is fulfilled. Then the precompactness of $``$ would be equivalent to the precompactness of any restriction $$_R=\left\{u\|_{B(0,R)}\|u\right\},$$ where $`R>0`$. Note that the condition (2.2) implies that the set $`_R`$ is bounded in $`L^2(B(0,R))`$. But then by the Sobolev–Kondrashov compactness theorem it is sufficient to establish a uniform $`L^2(B(0,R))`$-boundedness of the gradients of functions $`u_R`$ (for any fixed $`R>0`$). Since we assume that $`a_jL_{loc}^{\mathrm{}}`$, it is sufficient to establish that the magnetic gradients $`_au`$ are uniformly bounded in $`L^2(B(0,R))`$. This in turn follows from the definition of $``$ and uniform $`L^2`$-boundedness of the functions $`u`$ combined with the local boundedness of the potential $`V`$. Remark. The requirement $`uC_c^{\mathrm{}}(\text{}^n)`$ in the definition of $``$ (and in Lemmas 3.2 and 3.3) can be replaced by the requirement $`u\mathrm{Lip}_c(\text{}^n)`$ (the set of all Lipschitz functions with compact support in $`\text{}^n`$) because the space $`\mathrm{Lip}_c(\text{}^n)`$ is intermediate between $`C_c^{\mathrm{}}(\text{}^n)`$ and $`\mathrm{Dom}\left(H_{a,V}^{1/2}\right)`$. Now let us take a covering of $`\text{}^n`$ by balls $`B(x_k,r)`$, $`k=1,2,\mathrm{}`$, of fixed radius $`r>0`$, such that this covering has a finite multiplicity. Now take a partition of unity on $`\text{}^n`$ consisting of functions $`e_kC_c^{\mathrm{}}(\text{}^n)`$, $`k=1,2,\mathrm{}`$, such that $`0e_k1`$, $`\text{supp}e_kB(x_k,r)`$, and $$\underset{k=1}{\overset{\mathrm{}}{}}e_k^2=1,$$ $$\|e_k\|C,k=1,2,\mathrm{},$$ where $`C`$ does not depend on $`k`$. The main tool in proving the localization theorem is the following IMS localization formula (3.5) $$H_{a,V}=\underset{k=1}{\overset{\mathrm{}}{}}J_kH_{a,V}J_k\underset{k=1}{\overset{\mathrm{}}{}}\|e_k\|^2,$$ where $`J_k`$ is the multiplication operator by $`e_k`$ in $`L^2(\text{}^n)`$. Proofs of different versions of this formula can be found in , (Sect. 3.1), . Formally only the case $`a=0`$ is treated in , though the proof works with arbitrary $`a`$. Much more general case (second order differential operators on manifolds) is considered in , Sect.3. Now we will fix an operator $`H_{a,V}`$ and denote for brevity $`\lambda (B(x,r))=\lambda (B(x,r);H_{a,V})`$, $`\mu (B(x,r))=\mu (B(x,r);H_{a,V})`$. Proof of Theorem 1.1. $`(a)(b).`$ Let us assume that $`(a)`$ is satisfied and fix an arbitrary $`r>0`$. According to Lemma 3.3 for any $`\epsilon >0`$ there exists $`R>0`$ such that $`\|x\|>R`$ implies that $$_{B(x,r)}\|u(x)\|^2𝑑x<\epsilon ,$$ as soon as $`uC_c^{\mathrm{}}(B(x,r))`$ and $`h_{a,V}(u,u)1`$. It follows that $`\lambda (B(x,r))1/\epsilon `$, which implies $`(b)`$. Clearly $`(b)(c)`$. $`(c)(d)`$. Let us fix $`r>0`$ such that $`(c)`$ is satisfied and choose a covering of $`\text{}^n`$ by the balls $`B(x_k,r)`$, $`k=1,2,\mathrm{},`$, and a partition of unity $`\{e_k\|k=1,2,\mathrm{},\}`$ with $`\text{supp}e_kB(x_k,r)`$ and with the properties formulated above. Then for any $`uC_c^{\mathrm{}}(\text{}^n)`$ we obtain from (3.5): $$h_{a,V}(u,u)=\underset{k}{}h_{a,V}(e_ku,e_ku)\underset{k}{}\|e_k\|^2\|u\|^2(\mathrm{\Lambda }u,u),$$ where $$\mathrm{\Lambda }(x)=\underset{k}{}\lambda (B(x_k,r))e_k^2(x)\underset{k}{}\|e_k\|^2.$$ The first sum tends to $`+\mathrm{}`$ as $`x\mathrm{}`$ due to the condition $`(c)`$, whereas the second sum is bounded. This implies that $`\mathrm{\Lambda }(x)+\mathrm{}`$ as $`x\mathrm{}`$, hence $`(d)`$ is fulfilled. The implication $`(d)(e)`$ is obvious. $`(e)(a).`$ Assume that $`(e)`$ is satisfied with the function $`\mathrm{\Lambda }(x)`$. Let us take $`u`$. Then $$_\text{}^n\mathrm{\Lambda }(x)\|u(x)\|^2𝑑x1.$$ Therefore, $$_{\|x\|R}\|u(x)\|^2𝑑x\left(\underset{\{x\|\|x\|R\}}{inf}\{\mathrm{\Lambda }(x)\}\right)^10\text{as}R\mathrm{},$$ so $`(a)`$ follows from Lemma 3.3. Now we assume that $`V0`$ and proceed to some preparatory material which is needed to prove Theorem 1.2. Denote temporarily $`B_r=B(0,r)\text{}^n`$ and define $$\psi =\psi _{L^2(B_r)}=\left(_{B_r}\|\psi \|^2𝑑x\right)^{1/2},\psi _t=\psi _{L^2(B_{tr})},$$ where $`0<t1`$. Similarly define $`\psi =\psi _{L^2(B_r)}=\left({\displaystyle _{B_r}}\|\psi \|^2𝑑x\right)^{1/2},`$ $`\psi _t=\psi _{L^2(B_{tr})}.`$ ###### Lemma 3.4 The following estimates hold true for any magnetic potential $`a`$: (3.6) $$\psi t^{n/2}\psi _t+2^{n+1}r(1t)_a\psi ,t[1/2,1];$$ (3.7) $$\psi ^22t^n\psi _t^2+2^{2n+3}r^2(1t)^2_a\psi ^2,t[1/2,1].$$ In these estimates we assume that $`\psi \mathrm{Lip}(B_r)`$. Proof. With $`a=0`$ this Lemma was proved in (for cubes instead of balls) and in (see Lemma 2.8 there). Applying this particular case to $`\|\psi \|`$ and using Lemma 3.1, we obtain the desired result. The following Lemma for the case $`a=0`$ was proved in (for cubes) and (see Lemma 2.9 there): ###### Lemma 3.5 Let us assume that $`V0`$. Then (3.8) $$\mu (B(x,r))\lambda (B(x,r))C_1\mu (B(x,r))+C_2r^2,$$ for any $`x\text{}^n`$, and any $`r>0`$. Here $`C_1`$ and $`C_2`$ depend only on $`n`$; for example we can take $`C_1=2^{n+3}(1+2^{2n+6})`$ and $`C_2=2^{n+7}`$. Proof. The proof given in works in our case (with arbitrary $`a`$) if we replace $``$ by $`_a`$, use the Leibniz rule (2.5) and apply Lemma 3.4 above. Proof of Theorem 1.2. Equivalence of $`(b)`$ (respectively $`(c)`$) from Theorem 1.1 to $`(f)`$ (resp. $`(g)`$) from Theorem 1.2 follows from Lemma 3.5 provided we additionally assume that $`V0`$. Proof of Theorem 1.3. Without loss of generality we can assume that $`V0`$. The necessity part of the theorem follows from Theorem 1.1. (If $`\sigma =\sigma _d`$ then we can even make $`\mathrm{\Lambda }(x)+\mathrm{}`$ as $`x\mathrm{}`$.) Now let us assume that there exists $`\mathrm{\Lambda }(x)`$ which is semi-bounded below, satisfies $`(M_c)`$ with some $`c>0`$, and $`H_{a,V}\mathrm{\Lambda }(x)`$, i.e. $$h_{a,V}(u,u)(\mathrm{\Lambda }\|u\|,\|u\|)$$ for any $`u\mathrm{Lip}_c(\text{}^n)`$. Using Lemma 3.1 we also obtain $$h_{a,V}(u,u)h_{a,0}(u,u)h_{0,0}(\|u\|,\|u\|).$$ Adding these two inequalities we get $$2h_{a,V}(u,u)h_{0,\mathrm{\Lambda }}(\|u\|,\|u\|).$$ Let us introduce the set $$\stackrel{~}{}=\{u\|u\mathrm{Lip}_c(\text{}^n),h_{0,\mathrm{\Lambda }}(u,u)1\},$$ which has the small tails property (3.4) due to the Molchanov theorem and the Small Tails Lemma 3.3. With $``$ defined by (3.2) we see that the map $`u\|u\|`$ maps $``$ into $`\sqrt{2}\stackrel{~}{}`$. Hence $``$ also has the small tails property and we conclude from Lemma 3.3 that $`H_{a,V}`$ has a discrete spectrum. Proof of Theorem 1.5. We immediately conclude from the conditions that for any function $`u\mathrm{Lip}_c(\text{}^n)`$ $`h_{a,V}(u,u)`$ $`=h_{a,0}(u,u)+(Vu,u)`$ $`=(1\delta )h_{a,0}(u,u)+\delta h_{a,0}(u,u)+(V\|u\|,\|u\|)`$ $`(1\delta )h_{0,0}(\|u\|,\|u\|)+((\delta \mathrm{\Lambda }+V)\|u\|,\|u\|)`$ $`=(1\delta )\left[h_{0,0}(\|u\|,\|u\|)+((1\delta )^1(V+\delta \mathrm{\Lambda })\|u\|,\|u\|)\right].`$ Now the same argument based on the Small Tails Lemma 3.3, as in the proof of Theorem 1.3, ends the proof of Theorem 1.5. ## 4 Bounded magnetic field perturbations The main goal of this section is the proof of the following ###### Theorem 4.1 Suppose we are given an electric potential $`VL_{\text{loc}}^{\mathrm{}}(\text{}^n)`$ which is semi-bounded below, and two magnetic potentials $`a,\stackrel{~}{a}\mathrm{Lip}_{\text{loc}}(\text{}^n)`$ such that for the magnetic field $`\stackrel{~}{B}`$, associated with $`\stackrel{~}{a}`$, we have $`\stackrel{~}{B}\mathrm{Lip}_{\text{loc}}(\text{}^n)`$ and (4.1) $$\|\stackrel{~}{B}(x)\|C,x\text{}^n.$$ Then $`H_{a+\stackrel{~}{a},V}`$ is essentially self-adjoint, semi-bounded below and has a discrete spectrum if and only if this is true for $`H_{a,V}`$. We will start with the following version of Poincaré Lemma which is similar to the one used by A. Iwatsuka , Proposition 3.2. ###### Lemma 4.2 Let $`B=_{j<k}B_{jk}dx^jdx^k`$ be a closed 2-form in $`B(x_0,r)\text{}^n`$ with $`B_{jk}\mathrm{Lip}(B(x_0,r))`$ and $`B_{jk}_{\mathrm{}}C`$. Then there exists a 1-form $`a=_ja_jdx^j`$ with $`da=B`$, such that $`a_j\mathrm{Lip}(B(x_0,r))`$ and $`a_j_{\mathrm{}}nC`$. Here $`_{\mathrm{}}`$ means the $`L^{\mathrm{}}`$-norm on $`B(x_0,r)`$. Proof. We can obviously assume that $`x_0=0`$. Then we can produce $`a_j`$ by the following explicit formulas (see e.g. , p. 155–156): $$a_j(x)=\underset{k=1}{\overset{n}{}}x_k_0^1tB_{kj}(tx)𝑑t,$$ and all necessary estimates obviously follow. Proof of Theorem 4.1. Let us assume that $`H_{a+\stackrel{~}{a},V}`$ is essentially self-adjoint, semi-bounded below and has a discrete spectrum. We will prove that the same holds for holds for $`H_{a,V}`$. Clearly this is sufficient to prove the theorem. We can also assume that $`V0`$. Let us choose an arbitrary ball $`B(x_0,r)`$. We would like to estimate the bottom of the Dirichlet spectrum of $`H_{a+\stackrel{~}{a},V}`$ which we denote, as before, by $`\lambda (B(x_0,r);H_{a+\stackrel{~}{a},V})`$. Using the gauge invariance we can then arbitrarily change $`\stackrel{~}{a}`$ in the ball $`B(x_0,r)`$ as soon as the relation $`d\stackrel{~}{a}=\stackrel{~}{B}`$ is preserved. Therefore we can use Lemma 4.2 and assume that $`\stackrel{~}{a}_{\mathrm{}}nrC`$, where $`C`$ is the constant in (4.1) and $$\stackrel{~}{a}_{\mathrm{}}=\underset{j}{\mathrm{max}}\stackrel{~}{a}_j_{\mathrm{}}.$$ Denote $$P_j^{}=\frac{1}{i}\frac{}{x^j}+a_j+\stackrel{~}{a}_j,P_j=\frac{1}{i}\frac{}{x^j}+a_j.$$ Then for any $`u\mathrm{Lip}_c(B(x_0,r))`$ we have $$P_j^{}u^2=P_ju+\stackrel{~}{a}_ju^2=P_ju^2+\stackrel{~}{a}_ju^2+2\text{R}e(P_ju,\stackrel{~}{a}_ju),$$ hence $$h_{a+\stackrel{~}{a},0}(u,u)=\underset{j=1}{\overset{n}{}}P_j^{}u^2h_{a,0}(u,u)+n\stackrel{~}{a}_{\mathrm{}}^2u^2+2\underset{j=1}{\overset{n}{}}\text{R}e(P_ju,\stackrel{~}{a}_ju).$$ We have for any $`\epsilon >0`$ $`2\text{R}e(P_ju,\stackrel{~}{a}_ju)2P_ju\stackrel{~}{a}_ju`$ $`\epsilon P_ju^2+{\displaystyle \frac{1}{\epsilon }}\stackrel{~}{a}_ju^2\epsilon P_ju^2+{\displaystyle \frac{1}{\epsilon }}\stackrel{~}{a}_{\mathrm{}}^2u^2.`$ Combining this with the previous estimate we obtain: $$h_{a+\stackrel{~}{a},0}(u,u)(1+\epsilon )h_{a,0}(u,u)+n\left(1+\frac{1}{\epsilon }\right)\stackrel{~}{a}_{\mathrm{}}^2(u,u).$$ Now adding $`(Vu,u)`$ to the left hand side of this inequality and $`(1+\epsilon )(Vu,u)`$ to the right hand side, we obtain the following operator inequality (which holds in the sense of quadratic forms on functions $`uC_c^{\mathrm{}}(B(x_0,r))`$): $`H_{a+\stackrel{~}{a},V}(1+\epsilon )H_{a,V}+n\left(1+{\displaystyle \frac{1}{\epsilon }}\right)\stackrel{~}{a}_{\mathrm{}}^2`$ $`=(1+\epsilon )\left(H_{a,V}+{\displaystyle \frac{n}{\epsilon }}\stackrel{~}{a}_{\mathrm{}}^2\right)(1+\epsilon )\left(H_{a,V}+{\displaystyle \frac{n^3r^2C^2}{\epsilon }}\right).`$ It follows that $$\lambda (B(x_0,r);H_{a+\stackrel{~}{a},V})(1+\epsilon )\left(\lambda (B(x_0,r);H_{a,V})+\frac{n^3r^2C^2}{\epsilon }\right).$$ Using the IMS localization formula (3.5), we easily conclude that $`H_{a,V}`$ is semi-bounded below on $`C_c^{\mathrm{}}(\text{}^n)`$, hence essentially self-adjoint due to Theorem 2.1. Now the discreteness of the spectrum for $`H_{a,V}`$ immediately follows from the localization Theorem 1.1. ###### Corollary 4.3 Let $`a`$ be a magnetic potential such that $`a\mathrm{Lip}_{\text{loc}}(\text{}^n)`$ and $`B=da`$ is also in $`\mathrm{Lip}_{\text{loc}}(\text{}^n)`$ and bounded. Let also $`VL_{\text{loc}}^{\mathrm{}}(\text{}^n)`$ be a bounded below electric potential. Then $`H_{a,V}`$has a discrete spectrum if and only if $`V`$ satisfies the Molchanov condition $`(M_c)`$ with some $`c>0`$. Proof. Theorem 4.1 implies in our case that $`H_{a,V}`$ has a discrete spectrum if and only if this is true for $`H_{0,V}`$. By the Molchanov theorem the later is equivalent to the fulfillment of $`(M_c)`$ (for $`V`$) with some $`c>0`$. ###### Corollary 4.4 Let $`a`$ be a magnetic potential which corresponds to a constant magnetic field. Let also $`VL_{\text{loc}}^{\mathrm{}}(\text{}^n)`$ be a bounded below electric potential. Then $`H_{a,V}`$ has a discrete spectrum if and only if $`V`$ satisfies the Molchanov condition $`(M_c)`$ with some $`c>0`$. ## 5 Sufficient conditions ### 5.1 Case $`n=2`$ We will start by demonstrating the uncertainty principle argument in the proof of the following well known Lemma (see e.g. ): ###### Lemma 5.1 Assume that $`n=2`$. Then (5.1) $$H_{a,0}B(x)\text{and}H_{a,0}B(x),$$ where $`B`$ as usual denotes the magnetic field produced by the magnetic potential $`a`$. The inequalities (5.1) hold in the sense of operator inequalities (i.e. inequalities of the quadratic forms) on $`C_c^{\mathrm{}}(\text{}^2)`$. Proof. We have $`H_{a,0}=P_1^2+P_2^2`$ with $`[P_1,P_2]=iB`$. (See notations in Sect. 1.1.) Hence for any $`uC_c^{\mathrm{}}(\text{}^2)`$ $`(Bu,u)\|(Bu,u)\|`$ $`=\|((P_1P_2P_2P_1)u,u)\|=\|(P_2u,P_1u)(P_1u,P_2u)\|`$ $`=\|2\text{Im}(P_1u,P_2u)\|2P_1uP_2uP_1u^2+P_2u^2`$ $`=(H_{a,0}u,u).`$ The second inequality in (5.1) follows from the first one (by changing enumeration of coordinates). Proof of Theorem 1.6. The condition $`\|B(x)\|\mathrm{}`$ means that either $`B+\mathrm{}`$ or $`B+\mathrm{}`$. In any of these cases the condition $`(d)`$ of Theorem 1.1 is satisfied. ###### Corollary 5.2 For $`n=2`$ and any $`\delta [1,1]`$ we have (5.2) $$H_{a,V}V+\delta B$$ Proof. Using the decomposition $`H_{a,V}=H_{a,0}+V`$ and Lemma 5.1 we obtain (5.3) $$H_{a,V}B+V\text{and}H_{a,V}B+V.$$ Multiplying the first inequality by $`\kappa [0,1]`$, the second by $`1\kappa `$ and adding we obtain (5.2) after denoting $`12\kappa `$ by $`\delta `$. ###### Corollary 5.3 Assume that $`n=2`$ and there exists $`\delta [1,1]`$ such that (5.4) $$V(x)+\delta B(x)+\mathrm{}\text{as}x\mathrm{}.$$ Then $`H_{a,V}`$ is essentially self-adjoint, semi-bounded below and has a discrete spectrum. Proof. Using Corollary 5.2 we see that the condition $`(d)`$ of Theorem 1.1 is fulfilled. $`\mathrm{}`$ ###### Corollary 5.4 Assume that $`n=2`$ and $`\mathrm{\Omega }\text{}^2`$ is an open set. Then for any $`\delta [1,1]`$ $$\lambda (\mathrm{\Omega };H_{a,V})\underset{x\mathrm{\Omega }}{inf}\{V(x)+\delta B(x)\}.$$ Proof. We can restrict the estimates given in Lemma 5.1 and Corollary 5.2 to the functions from $`C_c^{\mathrm{}}(\mathrm{\Omega })`$ for any $`\mathrm{\Omega }\text{}^2`$. Proof of Theorem 1.9. The result immediately follows from Lemma 5.1 and Theorem 1.5 (with $`\mathrm{\Lambda }(x)B(x)`$). Remark. V. Ivrii noticed that if we take any $`\delta [1,1]`$ then the result of the Corollary 5.3 does not hold any more. Without loss of generality we can assume that $`\delta >1`$. More precisely, we have ###### Proposition 5.5 (V. Ivrii) There exists a magnetic Schrödinger operator $`H_{a,V}`$ with $`C^{\mathrm{}}`$ potentials $`a,V`$ in $`\text{}^2`$ such that $`H_{a,V}0`$, the spectrum of $`H_{a,V}`$ is not discrete, but for any $`\delta >1`$ we have $`V(x)+\delta B(x)+\mathrm{}`$ as $`x\mathrm{}`$. Proof. It is well known from the Landau calculation that the bottom of the spectrum of the operator $`H_{a,0}`$ with constant magnetic field $`B`$ is precisely $`\|B\|`$ (see e.g. Sect.6.1.1. in for a more general calculation). Let us take a sequence of constants $`B_j`$, $`j=1,2,\mathrm{}`$, $`B_j>0`$, $`B_j+\mathrm{}`$ as $`j\mathrm{}`$. Then define $`V_j=B_j`$ and consider the Schrödinger operators $`H_j=H_{a_j,V_j}`$, where $`a_j`$ is a magnetic potential corresponding to the constant magnetic field $`B_j`$. Then the bottom of the spectrum of $`H_j`$ in $`L^2(\text{}^2)`$ is $`0`$ for all $`j`$. Note that the bottom of the spectrum of $`H_j`$ in $`L^2(\text{}^2)`$ can be defined as the bottom of the Dirichlet spectrum (see (1.16)) which involves only functions with a compact support. Therefore for any $`\epsilon >0`$ and any $`x\text{}^2`$ we can find $`R_j>0`$ such that $`\lambda (B(x,R_j);H_j)<\epsilon `$. In fact this does not depend either on the choice of the potential $`a_j`$, or on the center point $`x`$ due to the gauge invariance. Let us fix $`\epsilon >0`$ (e.g. take $`\epsilon =1`$). Let us choose a sequence of points $`x_j`$, $`j=1,2,\mathrm{}`$, such that the balls $`B(x_j,R_j+1)`$ are disjoint. Let us construct a function $`BC^{\mathrm{}}(\text{}^2)`$ such that $`B(x)=B_j`$ on $`B(x_j,R_j)`$ and $`B(x)+\mathrm{}`$ as $`x\mathrm{}`$. By the Poincaré Lemma we can find a magnetic potential $`a=a_kdx^k`$ in $`\text{}^2`$ with $`a_kC^{\mathrm{}}(\text{}^2)`$, $`k=1,2`$, such that the corresponding magnetic field is $`B(x)dx^1dx^2`$ (or $`B(x)`$). Define also $`V(x)=B(x)`$ and consider the magnetic Schrödinger operator $`H_{a,V}`$ with $`a`$ and $`V`$ as constructed above. It follows from Corollary 5.2 (with $`\delta =1`$) that $`H_{a,V}0`$. On the other hand it is clear from the variational principle for the Dirichlet spectrum that there are infinitely many points of the spectrum of $`H_{a,V}`$ below $`2\epsilon `$. Therefore the spectrum of $`H_{a,V}`$ is not discrete. At the same time for any fixed $`\delta >1`$ we have $`V(x)+\delta B(x)=(\delta 1)B(x)+\mathrm{}`$ as $`x\mathrm{}`$. $`\mathrm{}`$ ### 5.2 Case $`n3`$ The case $`n3`$ is substantially more complicated than the case $`n=2`$, in particular because no growth condition on $`\|B\|`$ would suffice for a reasonable estimate below for $`H_{a,0}`$ as was shown by A. Iwatsuka . We will establish however that appropriate regularity conditions on $`B`$, such as the ones imposed in , , , can be incorporated in the growth conditions for suitable effective potentials, so that results similar to Theorem 1.9 hold. #### 5.2.1 The Iwatsuka identity We will start with the following Lemma ((6.2) on page 370 in ): ###### Lemma 5.6 (Iwatsuka identity) Assume that we are given an open set $`\mathrm{\Omega }\text{}^n`$, a magnetic potential $`a`$ $`(`$with components from $`\mathrm{Lip}_{\text{loc}}(\mathrm{\Omega })`$$`)`$ and a set of real-valued functions $`A_{jk}\mathrm{Lip}_{\text{loc}}(\mathrm{\Omega })`$, $`j,k=1,\mathrm{},n`$, $`A_{kj}=A_{jk}`$. Then (5.5) $$2\underset{k<j}{}\mathrm{Im}(A_{kj}P_ku,P_ju)=(\left[\underset{k<j}{}A_{kj}B_{kj}\right]u,u)+\underset{k,j}{}(\frac{A_{kj}}{x^k}P_ju,u),$$ for any $`uW_{\text{comp}}^{1,2}(\mathrm{\Omega })`$, i.e. $`uL^2(\mathrm{\Omega })`$ such that $`u`$ has a compact support in $`\mathrm{\Omega }`$ and $`u(L^2(\mathrm{\Omega }))^n`$ $`(`$in the sense of distributions$`)`$. Proof. We reproduce the proof for the sake of completeness and also because we have a different sign convention in the definition of $`a`$ compared with . An obvious approximation argument shows that it is sufficient to consider $`uC_c^{\mathrm{}}(\text{}^n)`$. Then using integration by parts and the commutation relations (1.11) we obtain $`2{\displaystyle \underset{k<j}{}}\mathrm{Im}(A_{kj}P_ku,P_ju)`$ $`=`$ $`{\displaystyle \frac{1}{i}}{\displaystyle \underset{k<j}{}}\left\{(A_{kj}P_ku,P_ju)(A_{kj}P_ju,P_ku)\right\}`$ $`=`$ $`{\displaystyle \frac{1}{i}}{\displaystyle \underset{k<j}{}}((P_jA_{kj}P_kP_kA_{kj}P_j)u,u))`$ $`=`$ $`{\displaystyle \frac{1}{i}}{\displaystyle \underset{k<j}{}}\left\{(A_{kj}[P_j,P_k]u,u)+{\displaystyle \frac{1}{i}}({\displaystyle \frac{A_{kj}}{x^j}}P_ku,u){\displaystyle \frac{1}{i}}({\displaystyle \frac{A_{kj}}{x^k}}P_ju,u)\right\}`$ $`=`$ $`(\left[{\displaystyle \underset{k<j}{}}A_{kj}B_{kj}\right]u,u)+{\displaystyle \underset{k,j}{}}({\displaystyle \frac{A_{kj}}{x^k}}P_ju,u).\text{ }\text{ }\text{ }\text{ }`$ It is natural to consider the functions $`A_{jk}=A_{kj}`$ as coefficients of a 2-form which we will call a dual field or a dual form. The Iwatsuka identity (5.5) plays a very important role in the arguments below. As A. Iwatsuka noticed in , by choosing different dual fields $`A_{jk}`$ we can obtain different estimates leading to various sufficient conditions for the discreteness of spectrum of the operator $`H_{a,0}`$. We will develop this idea further by incorporating the electric potential into the picture. In this way different effective potentials $`V_{\text{eff}}`$ emerge such that the Molchanov condition for $`V_{\text{eff}}`$ implies that $`H_{a,V}`$ is essentially self-adjoint, semi-bounded below and has a discrete spectrum. We will see that the results of J. Avron, I. Herbst and B. Simon , A. Dufresnoy and A. Iwatsuka about the discreteness of spectrum can be extended to the case when both electric and magnetic fields contribute to the localization of the quantum particle. #### 5.2.2 First choice The first choice of the dual field which we will discuss, is related to the Dufresnoy sufficiency result and will lead to its generalization which takes into account the electric field. In the situation of , considering the operator $`H_{a,0}`$, we can take (5.6) $$A_{jk}(x)=\frac{B_{jk}(x)}{\|B(x)\|},$$ where $`\|B\|`$ is defined as in (1.10). (Note, however, that A. Dufresnoy used different arguments in .) To be able to use this choice we need to know that $`B(x)0`$ for large $`\|x\|`$ (i.e. if $`\|x\|R`$ where $`R>0`$ is sufficiently large. This was irrelevant in the situation when $`\|B(x)\|\mathrm{}`$ as $`x\mathrm{}`$ and $`V=0`$ which was considered in . In a more general situation which will be considered here we do not want to impose any a priori growth or non-vanishing requirement on $`B`$. By this reason we will smooth down the field (5.6) at the places where $`\|B\|`$ is small and use the dual field given by (1.24). We will use the following notation: $$\|B\|=\left(\underset{j<k}{}\|B_{jk}\|^2\right)^{1/2}.$$ Similarly we define $`\|\beta \|`$ and $`\|A\|`$ below. Denote also $`\beta _{jk}=B_{jk}/\|B\|`$. Differentiating $`\beta _{jk}`$ gives (5.7) $$\beta _{jk}=\|B\|^1B_{jk}\|B\|^2B_{jk}\|B\|.$$ Using the inequality $`\|\|B\|\|\|B\|`$, we obtain (5.8) $$\|\beta _{jk}\|\|B\|^1\|B_{jk}\|+\|B\|^2\|B_{jk}\|\|B\|,$$ hence by the triangle inequality (5.9) $$\|\beta \|2\|B\|^1\|B\|.$$ Differentiating (1.24) we obtain (5.10) $$A_{jk}=\chi (\|B\|)\beta _{jk}+\beta _{jk}\chi ^{}(\|B\|)\|B\|,$$ hence $$\|A_{jk}\|\chi (\|B\|)\|\beta _{jk}\|+\chi ^{}(\|B\|)\|\beta _{jk}\|\|B\|.$$ Therefore for $`1/2\|B\|1`$ we have (5.11) $$\|A_{jk}\|\|\beta _{jk}\|+2\|\beta _{jk}\|\|B\|$$ and (5.12) $$\|A\|\|\beta \|+2\|B\|2\|B\|^1\|B\|+2\|B\|6\|B\|.$$ Since $`A=0`$ for $`\|B\|1/2`$, the estimate (5.12) holds if $`\|B\|1`$. Now introducing a majorant function (5.13) $$M_B(x)=\{\begin{array}{cc}\|\beta _{jk}\|\hfill & \text{if }\|B\|1\text{,}\hfill \\ 6\|B\|\hfill & \text{if }\|B\|<1,\hfill \end{array}$$ we see that (5.14) $$\|A(x)\|M_B(x)\text{for all}x\text{}^n.$$ Note also that $$\|B\|^12(1+\|B\|)^1\text{if }\|B\|1\text{,}$$ and $$\|B\|2(1+\|B\|)^1\|B\|\text{if }\|B\|1.$$ These estimates together imply that $$M_B(x)12(1+\|B\|)^1\|B\|$$ for all $`x`$, hence (5.14) gives (5.15) $$\|A\|12(1+\|B\|)^1\|B\|$$ for all $`x\text{}^n`$. Now we are ready for Proof of Theorem 1.12. Let us use the Iwatsuka identity (5.5) with $`A_{kj}`$ given by (1.24). Clearly $$\underset{k<j}{}A_{kj}B_{kj}=\chi (\|B\|)\|B\|,$$ and the first term in the right hand side of (5.5) becomes $`(\chi (\|B\|)\|B\|u,u)`$. Since $$\chi (\|B\|)\|B\|\|B\|1,\text{and}\|A_{jk}\|1,$$ we obtain, using (1.25), that for any $`u\mathrm{Lip}_c(\text{})`$: $`(\|B\|u,u)`$ $`2{\displaystyle \underset{k<j}{}}(\|A_{kj}\|\|P_ku\|,\|P_ju\|)+{\displaystyle \underset{k,j}{}}(\left\|{\displaystyle \frac{A_{kj}}{x^k}}\right\|\|P_ju\|,\|u\|)+(u,u)`$ $`(n1){\displaystyle \underset{k=1}{\overset{n}{}}}P_ku^2+{\displaystyle \underset{j=1}{\overset{n}{}}}(\|P_ju\|,X\|u\|)+(u,u)`$ $`=(n1)h_{a,0}(u,u)+{\displaystyle \underset{j=1}{\overset{n}{}}}(\|P_ju\|,X\|u\|)+(u,u).`$ Choosing an arbitrary $`\epsilon >0`$, we see that the middle term in the right hand side is estimated by $$\epsilon \underset{j=1}{\overset{n}{}}P_j^2+\frac{n}{4\epsilon }(X^2u,u)=\epsilon h_{a,0}(u,u)+\frac{n}{4\epsilon }(X^2u,u).$$ Therefore we obtain $$(\|B\|u,u)(n1+\epsilon )h_{a,0}(u,u)+\frac{n}{4\epsilon }(X^2u,u)+(u,u),$$ and $$h_{a,0}(u,u)\frac{1}{n1+\epsilon }(\|B\|u,u)\frac{n}{4\epsilon (n1+\epsilon )}(X^2u,u)\frac{1}{n1+\epsilon }(u,u).$$ Now we can apply Theorem 1.5 with $$\mathrm{\Lambda }(x)=\frac{1}{n1+\epsilon }\|B\|\frac{n}{4\epsilon (n1+\epsilon )}X^2\frac{1}{n1+\epsilon },$$ which immediately leads to the desired result. Taking a specific majorant $`X(x)`$ in Theorem 1.12 we can make the result more specific. In this way we obtain for example the following ###### Corollary 5.7 The result of Theorem 1.12 holds if we replace the majorant function $`X`$ in the definition of the effective potential (1.26) either by $`\sqrt{n1}M_B`$ or by $`12\sqrt{n1}(1+\|B\|^1)\|B\|`$. Proof. It is sufficient to notice that $$\underset{k=1}{\overset{n}{}}\left\|\frac{A_{kj}}{x^k}\right\|\sqrt{n1}\|A\|,$$ due to the Cauchy–Schwarz inequality. We will need a notation for domination between two real-valued functions $`f,g`$ on $`S\text{}^n`$. Namely, we will write (5.16) $$fg\text{or}gf$$ on $`S`$ if for any $`\epsilon >0`$ there exists $`C(\epsilon )>0`$ such that (5.17) $$f(x)\epsilon g(x)+C(\epsilon )\text{for all}xS.$$ It is easy to see that $`fg`$ and $`gh`$ imply $`fh`$. If $`f(x)0`$, $`g(x)>0`$ for all $`xS`$, $`f,g`$ are locally bounded and $`g(x)+\mathrm{}`$ as $`x\mathrm{}`$, then the relation (5.16) is equivalent to $$f(x)=o(g(x))\text{as}x\mathrm{},$$ i.e. $`f(x)/g(x)0`$ as $`x\mathrm{}`$. Note also that if $`g`$ is semi-bounded below on $`S`$, then (5.18) $$fgf\|g\|f1+\|g\|.$$ The proof immediately follows from the implication $$gC\|g\|g+2C.$$ Now we are ready for the formulation of an important corollary of Theorem 1.12. Note that though this theorem does not explicitly include any regularity requirements on $`B`$, they are implicit in the requirements on the effective potential. We will make the result more explicit (though weaker) by invoking some explicit domination requirements on the majorant $`X`$. ###### Corollary 5.8 Let us assume that there exists $`\delta [0,1)`$ such that the effective potential (5.19) $$V_{\text{eff}}^{(\delta )}(x)=V(x)+\frac{\delta }{n1}\|B(x)\|$$ is semi-bounded below and satisfies the Molchanov condition $`(M_c)`$ with some $`c>0`$ (in particular, this holds if $`V_{\text{eff}}^{(\delta )}(x)+\mathrm{}`$ as $`x\mathrm{}`$). In addition assume that the square of the majorant function $`X(x)`$ from (1.25) is dominated by $`V_{\text{eff}}^{(\delta )}`$: (5.20) $$X^2V_{\text{eff}}^{(\delta )}\text{on}\text{}^n.$$ Then the operator $`H_{a,V}`$ is essentially self-adjoint, semi-bounded below and has a discrete spectrum. Remark. Since $`V_{\text{eff}}^{(\delta )}`$ is bounded below and locally bounded, replacing $`V_{\text{eff}}^{(\delta )}`$ by $`1+\|V_{\text{eff}}^{(\delta )}\|`$ in (5.20) leads to an equivalent relation. Proof of Corollary 5.8. Let us choose an arbitrary $`\epsilon >0`$. The condition (5.20) means that for any $`\kappa >0`$ there exists $`C(\kappa )>0`$ such that (5.21) $$X^2(x)\kappa V_{\text{eff}}^{(\delta )}(x)+C(\kappa ),x\text{}^n.$$ Now for any $`\epsilon >0`$ and $`\delta [0,1)`$ we can find $`\epsilon >0`$ and $`\delta ^{}[0,1)`$ such that $$\frac{\delta ^{}}{n1+\epsilon }=\frac{\delta }{n1},$$ so that $$V_{\text{eff}}^{(\delta ^{},\epsilon )}(x)=V_{\text{eff}}^{(\delta )}(x)\frac{n\delta ^{}}{4\epsilon (n1+\epsilon )}X^2(x)$$ for $`V_{\text{eff}}^{(\delta ,\epsilon )}`$ as in (1.26). It follows from (5.21) that $$V_{\text{eff}}^{(\delta ^{},\epsilon )}(x)(1\kappa )V_{\text{eff}}^{(\delta )}C(\delta ,\epsilon ,\kappa ),$$ hence $`V_{\text{eff}}^{(\delta ^{},\epsilon )}`$ is also semi-bounded below and satisfies $`(M_c)`$ with the same $`c>0`$ as for $`V_{\text{eff}}^{(\delta )}`$. Therefore the desired result follows from Theorem 1.12. Now we will replace the conditions on the majorant $`X`$ in Corollary 5.8 by an explicit choice of $`X`$ which is of A. Dufresnoy type . In this way we get a theorem which improves the sufficiency result of adding electric field and capacity to the picture. ###### Theorem 5.9 Let us assume that there exists $`\delta [0,1)`$ such that the effective potential (5.19) is semi-bounded below and satisfies the Molchanov condition $`(M_c)`$ with some $`c>0`$ (in particular, this holds if $`V_{\text{eff}}^{(\delta )}(x)+\mathrm{}`$ as $`x\mathrm{}`$). In addition assume that one of the following conditions $`(a)`$, $`(b)`$ is satisfied: $`(`$a$`)`$ $`M_B^2V_{\text{eff}}^{(\delta )}`$, or, in other words, $`\|\beta \|^2V_{\text{eff}}^{(\delta )}\text{on}\{x\|\|B(x)\|1\},`$ and $`\|B\|^2V_{\text{eff}}^{(\delta )}(x)\text{on}\{x\|\|B(x)\|<1\};`$ $`(`$b$`)`$ $`(1+\|B\|)^2\|B\|^2V_{\text{eff}}^{(\delta )}`$ on $`\text{}^n`$. Then the operator $`H_{a,V}`$ is essentially self-adjoint, semi-bounded below and has a discrete spectrum. Proof. According to (5.14) and (5.15) the conditions $`(a)`$ or $`(b)`$ imply that we can take the majorant $$X(x)=M_B(x)\text{or}12(1+\|B(x)\|)^1\|B(x)\|$$ respectively, and apply Corollary 5.8. The following even more explicit result also improves the sufficiency result by A. Dufresnoy . ###### Theorem 5.10 Let us assume that $`B_{jk}\mathrm{Lip}(\text{}^n)`$ for all $`j,k`$, and the following conditions are satisfied: (5.22) $$\|B\|C\text{if}\|B\|1;$$ (5.23) $$\|\beta \|=o(\|B\|^{1/2})\text{as}\|B\|\mathrm{}.$$ Assume also that there exists $`\delta [0,1)`$ such that the effective potential (5.19) is semi-bounded below and satisfies the Molchanov condition $`(M_c)`$ with some $`c>0`$. Then the operator $`H_{a,V}`$ is essentially self-adjoint, semi-bounded below and has a discrete spectrum. Proof. The result easily follows from Theorem 5.9 if $`V0`$ (or if $`V`$ is semi-bounded below). In the general case it can be deduced from Theorem 1.12 or Corollary 5.7 by the same argument which was used in the proof of Corollary 5.8, except that in this case instead of (5.21) we should estimate $`X^2(x)`$ by $`\kappa \|B(x)\|+C(\kappa )`$. Proof of Theorem 1.13. To deduce this theorem from Theorem 5.10 it is sufficient to show that the condition (5.23) imposed on $`\beta `$ in Theorem 5.10, follows from the condition $`(B_{3/2}^0)`$ which is imposed on $`B`$ and means that $$\|B\|=o\left((1+\|B\|)^{3/2}\right)\text{as}\|B\|\mathrm{}.$$ This implication immediately follows from the estimate (5.9). Remark 1. It is also possible to prove Theorem 1.13 and all previous theorems of this subsection (except the ones with conditions which explicitly include $`\beta _{jk}`$) by the following choice of the dual field: $$A_{jk}(x)=\frac{B_{jk}(x)}{B(x)},$$ where $`B=(1+\|B\|^2)^{1/2}`$. This choice leads to the arguments and estimates which are similar to the ones used above in the proof of Theorem 1.12. Remark 2. Sufficient measure conditions similar to Corollaries 1.15 and 1.16 hold for all types of effective potentials discussed above. #### 5.2.3 Second choice Now we will discuss the choice of the dual field $`A_{jk}`$ which is associated with the field suggested by A. Iwatsuka : $$A_{jk}(x)=\frac{B_{jk}(x)}{\|B(x)\|^2}.$$ It is proved in that this choice leads to a weakest regularity condition on $`B`$ which guarantees the discreteness of spectrum for $`H_{a,0}`$ provided $`\|B(x)\|\mathrm{}`$ as $`x\mathrm{}`$. We will improve the result of by adding a term taking into account the electric field. Hence by the same reason as above we will slightly modify this field as follows: (5.24) $$A_{jk}(x)=\frac{B_{jk}(x)}{B(x)^2}.$$ We will assume that $`B_{jk}\mathrm{Lip}_{\text{loc}}(\text{}^n)`$ for all $`j,k`$. A. Iwatsuka assumes also that the condition $`(B_2^0)`$ from Sect.1.4 holds. Choosing an arbitrarily small $`r>0`$ denote (5.25) $$\epsilon _x=\underset{yB(x,r)}{sup}\frac{1+\|B(y)\|}{B(y)^2}.$$ The conditions (1.23) and $`(B_2^0)`$ together are equivalent to the relation (5.26) $$\epsilon _x0\text{as}x\mathrm{}.$$ It is proved in that this implies that $`H_{a,0}`$ has a discrete spectrum. We will not a priori require either (1.23) or $`(B_2^0)`$ to be satisfied, but we will include $`\epsilon _x`$ above into an effective potential so that possible violation of (5.26) is compensated by the electric field. More precisely, we will prove ###### Theorem 5.11 Let us assume that $`B_{jk}\mathrm{Lip}(\text{}^n)`$ for all $`j,k`$ and define an effective potential (5.27) $$V_{\text{eff}}^{(\delta )}(x)=V(x)+\frac{\delta }{n1}\epsilon _x^{1/2}(1\epsilon _x).$$ If there exists $`\delta [0,1)`$ such that $`V_{\text{eff}}^{(\delta )}`$ is semi-bounded below and satisfies the Molchanov condition $`(M_c)`$ with a positive $`c>0`$, then $`H_{a,V}`$ is essentially self-adjoint, semi-bounded below and has a discrete spectrum. Proof. Note first that due to the triangle inequality $$\|B\|B1+\|B\|,$$ so $`B`$ can be replaced by $`\|B\|`$ (and vice versa) in the domination relations. We will again use the Iwatsuka identity (5.5). With the choice (5.24) we have (5.28) $$\underset{k<j}{}A_{kj}B_{kj}=\frac{\|B\|^2}{B^2}=\frac{B^21}{B^2}=1\frac{1}{B^2}.$$ Let us fix $`r>0`$ and choose $`x\text{}^n`$. Then using the inequality $$\|A_{kj}(y)\|B(y)^1\epsilon _x^{1/2},yB(x,r),$$ we obtain for any $`uC_c^{\mathrm{}}(B(x,r))`$ for the left hand side of (5.5): (5.29) $`2\underset{k<j}{{\displaystyle }}\text{Im}(A_{kj}P_ku,P_ju)2\underset{k<j}{{\displaystyle }}(\|A_{kj}\|\|P_ku\|,\|P_ju\|)`$ $`2\epsilon _x^{1/2}\underset{k<j}{{\displaystyle }}(\|P_ku\|,\|P_ju\|)(n1)\epsilon _x^{1/2}h_{a,0}(u,u).`$ Now let us estimate the last term in (5.5). Using the inequality $`\|\|B\|\|\|B\|`$, we obtain $`\|A_{kj}\|`$ $`=\left\|B^2B_{kj}2B_{kj}B^4\|B\|\|B\|\right\|`$ $`B^2\|B_{kj}\|+2\|B_{kj}\|B^4\|B\|\|B\|`$ $`B^2\|B_{kj}\|+2\|B_{kj}\|B^3\|B\|.`$ Therefore by the triangle inequality $$\|A\|3B^2\|B\|.$$ Now using the Cauchy-Schwarz inequality we find: $$\underset{k}{}\left\|\frac{A_{kj}}{x^k}\right\|\sqrt{n1}\|A\|3\sqrt{n1}B^2\|B\|.$$ Hence we can estimate the last term in (5.5) as follows: $`\left\|{\displaystyle \underset{k,j}{}}({\displaystyle \frac{A_{kj}}{x^k}}P_ju,u)\right\|3\epsilon _x\sqrt{n1}{\displaystyle \underset{j}{}}(\|P_ju\|,\|u\|)`$ $`3\epsilon _x\sqrt{n1}\left(\kappa _x{\displaystyle \underset{j}{}}P_ju^2+{\displaystyle \frac{n}{4\kappa _x}}u^2\right)`$ $`=3\epsilon _x\sqrt{n1}\left(\kappa _xh_{a,0}(u,u)+{\displaystyle \frac{n}{4\kappa _x}}u^2\right),`$ where $`uC_c^{\mathrm{}}(B(x,r))`$ and $`\kappa _x>0`$ is arbitrary. Let us fix $`\epsilon >0`$ and then choose $`\kappa _x`$ so that $$3\epsilon _x\sqrt{n1}\kappa _x=\epsilon \epsilon _x^{1/2}.$$ This leads to the estimate (5.30) $$\left\|\underset{k,j}{}(\frac{A_{kj}}{x^k}P_ju,u)\right\|\epsilon \epsilon _x^{1/2}h_{a,0}(u,u)+\frac{9n(n1)}{4\epsilon }\epsilon _x^{3/2}u^2.$$ Now from the identities (5.5), (5.28) and from the estimates (5.29), (5.30) we obtain for any $`uC_c^{\mathrm{}}(B(x,r))`$ $$(n1+\epsilon )h_{a,0}(u,u)\left[(1\epsilon _x)\epsilon _x^{1/2}\frac{9n(n1)\epsilon _x}{4\epsilon }\right](u,u).$$ (Here we used the obvious estimate $`B(y)^2\epsilon _x`$ if $`yB(x,r)`$.) The last estimate gives a lower bound for the bottom of the Dirichlet spectrum of $`H_{a,0}`$ on any ball $`B(x,r^{})`$ with $`r^{}<r`$: $$\lambda (B(x,r^{});H_{a,0})(n1+\epsilon )^1\left[\epsilon _x^{1/2}(1\epsilon _x)\frac{9n(n1)\epsilon _x}{4\epsilon }\right].$$ If $`\epsilon _x<1`$, we can replace $`\epsilon _x`$ by 1 in the last term to get (5.31) $$\lambda (B(x,r^{});H_{a,0})(n1+\epsilon )^1\left[\epsilon _x^{1/2}(1\epsilon _x)\frac{9n(n1)}{4\epsilon }\right].$$ This also obviously holds if $`\epsilon _x1`$. Now we can apply the IMS localization formula argument (see proof of Theorem 1.1) by considering a finite multiplicity covering of $`\text{}^n`$ by balls of radius $`r/2`$, to conclude that $`H_{a,0}\mathrm{\Lambda }(x)`$ with $$\mathrm{\Lambda }(x)=(n1+\epsilon )^1\epsilon _x^{1/2}(1\epsilon _x)C,$$ where $`C=C(r,n,\epsilon )`$. Then the statement of the Theorem follows from Theorem 1.5. ###### Corollary 5.12 Let us assume that $`B`$ satisfies the Iwatsuka conditions (1.23) and $`(B_2^0)`$ (or, equivalently, $`\epsilon _x0`$ as $`x\mathrm{}`$). Then the statement of the Theorem 5.11 holds with the effective potential $$V_{\text{eff}}^{(\delta )}(x)=V(x)+\frac{\delta }{n1}\epsilon _x^{1/2}.$$ Remark 1. If $`V=0`$ then this Corollary immediately implies the Iwatsuka result about the sufficient condition of the discreteness of the spectrum because then the Molchanov condition $`(M_c)`$ for the function $`x\epsilon _x^{1/2}`$ is fulfilled automatically. Remark 2. If $`\|B(x)\|\mathrm{}`$ as $`x\mathrm{}`$ and $`\|B(x)\|`$ varies sufficiently slowly (e.g. if $`\|\|B\|\|\|B\|`$, which in turn holds e.g. if $`B`$ satisfies the condition $`(B_1^0)`$), then $`\epsilon _x^{1/2}`$ becomes equivalent to $`\|B(x)\|`$, so the results of Sect.5.2.3 agree with the results of Sect.5.2.2. #### 5.2.4 Third choice Here we will discuss the choice of the dual field $`A_{jk}`$ which leads to a generalization of the result by J. Avron, I. Herbst and B. Simon . In this choice $`A_{jk}`$ are constants, though we argue on balls of a fixed radius $`r>0`$ and these constants may also depend on the choice of the ball. The advantage of the result obtained in this way is that no local smoothness of $`B`$ is required. Even the continuity of $`B`$ is not needed, though we will still maintain the requirement $`a_j\mathrm{Lip}_{\text{loc}}(\text{}^n)`$, hence $`B_{jk}L_{\text{loc}}^{\mathrm{}}(\text{}^n)`$. Let us choose a finite multiplicity covering of $`\text{}^n`$ by balls $`B(\gamma ,r)`$, $`\gamma \mathrm{\Gamma }`$. (For example $`\mathrm{\Gamma }`$ may be an appropriate lattice in $`\text{}^n`$.) Then for any $`\gamma \mathrm{\Gamma }`$ chose $`A_{kj}(\gamma )=A_{jk}(\gamma )`$ such that $`\|A(\gamma )\|=1`$ and denote (5.32) $$r_\gamma =\underset{yB(\gamma ,r)}{inf}\underset{k<j}{}A_{kj}(\gamma )B_{kj}(y).$$ We are interested to make the numbers $`r_\gamma `$ as big as possible. Clearly $$r_\gamma \underset{yB(\gamma ,r)}{sup}\|B(y)\|,$$ and this inequality is close to equality if $`B`$ is almost constant in $`B(\gamma ,r)`$ and we have chosen $`A_{kj}(\gamma )=B_{kj}(\gamma )/\|B_{kj}(\gamma )\|`$. It is proved in that if we can choose $`A_{kj}(\gamma )`$ so that (5.33) $$r_\gamma \mathrm{}\text{as}\gamma \mathrm{},$$ then the operator $`H_{a,0}`$ has a discrete spectrum. We will improve this result by adding an electric field into the picture. Denote (5.34) $$r(x)=\mathrm{min}\{r_\gamma \|x\gamma \|r\}.$$ ###### Theorem 5.13 Denote (5.35) $$V_{\text{eff}}^{(\delta )}(x)=V(x)+\frac{\delta }{n1}r(x),$$ and assume that there exists $`\delta [0,1)`$ such that the effective potential (5.35) is semi-bounded below and satisfies the Molchanov condition $`(M_c)`$ with some $`c>0`$. Then the operator $`H_{a,V}`$ is essentially self-adjoint, semi-bounded below and has a discrete spectrum. Proof. Let us use the Iwatsuka identity (5.5) for $`uC_c^{\mathrm{}}(B(\gamma ,r))`$ and with $`A_{kj}(x)=A_{kj}(\gamma )`$. The last term in (5.5) vanishes. The first term in the right hand side of (5.5) is estimated from below by $`r_\gamma (u,u)`$ and the left hand side is estimated from above by $`(n1)h_{a,0}(u,u)`$ (see arguments given in the proof of Theorem 1.12). Therefore we obtain: (5.36) $$h_{a,0}(u,u)\frac{r_\gamma }{n1}(u,u),uC_c^{\mathrm{}}(B(\gamma ,r)).$$ Now using the IMS localization formula as in the proof of Theorem 5.11 we conclude that $$H_{a,0}\frac{r(x)}{n1}C,$$ and it remains to apply Theorem 1.5. ## 6 Operators on manifolds Let $`(M,g)`$ be a Riemannian manifold (i.e. $`M`$ is a $`C^{\mathrm{}}`$-manifold, $`(g_{jk})`$ is a Riemannian metric on $`M`$), $`dimM=n`$. We will always assume that $`M`$ is connected. We will also assume that we are given a positive smooth measure $`d\mu `$, i.e. a measure which has a $`C^{\mathrm{}}`$ positive density $`\rho (x)`$ with respect to the Lebesgue measure $`dx=dx^1\mathrm{}dx^n`$ in any local coordinates $`x^1,\mathrm{},x^n`$, so we will write $`d\mu =\rho (x)dx`$. This measure may be completely independent of the Riemannian metric, but may of course coincide with the canonical measure $`d\mu _g`$ induced by the metric (in this case $`\rho =\sqrt{g}`$ where $`g=det(g_{jk})`$, so locally $`d\mu _g=\sqrt{g}dx`$). Denote $`\mathrm{\Lambda }_{(k)}^p(M)`$ the set of all $`k`$-smooth (i.e. of the class $`C^k`$) complex-valued $`p`$-forms on $`M`$. We will write $`\mathrm{\Lambda }^p(M)`$ instead of $`\mathrm{\Lambda }_{(\mathrm{})}^p(M)`$. A magnetic potential is a real-valued 1-form $`a\mathrm{\Lambda }_{(1)}^1(M)`$. So in local coordinates $`x^1,\mathrm{},x^n`$ it can be written as in (1.6) where $`a_j=a_j(x)`$ are real-valued $`C^1`$-functions of the local coordinates. The usual differential can be considered as a first order differential operator $$d:C^{\mathrm{}}(M)\mathrm{\Lambda }^1(M).$$ The deformed differential (compare (1.12)) $$d_a:C^{\mathrm{}}(M)\mathrm{\Lambda }_{(1)}^1(M),udu+iua,$$ is also well defined. The Riemannian metric $`(g_{jk})`$ and the measure $`d\mu `$ induce an inner product in the spaces of smooth forms with compact support in a standard way. The corresponding completions are Hilbert spaces which we will denote $`L^2(M)`$ for functions and $`L^2\mathrm{\Lambda }^1(M)`$ for 1-forms. These spaces depend on the choice of the metric $`(g_{jk})`$ and the measure $`d\mu `$. However we will skip this dependence in the notations of the spaces for simplicity of notations. The corresponding local spaces will be denoted $`L_{loc}^2(M)`$ and $`L_{loc}^2\mathrm{\Lambda }^1(M)`$ respectively. These spaces do not depend on the metric or measure. Formally adjoint operators to the differential operators with sufficiently smooth coefficients are well defined through the inner products above. In particular, we have an operator $$d_a^{}:\mathrm{\Lambda }_{(1)}^1(M)C(M),$$ defined by the identity $$(d_au,\omega )=(u,d_a^{}\omega ),uC_c^{\mathrm{}}(M),\omega \mathrm{\Lambda }_{(1)}^1(M).$$ (Here $`C_c^{\mathrm{}}(M)`$ is the set of all $`C^{\mathrm{}}`$ functions with compact support on $`M`$.) Therefore we can define the magnetic Laplacian $`\mathrm{\Delta }_a`$ (with the potential $`a`$) by the formula $$\mathrm{\Delta }_a=d_a^{}d_a:C^{\mathrm{}}(M)C(M).$$ Now the magnetic Schrödinger operator on $`M`$ is defined as (6.1) $$H=H_{a,V}=\mathrm{\Delta }_a+V,$$ where $`VL_{loc}^{\mathrm{}}(M)`$, i.e. $`V`$ is a locally bounded measurable function which is called the electric potential. We will always assume $`V`$ to be real-valued. Then $`H`$ becomes a symmetric operator in $`L^2(M)`$ if we consider it on the domain $`C_c^{\mathrm{}}(M)`$. Note that for $`a=0`$ the operator $`\mathrm{\Delta }_a`$ becomes a generalized Laplace-Beltrami operator $`\mathrm{\Delta }`$ on scalar functions on $`M`$ and it can be locally written in the form (6.2) $$\mathrm{\Delta }u=\frac{1}{\rho }\frac{}{x^j}(\rho g^{jk}\frac{u}{x^k}).$$ The following expressions for $`H_{a,V}`$ are also useful (): (6.3) $$H_{a,V}u=\mathrm{\Delta }u2ia,du+(id^{}a+\|a\|^2)u+Vu,$$ and in local coordinates (6.4) $$H_{a,V}u=\frac{1}{\rho }\left(\frac{}{x^j}+ia_j\right)\left[\rho g^{jk}\left(\frac{}{x^k}+ia_k\right)u\right]+Vu,$$ or in slightly different notations (6.5) $$H_{a,V}u=\frac{1}{\rho }P_j[\rho g^{jk}P_ku]+Vu=\frac{1}{\rho }(D_j+a_j)[\rho g^{jk}(D_k+a_k)u]+Vu,$$ where $`D_j=i/x_j`$. Now we need a condition on the Riemannian manifold $`(M,g)`$ which would allow extending the results above to a more general context. This is the condition of bounded geometry which means that the injectivity radius $`r_{inj}`$ is positive and all covariant derivatives of the curvature tensor $`R`$ are bounded: $`(a)`$ $`r_{inj}>0;`$ $`(b)`$ $`\|^mR\|C_m;m=0,1,\mathrm{}.`$ Here the norm $`\|\|`$ of tensors $`^mR`$ is measured with respect to the given Riemannian metric $`g`$. (See more details about these conditions and their use in .) We will also impose a bounded geometry condition on the measure $`d\mu `$. We will say that the triple $`(M,g,d\mu )`$ has bounded geometry if $`(M,g)`$ is a manifold of bounded geometry, and for a small $`r>0`$ in local geodesic coordinates in any ball $`B(x,r)`$ we have $`d\mu =\rho (x)dx`$ where $`\rho \epsilon >0`$ and for any multiindex $`\alpha `$ we have $`\|^\alpha \rho \|C_\alpha `$ with the constants $`\epsilon `$, $`C_\alpha `$ independent of $`x`$. In particular this automatically holds for the Riemannian measure $`d\mu =\sqrt{g}dx`$ if $`(M,g)`$ has bounded geometry. The requirement on the measure was not needed for $`\text{}^n`$ which reflects the fact that our methods do not work as well for manifolds as they do for $`\text{}^n`$. Some conditions at infinity are needed to guarantee that the operator $`H_{a,V}`$ is essentially self-adjoint in $`L^2(M)`$ \- see e.g. M. Shubin and references there for such conditions. The essential self-adjointness result by A. Iwatsuka can also be extended to manifolds of bounded geometry. The result which is of particular importance for us is essential self-adjointness of any semi-bounded below magnetic Schrödinger operator on any complete Riemannian manifold (see ), in particular on any manifold of bounded geometry. For a triple $`(M,g,d\mu )`$ of bounded geometry we can define a capacity of a compact set $`FB(x,r)`$ for a small $`r>0`$ by use of geodesic coordinates or by use of norms induced by the metric $`g`$ and the measure $`d\mu `$. Using geodesic coordinates for different balls $`B(x,r)`$ and $`B(x^{},r)`$ to measure $`\mathrm{cap}(F)`$ for $`FB(x,r)B(x^{},r)`$ leads to equivalent results, so it does not affect our formulations. Using Riemannian norms of tensors in the definition of capacity also leads to an equivalent result. After these explanations all results formulated above make sense and can be extended to the triples of bounded geometry with minor changes (some constants depending on geometry appear in the formulations). In case when $`a=0`$ and with the Riemannian measure $`d\mu `$ this was done by the authors in . We will give examples of such extensions in Sect.6. Let $`(M,g,d\mu )`$ be a triple with bounded geometry. Let us chose a ball $`B(x_0,r)`$ of a fixed sufficiently small radius $`r>0`$, and let $`x^1,\mathrm{},x^n`$ be local geodesic coordinates in this ball. ###### Lemma 6.1 Under the conditions above there exists $`C=C(M,g,d\mu )`$ such that (6.6) $$C^1\underset{j=0}{\overset{n}{}}P_ju_0^2h_{a,0}(u,u)C\underset{j=0}{\overset{n}{}}P_ju_0^2,uC_c^{\mathrm{}}(B(x_0,r)),$$ where $`P_j`$ are defined by (1.2) in the geodesic coordinates, and $`_0`$ means the norm in $`L^2(B(x_0,r);dx^1\mathrm{}dx^n)`$. Proof. The result immediately follows from the local presentation (6.5) for $`H_{a,V}`$ and from the bounded geometry requirements. $`\mathrm{}`$ This Lemma allows an easy extension of all local estimates in the balls of a fixed small radius $`r>0`$ to the case of the triples $`(M,g,d\mu )`$ of bounded geometry. Then we need to use extensions of the localization results from Sect.3. They can indeed be extended due to Gromov’s observation on the existence of good coverings and the manifolds extension of the IMS localization formula (3.5) – see the details in . Now we can formulate the simplest result for 2-dimensional manifolds. ###### Theorem 6.2 Assume that we have a bounded geometry triple $`(M,g,d\mu )`$ with $`n=dimM=2`$. Then there exists $`\delta _0>0`$ depending on the triple $`(M,g,d\mu )`$, such that the following holds. Let $`H_{a,V}`$ be a magnetic Schrödinger operator on $`M`$ and there exists $`\delta (\delta _0,\delta _0)`$ such that the effective potential given by (1.21) with $`B`$ identified with the ratio $`B/d\mu `$, is semi-bounded below and satisfies the Molchanov condition $`(M_c)`$ with some $`c>0`$. Then $`H_{a,V}`$ is essentially self-adjoint, semi-bounded below and has a discrete spectrum. Proof. The estimate $`H_{a,0}\delta B`$ follows from Lemmas 6.1 and 5.1 on balls $`B(x,r)`$. Then the IMS-localization formula (3.5) used for a finite multiplicity covering of $`M`$ by such balls (of the same radius) leads to a global estimate of the form $$H_{a,0}\delta BC,$$ where $`C=C(M,g,d\mu )`$. Now the desired result immediately follows from the manifold version of Theorem 1.5. $`\mathrm{}`$ Let us formulate an extension of Theorem 1.12 to manifolds. Let us define a smoothed direction of the magnetic field as a 2-form (or a skew-symmetric (0,2)-tensor): (6.7) $$A(x)=\chi (\|B(x)\|)\frac{B(x)}{\|B(x)\|},$$ where $`\chi \mathrm{Lip}([0,\mathrm{}))`$, $`\chi (r)=0`$ if $`r1/2`$, $`\chi (r)=1`$ if $`r1`$, $`\chi (r)=2r1`$ if $`1/2r1`$, so $`0\chi (r)1`$ and $`\|\chi ^{}(r)\|2`$ for all $`r`$. The norm $`\|B\|`$ is measured by the use of the given Riemannian metric $`g`$. ###### Theorem 6.3 Let us assume that $`B\mathrm{Lip}_{loc}(\text{}^n)`$, $`A`$ is defined by (6.7), and a positive measurable function $`X(x)`$ in $`\text{}^n`$ satisfies (6.8) $$\|A(x)\|X(x),xM,$$ where $``$ means the covariant derivative. Then there exist constants $`\delta _0>0`$ and $`C_0>0`$ depending only on the triple $`(M,g,d\mu )`$, such that the following is true. If there exists $`\delta (\delta _0,\delta _0)`$ such that the effective potential (6.9) $$V_{\text{eff}}(x)=V(x)+\delta \|B(x)\|C_0\delta X^2(x)$$ is bounded below and satisfies the Molchanov condition $`(M_c)`$ with some $`c>0`$ (in particular, this holds if $`V_{\text{eff}}(x)+\mathrm{}`$ as $`x\mathrm{}`$), then $`H_{a,V}`$ is essentially self-adjoint, semi-bounded below and has a discrete spectrum. In case $`B=0`$ and $`V`$ semi-bounded below the condition on $`V`$ in this theorem becomes necessary and sufficient due to the extension of Molchanov theorem given in . Proof of Theorem 6.3 is similar to the proof of Theorem 6.2. Other results from Sect.5 have similar extensions as well. ## 7 Other results A review of other results on the discreteness of spectrum of the Schrödinger operators and related topics can be found in . Here we will restrict ourselves to a few specific remarks concerning magnetic Schrödinger operators. Cwickel-Lieb-Rozenblum (CLR) type estimates for the number of negative eigenvalues for magnetic Schrödinger operators have been proved by J. Avron, I. Herbst and B. Simon (see also ), though with the right hand side independent of the magnetic field. The proof used heat kernel estimates based on the Feynman–Kac formula as in the paper by E. Lieb , and also the diamagnetic inequality (see , or , Sect.1.3). An analytic proof was provided by M. Melgaard and G. Rozenblum . The CLR estimates used for $`H_{a,V}M`$ for arbitrary $`M>0`$ obviously imply sufficient conditions for the discreteness of spectrum (namely, the finiteness of the right hand sides of these estimates for all $`M`$). However, the above mentioned results still do not allow to take into account any interaction between the electric and magnetic fields. Under some stronger conditions on the fields it is possible to obtain even asymptotics for the counting function $`N(\lambda ;H_{a,V})`$ for the eigenvalues of $`H_{a,V}`$. One of the first results of this kind is due to Y. Colin de Verdière . Numerous further results on such asymptotics can be found in (or extracted from) the book of V. Ivrii (see also references there). The Lieb–Thirring inequalities give explicit estimates for sums of powers of the negative eigenvalues, or, in other words, for the $`l^p`$ norms of the sequence of these eigenvalues. If $`p=\mathrm{}`$ this means an estimate for the number of negative eigenvalues as in the case of the CRL estimate. Under some conditions Lieb–Thirring type inequalities were obtained for $`H_{a,V}`$ by L. Erdös and for a similar Pauli operator by A.V. Sobolev . A Feynman type estimate for $`\mathrm{Tr}(\mathrm{exp}(tH_{a,V}))`$ in explicit purely classical terms was obtained by J.M. Combes, R. Schrader and R. Seiler . There exists a useful interaction between capacities and the Feynman–Kac formula. It was discussed e.g. in the books by I. Chavel , K. Ito and H. McKean , M. Kac and B. Simon . M. Kac and J.-M. Luttinger noticed an interesting relation between the scattering length and capacity.
warning/0007/math0007083.html	ar5iv	text	# References Some Applications of Localization to Enumerative Problems Aaron Bertram<sup>1</sup><sup>1</sup>1Supported in part by NSF Research Grant DMS-9970412 Dedicated to Bill Fulton on the occasion of his 60th birthday 1. Introduction. A problem in enumerative geometry frequently boils down to the computation of an integral on a moduli space. We have intersection theory (with Fulton’s wonderful Intersection Theory as a prime reference) to thank for allowing us to make rigorous sense of such integrals, but for their computations we often need to look elsewhere. If a torus lurks in the background, acting on the moduli space, then the Atiyah-Bott localization theorem allows one to express equivariant cohomology classes on the moduli space in terms of their “residues” living on the connected components of the locus of fixed points (i.e. the fixed submanifolds). This can be very useful for computations, particularly when the fixed submanifolds are points. We will use localization in a different way. Here, the moduli space itself will be a fixed submanifold for a torus action on a larger ambient space. Localization is applied in this context to relate residues on the moduli space to residues on simpler spaces by means of suitable equivariant maps. This point of view can lead immediately to remarkably simple derivations of some complicated-looking formulas, for example when applied to: (a) Schubert calculus on the flag manifold, and (b) Gromov-Witten invariants of rational curves. In (a) the partial flag manifold Fl$`(1,2,\mathrm{},m,n)`$ is realized as a fixed submanifold of a blown-up projective space $`𝐏(\text{Hom}(W,V))`$ where $`W`$ and $`V`$ are vector spaces of ranks $`m`$ and $`n`$ respectively, and all torus actions come from the “standard” torus action of $`(𝐂^{})^m`$ on $`W^{}`$. The full locus of fixed points is a disjoint union of $`m!`$ fixed submanifolds in this setting, each isomorphic to the partial flag manifold, but with different (equivariant) Euler classes. For this warm-up application, we will simply list the results of Jian Kong , where residues on the flag manifold are computed, resulting in particular in some new methods for computing Schubert calculus on the Grassmannian $`G(m,n)`$. It would be quite interesting to compare this with other methods (e.g. Gröbner bases) for making such computations. In (b), the Kontsevich-Manin moduli space of stable maps $`\overline{M}_{0,m}(X,\beta )`$ of rational curves with image homology class $`\beta `$ is realized as a fixed submanifold of the “graph space” $`\overline{G}_{0,m}(X,\beta ):=\overline{M}_{0,0}(X\times (𝐏^1)^m,(\beta ,1^m))`$ again with a standard torus action. The main applications take place in this setting. When $`m=1`$ we investigate the $`J`$-functions introduced by Givental in his generalization of the enumerative side of mirror symmetry to arbitrary projective manifolds (see ). The $`J`$-function is a polynomial associated to a complex projective variety $`X`$ and ample system of nef divisors which encodes all the one-point Gromov-Witten invariants. The coefficients of the $`J`$-function are push-forwards of residues, and our point of view on residues leads to a simple proof of the multiplicativity of the $`J`$-functions. Our point of view also leads to a non-obvious property of the $`J`$-function under push-forward. The $`J`$-function of projective space is computed in this context as an immediate consequence of the existence of a nice “linear” approximation to the graph space. Following Givental’s proof of the enumerative mirror conjecture for complete intersections in toric varieties, Kim was led to the formulation of a “quantum Lefschetz principle” relating the $`J`$-function for $`X`$ with $`J`$-functions for very ample divisor classes in $`X`$ . This has recently been proved by Y.P. Lee in the general case building on the proof in of the case $`X=𝐏^n`$, which we briefly discuss here. When $`m>1`$ there are many other fixed submanifolds in the graph space besides $`\overline{M}_{0,m}(X,\beta )`$, but they all are built out of Kontsevich-Manin spaces involving smaller $`m`$’s and/or smaller $`\beta `$’s. This has been exploited in joint work with Holger Kley to produce recursive formulas for $`m`$-point Gromov-Witten invariants, and in particular to prove that when the cohomology is generated by divisor classes, the $`m`$-point Gromov-Witten invariants can be “reconstructed” from one-point Gromov-Witten invariants. We will give the formula and an outline of the proof of reconstruction in the two-point case as a final application of localization. Another proof of reconstruction has been achieved with very different techniques and different formulas by Lee and Pandharipande . As a direct consequence of reconstruction, the small quantum cohomology of Fano complete intersections in $`𝐏^n`$, or indeed any toric variety, can be explicitly computed, since the one-point invariants are computed from the quantum Lefschetz principle. As another consequence, the quantum cohomology of products are determined by reconstruction, since the $`J`$-functions multiply. 2. Localization. When a torus $`T=(𝐂^{})^m`$ acts on a compact complex manifold $`M`$, the fixed submanifolds $`FM`$ are closed and embedded (of varying dimensions). There is an equivariant cohomology space H$`{}_{T}{}^{}(M,𝐐)`$ which is naturally a module over the cohomology of the classifying space H$`{}_{}{}^{}(BT,𝐐)𝐐[t_1,\mathrm{},t_m]`$. If $`E`$ is a linearized vector bundle over $`M`$, then there are equivariant Chern classes $`c_d^T(E)`$ taking values in H$`{}_{T}{}^{2d}(M,𝐐)`$, and in particular, the normal bundles $`N_{F/M}`$ to the fixed loci are canonically linearized (for the trivial action of $`T`$ on $`F`$) and yield equivariant Euler classes: $$ϵ_T(F/M)\text{H}^{}(F,𝐐)_𝐐𝐐[t_1,\mathrm{},t_m]\text{H}_T^{}(F,𝐐)$$ which are the top equivariant Chern classes of the normal bundles. The Atiyah-Bott localization theorem states that these Euler classes are invertible in H$`{}_{}{}^{}(F,𝐐)_𝐐𝐐(t_1,\mathrm{},t_m)`$ and one can recover an equivariant Chern class $`\alpha \text{H}_T^{}(M,𝐐)`$ uniquely (modulo torsion) as a sum of residues: $$\underset{F}{}i_{}\frac{i^{}\alpha }{ϵ_T(F/M)}$$ where $`i^{}`$ and $`i_{}`$ are the equivariant pull-back and push-forward associated to the equivariant inclusion $`i:FM`$. It follows from the uniqueness that taking residues is functorial. That is, if $`\mathrm{\Phi }:MM^{}`$ is an equivariant map and $`j:F^{}M^{}`$ is the inclusion of a component of the fixed submanifold, then: $$()\underset{F\mathrm{\Phi }^1(F^{})}{}\mathrm{\Phi }\|_{F}^{}{}_{}{}^{}\frac{i^{}\alpha }{ϵ_T(F/M)}=\frac{j^{}f_{}\alpha }{ϵ_T(F^{}/M^{})}$$ where the sum is over the components $`X`$ of the fixed locus that are contained in $`\mathrm{\Phi }^1(F^{})`$ and $`\alpha `$ is any equivariant cohomology class on $`M`$ (see or ). Thus if we are asked to integrate a cohomology class $`\gamma `$ on a compact complex manifold $`F`$, and if $`F`$ happens to be isomorphic to a component of the fixed locus of an action of $`T`$ on $`M`$ as above, then the formula above expresses residues at $`F`$ in terms of residues at $`F^{}`$ and at the other fixed loci contained in $`\mathrm{\Phi }^1(F^{})`$. If $`\gamma `$ can be expressed in terms of residues of equivariant cohomology classes, then this formula yields a relation among integrals of cohomology classes related to $`\gamma `$. This will be our point of view throughout the rest of this paper. 3. Flag Manifolds and Grassmannians. The partial flag manifold: $$\text{Fl}(1,2,\mathrm{},m,V)=\{V_1V_2\mathrm{}V_mV𝐂^n\|V_r𝐂^r\}$$ is a component of the fixed-point locus of an action of $`T`$ on $`M`$. In this case, $`M`$ is the blow-up of $`𝐏(\text{Hom}(W,V))`$ along: $$Z_1𝐏(W^{})\times 𝐏(V)Z_2\mathrm{}Z_{m1}𝐏(\text{Hom}(W,V))$$ where $`W𝐂^m`$ and $`Z_r`$ is the locus of maps of rank $`r`$. That is, $`M`$ is obtained by blowing up along $`Z_1`$, followed by the proper transform of $`Z_2`$, followed by the proper transform of $`Z_3`$, etc. If we choose a basis $`e_1,\mathrm{},e_m`$ of $`W`$ and let $`T`$ act on the dual space $`W^{}`$with weights $`(t_1,\mathrm{},t_m)`$, then this induces an action of $`T`$ on $`M`$, and the following are checked in : $``$ The intersection of the $`m1`$ exceptional divisors on $`M`$ is: $$E_1\mathrm{}E_{m1}\text{Fl}(1,2,\mathrm{},W^{})\times \text{Fl}(1,2,\mathrm{},m,V)$$ $``$ The fixed-point loci for the action of $`T`$ on $`M`$ are all contained in this intersection and correspond via the isomorphism above to: $$\mathrm{\Lambda }_I\times \text{Fl}(1,2,\mathrm{},m,V)$$ where $`\mathrm{\Lambda }_I`$ are the (isolated) fixed points of the action of $`T`$ on Fl$`(1,2,\mathrm{},W^{})`$, indexed by the permutations of $`m`$ letters so that the permutation $`(i_1,\mathrm{},i_m)`$ corresponds to the flag: $$\mathrm{\Lambda }_{(i_1,\mathrm{},i_m)}=\{x_{i_1}x_{i_1},x_{i_2}\mathrm{}\}$$ $``$ Let $`\zeta _i`$ be the relative hyperplane class for the projection: $$\text{Fl}(1,2,\mathrm{},i+1,V)\text{Fl}(1,2,\mathrm{},i,V)$$ pulled back to Fl$`(1,\mathrm{},m,V)`$. Then the equivariant Euler class to the fixed locus $`F_I=\mathrm{\Lambda }_I\times \text{Fl}(1,\mathrm{},m,V)`$ is: $$ϵ_T(F_I/M)=\underset{1j<km}{}(t_{i_k}t_{i_j})\underset{s=1}{\overset{m1}{}}(t_{i_{s+1}}t_{i_s}\zeta _s)$$ We are therefore in a position to apply the formula $`()`$ to the diagram: $$\begin{array}{ccc}M& \stackrel{\mathrm{\Phi }}{}& M^{}=𝐏(\text{Hom}(W,V))\\ & & \\ F_I& & \end{array}$$ On the right side, each fixed locus belongs to $`Z_1𝐏(\text{Hom}(W,V))`$ as $`F_i^{}=p_i\times 𝐏(V)`$ where $`p_i𝐏(W^{})`$ is the fixed point $`x_i`$. In that case, one computes: $$ϵ_T(F_i^{}/M^{})=\underset{si}{}(h+t_st_i)^n$$ where $`h`$ is the hyperplane class on $`𝐏(V)`$. An $`F_I`$ belongs to $`\mathrm{\Phi }^1(F_i^{})`$ exactly when $`I`$ is of the form $`(i,i_2,\mathrm{}.,i_m)`$. In that case, the induced map: $$\mathrm{\Phi }\|_{F_I}:\mathrm{\Lambda }_I\times \text{Fl}(1,\mathrm{},m,V)p_i\times 𝐏(V)$$ is the natural projection, which we will denote by $`\pi `$. Thus $`()`$ with $`\alpha =1`$ gives us the following interesting formula for Schubert calculus: Schubert Formula 1: $$\underset{\{I\|i_1=i\}}{}\pi _{}\left(\frac{1}{_{1j<km}(t_{i_k}t_{i_j})_{s=1}^{m1}(t_{i_{s+1}}t_{i_s}\zeta _s)}\right)=\frac{1}{_{si}(h+t_st_i)^n}$$ This formula encodes all the information about intersection numbers on the flag manifold of the form: $$_{\mathrm{Fl}(1,2,\mathrm{},m,V)}h^a\zeta _1^{a_1}\mathrm{}\zeta _{m1}^{a_{m1}}$$ Of course, the same intersection numbers could be obtained by applying the Grothendieck relation to each of powers of the $`\zeta _i`$. But there is a second formula which is much more interesting, involving cohomology classes pulled back from the Grassmannian under: $$\rho :\text{Fl}(1,2,\mathrm{},m,V)G(m,V)$$ Recall that such a cohomology class is a symmetric polynomial: $$\tau (q_1,\mathrm{},q_m)$$ in the Chern roots $`q_i`$ of the universal subbundle $`SV𝒪_{G(m,n)}`$. The main theorem of Kong’s thesis is the following: Schubert Formula 2: $$\underset{\{I\|i_1=i\}}{}\pi _{}\left(\frac{\rho ^{}\tau (q_1,\mathrm{},q_m)}{_{1j<km}(t_{i_k}t_{i_j})_{s=1}^{m1}(t_{i_{s+1}}t_{i_s}\zeta _s)}\right)$$ $$=\frac{\tau (h+t_1t_i,\mathrm{}.,h+t_mt_i)}{_{si}(h+t_st_i)^n}+\text{irrelevant terms}$$ where the irrelevant terms are monomials in the $`t_i`$ which do not appear on the left side of the equation. Example: When $`m=2`$, set $`i=1`$ above, $`t_1=0`$ and $`t_2=t`$. Then: $$\pi _{}\frac{\rho ^{}\tau (q_1,q_2)}{t(t\zeta _1)}=\frac{\tau (h,h+t)}{(h+t)^n}+\text{irrelevant terms}$$ If we consider the coefficients of $`t^2`$ on both sides and integrate, we get the following new way of doing Schubert calculus on G$`(2,V)`$: $$_{G(2,V)}\tau (q_1,q_2)=_{\mathrm{Fl}(1,2,V)}\pi ^{}h\rho ^{}\tau =\text{coeff of }h^nt^2\text{ in}\frac{h\tau (h,h+t)}{(h+t)^n}$$ Kong proves this formula by finding a suitable equivariant class $`\alpha `$ on $`M`$ which restricts to the given $`\tau `$ on each of the fixed components $`F_I`$. This $`\tau `$ is well enough approximated by the pull-back of the corresponding equivariant class of a split bundle on $`M^{}`$ to give the formula. The example above for $`m=2`$ can be similarly worked out for $`m>2`$ with the main difference being that there are $`(m1)!`$ terms on the left which sum together to the attractive formula on the right. It can be shown that this suffices to compute Schubert calculus, and it seems that an analysis of the complexity of this computation ought to be done. Finally, there is no obstruction to carrying out this program when $`V`$ is replaced by a vector bundle over a base variety $`X`$. Kong also shows how the Chern classes of $`V`$ figure into this “relative” setting in . 4. Gromov-Witten Invariants of Rational Curves. We will describe the relevant Konstsevich-Manin spaces (and maps among them) only set-theoretically for simplicity. The interested reader may go to the literature (e.g. ) for rigorous constructions of the spaces and morphisms. A map $`f:CX`$ from an $`m`$-pointed rational curve is stable if: $``$ $`C`$ has only nodes as singularities, and the marked points are smooth. $``$ Every component of $`C`$ collapsed by $`f`$ has at least $`3`$ distinguished points, i.e. marked points and/or nodes. $$\overline{M}_{0,m}(X,\beta )$$ is the Kontsevich-Manin moduli space of isomorphism classes of stable maps with $`m`$ marked points and image homology class $`\beta `$. If $`X`$ is “convex” (e.g. a homogeneous space) then this moduli space is smooth as an orbifold, of the expected dimension. Otherwise, there is a “virtual class” on $`X`$ with “all the expected properties” (see ). There is always an injective morphism: $$\overline{M}_{0,m}(X,\beta )\overline{G}_{0,m}(X,\beta )=\overline{M}_{0,0}(X\times (𝐏^1)^m,(\beta ,1^m))$$ where the latter space is the “graph space” associated to the former. Given a stable map $`f:CX`$ and points $`p_1,\mathrm{},p_mC`$, we obtain the image of $`[f]`$ in the graph space by attaching a copy of $`𝐏^1`$ to each of the points, gluing $`p_iC`$ to $`0𝐏^1`$, and collapsing each $`𝐏^1`$ to construct the resulting stable map $`g:C𝐏^1X`$. It is convenient to number the $`𝐏^1`$’s, so $`𝐏_i^1=𝐏(W_i)`$ is the particular $`𝐏^1`$ which we attach to $`p_i`$. The actions of $`𝐂^{}`$ on the dual spaces $`W_i^{}`$ with weights $`(0,t_i)`$ give a natural action of the torus $`T`$ on the product of the $`𝐏^1`$’s and hence on the graph space above. Moreover, the $`m`$-pointed Kontsevich-Manin space is one of the components of the fixed-locus for the torus action. One computes, using for example : $$ϵ_T(\overline{M}_{0,m}(X,\beta )/\overline{G}_{0,m}(X,\beta ))=\underset{i=1}{\overset{m}{}}t_i(t_i\psi _i)$$ where the $`\psi _i`$ are the “gravitational descendants” $`\psi _i=c_1(\sigma _i^{}(\omega ))`$. Here $`\omega `$ is the relative dualizing sheaf of the universal curve $`𝒞`$ over $`\overline{M}_{0,m}(X,\beta )`$ and $`\sigma _i`$ is the section of $`𝒞`$ corresponding to the $`i`$th marked point. The Case $`m=1`$: Here we let $`t=t_1`$ and $`\psi =\psi _1`$. If $`H`$ is an ample divisor on $`X`$, then following Givental, we define: $$J_{X,H}(q)=1+\underset{\beta 0}{}e_{\beta }^{}{}_{}{}^{}\left(\frac{1}{t(t\psi )}\right)q^{\mathrm{deg}_\mathrm{H}(\beta )}$$ where $`e_\beta :\overline{M}_{0,1}(X,\beta )X`$ is the evaluation map $`e_\beta ([f])=f(p)`$. Since only a finite number of classes $`\beta `$ have a given degree against $`H`$, this sum makes sense. More generally, we will suppose $`H`$ is a system $`H=(H_1,\mathrm{},H_r)`$ of (linearly independent) nef divisors, and that some linear combination of the $`H_i`$ is ample. In that case, we define: $$J_{X,H}(q)=1+\underset{\beta 0}{}e_{\beta }^{}{}_{}{}^{}\left(\frac{1}{t(t\psi )}\right)q_1^{\mathrm{deg}_{H_1}(\beta )}\mathrm{}q_r^{\mathrm{deg}_{H_r}(\beta )}$$ The following “functorial” properties of the $`J`$-function are easily proved once we recognize that the coefficients are push-forwards of residues. Product Formula: Suppose $`X`$ and $`X^{}`$ are simply connected projective manifolds (so the curve classes on $`X\times X^{}`$ are all of the form $`(\beta ,\beta ^{})`$) and $`H`$ and $`H^{}`$ are ample systems of divisors, as above. Then: $$J_{X\times X^{},(\pi _1^{}H,\pi _2^{}H^{})}(q,q^{})=\pi _1^{}J_{X,H}(q)\pi _2^{}J_{X^{},H^{}}(q^{})$$ Proof: Kontsevich-Manin spaces are functorial, in the sense that a map $`f:XY`$ gives rise to maps: $$f_{0,m}:\overline{M}_{0,m}(X,\beta )\overline{M}_{0,m}(Y,f_{}\beta )$$ and analogous compatible equivariant maps on the graph spaces. Thus the projection maps give rise to a diagram of “lifts” of the identity map: $$\begin{array}{ccc}\overline{G}_{0,1}(X\times X^{},(\beta ,\beta ^{}))& \stackrel{\mathrm{\Phi }}{}& \overline{G}_{0,1}(X,\beta )\times \overline{G}_{0,1}(X^{},\beta ^{})\\ & & \\ \overline{M}_{0,1}(X\times X^{},(\beta ,\beta ^{}))& \stackrel{\varphi }{}& \overline{M}_{0,1}(X,\beta )\times \overline{M}_{0,1}(X^{},\beta ^{})\\ & & \\ X\times X^{}& & X\times X^{}\end{array}$$ $`\mathrm{\Phi }`$ is birational when $`X`$ and $`X^{}`$ are convex (and “virtally birational” always) even though $`\varphi `$ is not birational (the two sides have different dimensions!). Thus $`\mathrm{\Phi }_{}1=1`$ and we may apply $`()`$ to the class $`1`$ to obtain: $$\varphi _{}\left(\frac{1}{t(t\psi )}\right)=\pi _1^{}\left(\frac{1}{t(t\psi )}\right)\pi _2^{}\left(\frac{1}{t(t\psi )}\right)$$ Further pushing forward to $`X\times X^{}`$ yields the desired product formula. Push-Forward Formula: Suppose $`f:XY`$ is given. Then there are equivariant classes $`f_{\beta }^{}{}_{}{}^{}1\mathrm{H}^{}(\overline{M}_{0,1}(Y,f_{}\beta ),𝐐)[t]`$ so that: $$f_{}J_{X,H}(q)=f_{}1+\underset{\beta 0}{}(e_{f_{}\beta })_{}\left(\frac{f_{\beta }^{}{}_{}{}^{}1}{t(t\psi )}\right)q_1^{\mathrm{deg}_{H_1}(\beta )}\mathrm{}q_r^{\mathrm{deg}_{H_r}(\beta )}$$ Proof: Here we consider the diagram of lifts of $`f`$: $$\begin{array}{ccc}\overline{G}_{0,1}(X,\beta )& \stackrel{\mathrm{\Phi }}{}& \overline{G}_{0,1}(Y,f_{}\beta )\\ & & j\\ \overline{M}_{0,1}(X,\beta )& \stackrel{\varphi }{}& M_{0,1}(Y,f_{}\beta )\\ & & \\ X& \stackrel{f}{}& Y\end{array}$$ and note that applying $`()`$ to the class $`1`$ again, we get: $$f_{}e_{\beta }^{}{}_{}{}^{}\left(\frac{1}{t(t\psi )}\right)=(e_{f_{}\beta })_{}\varphi _{}\left(\frac{1}{t(t\psi )}\right)=(e_{f_{}\beta })_{}\left(\frac{j^{}\mathrm{\Phi }_{}1}{t(t\psi )}\right)$$ hence the push-forward formula with $`f_{\beta }^{}{}_{}{}^{}1:=j^{}\mathrm{\Phi }_{}1`$. Remark: If $`f`$ is an embedding, then $`\varphi ^{}\psi =\psi `$, in which case the projection formula tells us that $`f_{\beta }^{}{}_{}{}^{}1=\varphi _{}1`$ is constant in $`t`$. It seems that in general, however, $`f_{\beta }^{}{}_{}{}^{}1`$ is not constant in $`t`$. It would be very interesting to compute it, for instance, in case $`f`$ is the inverse of a blow-up along a submanifold. The $`J`$-function of Projective Space: Let $`H`$ be the hyperplane class on $`𝐏^n`$. Then: $$J_{𝐏^n,H}(q)=\underset{d=0}{\overset{\mathrm{}}{}}e_{d}^{}{}_{}{}^{}\left(\frac{1}{_{k=1}^d(H+kt)^{n+1}}\right)q^d$$ Proof: $`\overline{G}_{0,1}(𝐏(V),d)`$ has a natural birational map to a “linear” space $`𝐏(\text{Hom}(\text{Sym}^d(W),V)`$, where $`W=W_1`$. A general element of the graph space is represented by a degree $`d`$ morphism $`f:𝐏^1𝐏^n`$ which maps to an $`n+1`$-tuple of degree $`d`$ polynomials $`(p_0(x,y):\mathrm{}:p_n(x,y))`$ with no common factors. When the curve underlying the stable map picks up extra components, then the $`n+1`$-tuple of polynomials picks up common factors. In particular, the image of $`\overline{M}_{0,1}(𝐏^n,d)`$ under this weighted blow-down is a copy of $`𝐏^n`$, embedded via: $$\{x^d\}\times 𝐏(V)𝐏(\text{Sym}^d(W^{}))\times 𝐏(V)𝐏(\text{Hom}(\text{Sym}^d(W),V)$$ Thus, we have the diagram: $$\begin{array}{ccc}\overline{G}_{0,1}(𝐏^n,d)& \stackrel{\mathrm{\Phi }}{}& 𝐏(\text{Hom}(\text{Sym}^d(W),V)\\ & & \\ i& & j\\ & & \\ \overline{M}_{0,1}(𝐏^n,d)& \stackrel{e_d}{}& 𝐏^n\end{array}$$ One computes (see ) $`ϵ_T(𝐏^n/𝐏(\text{Hom}(\text{Sym}^d(W),V)))=_{k=1}^d(H+kt)^{n+1}`$ so that $`()`$ now applies with the class $`1`$, giving us: $$e_{d}^{}{}_{}{}^{}\left(\frac{1}{t(t\psi )}\right)=\frac{1}{_{k=1}^d(H+kt)^{n+1}}$$ proving the formula. Quantum Lefschetz Hyperplane: We will limit ourselves to considering hypersurfaces in $`𝐏^n`$, as in . See for the general version. Let $`f:X𝐏^n`$ be a hypersurface of degree $`l`$, and let $`H`$ denote the hyperplane class, either on $`𝐏^n`$ or on $`X`$. Let: $$I_{X/𝐏^n,H}(q)=\underset{d=0}{\overset{\mathrm{}}{}}\frac{_{k=0}^{dl}(lH+kt)}{_{k=1}^d(H+kt)}q^d$$ (a) If $`l<n`$, then $`f_{}J_{X,H}(q)=I_{X/𝐏^n,H}(q)`$. (b) If $`l=n`$, then $`f_{}J_{X,H}(q)=e^{\frac{l}{t}q}I_{X/𝐏^n,H}(q)`$. (c) If $`l=n+1`$, then there are $`a(q),b(q)q𝐐[[q]]`$ so that: $$f_{}J_{X,H}(q)=e^{\frac{H}{t}a(q)+b(q)}I_{X/𝐏^n,H}(qa(q))$$ To prove this, one uses the diagram for $`𝐏^n`$ and observes that $$f_{}J_{X,H}(q)=lH+\underset{d>0}{}e_{d}^{}{}_{}{}^{}\left(\frac{j^{}\mathrm{\Phi }_{}[X]_T}{t(t\psi )}\right)q^d$$ where $`[X]_T`$ is the equivariant Chern class: $$[X]_T=c_{dl+1}^T(\pi _{}e_d^{}𝒪_{𝐏^n}(l))\text{H}_T^{}(\overline{G}_{0,1}(𝐏^n,d),𝐐)$$ Thus the proof of quantum Lefschetz amounts to a detailed analysis of the class $`j^{}\mathrm{\Phi }_{}[X]_TH^{}(𝐏^n,𝐐)[t]`$. This is obtained by decomposing $`[X]_T`$ along boundary strata of the graph space by means of intersection theory. In particular, the open stratum of the graph space contributes $`_{k=0}^{dl}(lH+kt)`$ via an approximation of $`\pi _{}e_d^{}𝒪_{𝐏^n}(l)`$ by $`\text{Sym}^{dl}(W^{})\mathrm{\Phi }^{}𝒪(l)`$ in much the same way that the graph space is approximated by $`𝐏(\text{Hom}(\text{Sym}^d(W),V))`$. In the case $`l<n`$, this is the only stratum which contributes to $`j^{}\mathrm{\Phi }_{}[X]_T`$, giving us (a). In the other cases, the boundary strata do contribute, but in a self-similar manner. When tallied up, these contributions give formulas (b) and (c) in the cases $`l=n`$ and $`l=n+1`$ respectively. It is unknown whether a more general “change of coordinates” analogous to (b) and (c) occurs in the general type cases $`l>n+1`$. The case $`m>1`$. Reconstruction: In , reconstruction theorems make use of the following diagrams of K-M spaces and graph spaces: $$\begin{array}{ccc}\overline{G}_{0,m}(X,\beta )& \stackrel{\mathrm{\Phi }}{}& \overline{G}_{0,m1}(X,\beta )\times \overline{M}_{0,0}(𝐏_1^1\times 𝐏_m^1,(1,1))\\ & & \\ i& & j\\ & & \\ \overline{M}_{0,m}(X,\beta )& \stackrel{\pi }{}& \overline{M}_{0,m1}(X,\beta )\end{array}$$ $`\mathrm{\Phi }`$ is derived, as in the product formula, from projections. $`\overline{M}_{0,m1}(X,\beta )`$ is included in the graph space in the ordinary way, and the inclusion $`j`$ is given by the additional inclusion of the point corresponding to the inclusion of the intersecting lines $`\{0\}\times 𝐏_m^1𝐏_1^1\times \{0\}`$ in $`𝐏_1^1\times 𝐏_m^1`$. The fixed loci contained in $`\mathrm{\Phi }^1(\overline{M}_{0,m1}(X,\beta ))`$, in addition to $`\overline{M}_{0,m}(X,\beta )`$ are isomorphic to one of the following: $$\overline{M}_{0,k+1}(X,\beta _1)\times _X\overline{M}_{0,mk}(X,\beta \beta _1))\text{or}\overline{M}_{0,m1}(X,\beta )$$ and the induced maps to $`\overline{M}_{0,m1}(X,\beta )`$ are the gluing maps to boundary divisors (see ) and the identity map, respectively. The equation $`()`$ now tells us that given an equivariant cohomology class $`\alpha `$ on $`\overline{G}_{0,m}(X,\beta )`$, there is a relation among the residues of $`\alpha `$ along the fixed loci listed above, as well as the residue of $`\mathrm{\Phi }_{}\alpha `$ along the fixed locus $`\overline{M}_{0,m1}(X,\beta )`$. So the question now becomes, how to find interesting equivariant classes $`\alpha `$ on the graph space? The only source we know of to produce good residue classes comes from the linear approximation to $`\overline{G}_{0,1}(𝐏^n,d)`$. Namely, suppose a morphism (not necessarily an embedding) $`f:X𝐏^n`$ is given. Then we can pull back equivariant cohomology classes via: $$\overline{G}_{0,m}(X,\beta )\overline{G}_{0,1}(X,\beta )\overline{G}_{0,1}(𝐏^n,d)𝐏(\text{Hom}(\text{Sym}^d(W),V))$$ After all the equivariant Euler classes are computed, recursive formulas are obtained. Thus in this context the necessity of considering cohomology classes generated by divisor classes springs from our inability to find useful equivariant classes not coming from the linear approximation spaces to $`\overline{G}_{0,1}(𝐏^n,d)`$. As an example of the reconstruction theorems we obtain, we include the most useful one, which, in case the cohomology of $`X`$ is generated by divisor classes, already suffices to express (small) quantum cohomology in terms of the $`J`$-function. Reconstruction Theorem for $`2`$-Point Invariants: Given $`f:X𝐏^n`$, let $`H`$ be the hyperplane class on $`𝐏^n`$ and on $`X`$, and define: $`F_\beta (t)=e_{\beta }^{}{}_{}{}^{}\left(\frac{1}{t(t\psi )}\right)`$ (these are the coefficients of $`J`$) and $`G_\beta (\gamma ,t)=e_{\beta }^{1}{}_{}{}^{}\left(\frac{e_{\beta }^{2}{}_{}{}^{}\gamma }{t\psi _2}\right)\text{for evaluation maps}e_\beta ^1,e_\beta ^2:M_{0,2}(X,\beta )X\text{and}`$ cohomology class $`\gamma \mathrm{H}^{}(X,𝐐)`$, extended by linearity in the first factor. Then the expression: $$G_\beta (H^a,t)+\left(\underset{\beta _1+\beta _2=\beta }{}G_{\beta _1}(F_{\beta _2}(t)(Hd_{\beta _2}t)^a,t)\right)+F_\beta (t)(Hd_\beta t)^a$$ is polynomial in $`t`$, where $`d_\beta `$ is the degree of $`f_{}\beta \mathrm{H}_2(𝐏^n,𝐙)`$. Since $`G_\beta (H^a,t)`$ is polynomial in $`t^1`$ (with no contant term), this formula expresses $`G_\beta (H^a,t)`$ in terms of coefficients of $`J`$ and $`G_\beta ^{}(H^a,t)`$ for smaller $`\beta ^{}`$. Hence it inductively determines $`G_\beta (H^a,t)`$ in terms of $`J`$. University of Utah, Salt Lake City, UT 84112 bertram@math.utah.edu
warning/0007/astro-ph0007061.html	ar5iv	text	# Synchrotron Radiation as the Source of Gamma-Ray Burst Spectra. ## 1 Introduction The spectrum of a gamma-ray burst (hereafter, GRB) is a crucial element in understanding the nature of the event. It can provide information about the burst energetics, the local magnetic field, particle distributions and acceleration mechanisms, and the overall expansion of the fireball. Unlike the extreme variation in the light curves, the spectra of GRBs are fairly homogeneous. The photon spectrum of most GRBs in the BATSE spectral energy range can be parameterized by a broken power law with low and high energy photon spectral indices, $`\alpha `$ and $`\beta `$, and a break energy, $`E_b`$ (e.g. Band et al., 1993). For the majority of bursts, $`\alpha >2`$ and $`\beta <2`$, in which case the $`\nu F_\nu `$ spectrum has a maximum at $`E_p`$ that represents where the burst emits most of its energy; for a Band spectrum, $`E_p=(2+\alpha )E_b/(\alpha \beta )`$. In most emission models, $`E_p`$ (or $`E_b`$) usually reflects some characteristic electron energy. In models of GRB emission with a power law energy spectrum for the radiating particles, it could signify the presence of a cutoff in the underlying particle distribution, or the energy above which all electrons “cool” - rapidly radiating all of their kinetic energy within the source (see Piran, 1999, for a review). The high energy photon index, $`\beta `$, usually reflects the steepness of the particle energy distribution. The value of the low energy photon index, $`\alpha `$, on the other hand varies significantly from model to model and can be due to several factors. Therefore, it has greater potential to distinguish between the different scenarios for GRB emission. However, despite years of availability of high quality spectral data from GRBs, the burst emission process remains ambiguous. It was suggested (e.g. Katz, 1994) that synchrotron emission is a likely source of radiation from GRBs, and later shown (Tavani, 1996) that an optically thin synchrotron spectrum from a power law distribution of relativistic electrons with a sharp minimum energy cutoff ($`N(\gamma )\gamma ^p,\gamma >\gamma _m`$, where $`\gamma `$ is the electron Lorentz factor) and with an isotropic pitch angle distribution provides a good fit to some bursts. However, some features seen in the low energy portion of GRB spectra can not be explained by this simple synchrotron model. This model predicts that $`E_p`$ is related to the minimum value of the electron Lorentz factor, $`\gamma _m`$ ($`E_p\gamma _m^2h\nu _B`$, where $`\nu _B=eB/2\pi m_ec`$ is the gyrofrequency in a magnetic field with mean perpendicular component $`B`$), so that the asymptotic value of the low energy photon index, $`\alpha `$, should be a constant value of $`2/3`$ (e.g., Pacholczyk, 1970). Hence, if our spectral fits are determining the asymptote accurately, the $`\alpha `$ distribution should be a very narrow distribution centered about $`2/3`$. Figure 1 shows the time averaged distributions of $`\alpha `$ and $`\beta `$ taken from Preece et al. (1999). The $`\alpha `$ distribution clearly appears to disagree with this simple synchrotron model. It has also been claimed (e.g. Crider et al., 1997) that spectral evolution of $`\alpha `$ (and $`E_p`$, for that matter) throughout a burst is inconsistent with the standard synchrotron scenario, at least in the context of a single (external) shock model. Consequentially, other models - usually involving Compton scattering (Brainerd, 1994, Liang & Kargatis, 1996) - are invoked to explain this “anomolous” spectral behavior. Another difficulty with the synchrotron model stems from theoretical considerations. For typical spectral peak energies $`E_p`$ few$`\times 100`$ keV, and “typical” values for the electron Lorentz factors $`\gamma 100`$ and bulk Lorentz factor $`\mathrm{\Gamma }100`$, magnetic fields can reach values of about $`10^7G`$ (such fields are also reached simply assuming equipartition between photon, electron and magnetic field energy density). As a result, the synchrotron lifetime in the observer’s frame ($`\tau _s=(\gamma /\dot{\gamma _s})(1+z)/\mathrm{\Gamma })10^6s`$, where $`\dot{\gamma _s}\gamma ^2B^2`$ is the synchrotron energy loss rate, $`\mathrm{\Gamma }`$ is the bulk Lorentz factor of the medium, and $`z`$ is the GRB redshift) is much shorter than the physical timescales relevant in the standard fireball model for GRBs (again, see Piran, 1999 for a review), as well as the detector integration time. Thus, one expects these radiative loss effects to be apparent in the GRB spectrum. For example, if the electrons are injected in a source region where those with energies $`\gamma >\gamma _\mathrm{c}`$ lose all their energy to radiation (while those with $`\gamma <\gamma _\mathrm{c}`$ escape the emission region before suffering significant losses of their kinetic energy), then one obtains the so-called cooling spectrum which has an additional break at a photon energy $`E_c\gamma _\mathrm{c}^2h\nu _B`$. Depending on the relative values of $`\gamma _m`$ and $`\gamma _c`$, one has a variety of values for the spectral indices. For example, if $`\gamma _\mathrm{c}>\gamma _{\mathrm{min}}`$ (and we have the same power law distribution of electron energy mentioned above) then one has an index $`\beta ^{}=(p+2)/2`$ above $`E_\mathrm{c}`$. In the case when $`\gamma _\mathrm{c}<\gamma _{\mathrm{min}}`$, all of the injected electrons lose their energy to radiation and we therefore have an index $`\alpha ^{}=3/2`$ between $`\gamma _\mathrm{c}`$ and $`\gamma _{\mathrm{min}}`$. The peak of the $`\nu F_\nu `$ spectrum, $`E_\mathrm{p}`$, could represent either the cooling or the minimum Lorentz factor break, depending on their relative values and their positions in the BATSE spectral window (see, e.g. , Sari, Piran & Narayan 1998). Hence, in a cooling spectrum, we expect either $`\alpha `$ or $`\beta `$ to be about -1.5, which is not borne out by the observations as seen in Figure 1. Ghissilini et al. (1999) use this argument to rule out synchrotron emission in favor of a quasi-thermal inverse Compton model. We agree that the simple cooling spectrum is ruled out by the observations. However, we also believe that this does not rule out synchrotron emission in general, and that all of these above mentioned discrepancies arise from incorrect simplifying assumptions in the models. Let us first consider the arguments based on the discrepancy with the cooling spectrum. We believe the reason for this discrepancy lies in the artificial separation between particle acceleration and radiation. In this model particles are accelerated in one region, and then either the acceleration processes automatically cease or the particles are injected into another region in which they radiate away all of their energy. We believe that such a model of the particle acceleration and radiation processes is an unrealistic situation. It is likely that both acceleration and cooling or radiation processes take place behind the shock continually throughout the emission episode. The acceleration rate $`R_{acc}`$ and the synchrotron loss rate $`R_{syn}=\gamma /\dot{\gamma _s}`$, where $`\dot{\gamma _s}`$ is defined in the previous paragraph, compete with each other in forming the instantaneous spectrum of the electrons which produces the observed synchrotron spectrum. Since the synchrotron loss increases rapidly with energy, its effect is to inhibit acceleration beyond some maximum value $`\gamma _{\mathrm{max}}`$ where the two rates become equal, $`R_{acc}(\gamma _{\mathrm{max}})=R_{syn}(\gamma _{\mathrm{max}})`$. Thus, we expect a power law electron spectrum for $`\gamma <\gamma _{\mathrm{max}}`$, with a relatively sharp cutoff outside $`\gamma _{\mathrm{max}}`$. In reality, the spectrum of the accelerated electrons could be more complex and may lead to the formation of a plateau just below $`\gamma _{max}`$ or to other features (e.g., see Petrosian & Donaghy, 2000). Considering observations up to the GeV range of some GRBs by the EGRET instrument on CGRO, $`\gamma _{\mathrm{max}}`$ must be large enough so that the energy of the corresponding synchrotron photons will far exceed the BATSE range. Thus in this model $`E_p`$ is related to $`\gamma _m`$, the minimum energy of the accelerated electrons, and not to any cooling break. We point out that the high energy photon indices for a simple cooling spectrum and the so-called “instantaneous” spectrum (in which cooling effects are negligible) are $`(p+2)/2`$ and $`(p+1)/2`$, respectively. Observed values of $`\beta 2.1`$ indicate values of $`p2.2`$ for the simple cooling model and $`3.2`$ for the instantaneous spectrum. Recent studies of particle acceleration via the Fermi mechanism in relativistic shocks determine a “universal” index of $`p2.2`$ (e.g., Gallant et al., 2000, Guthman et al., 2000), which is more consistent the values expected from the cooling model, and this has been used as an argument against the instantaneous spectrum. However, we note that $`\beta `$ is usually not very well constrained by the data and the $`1\sigma `$ errors are frequently of the order $`\pm 1/2`$, so that observations cannot distinguish between the cooling and instantaneous spectra. Moreover, the cooling spectrum itself is not consistent with the universal index for the significant fraction of bursts with $`\beta 2.5`$. In any case, this universal index is based on a very simplified model; particle acceleration in the context of internal shocks has not been carefully studied, and it is unclear whether the Fermi mechanism will even work at all (see, e.g. Lyutikov & Blackman, 2000 and Kirk et al., 2000, the latter of whom show that including a magnetic field at the acceleration site will steepen this universal index). In most of our analysis below, we will use the instantaneous synchrotron spectrum, and not the simple cooling spectrum; we will show that this instantaneous synchrotron spectrum can indeed explain the behavior seen in the BATSE spectral data. Another simplification in standard synchrotron models that leads to discrepancies between the theory and the observations is the assumption that the GRB emission region is optically thin. Under certain conditions, the medium may in fact become opaque, e.g. to synchrotron self-absorption (discussed in §2.1). This will produce a steep low energy cutoff in the photon spectrum. This is one possible explanation for those bursts with $`\alpha `$ values above the so-called “line of death”, $`\alpha >2/3`$ (Preece et al., 1998a) A third unrealistic simplification involves the electron spectrum at low energies. The asymptote of the optically thin photon spectrum $`\alpha =2/3`$ is obtained if one assumes that the electron spectrum cuts off sharply below $`\gamma _m`$. It is unlikely that the acceleration process will give rise to such a spectrum, which is subject to plasma instabilities. Particles in front of a relativistic shock - after crossing it - will acquire a Lorentz factor comparable to the bulk Lorentz factor of the shock. Subsequent interaction with the turbulent medium will accelerate them to higher energies (up to $`\gamma _{max}`$ described above), but the same processes will decelerate some to lower energies (see, e.g., Park & Petrosian, 1995, and references cited therein). Plasma instabilities (e.g. the two stream instability) can further smooth this peaked spectrum (and in some cases give rise to a nearly flat spectrum for $`\gamma <\gamma _m`$). Inclusion of such effects will clearly modify the behavior of the low energy synchrotron emission. A fourth assumption in the simple synchrotron model, which may not always be correct, is that the pitch angle distribution of the electrons is isotropic. An isotropic distribution is expected if the acceleration mechanism is more efficient in changing the pitch angle rather than the energy of the particles, which is commonly the case. However, in low density high magnetic field plasmas, the opposite may be true so that a highly anisotropic distribution of pitch angles is expected; as a result, the low energy synchrotron spectrum will differ from the usual simple model. Note that as long as the magnetic field lines have random orientations, the total emission in the rest frame will be isotropic. For details, see Epstein (1973). In what follows, we will show that incorporating all of these effects can lead to a wide variety of low energy GRB spectral behavior. However, there is another important cause of the discrepancy between the predicted and observed spectral behavior, which arises as a result of the finite bandwidth of the instrument and the spectral fitting procedure. One cannot assume that the spectral fits to phenomenological models are able to accurately determine the asymptotic values of the spectral indices. This is because the actual spectra do not show sharp breaks at $`E_p`$ (or $`E_b`$) from the high energy spectral index to the low energy spectral index - rather, there is a smooth transition between the two indices. The fitted values of the low energy (and high energy, for that matter) spectral index depends on how far the spectral window extends below (or above) $`E_p`$. We must take these effects into consideration when testing any model. In this paper we show that when these procedural effects as well as more realistic synchrotron models are incorporated, synchrotron radiation can accomodate both the shape of the distributions and temporal evolution of GRB spectral parameters, and the discrepancies discussed above are resolved. Of course this does not prove the synchrotron model is the unique emission mechanism for GRBs, but merely shows that it is consistent with the existing data. In §2, we discuss the various spectral shapes obtained from a general form for synchrotron emission, allowing for the possibility for self-absorption, a smooth (instead of sharp) cutoff to the electron energy distribution, and a small pitch angle distribution. In §3, we consider the instrumental and fitting effects and point out a correlation that exists between $`\alpha `$ and $`E_p`$ as determined by the Band spectrum (Band, 1993). In §4, we use this relationship and knowledge of the $`E_p`$ distribution to determine the expected $`\alpha `$ distribution, and compare this with the observed distribution. We also discuss the presence of an absorption cutoff and evidence of small pitch angle scatering observable by BATSE in some GRBs to explain particularly those bursts with $`\alpha >2/3`$. These results also predict a certain relationship between $`\alpha `$ and $`E_p`$ during the temporal evolution of a GRB. In §5, we give examples of spectral evolution in GRBs and discuss whether these are consistent with synchrotron emission models. A summary of our conclusions are given in §6. ## 2 Synchrotron Spectra Synchrotron radiation will occur when relativistically charged particles have a component of their velocity perpendicular to a local magnetic field (i.e. a non-zero pitch angle). The importance of synchrotron emission (compared to, say, inverse Compton or brehmstrahlung radiation) depends - among other things - on the strength of the magnetic field. Given the relativistic nature of GRBs and probable physical conditions (e.g. significant magnetic fields), it is likely that synchrotron radiation plays a role in the emission from GRBs. The general form for an instantaneous synchrotron energy spectrum from electrons with a (homogeneous) power law distribution of Lorentz factors with a sharp cutoff, $`N(\gamma )=N_o\gamma ^p\mathrm{\Theta }(\gamma \gamma _m)`$ (where $`\mathrm{\Theta }`$ is the Heaviside step function), and with an isotropic pitch angle distribution, is given by (e.g. Pacholczyk, 1970) $$F_\nu =𝒜\nu ^{5/2}[\frac{I_1}{I_2}]\times [1.0\mathrm{exp}[Q\nu ^{(p+4)/2}I_2]]$$ (1) $$I_1=_0^{\nu /\nu _m}𝑑xx^{(p1)/2}_x^{\mathrm{}}K_{5/3}(z)𝑑z$$ (2) $$I_2=_0^{\nu /\nu _m}𝑑xx^{p/2}_x^{\mathrm{}}K_{5/3}(z)𝑑z,$$ (3) Here, we have assumed that the electrons are extremely relativistic, $`\gamma _m1`$, and that the magnetic field at the source is randomly oriented, or that the emission is isotropic. The coefficient $`𝒜`$ is the normalization and contains factors of the magnetic field, $`B`$, bulk Lorentz factor, $`\mathrm{\Gamma }`$, and electron number, $`N_o`$. The integrand $`K_{5/3}(z)`$ is a modified Bessel function of order $`5/3`$. The parameter $`\nu _m=((3/2)\mathrm{\Gamma }\gamma _m^2\nu _B)`$, and $`Q`$ is proportional to the optical depth (to self-absorption) of the medium; for $`\nu \nu _m`$, the photon spectrum is self-absorbed at $`\nu <\nu _aQ^{2/(p+4)}`$. The high energy optically thin asymptotic behavior is the usual $`F_\nu \nu ^{(p1)/2}`$. The low energy asymptotic forms of the function depend on the relative values of $`\nu _m`$ and $`\nu _a`$: $`F_\nu \nu ^{5/2}`$ for $`\nu _m<\nu \nu _a`$, but is much flatter - $`F_\nu \nu ^{1/3}`$ \- for $`\nu _a<\nu <\nu _m`$. For very low frequencies, $`\nu \mathrm{min}[\nu _a,\nu _m]`$, $`F_\nu \nu ^2`$. We point out that the photon spectral index $`\alpha `$ is $`d`$log($`F_\nu )/d`$log$`(\nu )1`$. Figure 2 shows the many different types of low energy behavior one can obtain from synchrotron emission. The dot-dashed line is the usual simple optically thin spectrum, and the solid line shows a self-absorbed spectrum for $`\nu _a>\nu _m`$. The dotted and short dashed lines show optically thin spectra for a gradual (rather than sharp) cutoff to the electron distribution, discussed in §2.2. The long dashed line shows a spectrum in the case of small pitch angle scattering, discussed in §2.3. Since we are focusing on the low energy spectral index $`\alpha `$, we have kept $`\beta `$ constant and normalized the spectra at their peaks to a representative value within the BATSE window approximated by the vertical dotted lines. ### 2.1 Synchrotron Self Absorption Most treatment of synchrotron radiation from GRBs has been in the optically thin case (see, however, Papathanassiou, 1998). This is because synchrotron self-absorption requires “extreme” physical conditions in a GRB - particularly, high magnetic fields and a high column density of electrons. For example, for $`\nu _a<\nu _m`$, the optical depth to synchrotron self absorption is $`\tau (l/10^{13}cm)(n/10^8cm^3)(B/10^8G)^{2/3}(\gamma _m/50)^{8/3}(\mathrm{\Gamma }/10^3)^3(\nu _{obs}/10^{19}Hz)^{5/3}`$, where $`l`$ and $`n`$ are the path length and particle density in the co-moving frame, $`\nu _{obs}`$ is the absorption frequency in the observer’s frame (note that this frequency falls within BATSE’s spectral window), and we have assumed an electron energy distribution index $`p=2`$. For a more detailed discussion of absorption frequencies in various regimes, see Granot et al., 2000. An example of a self-absorbed spectrum is given by the solid line if Figure 2. As evident from the above expression, optical depths of order unity can be achieved within the BATSE spectral range with somewhat extreme values of the physical parameters, most notably the value of the magnetic field. However, such high values of the magnetic field can be reached simply through equipartition in internal shocks (see, e.g., Piran, 1999). In general, we understand far too little about the generation of magnetic fields and the hydrodynamical conditions relevant for GRBs to rule out these extreme physical parameters. We also point out that the required values of $`n`$, $`l`$, and $`\gamma _m`$ lead to a Compton Y parameter of order unity or greater, thereby reducing the efficiency of synchrotron radiation in the observed range (Kumar, 1999, Piran, 1999). In this paper, we do not focus on how the conditions required for observable synchrotron self-absorption are achieved, but rather investigate the consequences if they indeed are reached. In any case, as shown in Figure 1 above and as we will see below, a self-absorbed spectrum may be applicable to only a small fraction of the GRB population. To test how well the optically thin and thick synchrotron spectra fit the data, we fit 11 bursts with 256 channel energy resolution to the synchrotron spectral form in equation (1). In six cases, the optically thin spectrum fit the bursts well. But for five bursts, which have a low energy photon index $`\alpha >2/3`$ (as determined by fitting a Band spectrum, Preece et al., 1999), the optically thin spectrum does not provide a good fit. For these bursts, we find that when including a parameter characterizing the optical depth ($`Q`$ in equation 1), the fits are improved significantly over those to the optically thin spectra. Crider & Liang, 1999, also show self-absorption is consistent with GRB 970111, another burst with a hard low energy spectral index ($`\alpha >2/3`$). It should be noted, however, that in all of these bursts, $`\nu _m>\nu _a`$, and $`\nu _a`$ was close to the edge of the BATSE window (in fact, the presence of only a few bursts with $`\alpha =1`$ or $`3/2`$ (see Figure 1) indicates that in general $`h\nu _a50`$ keV). We discuss the implications this has on the observed $`\alpha `$ distribution below. Figures 3a and 3b show the spectral fits for 2 GRBs in our sample (burst triggers 3893 and 3253). An optically thin spectrum is the best fit to 3253, while a self-absorbed spectrum is required for 3893 - the rollover at low energies expected from synchrotron self-absorption is evident, and clearly cannot be accomodated by an optically thin spectrum. \[Note that these are counts spectra and not photon spectra, so the bumps in the data and the model are a result of the convolution of the impinging photon spectrum and the model convolved with the detector response matrix.\] ### 2.2 The Electron Energy Spectrum In most models of synchrotron emission, the electron distribution is modeled by a power law with with spectral index $`p`$ (as done above). Since the high energy index $`\beta =(p+1)/2`$ (or -(p+2)/2 for a simple cooling model), and the majority of bursts have $`\beta >1.5`$ (see Figure 1), this means that the index $`p>2`$. As a result, a cutoff to the power law distribution at some minimum energy $`\gamma _m`$ must be imposed to prevent divergence of the accelerated electron energy distribution at low energies. However, as discussed in the introduction, such a sharp cutoff is not a realistic - even stable - scenario. Instead, one expects to produce a gradual cutoff to the electron distribution (perhaps due to such effects as the two stream instability). In general, particle acceleration in relativistic shocks is not well understood, and could produce a range of low energy power law tails (among more complicated behavior) to the electron distribution, or even a nearly flat low energy electron spectrum. In order to account for this range of possible low energy behavior, we characterize the electron distribution by the following equation: $$N(\gamma )=N_o\frac{(\gamma /\gamma _m)^q}{1+(\gamma /\gamma _m)^{p+q}}$$ (4) where, now, $`\gamma _m`$ is some critical energy below which the electron distribution changes from $`p`$ to $`q`$ . For $`\gamma \gamma _m`$, $`N(\gamma )\gamma ^p`$, while for $`\gamma \gamma _m`$, $`N(\gamma )\gamma ^q`$. Hence, $`q`$ characterizes the “smoothness” of the cutoff (note for $`q\mathrm{}`$, $`\gamma _m`$ becomes a sharp cutoff to the distribution as defined in the previous section). For this new electron distribution, an optically thin synchrotron spectrum for an isotropic pitch angle distribution takes the form: $$F_\nu =𝒞(\nu /\nu _m)^{(q+1)/2}_0^{\mathrm{}}𝑑x\frac{x^{(q+1)/2}}{1+((\nu /\nu _m)^{(q+p)/2}x^{(p+q)/2})}_x^{\mathrm{}}K_{5/3}(z)𝑑z$$ (5) where, again, $`\nu _m=(3/2)\mathrm{\Gamma }\gamma _m^2\nu _B`$, and $`𝒞`$ is the normalization. Note that for $`q<1/3`$, the asymptotic photon index $`\alpha =\frac{q+1}{2}1<2/3`$, but for $`q>1/3`$, this index $`\alpha =2/3`$. This simply means that for $`q<1/3`$, the superposition of the photon spectra from individual electrons below $`\gamma _m`$ becomes powerful enough to change the low energy behavior. That is, there is a “competition” between $`\nu ^{1/3}`$ (the low energy photon spectrum for an individual electron) and $`\nu ^{(q+1)/2}`$ (the superposed low energy behavior of the electron distribution below $`\gamma _m`$) - when $`(q+1)/2<1/3`$, the low energy behavior below $`\nu _m`$ is different than the usual $`\nu ^{1/3}`$. Nonetheless, even for $`q>1/3`$, the “smoothness” of the cutoff can change the spectrum of the emitted photons significantly, primarily making the transition from the high energy index $`\beta `$ to the low energy asymptotic value of $`\alpha =2/3`$ more gradual. Examples of spectra with $`q\mathrm{}`$ are shown in Figure 2 by the dotted ($`q=0`$) and short dashed lines ($`q=2`$). We point out that as the cutoff of the electron distribution is made smoother (q smaller), the peak of the spectrum shifts to lower energies, while the asymptotic value of the spectral index is approached at a slower rate. Figure 4 shows both of these effects. The solid line shows the peak of the $`F_\nu `$ spectrum, denoted as $`h\nu _{max}`$, as a function of $`q`$; the dotted line shows the ratio of $`h\nu _{max}`$ to $`h\nu _{asymp}`$, where $`h\nu _{asymp}`$ is defined as the energy in which the optically thin spectral index is within 5% of its low energy asymptotic value of $`1/3`$. ### 2.3 Small Pitch Angle Emission The usual analysis of synchrotron radiation assumes electrons are distributed isotropically in a either a uniform or randomly oriented magnetic field geometry (Pacholczyk, 1970). However, this is a simplifying assumption and sometimes one expects a non-isotropic pitch angle distribution, usually beamed along the field lines. If the beaming is strong so that most electrons have pitch angles $`\theta <1/\gamma `$ (the Lorentz factor of the electron), the shape of the synchrotron spectrum changes significantly (see Epstein, 1973, and Epstein & Petrosian, 1973). An efficient method of producing isotropic pitch angle distributions at relativistic energies is by plasma turbulence, which can scatter and change the energy of the electrons. For high density, low magnetic field plasmas the Alfvén phase velocity is less than the speed of light and therefore the speed of the particles (relativistic in our case). In this case, the pitch angle diffusion rate of the electrons interacting with plasma turbulence is much larger than the acceleration rate; consequently, the accelerated electrons will have an isotropic pitch angle distribution. However, for the low density, high magnetic field condition expected for the sources of GRBs the opposite is true. In this case the fluctuation in the electric field of the waves exceeds the fluctuation of the magnetic field so that the above situation is reversed (see e.g. Dung & Petrosian, 1994 or Pryadko & Petrosian 1997). Then the pitch angle distribution of the accelerated electrons could become highly anisotropic as required in the small pitch angle model. One consequence of a small pitch angle distribution is that the resulting spectrum at low frequencies is approximately: $`F_\nu \nu `$, for $`\nu <\nu _s`$, where $`\nu _s`$ characterizes the where small pitch angle scattering is important. Above this frequency, $`F_\nu `$ follows the usual synchrotron spectrum (again, we have assumed that $`\theta \gamma 1`$, where $`\theta `$ is the electron pitch angle. This leads to a single particle emissivity that goes as $`\nu `$ up to $`2\gamma \nu _B`$, and sharply cuts off after that; see Epstein, 1973, §II). An example of a synchrotron spectrum with a small pitch angle distribution in this regime is shown by the long dashed line in Figure 2. Note that this is an optically thin spectrum with $`\alpha =0`$; this is above the so-called “line of death” $`\alpha =2/3`$. It should be noted that a very similiar low energy spectrum is obtained when an isotropic distribution of electrons is embedded in a region with a very tangled magnetic field. If there exist magnetic field fluctuations with a correlation length that is less than the Larmor radius of the electrons, then there will be transverse deflections of the electrons with angles less than the relativistic beaming angle. Emission due to these fluctuations is very similiar to the small pitch angle case and therefore also produces a spectrum $`F_\nu \nu `$ at low energies (Medvedev, 2000). Note that the latter scenario imposes restrictions on the structure of the magnetic field, while the former (small pitch angle) model constrains the acceleration mechanism. ## 3 The Low Energy Asymptotic Behavior We have shown so far that synchrotron emission can result in low energy spectral index values of $`\alpha =`$ $`d`$log($`F_\nu )/d`$log$`(\nu )1=2/3,0,,1.5`$ (and possibly $`1.5`$ if “cooling” effects are important). The observed distribution shown in Figure 1, on the other hand, shows a broad and continuous distribution of $`\alpha `$ with about $`96`$% of bursts in the range $`1.5<\alpha <0`$. We propose that this dispersion is caused by the finite bandwidth of the BATSE instruments and arises from the fitting of a phenomenological model to the data. The theoretical values stated above represent the asymptotic logarithmic slopes far from the break or peak photon energy. How well a fitting algorithm can determine these asymptotic values depends on whether the asymptotic value of the spectrum is reached within the detector window. This depends both on how “quickly” the spectrum reaches its asymptote and the value of $`E_p`$ relative to the lower end of the detector window. In other words, we need to understand how well spectral fits can determine the low energy asymptote for the different cases presented above. For example, as $`E_p`$ moves to lower and lower energies, we get less and less of the low energy portion of the spectrum; in this case, our spectral fits probably will not be able to determine the asymptotic index and will measure a smaller (softer) value of $`\alpha `$. Preece et al. (1998a) pointed out this effect and attempt to minimize it by defining an effective low energy index, $`\alpha _{eff}`$, which is the slope of the spectrum at 25keV (the edge of the BATSE window). However, a correlation between $`\alpha _{eff}`$ and $`E_p`$ will still exist if the asymptotic value is not reached by 25keV. This difficulty becomes more severe the smoother the cutoff to the electron distribution, because the spectrum takes longer to reach its asymptotic index. To determine the extent of these effects, we simulated data from optically thin synchrotron models with different values of the parameters $`\nu _m`$ and $`q`$ (which determine the values of $`E_p`$ and $`\alpha `$ in a Band spectral fit). We have included spectra both for an isotropic pitch angle distribution and for a distribution of electrons having small pitch angles only. Note that since we take an optically thin spectra with $`q0`$, all of the spectra have a low energy asymptote of $`2/3`$ (isotropic pitch angle distribution) or $`0`$ (small pitch angle distribution). We normalize the spectra so that the peak photon flux in the range $`50300`$ keV is $`10\mathrm{p}\mathrm{h}/\mathrm{cm}^2/\mathrm{s}`$ (e.g. a fairly bright GRB). Our data points are then drawn from a Poisson distribution with a mean at the value given by the synchrotron photon spectrum at a particular frequency; that is, the data points are drawn from a Poisson distribution with a mean $`N_{\nu ,i}=(F_\nu /h\nu )_i`$, where $`i`$ indexes the data point. We then fit a Band spectrum to this data using a conservative estimate that BATSE is sensitive to all photons above $`10`$ keV. Figure 5 shows the value of $`\alpha `$ as a function of $`E_p`$, for different degrees of the smoothness of the electron energy distribution cutoff in the isotropic pitch angle case (three lower curves), and for an intermediate cutoff in the small pitch angle distribution case (upper curve). Not surprisingly, there is a strong correlation between the value of $`E_p`$ and the value of the index, $`\alpha `$. To make sure this isn’t purely an artifact of the Band function, we also tried a broken power law fit ($`A(E)E^\alpha `$, $`E<E_{break}`$, $`A(E)E^\beta `$, $`E>E_{break}`$, where $`A`$ is in photons/cm<sup>2</sup>/s/keV) to our simulated data and found a very similiar relationship (although the power law does not give as good of fits as the Band spectrum, so the relationship was noisier). Because GRBs have a relatively broad distribution of $`E_p`$, the above correlation between $`\alpha `$ and $`E_p`$ will lead to a dispersion in the distribution of $`\alpha `$; to determine the extent of this effect, we need the distribution of $`E_p`$’s. In our analysis, we use the distribution observed by BATSE, plotted in Figure 6 (solid line). Note that there has been some controversy over the dispersion in this distribution and we have shown (Lloyd & Petrosian, 1999) that indeed this distribution suffers from selection effects in the BATSE spectral window that tend narrow the it. These effects are evident in Figure 6, which also plots $`E_p`$ distributions observed by SMM (sensitive to higher energies than BATSE; Harris & Share, 1998) and GINGA (sensitive to lower energies than BATSE; Strohmeyer et al., 1998). The former extends the BATSE $`E_p`$ distribution on the upper end, while the latter extends it on the lower end. However, we are interested in the BATSE observed $`\alpha `$ distribution through its relationship with the BATSE observed $`E_p`$ distribution. Hence, the intrinsic $`E_p`$ distribution (with selection effects accounted for) - although relevant for understanding other physical aspects of the radiation processes - is not relevant for our discussion of the observed $`\alpha `$ distribution. ## 4 The Observed $`\alpha `$ Distribution From the observed distribution of $`E_p`$ (Figure 6, solid histogram), we can determine how the correlation between $`E_p`$ and $`\alpha `$ introduced by the fitting procedure (Figure 5) “smears” the distribution of $`\alpha `$ away from the expected narrow distribution around the physical asymptotes of $`2/3`$, $`0`$, etc.. We represent the correlation between $`\alpha `$ and $`E_p`$ by an approximate simple analytical function; log$`(E_p)=h(\alpha )`$, where the function $`h(\alpha )`$ depends on the specifics of the synchrotron model (see Figure 5). We then approximate the $`E_p`$ distribution, $`f(\mathrm{log}(E_p))`$ by a Gaussian in log$`(E_p)`$, with a mean and dispersion equal to those of the observed distribution. The distribution of $`\alpha `$ is then obtained from the relation $$g(\alpha )=f(\mathrm{log}(E_p))\frac{dh(\alpha )}{d\alpha }.$$ (6) Figure 7 shows the resultant $`\alpha `$ distributions in the case of an isotropic pitch angle distribution for a sharp ($`q=\mathrm{}`$, right solid curve), intermediate ($`q=2`$, middle short-dashed curve), and flat ($`q=0`$, left long-dashed curve) cutoff to the electron energy spectrum, as well as the small pitch angle distribution case for an intermediate cutoff (right dot-dashed line). The dotted histogram is the observed $`\alpha `$ distribution shown in Figure 1. The simulated curves have been arbitrarily normalized to the height of the observed distribution at their central values of $`\alpha `$. We note the following: 1) It appears that almost the entire range of the observed distribution can be covered by these optically thin synchrotron models. In particular, given a distribution in $`q`$, an instantaneous optically thin spectrum in the large pitch angle scattering regime can easily accomodate bursts below the “line of death” with $`\alpha 2/3`$, where most of the bursts are located. 2) The second conclusion is that the electron energy distribution below the turnover energy $`\gamma _m`$ must be falling off, or must be (at most) flat ($`q0`$). Otherwise, the optically thin model would predict too many bursts with $`\alpha `$ less than about $`1.5`$. 3) We have shown than an optically thin spectrum from a small pitch angle distribution and a smooth cutoff to the electron energy distribution can explain all bursts with $`2/3\alpha 0`$. 4) However, this is not necessarily the only explanation. Bursts with $`\alpha >2/3`$ all the way up to $`3/2`$ can also be due to absorption effects. As shown in §2, for some bursts synchrotron self-absorption is necessary to provide a good spectral fit to the data. However, we usually do not see the values $`\alpha =3/2`$ ($`\nu _m<\nu \nu _a`$) or $`\alpha =1.0`$ ($`\nu \mathrm{min}[\nu _m,\nu _a]`$) expected from self-absorption (there are, however, some bursts with such sharp breaks; see Preece et al., 1999). If self-absorption is important, $`h\nu _a`$ must be near the lower edge of the BATSE spectral window. In this case, we just begin to see the absorption cutoff and the steep asymptotic low energy slopes are not yet reached. Indeed, this is the case for the fits mentioned in §2. 5) The analysis above becomes more complicated when there are two or more spectral breaks in the detector bandpass (there is some evidence for two breaks in GRB spectra; see Strohmeyer et al., 1998). In this case, a fit to a phenomenological model with a single break will lead to averaging of the slopes above and below the break it did not fit. For example, if $`\nu _m`$ and $`\nu _a`$ are both present in the spectral window with $`\nu _m>\nu _a`$ and the fit places $`E_p\nu _m`$, then the Band function will not accomodate the additional absorption break. As a result, the low energy index ends up being a weighted average of the optically thin (-2/3) and optically thick (1) indices. If the fit places $`E_p\nu _a`$, then the high energy index $`\beta `$ will be an average of $`2/3`$ and $`(p+1)/2`$. This of course applies to any other two (or more) characteristic frequencies as well. 6) It should also be noted that we have included only a small amount of noise (from counting statistics) in our simulations - adding more noise will contribute to the spread in the simulated $`\alpha `$ distributions and strenghten our arguments. 7) Finally, we note that other effects such as inhomogeneities in the electron distribution (Granot et al., 2000) an other radiative transport effects (Grusinov & Meszaros, 2000, Dermer & Boettcher, 2000) can also produce spectra both below and above the values of $`\alpha =2/3`$. ## 5 Spectral Evolution Not only do emission models have to accomodate the shape of the observed distributions, but also the temporal behavior of the spectral parameters. The behavior of the spectral characteristics with time throughout a GRB can give us information about the environment of the local emission region and conceivably constrain the emission mechanism. If the above interpretation of the $`\alpha `$ distribution is correct, then we expect considerable correlation between the variation of the observed values of $`\alpha `$ and $`E_p`$ throughout a burst. Many studies have looked at time evolution of spectral parameters (e.g. Norris et al., 1986, Ford et al., 1995, Crider et al., 1997, Preece et al., 1998(b)). Ford et al. analyze the evolution of $`E_p`$ for 37 bright, long duration GRBs observed by BATSE. This study presents evidence of a general envelope of “hard-to-soft” evolution of $`E_p`$ throughout the duration of the burst. The behavior is explained in the context of an external shock model in which there is a gradual decline of average energy as more particles encounter the shock; the emission mechanism is unspecified. Crider et al. investigate the behavior of the low energy spectral index $`\alpha `$ for a sample of 30 BATSE GRBs. They find that 18 of these bursts show hard-to-soft evolution, while 12 exhibit so-called “tracking” behavior - when the evolution of $`\alpha `$ correlates with the burst time profile. All of these bursts also show a strong correlation between $`\alpha `$ and the peak energy, $`E_p`$, as a function of time. The authors attribute the hard to soft behavior is to Thomson thinning of a Comptonizing plasma. The “tracking” behavior is not explained. Recently, Preece et al. (1999) published a catalog of spectral data with high time resolution fitting. We present the time evolution of three of these bursts (Figures 8, 9, and 10), plotting all of the variables that parameterize the (Band) spectral fits. Starting in the upper left hand corner and going clockwise we plot the evolution of the normalization $`A(t)`$ (in units of photons/cm<sup>2</sup>/s/keV), $`\alpha `$, $`E_p`$, and the high energy photon index $`\beta `$. \[We point out that the time resolution of the spectral fitting is sometimes coarser than the real time variation of the burst (at least on the shortest detector timescale (64ms)). This will lead to some averaging effects which may weaken the correlation we expect. We have included the higher resolution time profiles (in arbitrary units) superposed on the plot of the normalization $`A(t)`$, as a reference.\] Given the expected correlation between $`\alpha `$ and $`E_p`$ from instrumental and fitting procedural effects discussed in the previous section, we expect evolution of $`\alpha `$ to some extent mimic the evolution of $`E_p`$ in time. Indeed this is what we see in Figure 8 - the parameter $`\alpha `$ “tracks” $`E_p`$ (and both of these track the flux, $`A`$). Note the values of $`E_p`$ and $`\alpha `$ are consistent with the case of a flat ($`q=0`$) cutoff to the electron distribution, but there may be indication of variation of $`q`$ as well. However, it is not always this simple. In an internal shock scenario, we can regard each pulse in the time profile as a separate emission episode; in this case, the internal parameters can vary depending on the physical conditions at each shock. In particular, a change in $`q`$ can create a change in $`\alpha `$ from pulse to pulse, independent of $`E_p`$ (although we might expect some anti-correlation with $`E_p`$ in this case - high $`q`$ gives a high $`\alpha `$, but decreases $`E_p`$ if all other paramters are the same). Figure 9 shows such behavior. The parameter $`\alpha `$ tracks the flux; $`E_p`$ varies on the same timescale as the flux, but these variations are superimposed in an envelope of hard-to-soft evolution. This can be explained by a sharpening of the cutoff of the electron distribution from peak to peak. That is, in the first peak, we have fairly high $`E_p`$ values ($`500`$ keV) and moderately low ($`1`$) values of $`\alpha `$. This is what we expect for a $`q0`$ or $`1`$ cutoff. In the second peak, $`E_p`$ is about $`300`$ keV, but the $`\alpha `$ value is around $`0.6`$ \- this is what we’d expect for a sharper cutoff ($`q5`$). In addition, we may see evidence of evolution of the opacity or pitch angle distribution from pulse to pulse. For example, in Figure 10, the parameter $`\alpha `$ appears to evolve from hard (above the “line of death”) to soft, while $`E_p`$ appears to tracks the normalization. In the early phase, $`\alpha >2/3`$ and $`E_p`$ is moderate ($`300`$ keV). At later times, $`E_p`$ is at fairly high values ($`500`$ keV), while $`\alpha `$ is fluctuating around the expected asymptote for optically thin emission, $`\alpha =2/3`$ (this is what we expect from the correlation discussed in this paper - when $`E_p`$ is high enough, the asymptote of the spectrum is reached inside the BATSE window). The overall behavior can be attributed either to a transition from a small pitch angle to isotropic regime, or to a transition from an optically thick (to synchrotron self-absorption) to optically thin regime. Note that we have not accounted for any biases determination of the high energy photon index $`\beta `$ might cause. This ultimately can affect the value of $`\alpha `$ or $`E_p`$. For most cases, $`\beta `$ remains fairly constant, with only minor fluctuations. Our purpose here is not to present a rigorous temporal analysis of the spectral parameters, but to show the diversity of spectral evolution from burst to burst, and to discuss both the procedural (i.e. spectral fitting) and physical reasons that produce this behavior. The variety of observed spectral evolution and its meaning in terms of synchrotron emission from internal shocks will be discussed in much more detail in an upcoming publication. ## 6 Summary and Conclusions Although synchrotron radiation has been suggested to be a viable emission process for GRBs (Katz, 1994, Tavani, 1996), there has since been much controversy over what radiation mechanism actually gives rise to the observed GRB spectra. Several studies (e.g., Ghissilini et al., 1999, Celotti & Ghissilini, 1999) have attempted to rule out synchrotron radiation as the source of GRB spectra based on certain discrepancies between the model’s predictions for the low energy spectral behavior and the observed data. In particular, the simple model of synchrotron radiation in an optically thin medium from a power law distribution of electrons with an isotropic pitch angle distribution (and neglecting radiation losses) predicts a value of the low energy photon index $`\alpha `$ of $`2/3`$. The observed distribution of $`\alpha `$ is clearly inconsistent with this predicted behavior. The predicted value of $`\alpha `$ expected in a a simple cooling spectrum is $`3/2`$, which is also inconsistent with the data. Furthermore it has been suggested (Crider et al., 1997) that the variation of $`\alpha `$ throughout the burst’s duration is inconsistent with what is expected from these synchrotron models. In this paper, we show that synchrotron radiation can in fact explain the observed spectral behavior of GRBs. We show that the source of the discrepancies described above are a result of two factors: a) the simplistic assumptions made in standard synchrotron models, b) assuming that the spectral fitting procedure accurately determines the asymptotic values of the spectral indices, without accounting for the effects of the detector bandwidth in the fitting procedure. \[The latter effect must be accounted for not only when comparing to synchrotron models, but all other emission models as well.\] We focus on the so-called “instantaneous” synchrotron spectra (in which radiation losses are not evident). As mentioned above, the simple cooling spectrum is not consistent with observations. A cooling spectrum is obtained when electrons are injected in a magnetized region, where they then lose all of their energy to radiation; it is not clear that this is a viable scenario for GRBs. Athough the “cooling time” (the timescale over which electrons lose their energy to radiation) is much shorter than other timescales involved in the fireball model (e.g. the hydrodynamical timescale), it is probably unrealistic to separate the acceleration and radiation loss processes at the site of the GRB. To sustain the burst, there should be continual acceleration of the emitting particles with a rate equal or greater than the synchrotron loss rate. In this case, particles can still lose all of their energy to radiation quickly, but we will have the instantaneous spectrum as a result of the competing acceleration and loss processes. \[Note that because electrons are losing all of their energy to radiation, the efficiency (i.e. the ratio of energy radiated to the total energy of the GRB) in this model is the same as calculated using the simple cooling model - about $`1`$% (see, e.g., Kumar, 1999).\] We believe that this is an attainable scenario, although proof of this is beyond the aim of this paper. To investigate whether the instantaneous synchrotron models are consistent with the existent data, we relax four important assumptions used in the usual analysis with synchrotron models: 1) We allow for the possibility of synchrotron self-absorption, which produces a hard low energy photon index of either $`1`$ or $`3/2`$, depending on the relative values of the synchrotron self-absorption frequency and the minimum electron frequency. 2) We allow for a more realistic smooth (rather than sharp) cutoff to the electron distribution, which tends to soften the low energy spectral behavior. 3) We include effects of small pitch angle scattering, which leads to a value of $`\alpha `$ of $`0`$. 4) We account for the fact that as the break energy of the GRB spectrum approaches the lower edge of the BATSE spectral window, the low energy spectral index will become softer (the expected asymptote of $`2/3`$ is not reached). Items 1-3 above produce physically different low energy spectral behavior. Item 4 is a procedural and instrumental effect which leads to additional dispersion of the $`\alpha `$ distribution. We show that including all of these effects can explain the observed distribution of $`\alpha `$’s. We also infer from our analysis that the electron energy distribution must flatten or decline at low energies; otherwise, we would see many more bursts with $`\alpha <3/2`$. GRBs whose $`\alpha `$ values lie above the “line of death” ($`\alpha =2/3`$) may be explained either by emission from electrons with small pitch angles or absorption processes (e.g. synchrotron self absorption). Finally, we present some examples of the temporal behavior of GRB spectral parameters. We show that some of the temporal evolution is explained by the expected correlation between $`\alpha `$ and $`E_p`$, that results from the detector bandwidth and the fitting process. Other temporal variations in the spectral parameters suggest changes in the physical conditions at the GRB source from pulse to pulse (where each pulse may represent independent emission episodes, for example, in an internal shocks model). A detailed study of the variety of spectral evolution and its consistency with synchrotron emission from internal shocks is the subject of an upcoming publication. In this paper, we have dealt only with the low energy spectral index $`\alpha `$. Similiar discrepancies and limitations exist for the high energy spectral index $`\beta `$. However, the value of this index is not as well determined or reliable because the signal-to-noise decreases rapidly with energy, to an extent that it is questionable whether we even have a simple power law spectrum above $`E_p`$. However, there is some reliable data for bright bursts, and we shall explore the effects of modifications described here for $`\beta `$ in a future publication. Acknowledgements: This work was supported by the CGRO Guest Investigator Program and by the Stanford McMicking Fellowship. We would like to thank Rob Preece for many useful discussions, as well as providing and helping with the some of the BATSE data used in this analysis. We would also like to thank Pawan Kumar for helpful comments. Finally, we are greatly indebted to the referee for an extremely careful reading of the manuscript, and helpful comments which much improved this paper.
warning/0007/cond-mat0007382.html	ar5iv	text	# On time’s arrow in Ehrenfest models with reversible deterministic dynamics ## I Introduction Physical laws of motion are time reversible. As a consequence, from a movie of a few interacting particles one cannot distinguish whether the movie is running forward or backward. This seems to be different when we observe a very large number of particles. Macroscopic properties like the density of particles show a direction of time. If the milk in the coffee shrinks to a single drop we know that the movie is running backward. More than a hundred years ago Ludwig Boltzmann gave convincing arguments for the irreversible macroscopic behaviour of particles moving with reversible microscopic laws. The fraction of initial states which leads to reversible macroscopic behaviour is so extremely small that we will never be able to observe it. With probability one, we observe a time’s arrow in macroscopic properties. Hence, through the selection of the initial state, the deterministic motion receives probabilistic elements. In fact describing a large system of interacting particles by either a stochastic motion (Boltzmann transport equations) or ensemble theory (statistical mechanics) turned out to be very successful. However, the success of Boltzmann’s ideas did not suppress the discussion about the foundations of irreversible properties, even hundred years later . For example, there was a round table discussion on this subject on the STATPHYS 98 conference , where Prigogine introduced novel microscopic laws which are irreversible with time. One reason for this ongoing discussion is the absence of rigorous mathematical proofs of irreversible properties in the thermodynamical limit. Furthermore, solvable models where irreversible macroscopic properties can be well defined and investigated almost do not exist, to our knowledge. Either one studies stochastic models like the famous one of Paul and Tanja Ehrenfest , or disorder averages as for the Kac ring , or ensembles of chaotic noninteracting particles as in the Lorentz gas . However, a single trajectory of interacting particles is difficult to investigate in this context. Numerical simulations of the equations of motions do not help much in this context, either. The algorithms can be formulated in a time reversible way and the decay of Boltzmann’s H–function has been demonstrated . However, the state space is too large to study any details, for instance to estimate the fraction of untypical inital states or to find the time scale of the Poincaré return time. One should keep in mind that we want to understand the property of a single trajectory for a short time. Hence, ensemble averages do not give a basic explanation of irreversible properties, since they contain an average over infinitely many trajectories. Ergodic theory does not help either, since it needs time averages over infinitely large times (at least up to the Poincaré return time) whereas real physical systems or computer simulations run over time scales which are extemely short compared to the return time. In this paper we introduce a model with deterministic time reversible dynamics which can be analysed in detail. It is a deterministic extension of the Ehrenfest model discussed in Section IV. The dynamics is ergodic, each state is visited once during a cycle. The Poincaré return time is known exactly and the Boltzmann and Gibbs entropies can be defined precisely. From numerical calculations we obtain the distribution of first passage times the system needs to go from a non equilibrium state to a state with largest Boltzmann entropy as well as in the reverse direction. The scaling of average first passage times with system size is calculated. Finally we compare the results of the deterministic model with the corresponding ones of the stochastic version of the Ehrenfest model. ## II Deterministic Ehrenfest model: shift register generator The microscopic state $`x`$ of our model has $`N`$ binary variables $`x_i\{0,1\};i=1,\mathrm{},N`$, $$x(t)=(x_1,x_2,\mathrm{},x_N).$$ (1) The time $`t`$ is discrete. At each time step a new bit $`x_0`$ is generated and replaces $`x_1`$, whereas the old bits $`x_1,\mathrm{},x_{N1}`$ are moved one step to the right: $`x_1`$ $`x_2`$ (2) $`\mathrm{}`$ (3) $`x_{N1}`$ $`x_N`$ (4) $`x_0`$ $`x_1.`$ (5) The last bit $`x_N`$ is deleted. The new bit $`x_0`$ is constructed using the theory of primitive polynomials modulo two which is also used to generate pseudo random numbers . If $`x_0`$ is the sum modulo two of a few bits $`x_i`$ at certain positions $`i`$, then any initial state $`x(0)`$ runs through all possible states $`x`$ except the state zero $`(0,0,\mathrm{},0)`$. For instance for $`N=97`$ one finds from the table of Ref. : $$x_0=(x_{97}+x_6)mod\mathrm{\hspace{0.33em}2}.$$ (6) In our investigation we use such sequences with maximal cycle length, taking the generators from for $`N<100`$ and from for higher $`N`$. For the macroscopic property $`M(x)`$, which we will investigate, we take the number of $`x_i=1`$, $$M(x)=\underset{i=0}{\overset{N}{}}x_i.$$ (7) Note that at each time step $`M`$ changes by $`\pm 1`$ at most. The original Ehrenfest model considered $`N`$ balls in two urns. $`x_i=0(1)`$ means that ball $`i`$ is in the left (right) urn. $`M`$ is the number of balls in the right urn. In the Ehrenfest model the balls were chosen randomly. Here we define a deterministic rule to move the particles between the two urns. The equations of motion (2) and (6) are time reversible; for instance from (6) one finds $$x_{97}=(x_0+x_6)mod\mathrm{\hspace{0.33em}2}.$$ (8) Hence, from the first and seventh bit of state $`x(t+1)`$ one calculates the 97<sup>th</sup> bit of $`x(t)`$, and a shift to the left gives the rest of $`x(t)`$. We have now defined a model where time and states are discrete. It has the following properties: 1. The equations of motion are deterministic and time reversible. 2. Each initial state returns to itself after $`2^N1`$ time steps; i.e. the Poincaré return time is $`T_R=2^N1`$. 3. The Boltzmann entropy $`S_B`$ is given by the number of microstates $`x`$ which have the macroscopic property $`M`$. Since the distribution of macrostates is the binomial distribution (with the exception of the zero state), $$p(M(x))=\left(\genfrac{}{}{0pt}{}{N}{M(x)}\right)/(2^N1),$$ (9) one finds for the entropy in dimensionless units: $$S_B(x)=\mathrm{ln}\left(\genfrac{}{}{0pt}{}{N}{M(x)}\right).$$ (10) 4. Since each state $`x`$ is visited once during each cycle, it has an identical statistical weight in a Gibbs ensemble. Hence, the Gibbs entropy $`S_G`$ is given by $$S_G=\mathrm{ln}(2^N1)N\mathrm{ln}2.$$ (11) 5. The system is ergodic. When the system is observed over the return time $`T_R`$, the time average $`\overline{M}`$ agrees with ensemble average $`M`$: $$\overline{M}=\frac{1}{T_R}\underset{t>0}{\overset{T_R}{}}M[x(t)]=M=\frac{N}{2}.$$ (12) 6. The most probable value of $`M`$ agrees with time and ensemble average, $$M_{mp}=\overline{M}=M=\frac{N}{2}.$$ (13) There are $$\left(\genfrac{}{}{0pt}{}{N}{N/2}\right)\sqrt{\frac{2}{\pi }}\frac{2^N}{\sqrt{N}}$$ (14) states $`x`$ whith $`M(x)=M_{mp}`$. We call these states equilibrium states. During a cycle, the average time interval between two consecutive equilibrium states increases with $`\sqrt{N}`$ according to Eq. (14). ## III Macroscopic dynamics We investigate the two following questions: 1. When the system starts from an initial state $`x`$ which is far away from equilibrium, for instance with $`M(x)=N/4`$, how long does it take to reach an equilibrium (= most probable) state? 2. When the system starts from an equilibrium state, $`M(x)=N/2`$, how long does it take to reach a non-equilibrium state with $`M=N/4`$? We call these two time intervals first passage time $`T_{eq}`$ and $`T_{neq}`$, respectively. Both of these times depend on the special choice of the initial state, therefore we obtain a distribution of first passage times. Fig. 1 shows the result for $`N=24`$ which is obtained from simulations of all $`2^{24}1`$ states. The time $`T_{eq}`$ to reach an equlibrium state has a sharp peak around $`T_{eq}N`$. In contrast, the time to leave equilibrium has a broad distribution extending from the minimal possible time $`T_{neq}=N/4`$ to $`T_{neq}=5079`$. Since the system is time reversible, the distribution of equilibrium times should be part of the distribution of non-equilibrium times $`T_{neq}`$. In fact, we observe the corresponding peak in Fig. 1. But surprisingly there are more equilibrium states which reach non equilibrium ones for the first time in short than in long times. The distribution of $`T_{neq}`$ has its maximum for values of the order of $`N`$, while the tail can be fitted to an exponential distribution $$P(t)e^{t/\tau },$$ (15) where $`\tau `$ is of the order of the average non-equilibrium time $`T_{neq}`$, averaged over all initial equilibrium states. In Figs. 2 and 3 it is shown how the average value of these times scales with system size $`N`$. For small systems $`N<32`$ we average over all possible initial states whereas for large systems we average over several hundred randomly selected initial states. We obtain a good fit with $`T_{eq}`$ $`=`$ $`1.30N12.3`$ (16) $`T_{neq}`$ $`=`$ $`15.0\mathrm{exp}(0.146N).`$ (17) Hence, the time to reach equilibrium is short, it scales with system size $`N`$, whereas the time to reach a state far from equilibrium is very long and increases exponentially with $`N`$. A lower bound on $`T_{neq}`$ can be found for any system whose distribution of macrostates is the binomial distribution: take a sequence of $`\left(\genfrac{}{}{0pt}{}{N}{N/2}\right)/\left(\genfrac{}{}{0pt}{}{N}{N/4}\right)`$ equilibrium states followed by one nonequilibrium state with $`M=N/4`$, repeated as often as necessary. Incidentally, this lower bound differs only by a factor of $`\sqrt{2/(\pi N)}`$ from the average return time of the non-equilibrium state. Choosing an initial occupation $`M_0=mN`$, the lower bound can be approximated for large $`N`$ using Stirling’s Formula: $$T_{neq}\sqrt{2m(1m)}(2m^m(1m)^{1m})^N.$$ (18) For $`M_0=N/4`$, this gives $$T_{neq}\sqrt{3}(3^{3/4}/2)^N\sqrt{3}\mathrm{exp}(0.131N);$$ (19) the exponent is close to the fit given in Eq. (16). Together with the upper bound $`T_{neq}2^N1`$ one obtains exponential scaling with $`N`$. We find that usually $`T_{neq}`$ is proportional to the return time of the non-equilibrium state, with a proportionality constant that depends on the details of the generator (for example, if a two-tap or four-tap register is chosen, see Fig. 3). Fig. 4 shows the time dependence of the Boltzmann entropy $`S_B(x(t))`$, averaged over initial states with $`M=N/4`$ (non-equilibrium) and $`M=N/2`$ (equilibrium), respectively. The non-equilibrium entropy decays to its equilibrium value $`S_BS_G`$ in $`N`$ time steps, since every bit $`x_i`$ is visited and “randomized” exactly every $`N`$th time step. Starting from equilibrium states, the entropy stays constant. Note that after the Poincaré return time $`T_R`$, $`S_B`$ has to return to its initial value. But the time to reach equilibrium is of the order of the microscopic time, whereas the return time $`T_R`$ increases exponentially with system size. ## IV Stochastic Ehrenfest model The original Ehrenfest model describes a stochastic process: $`N`$ balls are distributed among two urns. At every time step, a ball is picked at random and moved to the respective other urn; $`M`$ is the number of balls in one of the urns. This process can be mapped to a random walk with drift . $`M`$ is the position of the walker, who at each time step performs a step $`\mathrm{\Delta }M=\pm 1`$ with probability $`P(MM+1)`$ $`=`$ $`(NM)/M`$ (20) $`P(MM1)`$ $`=`$ $`M/N.`$ (21) Hence, there is a drift towards the center $`M=N/2`$. To our knowledge, analytic calculations of first passage times exist only between equilibrium states , whereas we are interested in the decay to or from non equilibrium states. We have calculated these times from numerical iteration of the corresponding transition matrices using absorbing states . The average times for different system sizes are shown in Figs. 2 and 3 in comparison with the corresponding deterministic model; the distribution of times for $`N=24`$ is shown in Fig. 5. Again, the decay time $`T_{eq}`$ to equilibrium scales with $`N`$ (although with more pronounced finite-size effects) while the time to reach a non equilibrium state from an equilibrium one increases exponentially with system size. Indeed the stochastic model has similar properties as the deterministic one. ## V Hierarchical dynamics There are several deterministic algorithms which change only one bit per time step and run through all of the $`2^N`$ configurations of the bit sequence $`x`$. One of them is the Gray Code : Find the largest power $`2^j`$ by which the discrete time $`t`$ is divisible and change bit $`x_{j+1}`$. The quantity $`M`$ changes by $`\pm 1`$ only, and the system runs through all of the $`2^N`$ possible states. However, now the access to the different variables $`x_j`$ is hierarchical; it takes $`2^{j+1}`$ steps before site $`j`$ is visited again. This hierarchical structure of microscopic dynamics leads to slow relaxation times of the macroscopic property $`M`$, as shown in Fig. 6. The time to reach an equilibrium state from an initial non-equilibrium one scales exponentially with system size $`N`$. There is a corresponding stochastic version of the Gray code: * Start with $`j=1;`$ * with probability $`1/2`$, increase $`j`$ by 1; else flip $`x_j`$, increase $`t`$ by one and return to the first step; * if $`j>N`$, increase $`t`$ and return to the first step; else repeat the second step. The ensemble average of $`M`$ decays as $$M=N/2+\underset{j=1}{\overset{N}{}}(M_0/N1/2)\mathrm{exp}(t/2^{j1}).$$ (22) The simulations of the stochastic model agree with Eq. (22); however, in the deterministic Gray code, $`M`$ shows structures caused by ”almost returns” after times that are multiples of high powers of $`2`$ (see Fig. 6). The first passage times of the stochastic model show a similar exponential scaling as the corresponding times of the deterministic model (see Fig. 7). It is possible to have a quadratic dependence on $`N`$ of the decay time of $`M(t)`$ by the following stochastic algorithm: Draw a uniform random number $`r[0,1]`$ and flip $`x_j`$ with the largest index $`j<1/r`$, if $`jN`$. The probability to flip $`x_j`$ is of order $`1/j^2`$. The average value of $`M(t)`$ decays as $`M`$ $``$ $`N/2+(M_0/N1/2)[(N+1)\mathrm{exp}({\displaystyle \frac{2t}{(N+1)^2}})`$ (24) $`\mathrm{exp}(2t)+\sqrt{2\pi t}(\text{erf}({\displaystyle \frac{\sqrt{2t}}{N+1}})\text{erf}(\sqrt{2t}))].`$ The long term behaviour is determined by the longest decay constant, which is $`2/(N+1)^2`$. We have also calculated the first passage time to equilibrium, and find that it scales with $`N^2`$. ## VI Summary A simple model with time reversible deterministic dynamics is investigated. The Poincaré return time is $`2^N1`$, where $`N`$ is the size of the system; each initial state returns to itself after $`2^N1`$ time steps. We define a macroscopic quantity $`M(t)`$ and study its time dependence. When the system starts from an equilibrium state (state with most probable value of $`M`$) it takes on average an exponentially large time before it reaches a state far from equilibrium. On the other side, when the system starts far from equilibrium it reaches an equilibrium state after a period of the order of the microscopic time. The same is true for the corresponding Boltzmann entropy. This behaviour explains the time’s arrow in macroscopic properties for a system with time reversible dynamics. We have calculated the full distributions of these first passage times and compared them to the ones of the corresponding stochastic Ehrenfest model. Surprisingly the distribution of times the system needs to reach non-equilibrium for the first time reaches a maximum at times of order $`N`$ and monotonically decreases for larger times. However, the integral of $`T_{neq}`$ over short times is small compared to the total number of equilibrium states. Therefore, on average the time to leave equilibrium scales exponentially with system size. If our computer needs $`10^6`$ steps per second, Eq. (19) gives the age of the universe for $`N400`$. Even for our simple model the return to non-equilibrium cannot be observed in reasonable time for $`N`$ significantly larger than 70. All these results hold only for a model with identical microscopic times for all of the local variables. We have also studied models with hierarchical or power law access times. In these cases the relaxation to equilibrium is slowed down. The time to reach equilibrium scales exponentially or with a power of system size, depending on the distribution of microscopic times. ## VII Acknowledgement All authors are grateful for financial support by the German-Israeli Foundation. This work also benefitted from a conference at the Max-Planck Institut für Physik komplexer Systeme, Dresden.
warning/0007/hep-th0007173.html	ar5iv	text	# 1 Introduction ## 1 Introduction The construction of string theory models which at low energies reproduce the basic features of the standard model of elementary particles (or some extension thereof) is a non-trivial task from the physical point of view. Interestingly, it also involves beautiful mathematics, and leads to a rich interplay between the physical and mathematical points of view. To quote a traditional example, the construction of four-dimensional $`𝒩=1`$ supersymmetric heterotic string vacua involves the study of stable $`G`$-bundles on Calabi-Yau threefolds, with $`G`$ a subgroup of $`E_8\times E_8`$ or $`SO(32)`$ (for recent results on this approach, see ). Recent developments in string theory have led to new kinds of string theory vacua with potential phenomenological interest. A particular class corresponds to compactification on a Calabi-Yau threefold $`X`$ with gauge bundles defined on subvarieties of the internal space (times non-compact spacetime). In physical terms, the compactification includes a set of D-branes , which are dynamical extended objects partially wrapped on the internal manifold, and filling non-compact spacetime. Their dynamics is controlled by a gauge theory defined on their world-volume. The rank of the gauge bundle they carry is usually referred to as the number of D-branes. In the following we center on the simplest case of D3-branes in type IIB string theory, where each D-brane spans four-dimensional spacetime times a point $`PX`$. The properties of the corresponding gauge theory sector in the low-energy theory are then determined in terms of the local geometry of $`X`$ around $`P`$. For a set of D3-branes sitting at a smooth point in $`X`$, the low-energy gauge field theory on the world-volume has $`𝒩=4`$ supersymmetry, and is therefore non-chiral and phenomenologically uninteresting. Chiral gauge sectors arise when D3-branes sit at singular points in $`X`$. Singularities preserving at least $`𝒩=1`$ supersymmetry, and at a finite distance in Calabi-Yau moduli space must be Gorenstein canonical singularities. In Section 2 we discuss several systems of branes at singular points in non-compact Calabi-Yau spaces. In section 2.1 we center on the simple case of threefold quotient singularities, which leads to an implementation of the McKay correspondence in string theory. In section 2.2. and 2.3 we point out several generalizations suggested by string theory. In section 3 we discuss how an extremely simple $`𝐙_\mathrm{𝟑}`$ quotient singularity leads to gauge groups and particle contents remarkably similar to those of the (minimal supersymmetric) standard model. I am grateful to G. Aldazabal, A. Hanany, L. E. Ibáñez, J. Park, F. Quevedo and R. Rabadán for collaboration and useful discussion on these issues. I also thank M. González for encouragement and support. ## 2 Branes at singularities ### 2.1 Branes at orbifold singularities Let $`\mathrm{\Gamma }`$ be a discrete group of $`SU(3)`$, and $`\{𝐫_i\}`$ the set of its unitary irreducible representations. We want to consider a set of D3-branes at the origin of $`𝐂^\mathrm{𝟑}/\mathrm{\Gamma }`$, where $`\mathrm{\Gamma }`$ acts on $`𝐂^3`$ through a three-dimensional representation $`_\mathrm{𝟑}`$. Consider first a set of $`N`$ D3-branes at the origin in $`𝐂^3`$, labeled by an index $`a=1,\mathrm{},N`$ referred to as Chan-Paton index. Quantization of open strings with endpoints of the D3-branes leads to a set of dynamical fields propagating on the D3-brane world-volume. In terms of $`𝒩=1`$ supersymmetry multiplets, the corresponding gauge field theory contains a set of vector multiplets (each containing one gauge field and one complex fermion) with gauge group $`U(N)`$, and three chiral multiplets $`\mathrm{\Phi }^a`$, $`a=1,2,3`$ (each containing one complex scalar and one complex fermion). The latter transform in the adjoint representation of $`U(N)`$ and form a triplet under the $`SU(3)`$ action on $`𝐂^\mathrm{𝟑}`$. The interactions are encoded in the superpotential function $`W=ϵ_{abc}\mathrm{tr}(\mathrm{\Phi }^a\mathrm{\Phi }^b\mathrm{\Phi }^c)`$, where $`ϵ_{abc}`$ correspond to the components of the $`SU(3)`$ invariant tensor. Following (see for generalizations, and for related discussions) the field theory on D3-branes at the origin in $`𝐂^3/\mathrm{\Gamma }`$ is obtained from the above field theory associated to $`N`$ D3-branes in flat space, by keeping the states which are invariant under the combined action of $`\mathrm{\Gamma }`$ on $`𝐂^3`$ (as determined by $`_\mathrm{𝟑}`$) and on the space of Chan-Paton indices (through a $`N`$-dimensional representation $``$, with decomposition $`=_iN_i𝐫_i`$). Following , we regard fields in the adjoint of $`U(N)`$ as $`\mathrm{Hom}(𝐂^N,𝐂^N)`$. The projection on the $`𝒩=1`$ vector multiplets leaves the following fields $`\mathrm{Hom}(𝐂^N,𝐂^N)^\mathrm{\Gamma }={\displaystyle \underset{i}{}}\mathrm{Hom}(𝐂^{N_i},𝐂^{N_i})`$ (1) corresponding to a gauge group $`_iU(N_i)`$. The projection on the $`SU(3)`$ triplet of $`𝒩=1`$ chiral supermultiplets leaves the following fields $`(_\mathrm{𝟑}\mathrm{Hom}(𝐂^N,𝐂^N))^\mathrm{\Gamma }={\displaystyle \underset{i,j}{}}a_{ij}^\mathrm{𝟑}\mathrm{Hom}(𝐂^{N_i},𝐂^{N_j})`$ (2) where $`a_{ij}^\mathrm{𝟑}`$ are defined by $`_\mathrm{𝟑}𝐫_i=_ja_{ij}^\mathrm{𝟑}𝐫_j`$. Hence we obtain $`a_{ij}^\mathrm{𝟑}`$ $`𝒩=1`$ chiral multiplets transforming in the representation $`(N_i,\overline{N}_j)`$ of the gauge group. The superpotential is obtained by restricting the above one to the surviving fields. The field content on the D3-brane world-volume can be encoded in a quiver diagram, where the $`i^{th}`$ node represents the $`U(N_i)`$ factor in the gauge group, and $`a_{ij}^\mathrm{𝟑}`$ oriented arrows from the $`i^{th}`$ to the $`j^{th}`$ node correspond to the $`𝒩=1`$ chiral multiplets in the $`(N_i,\overline{N}_j)`$ representation. Finally, closed triangles of oriented arrows are associated to superpotential couplings of the corresponding chiral multiplets. The quiver for a $`𝐂^\mathrm{𝟑}/𝐙_\mathrm{𝟓}`$ singularity is depicted in Fig. 2a. Several interesting mathematical connections arise at this point. For instance, the quiver diagrams encoding the field theory content and interactions coincide with the McKay quivers associated to the singularity, and which are related to the homology of the resolved space by the McKay correspondence . The correspondence arises in the string theory context since branes giving rise to a specific gauge factor have the geometrical interpretation of higher-dimensional branes wrapped on homology cycles of the space (see for a detailed description). Hence, gauge theory data (the quiver diagram) are related to the homology of the ambient space. Also, if $``$ is chosen to be the regular representation of $`\mathrm{\Gamma }`$, the moduli space of vacua of the gauge theory corresponds to the space of possible locations of the D3-brane, which is isomorphic to the transverse space $`𝐂^3/\mathrm{\Gamma }`$. The construction of the moduli space amounts to performing a symplectic quotient in the subspace of fields subject to relations $`\frac{W}{\mathrm{\Phi }_i}=0`$ (F-term constraints) determined by the superpotential . It provides the string theory counterpart of the construction of $`𝐂^\mathrm{𝟑}/𝚪`$ as the moduli of representations of a quiver diagram with relations . In the particular case $`\gamma SU(2)`$, studied in , one recovers the hyperkähler quotient construction of ALE spaces . String theory also provides a description of the resolved spaces, by a suitable modification of the symplectic quotient due to a non-zero moment map (D-term). We conclude this section with some physical considerations. In string theory D-branes are sources of certain $`p`$-form gauge fields from the closed string sector, whose equations of motion may impose certain consistency conditions on the D-brane configuration. In the particular case of D3-branes at $`𝐂^3/\mathrm{\Gamma }`$ singularities, the equations of motion for fields located at the singularity impose the so-called twisted tadpole cancellation conditions, which for quotient singularities amount to the vanishing of the character of the representation $``$. This constraint is equivalent to the cancellation of non-abelian anomalies in the gauge field theory on the D3-branes world-volume . They also imply that the remaining mixed $`U(1)`$\- non-abelian anomalies have a factorized form and are cancelled by a version of the Green-Schwarz mechanism . The anomalous $`U(1)`$ factors become massive and disappear from the low-energy dynamics. ### 2.2 Generalizations #### 2.2.1 Non-orbifold singularities There is no simple recipe to obtain the field theory on the world-volume of stacks of D3-branes at a general singularity. However, the requirement that its moduli space should correspond to the space of possible locations of D-branes in the transverse space is enough to determine the field theory in the simple example of the conifold singularity $`X_{con}`$ , defined by the hypersurface $`x^2+y^2+z^2+w^2=0`$ in $`𝐂^4`$. A set of $`N`$ D3-branes at a conifold singularity yields a $`𝒩=1`$ supersymmetric field theory with gauge group $`U(N)\times U(N)`$ and chiral multiplets $`A_i`$, $`B_i`$, $`i=1,2`$ in the representations $`(N,\overline{N})`$ and $`(\overline{N},N)`$ under the gauge group, and in the representations $`(2,1)_{1/2}`$, $`(1,2)_{1/2}`$ under the $`SU(2)^2\times U(1)`$ symmetry group of $`X_{con}`$. The interactions are determined by a quartic superpotential $`W=ϵ^{ij}ϵ^{kl}\mathrm{tr}A_iB_jA_kB_l`$. This example provides a whole new family of models (see for a related discussion) corresponding to D3-branes at quotients of the conifold, $`X_{con}/\mathrm{\Gamma }`$, with $`\mathrm{\Gamma }`$ a subgroup of $`SU(2)\times SU(2)`$ in order to preserve $`𝒩=1`$ supersymmetry. The strategy is, as in section 2.1, to start with D3-branes at $`X_{con}`$, embed the action of $`\mathrm{\Gamma }`$ on the Chan-Paton indices, and keep only fields which are invariant under the combined geometrical and Chan-Paton action of $`\mathrm{\Gamma }`$. The resulting field theories can be encoded in a quiver diagram, with nodes and arrows corresponding to gauge factors and chiral multiplets, and superpotential terms correspond to closed polygons formed by four arrows. When $`\mathrm{\Gamma }`$ acts on the Chan-Paton indices in the regular representation, the moduli space of the gauge theory is a symplectic quotient subject to F-term constraints from the superpotential, as obtained in explicit examples , generalizing the orbifold result. The string theory construction also suggest these quiver diagrams also encode the homology of the resolution of $`X_{con}/\mathrm{\Gamma }`$. A general recipe to obtain the field theory on D3-branes at a general toric singularities $`X`$ was proposed in (see for a more general discussion). It is based on realizing $`X`$ as a partial resolution of a threefold quotient singularity $`𝐂^3/\mathrm{\Gamma }`$ for suitable $`\mathrm{\Gamma }SU(3)`$. The construction requires a precise identification of the effect of the resolution on the field theory, which is a specific Higgs mechanism in which several gauge factors break to their diagonal subgroup, and some chiral multiplets become massive and disappear from the light spectrum. The detailed map has been worked out in some explicit examples, e.g. in , and provides an explanation for the existence of quiver diagrams for non-orbifold singularities. They are obtained by joining nodes and deleting arrows in the quiver of the initial quotient singularity, as dictated by the Higgs mechanism in the field theory. The quiver version of the resolution of the $`𝐂^\mathrm{𝟑}/(𝐙_\mathrm{𝟐}\times 𝐙_\mathrm{𝟐})`$ singularity to the conifold is shown in Figure 1 . It would be interesting to obtain a more precise characterization of these operations on quiver diagrams. Another interesting question is related to the uniqueness of the field theory corresponding to a given singularity . In some cases, different field theories may have isomorphic moduli spaces. Mathematically, the representation moduli of the corresponding quivers with relations are isomorphic. From the physical point of view, the equality of moduli spaces may in some cases be a reflection of Seiberg duality, a non-trivial infrared equivalence of seemingly different field theories. #### 2.2.2 Orientifold projections Type IIB string theory is invariant under an operation $`\mathrm{\Omega }`$ which reverses the orientation on the world-sheet. Hence it is possible to consider modding out type IIB configurations by $`\mathrm{\Omega }`$, possibly accompanied by a geometric involution $`g`$, also leaving the theory invariant . The new configurations thus obtained are called orientifolds, and are characterized by the inclusion of non-orientable world-sheets in the string theory perturbative expansion. We are interested in studying D3-branes at orientifold threefold singularities, i.e. $`\mathrm{\Omega }g`$ quotients of the system of D3-branes at singularities in CY threefolds. The best studied case is again that of quotient singularities $`𝐂^3/\mathrm{\Gamma }`$, and for the sake of clarity here we center on $`\mathrm{\Gamma }=𝐙_𝐍`$ with the action of the generator $`\theta `$ of $`𝐙_N`$ on $`𝐂^3`$ represented by $`\gamma =\mathrm{diag}(e^{2\pi it_1/N},e^{2\pi it_2/N},e^{2\pi it_3/N})`$, with $`t_i`$ integers defined modulo $`N`$, and $`_{a=1}^3t_a=0\mathrm{mod}N`$. Let us embed the action of $`\theta `$ on the Chan-Paton indices by a diagonal matrix $`\gamma _{\theta ,3}`$ with $`N_k`$ entries $`e^{2\pi ik/N}`$. The field theory one obtains has vector multiplets with gauge group $`_{i=1}^NU(N_i)`$, and chiral multiplets $`\mathrm{\Phi }_{i,i+t_a}^a`$, $`a=1,\mathrm{},3`$, $`i=1,\mathrm{},N`$ in the representation $`(N_i,\overline{N}_{i+t_a})`$. The orientifold action may also be embedded in the space of Chan-Paton indices, through a matrix $`\gamma _{\mathrm{\Omega }g,3}`$. The representations of $`\mathrm{\Gamma }`$ and $`\mathrm{\Omega }g`$ are usually constrained from mutual consistency requirements . Several solutions to these conditions are known (quite exhaustively in the two-fold case ), but a complete classification is lacking. In the following we center on a concrete example where $`g`$ inverts all coordinates in $`𝐂^3`$, and exchanges closed string fields in oppositely twisted sectors, i.e. has a non-trivial action on the homology cycles shrunk at the singularity. We also choose $`\gamma _{\mathrm{\Omega }g,3}`$ symmetric and such that exchanges the eigenspaces of conjugate eigenvalues in $`\gamma _{\theta ,3}`$. The action of $`𝐙_N`$ on Chan-Paton indices is therefore constrained to form a real representation, $`N_k=N_k`$. The effect of the orientifold projection on the spectrum is as follows. Gauge factors associated to conjugate irreducible representations of $`𝐙_N`$ are identified, and unitary gauge factors associated to real representations are reduced to their maximal orthogonal subgroups. Correspondingly, the chiral multiplets $`\mathrm{\Phi }_{i,i+t_a}^a`$ and $`\mathrm{\Phi }_{it_a,i}^a`$ are identified. When $`i+t_a=i`$ the bifundamental field is projected down to a two-index antisymmetric representation of the unitary gauge factor after the projection. The quiver diagram for an orientifold of a $`𝐂^\mathrm{𝟑}/𝐙_\mathrm{𝟓}`$ singularity is depicted in Figure 2b. The effect on the quiver diagram of $`𝐂^3/\mathrm{\Gamma }`$ is an identification of nodes and arrows related by a $`𝐙_2`$ action, with a specific prescription for nodes and arrows which are mapped to themselves (it would be interesting to characterize these mappings more precisely in order to allow the classification of resulting models). One may wonder about the meaning of the resulting quiver. String theory suggests it encodes the information about the homology of the resolved orientifold singularity. In fact, the string theory counting of homology cycles in the resolved space (by counting of twisted sector blow-up moduli) gives roughly speaking half the number encountered before the orientifold projection, due to the non-trivial action of $`g`$ by exchanging oppositely twisted sectors. This agrees with the counting of nodes in the orientifolded quiver, which also gives half the number encountered before the orientifold projection. We conclude by pointing out that orientifolds of non-orbifold singularities can be constructed as partial resolutions of orientifold of quotient singularities. Preliminary results on some simple examples indicate that the effect of the orientifold projection on the quiver of the non-orbifold singularity is also a $`𝐙_\mathrm{𝟐}`$ involution, and that the orientifold singularity can be constructed as a symplectic quotient on the space of fields constrained by the F-terms conditions. ## 3 Particle physics In this section we discuss particular examples of singularities leading to interesting low-energy field theories, in that they resemble the structure of the (minimal supersymmetric) standard model (or some extension thereof) which we now review. Such models would provide phenomenologically interesting string theory vacua when embedded in a compact Calabi-Yau context. ### 3.1 Review of the (minimal supersymmetric) standard model All known non-gravitational interactions between elementary particles are described by the quantum field theory known as standard model. The simplest supersymmetric extension of this model contains vector multiplets with gauge group $`SU(3)\times SU(2)\times U(1)`$. It also contains a set of chiral multiplets transforming in three copies of the representation $`(3,2)_{1/6}+(\overline{3},1)_{2/3}+(\overline{3},1)_{1/3}+(1,2)_{1/2}+(1,1)_1,`$ (3) where subscripts denote $`U(1)`$ charges. Successful breaking of the electroweak interactions requires also at least one chiral multiplet in the representation $`(1,2)_{1/2}+(1,2)_{1/2}`$. Interactions are encoded in a set of complicated gauge invariant functions of these chiral multiplets, the superpotential and gauge kinetic functions, which are holomorphic, and the Kähler potential, which is real. We will skip their details since realistic phenomenology usually involves additional model-dependent assumptions, like additional global symmetries. ### 3.2 Realistic models from $`𝐙_\mathrm{𝟑}`$ singularities The replication of families is an intriguing feature of the above field theory. To reproduce it from branes at singularities, the case of $`𝐂^3/𝐙_3`$ with the $`𝐙_3`$ action defined by $`(z_1,z_2,z_3)(e^{2\pi i/3}z_1,e^{2\pi i/3z_2},e^{2\pi i/3}z_3)`$, is singled out, in that it leads to natural triplication. The two examples we are to consider are based on this singularity. The first example we consider forms a subsector of the model considered in . It is constructed by placing eleven D3-branes on the orientifold of the $`𝐂^3/𝐙_3`$ singularity, introduced in section 2.2.2. We choose $`\gamma _{\theta ,3}=\mathrm{diag}(1,e^{2\pi i/3}\mathrm{𝟏}_5,e^{2\pi i/3}\mathrm{𝟏}_5)`$. Following our rules above, we obtain vector multiplets with gauge group $`SU(5)`$ (the $`U(1)`$ factor disappears from the light spectrum as explained in section 2.1), and there are chiral multiplets in three copies of the representation $`5+\overline{10}`$. This spectrum resembles the structure of $`SU(5)`$ grand unified theories, which reproduce the spectrum of the minimal supersymmetric standard model when an additional field in the adjoint representation is present to break $`SU(5)`$ to the standard model group through the Higgs mechanism (an additional $`5+\overline{5}`$ pair is further required for electroweak symmetry breaking). Unfortunately, all Higgs fields are absent in our field theory, which therefore is suggestive but not truly realistic. Our second example is based on the $`𝐂^3/𝐙_3`$ orbifold (rather than orientifold) singularity <sup>2</sup><sup>2</sup>2This model has been studied in and is similar to a subsector of models in ., with Chan-Paton embedding $`\gamma _{\theta ,3}=\mathrm{diag}(\mathrm{𝟏}_\mathrm{𝟑},e^{2\pi i/3}\mathrm{𝟏}_2,`$ $`e^{2\pi i/3})`$. This embedding does not satisfy the tadpole cancellation condition stated in section 2.1 (the character of the representation of $`𝐙_3`$ is non-zero). The problem can be solved by introducing an additional set of D-branes, for instance D7<sub>a</sub>-branes, with $`a=1,2,3`$, wrapped on the complex surfaces $`z_a=0`$, which do not break further supersymmetries. For non-compact spaces these additional branes are non-dynamical, but they contribute additional fields in the four-dimensional world-volume of the D3-branes, arising from open strings stretched between D3- and D7-branes. Hence the corresponding gauge field theories correspond to extended quivers, where additional nodes correspond to the global symmetries (symmetries on the D7-branes), and additional arrows correspond to 3-7 and 7-3 string states. These quiver diagrams generalize to the threefold case those in . Their moduli spaces therefore provide a generalization of the Kronheimer-Nakajima construction of the moduli space of instantons on ALE spaces . The $`𝐙_3`$ acts on the space of D7-brane Chan-Paton indices, through matrices $`\gamma _{\theta ,7_a}`$. The modified tadpole cancellation conditions (which reduce to $`_a\mathrm{tr}\gamma _{\theta ,7_a}+3\mathrm{t}\mathrm{r}\gamma _{\theta ,3}=0`$) are satisfied for the very symmetric choice $`\mathrm{tr}\gamma _{\theta ,7_a}=\mathrm{tr}\gamma _{\theta ,3}`$, $`a=1,2,3`$. Choosing e.g. $`\gamma _{\theta ,7_a}=\mathrm{diag}(e^{2\pi i/3},e^{2\pi i/3}\mathrm{𝟏}_2)`$ we obtain the following spectrum $`\begin{array}{ccc}33\mathrm{strings}& \mathrm{Gauge}\mathrm{group}& SU(3)\times SU(2)\times U(1)\\ & \mathrm{Chiral}\mathrm{multiplets}& 3[(3,2)_{1/6}+(1,2)_{1/2}+(\overline{3},1)_{2/3}]\\ 37_a\mathrm{strings}& \mathrm{Chiral}\mathrm{multiplets}& (3,1)_{1/3}+2(1,2)_{1/2}+\\ a=1,2,3& & +2(\overline{3},1)_{1/3}+(1,1)_1\end{array}`$ (8) We have included the charges under the only non-anomalous linear combinations of the $`U(1)`$ factors in the original $`U(3)\times U(2)\times U(1)`$ gauge group. This $`U(1)`$ does not become massive and plays the role of hypercharge in the standard model we have just constructed. The quiver diagram of this gauge theory is shown in Figure 3. Notice how close to the spectrum in section 3.1 one can get using very simple singularities. In particular, the later model differs from the MSSM just in that it produces three Higgs pairs instead of one, and in fact constitutes one of the simplest semi-realistic string models ever built. Since such constructions would provide a rationale for the existence of three generations (one per complex transverse dimensions), and for hypercharge assignments (see for discussion of hypercharge in a other brane contexts), we believe these examples illustrate the phenomenological interest of string theory compactifications with branes at singularities.
warning/0007/astro-ph0007474.html	ar5iv	text	# Launching of jets and the vertical structure of accretion disks ## 1. Introduction Jets and other outflows are commonly observed from young stellar objects, interacting binary stars and active galactic nuclei. It is generally believed that the outflows are produced by the accretion disks in these systems. The widely differing properties of the central objects in these systems suggest that a mechanism may be at work that is largely independent of the nature of the central object (e.g. Livio 1997). A particularly promising mechanism for the acceleration and collimation of outflows involves a large-scale poloidal magnetic field that threads the disk. The influential model of Blandford & Payne (1982) established the significance of the angle of inclination, $`i`$, of the poloidal magnetic field lines to the vertical at the surface of the disk. When $`i>30^{}`$, matter that is just above the surface, being forced to corotate with the foot-point of the field line, is accelerated outwards along it by the centrifugal force. At greater distances, the flow may collimate into a jet by magnetic hoop stresses or poloidal collimation. A clear review of the physics of magnetocentrifugally driven outflows has been given by Spruit (1996). Several groups have performed axisymmetric numerical simulations of outflows using this mechanism (e.g. Ustyugova et al. 1995; Romanova et al. 1997; Ustyugova et al. 1999; Ouyed & Pudritz 1997a, 1997b, 1999; Krasnopolsky, Li, & Blandford 1999). Like the model of Blandford & Payne (1982), these calculations do not resolve the vertical structure of the disk but treat it as a boundary surface that loads mass at a specified rate on to the magnetic field lines. While these simulations have convincingly demonstrated some aspects of the magnetocentrifugal model of jet production, several fundamental issues remain to be resolved. In this paper, we will attempt to address two of these questions. First, what determines the rate of mass loss in the outflow? How does this depend on the strength and inclination of the magnetic field, or on other properties of the disk? Second, what is the long-term evolution of a large-scale magnetic field in a disk? Is it possible to assemble a magnetic configuration suitable for jet launching, or to maintain it against dissipation? A proper understanding of these issues is essential if we are to explain the conditions that regulate the production of astrophysical jets. The magnetocentrifugal mechanism can explain the acceleration of outflows even when the matter is ‘cold’ in the sense that the temperature is much less than the virial temperature or escape temperature. However, as mentioned by Blandford & Payne (1982) and calculated in detail by Ogilvie (1997) and by Ogilvie & Livio (1998; hereafter Paper I), some amount of thermal assistance is still required to overcome the potential barrier between the surface and the slow magnetosonic point (‘sonic point’). The barrier occurs because the field lines are not straight inside the disk and the angular velocity deviates slightly from the Keplerian value because of the radial Lorentz force associated with the bending of the field lines (Shu 1991; Wardle & Königl 1993). A proper calculation of the vertical structure of the disk, including these effects, is therefore required in order to determine the height of the potential barrier and the rate of mass loss in the outflow. Indeed, it is obvious from the well-known properties of transonic outflows that the mass loss rate depends on the physical conditions below the sonic point, and that the vertical structure must therefore be resolved. In Paper I we carried out such a calculation for disks that are rather strongly magnetized in the sense that the magnetorotational instability, which leads to turbulence in accretion disks (Balbus & Hawley 1998), is suppressed or nearly so. We assumed that the magnetic field enforces strict isorotation and showed that the potential barrier increases very steeply as the field is made stronger and the disk becomes more sub-Keplerian. This effect would suppress outflows from strongly magnetized disks unless an additional source of energy, such as coronal heating, were present (in accord with the suggestion made by Livio 1997). A disadvantage of the calculation in Paper I is that no explanation was given for the source of the effective viscosity of a disk in which the magnetorotational instability is suppressed. The solutions did not extend to field strengths much below the stability boundary. When the field is weaker than this, it cannot be expected to enforce isorotation, and the model requires some modification. In addition, the effect of the turbulence on the mean magnetic field should be modeled, most simply through the introduction of an effective magnetic diffusivity. However, the presence of a turbulent diffusivity may cause problems for the magnetocentrifugal mechanism. The analysis by Lubow, Papaloizou, & Pringle (1994a) suggests that, if the effective magnetic Prandtl number of the disk is of order unity, as might be expected, it is impossible to sustain a steady configuration with a significantly inclined field in a thin disk (see also Heyvaerts, Priest, & Bardou 1996; Reyes-Ruiz & Stepinski 1996). Although the accretion flow tends to drag magnetic flux inwards, the turbulent diffusivity expels flux faster if the inclination is large. The inclination angle in a steady state is then expected to be comparable to the angular semi-thickness $`H/r`$ of the disk. This constraint can be avoided if the accretion flow is due primarily to the loss of angular momentum in an outflow, and might also be relieved by a dynamo operating in the disk if special conditions are met (Campbell, Papaloizou, & Agapitou 1998). The purpose of the present paper is to explore a model for the vertical structure of magnetized disks. We will extend the analysis of Paper I to allow for the possibility of more weakly magnetized disks in which there is turbulence, and strict isorotation does not hold. We will also give a better treatment of the matching between the disk and the atmosphere by applying photospheric boundary conditions, allowing us to calculate the mass loss rate in the outflow explicitly. At the same time, we will solve for the rate of dragging of magnetic flux and thereby refine the previous estimates which may have been based on oversimple arguments. Related calculations of the vertical structure of magnetized disks, and of the disk-jet connection, have been made by Königl (1989), Wardle & Königl (1993), Li (1995, 1996), Ferreira (1997), Campbell (1999), Casse & Ferreira (2000), and Shalybkov & Rüdiger (2000). Our method and results are significantly different from all previous calculations and comparisons of the relevant issues will be made towards the end of this paper. ## 2. Context of the present calculation To set this calculation in its proper context, we consider here some of the wider issues relating to magnetized accretion disks and outflows. For a thin disk without a large-scale magnetic field, there is a clear division of the physical problem into two aspects. First, a model is required for the local vertical structure at any given radius $`r`$ and time $`t`$. Such a model takes the surface mass density $`\mathrm{\Sigma }`$ as a parameter and predicts quantities such as the vertically integrated viscous stress $`𝒢`$. The vertical structure may be assumed to be stationary on the dynamical time-scale. Second, the one-dimensional conservation equations for mass and angular momentum determine how $`\mathrm{\Sigma }(r,t)`$ evolves on the viscous time-scale, given the local relation between $`𝒢`$ and $`\mathrm{\Sigma }`$ (e.g. Lynden-Bell & Pringle 1974). For a disk with a large-scale poloidal magnetic field but no outflows (e.g. a polytrope with a force-free or vacuum exterior), the situation is similar but more complicated. Now the local vertical structure at a given radius depends not only on $`\mathrm{\Sigma }`$ but also on the vertical magnetic field $`B_z`$ and the inclination angle $`i`$. In addition to an evolutionary equation for $`\mathrm{\Sigma }`$, there is a one-dimensional conservation equation for magnetic flux, which determines how the flux function $`\psi (r,t)`$ evolves on the viscous time-scale. $`B_z`$ is simply related to the radial derivative of $`\psi `$. Finally, the force-free magnetic equilibrium in the exterior of the disk leads to a global relation between $`\psi `$ and the inclination angle $`i`$. The full problem therefore involves an integro-differential equation. This has been considered by Lubow et al. (1994a), although they did not examine the effect of the magnetic field on the vertical disk structure. When outflows occur along the magnetic field lines, further couplings exist. The exterior magnetic field is modified from the force-free solution by the inertial forces associated with the outflow. This effect depends on the amount of mass loading, and is therefore coupled to the launching problem. The exterior field is no longer purely poloidal but becomes significantly twisted beyond the Alfvén surface, although non-axisymmetric (e.g. kink) instabilities may set in here (e.g. Spruit, Foglizzo, & Stehle 1997). This part of the solution is established on the Alfvén travel time to the Alfvén surface, which is comparable to the dynamical time-scale of the disk. Finally, the mass loss and especially the angular momentum loss in the outflow feed back into the evolutionary equations for the disk. The full problem therefore involves three aspects: the local vertical structure of the disk and the launching problem; the global structure of the exterior magnetic field and the dynamics of the outflow along those field lines; and the evolution of mass, angular momentum, and magnetic flux on a longer time-scale. The first aspect is the subject of this paper, while the second aspect is well described by some of the numerical simulations referred to above. The third aspect has been considered implicitly in steady models such as that of Casse & Ferreira (2000), but the development of a complete evolutionary scheme for mass, angular momentum and magnetic flux in the general, non-steady case remains a challenge (see Lovelace, Newman, & Romanova 1997 for a simplified version of such a scheme).<sup>1</sup><sup>1</sup>1There have been many calculations simulating the time-dependent development of a jet from a numerically resolved disc. For example, Kudoh, Matsumoto, & Shibata (1998) set up a thick torus with a non-rotating corona, and introduced a vertical magnetic field. Owing to the lack of equilibrium in the initial condition, a rapid adjustment occurs, and Kudoh et al. followed the evolution for only a single orbit or so. Unfortunately there is little reason to expect the properties of such transient phenomena to agree with quasi-steady models of jets from thin discs. There has been some confusion in the literature regarding the interplay between these aspects of the problem. For example, studies of the launching problem have often imposed the condition that the poloidal magnetic flux should not migrate radially (e.g. Königl 1989; Li 1995). That is, there should be an instantaneous balance between inward advection of flux by the accretion flow, and outward transport due to turbulent, Ohmic, or ambipolar diffusion. This may not always be appropriate because the migration occurs on the viscous time-scale, whereas the solution of the launching problem need only be stationary on the dynamical time-scale. In Wardle & Königl (1993) the rate of flux migration was treated as a free parameter. Moreover, some of these studies appear to predict the value of the toroidal magnetic field at the surface of the disk, $`B_{\varphi \mathrm{s}}`$, even when the dynamics of the outflow beyond the sonic point has not been considered (e.g. Campbell 1999) or when there is no outflow (e.g. Shalybkov & Rüdiger 2000). This is unsatisfactory because $`B_{\varphi \mathrm{s}}`$ determines the external magnetic torque acting on the disk, which certainly depends on the dynamics of the outflow in the trans-Alfvénic region. These incorrectly specified studies may be motivated by the fact that the model of Blandford & Payne (1982) appears to require $`B_{\varphi \mathrm{s}}`$ to be determined by the disk (through their parameter $`\lambda `$). However, this is misleading and may be attributed to the fact that the Blandford & Payne model is underconstrained because it is missing a boundary condition at large distances from the disk (cf. Ostriker 1997). Indeed, solutions of their model generically behave unphysically at large distances without having passed through the modified fast magnetosonic point. In contrast, Krasnopolsky et al. (1999) have given a clear discussion of which quantities may or may not be specified as boundary conditions when treating the disk as a boundary surface. In their model $`B_{\varphi \mathrm{s}}`$ is determined by the outflow, not by the disk. In our local study, therefore, we will not impose the constraint of zero flux migration; rather, we will impose the value of $`B_{\varphi \mathrm{s}}`$ and determine the rate of flux migration. Usually we will specify $`B_{\varphi \mathrm{s}}=0`$, meaning that the outflow is absent or exerts only a weak torque. Efficient magnetocentrifugal outflows are expected to have $`\|B_{\varphi \mathrm{s}}\|\|B_z\|`$, with the toroidal field becoming comparable to the poloidal component only at greater distances from the source, where the Alfvén surface is located (Spruit 1996). ## 3. The evolution of magnetic flux As noted above, an important investigation of the evolution of the poloidal magnetic flux was made by Lubow et al. (1994a), who concluded that, if the disk is turbulent, with an effective magnetic Prandtl number of order unity, the accretion flow will be almost entirely ineffective in dragging in magnetic flux. This can be explained by noting that the effective magnetic Reynolds number of the accretion flow, based on the disk thickness $`H`$, is small, of order $`H/r`$ (Heyvaerts et al. 1996). Similar arguments suggest that a configuration in which the field lines are bent significantly from the vertical (e.g. to achieve $`i>30^{}`$) cannot be sustained on the viscous time-scale. These arguments are based on a kinematic analysis of the magnetic induction equation in which the radial velocity and magnetic diffusivity are prescribed quantities. They also depend on a simple order-of-magnitude treatment of the vertical structure of the magnetic field. In this paper we will present a numerical treatment of a set of equations, including the induction equation, in which the radial velocity is self-consistently determined. This is important because, in a jet-launching configuration, the magnetic field must become dynamically dominant above some height and will then control the radial velocity. We will find results that differ significantly from the estimates of Lubow et al. (1994) under some circumstances. Before presenting the numerical model, we therefore reconsider the problem of the induction equation from an analytical viewpoint. We assume that the mean poloidal magnetic field is axisymmetric and may be described by a flux function $`\psi (r,z,t)`$ such that $$B_r=\frac{1}{r}\frac{\psi }{z},B_z=\frac{1}{r}\frac{\psi }{r}.$$ (1) For a thin disk containing a bending poloidal field of dipolar symmetry, the flux function has the form (Ogilvie 1997) $$\psi =\psi _0(r,t)+\psi _1(r,z,t),$$ (2) where $`\|\psi _1\|\|\psi _0\|`$. Then $`B_r`$ and $`B_z`$ are comparable in magnitude if $`\|\psi _1/\psi _0\|=O(H/r)`$. The flux content of the disk is determined essentially by $`\psi _0`$, while $`\psi _1`$ allows for bending of the field lines within the disk. We assume that the mean magnetic field in a turbulent disk may be treated within the framework of mean-field electrodynamics (e.g. Moffatt 1978). In the absence of a mean-field dynamo effect (or in the absence of a toroidal field), the flux function then satisfies the mean induction equation, $$\frac{\psi }{t}+𝒖\psi =\eta r^2\left(\frac{1}{r^2}\psi \right),$$ (3) where $`𝒖`$ is the mean velocity and $`\eta `$ the turbulent magnetic diffusivity. In a thin disk the dominant terms are $$\frac{\psi _0}{t}+u_r\frac{\psi _0}{r}=\eta r\frac{}{r}\left(\frac{1}{r}\frac{\psi _0}{r}\right)+\eta \frac{^2\psi _1}{z^2}.$$ (4) The neglected terms are all smaller than those retained because $`\|\psi _1\|\|\psi _0\|`$ and $`\|u_r\|\|u_z\|`$. The first term on the right-hand side is also sometimes neglected (e.g. Lubow et al. 1994a). Although this equation appears to be an evolutionary equation for $`\psi _0`$ it should be observed that the equation is defined for all values of $`z`$ whereas $`\psi _0`$ is a function of $`r`$ and $`t`$ only. Moreover, the equation is not closed because of the term involving $`\psi _1`$. Following Lubow et al. (1994a), we might attempt to close the equation by vertical averaging. Conventionally in accretion-disk theory one uses density-weighted averages $`\overline{u}_r`$ and $`\overline{\eta }`$ defined by $$\mathrm{\Sigma }\overline{u}_r=\rho u_rdz,\mathrm{\Sigma }\overline{\eta }=\rho \eta dz,$$ (5) where $`\rho `$ is the density and $$\mathrm{\Sigma }=\rho dz$$ (6) is the surface density, the integrals being over the full vertical extent of the disk. The density-weighted average $`\overline{u}_r`$ is particularly appropriate because it is directly related to the mass accretion rate. The vertical average of the induction equation is then $$\frac{\psi _0}{t}+\overline{u}_r\frac{\psi _0}{r}=\overline{\eta }r\frac{}{r}\left(\frac{1}{r}\frac{\psi _0}{r}\right)\frac{r}{\mathrm{\Sigma }}\rho \eta \frac{B_r}{z}dz.$$ (7) Lubow et al. proceeded by approximating the last term as $`(r/H)\overline{\eta }B_{r\mathrm{s}}`$. This closes the equation because $`B_{r\mathrm{s}}`$, which is proportional to the vertically integrated toroidal electric current, can be globally related to the flux function by considering the force balance exterior to the disk. If the exterior region is treated either as an insulating vacuum, or as a force-free medium with vanishing $`B_\varphi `$, the exterior poloidal field is potential. This leads to a global relation of the form $$B_{r\mathrm{s}}=^1\psi _0,$$ (8) where $``$ is a certain linear integral operator (see also Ogilvie 1997). This order-of-magnitude estimate of the final term appears reasonable if the field lines bend mainly in the densest layers of the disc near the midplane. However, the term might be significantly overestimated if the bending occurs mainly in the upper layers of the disk where $`\rho \eta `$ is likely to be considerably smaller. Another way of averaging the induction equation reinforces this concern. If the equation is divided by $`\eta `$ and integrated with respect to $`z`$, we obtain $$\frac{\psi _0}{t}+u_{}\frac{\psi _0}{r}=\eta _{}r\frac{}{r}\left(\frac{1}{r}\frac{\psi _0}{r}\right)QB_{r\mathrm{s}},$$ (9) where $$u_{}(r,t)=\frac{u_r}{\eta }dz/\frac{1}{\eta }dz,$$ (10) $$\eta _{}(r,t)=dz/\frac{1}{\eta }dz,$$ (11) and $$Q(r,t)=2r/\frac{1}{\eta }dz.$$ (12) Now the precise shape of the field line has truly been eliminated in favor of $`B_{r\mathrm{s}}`$. This therefore appears to be the ‘correct’ way of averaging the equation. However, we are now faced with vertical averages of $`u_r`$ and $`\eta `$ weighted by $`1/\eta `$, not by $`\rho `$. Even if $`\eta `$ is independent of $`z`$, these could be significantly different from the density-weighted averages, leading to different conclusions about flux dragging. Furthermore, as mentioned above, the presence of the magnetic field will in general change the vertical profiles of $`u_r`$ and $`\eta `$. If the field is extremely weak so that the induction equation may be treated ‘kinematically’, then the $`1/\eta `$ averaging method is surely the correct one. However, the vertical profile of $`u_r`$, in particular, is determined by subtle effects (e.g. Kley & Lin 1992) and is easily distorted by even weak Lorentz forces. In this case we cannot obtain a closed equation because $`u_{}`$ cannot be determined from a knowledge of the mass accretion rate. The essential difficulty is that, while $`\overline{u}_r`$ is the relevant mean velocity for mass accretion, $`u_{}`$ is the appropriate mean for flux accretion. These could be significantly disparate and could even differ in sign. The above discussion shows that the only solution to this problem is to solve explicitly for the vertical structure of the disk including the physics that determines the vertical profile of the radial velocity. ## 4. Mathematical model We now consider the set of equations governing the local vertical structure of a disk containing a mean poloidal magnetic field. The equations of Paper I will be augmented by the inclusion of additional physical effects. ### 4.1. Basic equations for a thin disk The angular velocity is written as $$\mathrm{\Omega }=\mathrm{\Omega }_0(r)+\mathrm{\Omega }_1(r,z,t),$$ (13) where $$\mathrm{\Omega }_0=\left(\frac{GM}{r^3}\right)^{1/2}$$ (14) is the Keplerian value, and $`\|\mathrm{\Omega }_1\|\mathrm{\Omega }_0`$. The flux function is written as $$\psi =\psi _0(r,t)+\psi _1(r,z,t),$$ (15) where $`\|\psi _1\|\|\psi _0\|`$, and we approximate $$B_r=\frac{1}{r}\frac{\psi _1}{z},B_z\frac{1}{r}\frac{\psi _0}{r},$$ (16) so that $`B_z`$ is independent of $`z`$. The required equations may be approximated as follows. The radial component of the equation of motion is $$2\rho r\mathrm{\Omega }_0\mathrm{\Omega }_1=\frac{B_z}{\mu _0}\frac{B_r}{z},$$ (17) where $`\mu _0`$ is the permeability of free space. The azimuthal component is $$\frac{\rho u_r}{r}\frac{d}{dr}(r^2\mathrm{\Omega }_0)=\frac{B_z}{\mu _0}\frac{B_\varphi }{z}+\frac{1}{r^2}\frac{}{r}\left(\rho \nu r^3\frac{d\mathrm{\Omega }_0}{dr}\right),$$ (18) where $`\nu `$ is the kinematic viscosity. The vertical component is $$0=\rho \mathrm{\Omega }_0^2z\frac{}{z}\left(p+\frac{B_r^2}{2\mu _0}+\frac{B_\varphi ^2}{2\mu _0}\right),$$ (19) where $`p`$ is the pressure. The poloidal part of the induction equation is $$\frac{\psi _0}{t}+ru_rB_z=r\eta \left(\frac{B_z}{r}\frac{B_r}{z}\right).$$ (20) The toroidal part is $$0=r\left(B_r\frac{d\mathrm{\Omega }_0}{dr}+B_z\frac{\mathrm{\Omega }_1}{z}\right)+\frac{}{z}\left(\eta \frac{B_\varphi }{z}\right).$$ (21) The energy equation is $$\frac{F}{z}=\rho \nu \left(r\frac{d\mathrm{\Omega }_0}{dr}\right)^2+\frac{\eta }{\mu _0}\left[\left(\frac{B_r}{z}\right)^2+\left(\frac{B_\varphi }{z}\right)^2\right],$$ (22) where $$F=\frac{16\sigma T^3}{3\kappa \rho }\frac{T}{z}$$ (23) is the radiative energy flux, with $`\sigma `$ the Stefan-Boltzmann constant, $`T`$ the temperature, and $`\kappa `$ the Rosseland mean opacity. These equations must be supplemented by constitutive relations specifying the equation of state, the opacity, the viscosity, and the magnetic diffusivity. We will adopt the ideal-gas equation of state, $$p=\frac{k\rho T}{\mu m_\mathrm{H}},$$ (24) (where $`k`$ is the Boltzmann constant, $`\mu `$ the mean molecular mass, and $`m_\mathrm{H}`$ the mass of the hydrogen atom), and a generic power-law opacity, $$\kappa =C_\kappa \rho ^xT^y,$$ (25) where $`C_\kappa `$, $`x`$, and $`y`$ are constants. This includes the cases of Thomson scattering opacity ($`x=y=0`$) and Kramers opacity ($`x=1`$, $`y=7/2`$). ### 4.2. Turbulent viscosity and magnetic diffusivity The viscosity is less certain, and we will adopt the standard prescription $$\rho \nu =\frac{\alpha p}{\mathrm{\Omega }_0},$$ (26) where $`\alpha `$ is a dimensionless constant. For the magnetic diffusivity, we assume that the magnetic Prandtl number, $$\mathrm{Pm}=\frac{\nu }{\eta },$$ (27) is constant. We will further assume either that $`\alpha `$ is a fixed parameter (‘fixed alpha hypothesis’), or that $`\alpha `$ can adjust so as to keep the equilibrium at marginal magnetorotational stability, as explained in Section 4.8 below (‘marginal stability hypothesis’). ### 4.3. Neglected terms Several terms in the equations have been omitted on the grounds that the disk is thin and the solution should be stationary on the dynamical time-scale (although not necessarily on the viscous time-scale). The radial pressure gradient, the vertical variation of radial gravity, and the inertial terms associated with the meridional flow, have been neglected as usual. However, enough terms have been retained to determine the profile of radial velocity in the absence of a magnetic field. Generally, $`\mathrm{\Omega }_1`$ has been neglected relative to $`\mathrm{\Omega }_0`$, and $`\psi _1`$ relative to $`\psi _0`$, except where physically essential. It has also been assumed that the viscosity does not act on the shear components $`u_r/z`$ and $`r\mathrm{\Omega }_1/z`$. Such terms are expected to be relatively unimportant in most cases of interest, and it was found that the inclusion of these terms increases the order of the differential system and makes it difficult to obtain a solution. It is especially unclear how to connect a viscous disk to an inviscid atmosphere, when these terms are included, without introducing an artificial boundary layer. We did not include a dynamo alpha-effect in the equations for the mean magnetic field. Although there is no difficulty in principle in doing so, this effect is even less well understood than the turbulent viscosity and magnetic diffusivity, and we chose to avoid this additional complication in the present study. Finally, in the energy equation (22) it has been assumed that the heat generated in each annulus is radiated away locally and not advected through the disk. ### 4.4. Local and non-local effects The equations of Section 4.1 describe the local vertical structure of an accretion disk with a mean poloidal magnetic field. When supplemented with appropriate boundary conditions, as described below, they constitute a problem similar in type to a stellar structure calculation, although the details are of course very different. Obviously it is possible in principle to include a more detailed equation of state, accurate opacity tables, a more sophisticated approach to radiative transfer, and further possibilities such as convective energy transport. However, the principal uncertainties concern the viscosity and magnetic diffusivity. In fact, these equations are not strictly local in radius and time because some radial derivatives and one time-derivative remain. The case of $`d\mathrm{\Omega }_0/dr=3\mathrm{\Omega }_0/2r`$ is trivial. The terms $`\psi _0/t`$ and $`B_z/r`$ in the poloidal part of the induction equation are retained in order to determine whether the magnetic flux migrates inwards ($`\psi _0/t>0`$) or outwards as a result of the local disk solution. The gradient $`B_z/r`$ can affect this result because it contributes to the diffusion of flux. Therefore $`B_z/r`$ appears as a parameter of the local model and $`\psi _0/t`$ as an eigenvalue; note that both are independent of $`z`$. The case of $`(\rho \nu )/r`$ is more problematic. This viscous term is retained in the angular momentum equation because it partially determines the radial velocity, which in turn causes radial advection of magnetic flux and also affects the shape of the field lines, at least when the field is weak. This term may be approximated by arguing that $$\rho \nu \frac{\overline{\nu }\mathrm{\Sigma }}{H}f\left(\frac{z}{H}\right)$$ (28) in the neighborhood of the radius under consideration, where $$\overline{\nu }\mathrm{\Sigma }=_H^H\rho \nu 𝑑z$$ (29) is the vertically integrated dynamic viscosity, $`H`$ the semi-thickness, and $`f`$ an undetermined dimensionless function. Under this assumption, $$\frac{\mathrm{ln}(\rho \nu )}{\mathrm{ln}r}=\frac{\mathrm{ln}(\overline{\nu }\mathrm{\Sigma })}{\mathrm{ln}r}\left[1+\frac{\mathrm{ln}(\rho \nu )}{\mathrm{ln}z}\right]\frac{\mathrm{ln}H}{\mathrm{ln}r}.$$ (30) The vertical derivative is available as part of the local solution, while $`\mathrm{ln}(\overline{\nu }\mathrm{\Sigma })/\mathrm{ln}r`$ and $`\mathrm{ln}H/\mathrm{ln}r`$ appear as additional dimensionless parameters, which can be estimated from the well-known behavior of the steady, non-magnetized solution (Shakura & Sunyaev 1973). In the limit of a non-magnetized disk, this prescription predicts the radial velocity in the disk as $$u_r=\frac{3\nu }{2r}\left\{1+2\frac{\mathrm{ln}(\overline{\nu }\mathrm{\Sigma })}{\mathrm{ln}r}2\left[1+\frac{\mathrm{ln}(\rho \nu )}{\mathrm{ln}z}\right]\frac{\mathrm{ln}H}{\mathrm{ln}r}\right\}.$$ (31) ### 4.5. Boundary conditions We have derived a sixth-order system of nonlinear ordinary differential equations (ODEs) in $`z`$. The dependent variables may be taken as $`p`$, $`T`$, $`F`$, $`\mathrm{\Omega }_1`$, $`B_r`$, and $`B_\varphi `$. The variables $`\rho `$ and $`u_r`$ are determined algebraically from these. The solution should be symmetric about the mid-plane, with $$F=B_r=B_\varphi =0$$ (32) at $`z=0`$. The solution extends up to a photospheric surface $`z=H`$ at which it is matched to a simple atmospheric model, discussed below. The photospheric boundary conditions are $$F_\mathrm{s}=\sigma T_\mathrm{s}^4$$ (33) and $$\tau _\mathrm{s}=_H^{\mathrm{}}\kappa \rho 𝑑z=\frac{2}{3},$$ (34) where the subscript ‘s’ denotes a surface value. We also have $$B_r=B_{r\mathrm{s}},B_\varphi =B_{\varphi \mathrm{s}}$$ (35) there, where $`B_{r\mathrm{s}}`$ and $`B_{\varphi \mathrm{s}}`$ are assigned values which, physically, are determined by the solution of the global exterior outflow problem (not considered here). With these boundary conditions the equation of angular momentum conservation, obtained from a vertical integration of equation (18), has the form $$\frac{\dot{M}}{2\pi r^2}\frac{d}{dr}(r^2\mathrm{\Omega }_0)=\frac{2B_zB_{\varphi \mathrm{s}}}{\mu _0}+\frac{1}{r^2}\frac{}{r}\left(\overline{\nu }\mathrm{\Sigma }r^3\frac{d\mathrm{\Omega }_0}{dr}\right),$$ (36) where $$\dot{M}=2\pi r_H^H\rho u_r𝑑z$$ (37) is the mass accretion rate, not necessarily constant. This shows the contributions to angular momentum transport from magnetic and viscous torques. The solution is determined as follows. One guesses the values of the quantities $`H`$, $`F_\mathrm{s}`$, and $`\mathrm{\Omega }_{1\mathrm{s}}`$. This fixes the atmospheric model and determines $`p_\mathrm{s}`$ and $`T_\mathrm{s}`$. The quantities $`B_{r\mathrm{s}}`$ and $`B_{\varphi \mathrm{s}}`$ are given. One must further guess $`\psi _0/t`$ to start the downward integration. The three guessed quantities should be adjusted to match the three symmetry conditions on $`z=0`$. The main parameters of the model are then $`\mathrm{\Omega }_0`$, $`\mathrm{\Sigma }`$, $`B_z`$, $`B_{r\mathrm{s}}`$, $`B_{\varphi \mathrm{s}}`$, $`\alpha `$, and $`\mathrm{Pm}`$. The quantity $`\psi _0/t`$ is to be determined as an eigenvalue. ### 4.6. Atmospheric model In the simplest atmospheric model, $`T`$, $`F`$, $`B_r`$, and $`B_\varphi `$ are independent of $`z`$ between the photosphere and the sonic point, with $`T`$ $`=`$ $`T_\mathrm{s}=T_{\mathrm{eff}},`$ (38) $`F`$ $`=`$ $`F_\mathrm{s}=\sigma T_\mathrm{s}^4,`$ (39) $`B_r`$ $`=`$ $`B_{r\mathrm{s}}=B_z\mathrm{tan}i,`$ (40) $`B_\varphi `$ $`=`$ $`B_{\varphi \mathrm{s}},`$ (41) where $`i`$ is the inclination of the field lines to the vertical. The constancy of $`F`$ relies on there being little or no dissipation in the atmosphere, while the constancy of $`𝑩`$ relies on the atmosphere being magnetically dominated so that field line bending cannot be supported. That is, the plasma beta based on the poloidal magnetic field strength, $$\beta =\frac{2\mu _0p}{B_r^2+B_z^2},$$ (42) should be less than unity. In the atmosphere, all diffusive terms in the equations are neglected. If $`\psi _0/t0`$, there must be a slow drift orthogonal to the poloidal field lines, but otherwise the flow is constrained to follow the field. We know from Paper I that there will be a transonic outflow if $`i>30^{}`$, otherwise a modified hydrostatic atmosphere. As in Paper I, the angular velocity in the atmosphere is determined by isorotation, $$\mathrm{\Omega }_1=\mathrm{\Omega }_{1\mathrm{s}}+\frac{3\mathrm{\Omega }_0}{2r}(zH)\mathrm{tan}i.$$ (43) Again, this assumes that the atmosphere is magnetically dominated. The centrifugal-gravitational potential for matter constrained to follow the field is $$\mathrm{\Phi }^{\mathrm{cg}}=\frac{1}{2}(3\mathrm{tan}^2i1)\mathrm{\Omega }_0^2(zz_{\mathrm{sonic}})^2,$$ (44) where $$z_{\mathrm{sonic}}=H\left[1+(3\mathrm{tan}^2i1)^1\left(12\frac{\mathrm{\Omega }_{1\mathrm{s}}r}{\mathrm{\Omega }_0H}\mathrm{tan}i\right)\right]$$ (45) is the height of the sonic point in the case $`i>30^{}`$. The density scale-height $`h_\mathrm{s}`$ at the photosphere is given by $$\frac{c_\mathrm{s}^2}{h_\mathrm{s}}=\left(\frac{\mathrm{\Phi }^{\mathrm{cg}}}{z}\right)_\mathrm{s}=\mathrm{\Omega }_0^2H\left(12\frac{\mathrm{\Omega }_{1\mathrm{s}}r}{\mathrm{\Omega }_0H}\mathrm{tan}i\right),$$ (46) where $`c_\mathrm{s}`$ is the isothermal sound speed, given by $$c_\mathrm{s}^2=\frac{kT_\mathrm{s}}{\mu m_\mathrm{H}}.$$ (47) To a good approximation, the photospheric boundary condition (34) equates to $$\frac{2}{3}=C_\kappa \rho _\mathrm{s}^{1+x}T_\mathrm{s}^y_0^{\mathrm{}}\mathrm{exp}\left[(1+x)\left(\frac{zH}{h_\mathrm{s}}\right)\right]d(zH).$$ (48) This may be rearranged into the form $`\rho _\mathrm{s}^{1+x}T_\mathrm{s}^{1+y}`$ $`=`$ $`\left({\displaystyle \frac{1+x}{8\sigma }}\right)\left({\displaystyle \frac{\mu m_\mathrm{H}}{k}}\right)\left({\displaystyle \frac{16\sigma }{3C_\kappa }}\right)\mathrm{\Omega }_0^2H`$ (49) $`\times \left(12{\displaystyle \frac{\mathrm{\Omega }_{1\mathrm{s}}r}{\mathrm{\Omega }_0H}}\mathrm{tan}i\right).`$ Finally, in the case of a transonic outflow, the Mach number at the surface is determined from the equation (see Paper I) $$\frac{1}{2}(_\mathrm{s}^21)\mathrm{ln}_\mathrm{s}=\frac{\mathrm{\Delta }\mathrm{\Phi }}{c_\mathrm{s}^2},$$ (50) where $$\mathrm{\Delta }\mathrm{\Phi }=\frac{(\mathrm{\Omega }_0H2\mathrm{\Omega }_{1\mathrm{s}}r\mathrm{tan}i)^2}{2(3\mathrm{tan}^2i1)}$$ (51) is the potential barrier to the outflow. The vertical mass flux density in the outflow is then $$\dot{m}_\mathrm{w}=\rho u_z=_\mathrm{s}\rho _\mathrm{s}c_\mathrm{s}\mathrm{cos}i.$$ (52) ### 4.7. Non-dimensionalization We now rewrite the equations in a non-dimensional form suitable for numerical analysis. Given the local surface density $`\mathrm{\Sigma }`$, the angular velocity $`\mathrm{\Omega }_0`$, and the constants appearing in the constitutive relations, we identify $`U_H`$ $`=`$ $`\mathrm{\Sigma }^{(2+x)/(6+x2y)}\mathrm{\Omega }_0^{(52y)/(6+x2y)}`$ (53) $`\times \left({\displaystyle \frac{\mu m_\mathrm{H}}{k}}\right)^{(4y)/(6+x2y)}\left({\displaystyle \frac{16\sigma }{3C_\kappa }}\right)^{1/(6+x2y)}`$ as a natural unit for the semi-thickness of the disk. Natural units for other physical quantities follow according to $$U_\rho =\mathrm{\Sigma }U_H^1,U_p=\mathrm{\Sigma }\mathrm{\Omega }_0^2U_H,$$ (54) $$U_T=\mathrm{\Omega }_0^2\left(\frac{\mu m_\mathrm{H}}{k}\right)U_H^2,U_F=\mathrm{\Omega }_0U_HU_p,$$ (55) $$U_B=(\mu _0U_p)^{1/2}.$$ (56) Note that the above expression for $`U_H`$ can be obtained from the condition $$U_F=\frac{16\sigma U_T^3}{3C_\kappa U_\rho ^xU_T^yU_\rho }\frac{U_T}{U_H},$$ (57) which is a dimensional analysis of the definition (23) of the radiative flux. There are two small dimensionless parameters in the problem, $$ϵ=\frac{U_H}{r},\delta =\frac{U_F}{\sigma U_T^4}.$$ (58) Evidently $`ϵ`$ is a characteristic measure of the angular semi-thickness of the disk, while $`\delta `$ is an inverse measure of the total optical thickness. Non-dimensional variables are then introduced using the transformations $$z=z_{}U_H,H=H_{}U_H,$$ (59) $$\rho =\rho _{}U_\rho ,p=p_{}U_p$$ (60) $$T=T_{}U_T,F=F_{}U_F,$$ (61) $$\nu =\nu _{}\mathrm{\Omega }_0U_H^2,\eta =\eta _{}\mathrm{\Omega }_0U_H^2,$$ (62) $$u_r=u_r\mathrm{\Omega }_0U_H,\mathrm{\Omega }_1=\mathrm{\Omega }_1ϵ\mathrm{\Omega }_0,$$ (63) $$B_r=B_rU_B,B_\varphi =B_\varphi U_B,$$ (64) $$B_z=B_zU_B,\frac{\psi _0}{t}=\dot{\psi }_{}r\mathrm{\Omega }_0U_HU_B.$$ (65) The transformed equations are $$2\rho _{}\mathrm{\Omega }_1=B_zB_r^{},$$ (66) $$\frac{1}{2}\rho _{}u_r=B_zB_\varphi ^{}\frac{3}{4}ϵ(1+2D_{\nu \mathrm{\Sigma }})\rho _{}\nu _{}+\frac{3}{2}ϵD_H(\rho _{}\nu _{}z_{})^{},$$ (67) $$0=\rho _{}z_{}p_{}^{}B_rB_r^{}B_\varphi B_\varphi ^{},$$ (68) $$\dot{\psi }_{}+u_rB_z=\eta _{}\left(ϵD_BB_zB_r^{}\right),$$ (69) $$0=\frac{3}{2}B_r+B_z\mathrm{\Omega }_1^{}+(\eta _{}B_\varphi ^{})^{},$$ (70) $$F_{}^{}=\frac{9}{4}\rho _{}\nu _{}+\eta _{}\left[(B_r^{})^2+(B_\varphi ^{})^2\right],$$ (71) $$F_{}=\frac{T_{}^{3y}}{\rho _{}^{1+x}}T_{}^{},$$ (72) $$p_{}=\rho _{}T_{},$$ (73) where $$D_{\nu \mathrm{\Sigma }}=\frac{\mathrm{ln}(\overline{\nu }\mathrm{\Sigma })}{\mathrm{ln}r},D_H=\frac{\mathrm{ln}H}{\mathrm{ln}r},D_B=\frac{\mathrm{ln}B_z}{\mathrm{ln}r}$$ (74) are three dimensionless parameters, and the prime denotes differentiation with respect to $`z_{}`$. The dimensionless viscosity and magnetic diffusivity are given by $$\nu _{}=\mathrm{Pm}\eta _{}=\alpha T_{}.$$ (75) The mid-plane symmetry conditions are $$F_{}(0)=B_r(0)=B_\varphi (0)=0.$$ (76) The photospheric boundary conditions are $$\delta F_{}(H_{})=[T_{}(H_{})]^4$$ (77) and $$[\rho _{}(H_{})]^{1+x}[T_{}(H_{})]^{1+y}=\frac{1}{8}\delta (1+x)[H_{}2\mathrm{\Omega }_1(H_{})\mathrm{tan}i],$$ (78) together with $$B_r(H_{})=B_{r\mathrm{s}},B_\varphi (H_{})=B_{\varphi \mathrm{s}}.$$ (79) Finally, the definition of $`\mathrm{\Sigma }`$ implies the normalization condition $$2_0^H_{}\rho _{}𝑑z_{}=1.$$ (80) When the solution is found, the optical depth at the mid-plane can be determined from $$\tau _\mathrm{c}=\frac{2}{3}+_0^H\kappa \rho 𝑑z=\frac{2}{3}+\frac{16}{3\delta }_0^H_{}\rho _{}^{1+x}T_{}^y𝑑z_{}.$$ (81) The fact that the small parameter $`ϵ`$ cannot be fully scaled out of the equations indicates that a strictly self-consistent thin-disk asymptotic solution cannot be obtained when all the physical effects we have considered are included. The $`ϵ`$ terms in equation (67) are expected to be important only in the case of a weak magnetic field. The $`ϵ`$ term in equation (69) may be important when the field is nearly vertical. ### 4.8. Stability considerations Under the fixed alpha hypothesis the magnitudes of the turbulent viscosity and magnetic diffusivity are determined by the prescribed value of $`\alpha `$. A problem with this approach, as will be seen in the next section, is that equilibrium solutions can then be obtained that are magnetorotationally unstable even when the effect of the turbulent diffusivity on stability is taken into account. This suggests that the model may be physically inconsistent for these examples, because the ‘channel solution’ (Hawley & Balbus 1991) would continue to grow exponentially. Furthermore, there is a close connection between the stability of the equilibria and properties such as the shape of the field lines (Ogilvie 1998). We will find solutions that we believe to be unstable, in which the field lines bend more than once, and which have other undesirable properties. A resolution of these difficulties is suggested by the numerical simulations of magnetorotational turbulence, which indicate that the turbulence is much more vigorous in the presence of a mean poloidal magnetic field (Hawley, Gammie, & Balbus 1995), provided that the field is not so strong as to suppress the instability. Indeed, Stone et al. (1996) found it impossible to run a simulation in a stratified disk with a net vertical field, while Miller & Stone (2000) found a dramatic difference between simulations with no net field and those with a fairly weak uniform vertical field. In an effort to understand and model this behavior, we propose the following physical hypothesis for turbulent disks with a mean field: the value of $`\alpha `$ adjusts so that the equilibrium is marginally stable to the magnetorotational instability when the turbulent magnetic diffusivity is taken into account. We remark that a similar situation occurs in stellar convection, where an equivalent hypothesis can be used as a basis for the mixing-length theory. In a further example of this approach, Kippenhahn & Thomas (1978) modeled the outcome of a shear instability by imposing marginal stability according to the Richardson criterion. When the mean poloidal magnetic field is very weak, the instability favors small length scales and only a small value of $`\alpha `$ is required to suppress it. For stronger fields, the required value of $`\alpha `$ will be larger. For even stronger fields with $`\beta 1`$ on the mid-plane, the instability is suppressed or nearly so even without any diffusivity, and $`\alpha `$ will be small or zero. This variation of $`\alpha `$ with the strength of the mean poloidal field is qualitatively in agreement with that found in numerical simulations (Hawley, Gammie, & Balbus 1995; Brandenburg 1998), suggesting that our physical hypothesis is reasonable. However, our model will require this hypothesis to extend into a parameter range that cannot be (or at least has not been) reproduced in numerical simulations of stratified disks. Although a magnetorotationally unstable disk may contain numerous unstable modes, it has been argued by Ogilvie (1998) that, in ideal MHD, the last mode to be stabilized is an axisymmetric mode with vanishing radial wavenumber (i.e. $`/rr^1`$ rather than $`/rH^1`$), and is the first mode of odd symmetry about the mid-plane. Marginal stability of the annulus is then imposed by locating the marginal stability condition for this critical mode. We conjecture that this remains true when dissipation is included. From a consideration of the linearized equations in the limit of vanishing eigenfrequency, we find that such a mode should involve only horizontal motions and satisfy the horizontal components of the perturbed equation of motion, $$2\rho \mathrm{\Omega }_0\mathrm{\Delta }u_\varphi =\frac{B_z}{\mu _0}\frac{\mathrm{\Delta }B_r}{z},$$ (82) $$\frac{1}{2}\rho \mathrm{\Omega }_0\mathrm{\Delta }u_r=\frac{B_z}{\mu _0}\frac{\mathrm{\Delta }B_\varphi }{z},$$ (83) and the perturbed induction equation, $$0=B_z\frac{\mathrm{\Delta }u_r}{z}+\frac{}{z}\left(\eta \frac{\mathrm{\Delta }B_r}{z}\right),$$ (84) $$0=\frac{3}{2}\mathrm{\Omega }_0\mathrm{\Delta }B_r+B_z\frac{\mathrm{\Delta }u_\varphi }{z}+\frac{}{z}\left(\eta \frac{\mathrm{\Delta }B_\varphi }{z}\right),$$ (85) where $`\mathrm{\Delta }`$ denotes a linearized Eulerian perturbation. These equations should be compared with the unperturbed equations (17)–(18) and (20)–(21). Note that this problem is different from the marginal stability problem considered in Paper I, where the magnetic diffusivity was zero. The critical mode is expected to be the first mode of odd symmetry about the mid-plane (Ogilvie 1998). This should satisfy the symmetry conditions $`\mathrm{\Delta }u_r=\mathrm{\Delta }u_\varphi =0`$ on the mid-plane, and the photospheric boundary conditions $`\mathrm{\Delta }B_r=\mathrm{\Delta }B_\varphi =0`$ (since $`B_{r\mathrm{s}}`$ and $`B_{\varphi \mathrm{s}}`$ are fixed in the local model). For such a mode, the linearized equations may be combined into the second-order ODE $$\frac{}{z}\left[\left(\frac{B_z^2}{\mu _0\rho }+\frac{\mu _0\rho \eta ^2\mathrm{\Omega }_0^2}{B_z^2}\right)\frac{\mathrm{\Delta }B_r}{z}\right]+3\mathrm{\Omega }_0^2\mathrm{\Delta }B_r=0,$$ (86) or, in dimensionless form, $$\left[\left(\frac{B_z^2}{\rho _{}}+\frac{\rho _{}\eta _{}^2}{B_z^2}\right)\mathrm{\Delta }B_r^{}\right]^{}+3\mathrm{\Delta }B_r=0,$$ (87) with boundary conditions $$\mathrm{\Delta }B_r^{}(0)=\mathrm{\Delta }B_r(H_{})=0.$$ (88) We are now in a position to determine the value of $`\alpha `$ as follows. For any given value of $`\alpha `$ and suitable values of the other parameters, we might expect to find an equilibrium solution. However, the equilibrium will not in general possess a marginal mode satisfying the correct boundary conditions. By integrating the equations for a marginal mode simultaneously with those defining the equilibrium, we aim to tune $`\alpha `$ until a marginal mode is obtained. If this mode is the first mode of odd symmetry, we then believe we have determined the value of $`\alpha `$ corresponding to a marginally stable equilibrium. A crude estimate of this may be obtained by approximating the second derivative with respect to $`z_{}`$ as $`H_{}^2`$. This leads to the estimate $$\eta _{}^2\frac{B_z^2}{\rho _{}}\left(3H_{}^2\frac{B_z^2}{\rho _{}}\right),$$ (89) which has the properties described above. It might be argued that our marginal stability condition is too strong because the presence of one or two unstable modes of relatively long radial wavelength might not contribute significantly to turbulent transport. However, if the critical mode is not stabilized, it will lead to the growth of the fundamental ‘channel solution’, which is very efficient in transporting angular momentum and probably hydrodynamically unstable in three dimensions, leading to enhanced turbulence. Admittedly, the description of the disk under these circumstances is a matter of some uncertainty. ## 5. Numerical investigation ### 5.1. Numerical method The dimensionless ODEs are integrated from the photosphere $`z_{}=H_{}`$ to the mid-plane $`z_{}=0`$. The dependent variables are $`p_{}`$, $`T_{}`$, $`F_{}`$, $`\mathrm{\Omega }_1`$, $`B_r`$, and $`B_\varphi `$. The values of $`H_{}`$, $`F_{}(H_{})`$, $`\mathrm{\Omega }_1(H_{})`$, and $`\dot{\psi }_{}`$ are guessed and then adjusted by Newton-Raphson iteration to match the symmetry conditions on the mid-plane. In the marginal stability model, the equations for the marginal mode are integrated simultaneously. An arbitrary normalization $`\mathrm{\Delta }B_r^{}(H_{})=1`$ of the linear problem is applied, and the value of $`\alpha `$ is also guessed and adjusted by Newton-Raphson iteration to match the symmetry condition for the mode on the mid-plane. To obtain the derivatives $`B_\varphi ^{}`$ and $`\mathrm{\Omega }_1^{}`$ some algebra is required. Eliminating $`u_r`$, $`p_{}^{}`$, and $`B_r^{}`$ we obtain $`B_\varphi ^{}={\displaystyle \frac{\rho _{}}{2\mathrm{P}\mathrm{m}B_z^2(2B_z3ϵD_H\alpha B_\varphi z_{})}}`$ $`\times \{2\mathrm{P}\mathrm{m}B_z\dot{\psi }_{}+4\alpha p_{}\mathrm{\Omega }_1`$ $`+ϵ\left[2D_B+3\mathrm{P}\mathrm{m}(1+2D_{\nu \mathrm{\Sigma }}2D_H)\right]\alpha B_z^2T_{}`$ $`+6\mathrm{P}\mathrm{m}ϵD_H\alpha B_z(B_zz_{}^22B_r\mathrm{\Omega }_1z_{})\}.`$ (90) By multiplying this expression by $`\eta _{}`$ and differentiating, we find $`\mathrm{\Omega }_1^{}`$ from equation (70). ### 5.2. Unmagnetized solution In the absence of a mean magnetic field, a solution is obtained by omitting the induction equation and setting $`𝑩=\mathrm{𝟎}`$ elsewhere. This also implies $`\mathrm{\Omega }_1=0`$. The form of the solution for $`p_{}`$, $`T_{}`$, and $`F_{}`$ depends only on the dimensionless parameter $`\delta `$ (for a given opacity law), although the variables also have power-law scalings with $`\alpha `$. If $`u_r`$ is required, there are further dependences on $`D_{\nu \mathrm{\Sigma }}`$ and $`D_H`$, and the value of $`u_r`$ is proportional to $`ϵ`$ as a consequence of the scalings we have adopted for the general problem. We focus on the case of Thomson scattering opacity ($`x=y=0`$, $`C_\kappa 0.33\mathrm{cm}^2\mathrm{g}^1`$). In a steady disk, far from the inner edge, we may then set $`D_{\nu \mathrm{\Sigma }}=0`$ and $`D_H=21/20`$ (Shakura & Sunyaev 1973). We also take $`\alpha =0.1`$. For the purposes of illustration, we consider a location at 1000 Schwarzschild radii from a black hole of mass $`10M_{}`$, i.e. $`r=2.95\times 10^9\mathrm{cm}`$. For a surface density $`\mathrm{\Sigma }=10^4\mathrm{g}\mathrm{cm}^2`$, and assuming $`\mu =0.6`$, we find illustrative values $`U_H=6.34\times 10^7\mathrm{cm}`$, $`U_\rho =1.58\times 10^4\mathrm{g}\mathrm{cm}^3`$, $`U_p=3.27\times 10^{10}\mathrm{dyn}`$, $`U_T=1.50\times 10^6\mathrm{K}`$, $`U_F=4.70\times 10^{17}\mathrm{erg}\mathrm{cm}^2\mathrm{s}^1`$, and $`U_B=6.41\times 10^5\mathrm{G}`$. Then $`ϵ=2.15\times 10^2`$ and $`\delta =1.62\times 10^3`$. The numerically determined unmagnetized solution is shown in Fig. 1. It has a dimensionless photospheric height of $`H_{}=1.70`$ and an optical depth at the mid-plane of $`\tau _\mathrm{c}=\frac{2}{3}+\frac{8}{3}\delta ^1`$. Thus our illustrative values correspond to $`H/r=0.0364`$ and $`\tau _\mathrm{c}=1650`$, representing a geometrically thin and optically thick disk. The illustrative accretion rate is $`\dot{M}=3.55\times 10^{18}\mathrm{g}\mathrm{s}^1`$, which is less than the Eddington accretion rate $`1.69\times 10^{19}\mathrm{g}\mathrm{s}^1`$ for an accretion efficiency of $`0.1`$. The profile of radial velocity is of particular interest, since this will affect the advection of the magnetic field. The radial velocity is positive on the mid-plane and becomes negative at larger $`z`$, the density-weighted average being of course negative. This result has been found previously by several authors (e.g. Urpin 1984; Kley & Lin 1992, where further analysis of this phenomenon may be found). In the present model it results from the vertical profile of the viscosity, specifically from the $`\mathrm{ln}H/\mathrm{ln}r`$ term in equation (31). ### 5.3. Magnetized solutions under the fixed alpha hypothesis We first adopt the fixed alpha hypothesis and search for magnetized solutions, taking parameter values $`\alpha =0.1`$, $`B_{\varphi \mathrm{s}}=0`$, $`\mathrm{Pm}=1`$, $`D_B=0`$ in addition to $`x=y=0`$, $`D_{\nu \mathrm{\Sigma }}=0`$, $`D_H=21/20`$, and the illustrative values for $`ϵ`$ and $`\delta `$. There are two limits in which physically acceptable solutions are easily obtained over a wide range of inclination angles: very weak fields, and strong fields. We illustrate two such cases, with $`i=45^{}`$, in Figs 2 and 3. In the very weakly magnetized solutions, the magnetic field behaves kinematically and the equilibrium structure of the unmagnetized solution is not significantly distorted (e.g. the radial velocity profile is almost indistinguishable from that in the unmagnetized solution). The poloidal field lines bend smoothly according to a balance between advection and diffusion. Isorotation is not enforced (there is only a minuscule deviation from Keplerian rotation) and the field is strongly wound up ($`\|B_\varphi \|>\|B_z\|`$) in the example shown, with $`i=45^{}`$. The poloidal field lines continue to bend as the photosphere is approached. This is permitted because $`\beta `$ is very large throughout. Our atmospheric model, which assumed $`\beta <1`$, is not relevant here and certainly the outflow solution should be disregarded. It is not clear how to treat the atmosphere in a case such a this, because it is not understood whether the radial flow and magnetic diffusivity continue above the photosphere. By setting $`B_{\varphi \mathrm{s}}=0`$, we have ensured that the radial mass flux is caused by viscosity alone (cf. eq. ). As expected, the radial velocity causes the flux to migrate inwards for almost vertical fields, but outwards for larger inclination angles. With $`\mathrm{Pm}=1`$ and $`D_B=0`$, one finds the intermediate case $`\dot{\psi }_{}=0`$ at $`i=0.485=27.8^{}`$. This is much larger than $`H/r=0.0364`$, indicating that our concerns expressed in Section 3 were well founded. Alternatively, one can prevent the flux from being expelled (i.e. achieve $`\dot{\psi }_{}=0`$) when $`i=45^{}`$ by increasing $`\mathrm{Pm}`$ from $`1`$ to $`1.90`$. Increasing $`D_B`$ does not change $`\dot{\psi }_{}`$ significantly. For strong fields, we recover solutions very similar to those we obtained in Paper I. In the example shown in Fig. 3, the disk is significantly compressed by the Lorentz force (compare the plot of $`p_{}`$ with the unmagnetized case), and is hotter (cf. $`T_{}`$) and more luminous (cf. $`F_{}`$) than its unmagnetized counterpart. Most of the bending of the poloidal field lines occurs near the mid-plane where $`\beta `$ is close to unity. The field enforces isorotation, resulting in sub-Keplerian rotation, and very little toroidal field is produced. A feature of the solution is that the radial velocity, although everywhere subsonic, is rather large and non-uniform in direction. The reason for the profile of $`u_r`$ is that the flux is being expelled rapidly ($`\dot{\psi }_{}=0.551`$) by turbulent diffusion. In the upper layers, where $`\beta `$ is small, the fluid must follow the field and therefore flows rapidly outwards. To achieve the net accretion rate imposed by the vertically integrated viscous stress, a strong inflow is then required close to the mid-plane. Such a profile might even be considered advantageous for jet launching, since the outflow itself must involve a positive radial velocity. With $`B_z=2`$, as in this example, and assuming $`\mathrm{Pm}=1`$, zero flux migration is found at $`i=0.0346=1.98^{}`$. This is now close to $`H/r=0.0364`$ as expected from simple arguments, probably because the field-line bending occurs close to the mid-plane. On the other hand, these considerations may be inconsistent because the example in Fig. 3 is magnetorotationally stable and therefore unlikely to be turbulent. Of course, if there is no turbulent magnetic diffusivity, the origin of the viscosity should also be questioned. For a range of intermediate field strengths, the situation is more complicated. Solutions are not found in which the field lines bend in the simple way seen in Figs 2 and 3. Instead, branches of solutions appear in which the field lines bend several times as they pass through the disk. An example is shown in Fig. 4. According to Ogilvie (1998, Theorem 1) such a solution would be magnetorotationally unstable in ideal MHD. The multiple bending is closely related to the appearance of the ‘channel solution’ which is the first stage of the magnetorotational instability for a vertical field. We conjecture that such multiple-bending solutions are also unstable when the turbulent magnetic diffusivity is taken into account (in both the equilibrium and perturbation equations), although this has not been proven. If correct, this would imply that solutions of this type are physically inconsistent and ought to be disregarded, which in turn would suggest that no steady solution is possible for these intermediate field strengths. However, our alternative proposal, the marginal stability hypothesis, allows for a possible solution of this difficulty. ### 5.4. Magnetized solutions under the marginal stability hypothesis Under the marginal stability hypothesis, the strength of the turbulence (quantified through the parameter $`\alpha `$) is just sufficient to bring the equilibrium to marginal magnetorotational stability. This enables us to find single-bending solutions in a continuous range of field strengths from very small values up to the strength at which the instability is suppressed even without a turbulent diffusivity. We adopt the same parameters as in the previous section, except that $`\alpha `$ is now to be determined as part of the solution. Fig. 5 shows how the required value of $`\alpha `$ varies with $`B_z`$ for solutions with vertical fields ($`i=0`$). This has the general form expected from the crude estimate, equation (89). However, it is somewhat worrying that values of $`\alpha `$ in excess of unity may be required. It is often argued that $`\alpha `$ should not exceed unity because the turbulence would then have to be supersonic, or to have a correlation length greater than the disk thickness (e.g. Pringle 1981). It is not clear whether these constraints necessarily apply to magnetorotational turbulence, which is dominated by anisotropic Maxwell stresses whose correlation length in the azimuthal direction could exceed $`H`$ (e.g. Armitage 1998), but which could be limited instead by magnetic buoyancy. Moreover, the effective transport coefficients may themselves be anisotropic in the presence of a strong mean field (e.g. Matthaeus et al. 1998) and the magnetic Prandtl number may also differ from unity. Indeed, it is not certain whether the effects of the turbulence can be described adequately in terms of an effective diffusivity; it may be that coherent magnetic structures are formed. Nevertheless, the question remains as to whether the turbulence can truly reach a level sufficient to achieve marginal stability in the presence of a significant mean field. The numerical simulations by Hawley et al. (1995) are in qualitative agreement with Fig. 5 but suggest that, when the field strength is just below the stability boundary for a given computational domain, the channel solution may continue to grow without degenerating into turbulence. However, simulations of stratified models with fairly strong mean vertical fields do not appear to have been successful. An example solution calculated under the marginal stability hypothesis is shown in Fig. 6. This has $`\alpha =3.81`$ and should be contrasted with the solution in Fig. 4 which has the same vertical field strength but a fixed $`\alpha =0.1`$. The marginally stable solution has field lines with a single bend which become straight not far below the photosphere. The pressure and temperature decline smoothly and monotonically with increasing height. The field is not significantly wound up ($`\|B_\varphi \|<\|B_z\|`$). ### 5.5. Mass loss rates for jet-launching disks The principal aim of this paper was to determine how the mass loss rate in the outflow, $`\dot{m}_\mathrm{w}`$, varies with the strength and inclination of the magnetic field. In Fig. 7 we show the result of this calculation under the fixed alpha hypothesis with $`\alpha =0.1`$ and other parameters as given in Section 5.3. Note that the reciprocal of $`\dot{m}_\mathrm{w}/(\mathrm{\Sigma }\mathrm{\Omega })`$ is approximately the number of orbits in which the disk would be evaporated if not replenished by the accretion flow. In Fig. 8 we plot the quantity $`4\pi r^2\dot{m}_\mathrm{w}/\dot{M}`$ as a function of $`B_z`$ for solutions with $`i=45^{}`$. This alternative dimensionless measure of the outflow rate is the local mass loss rate per unit logarithmic interval in radius, divided by the accretion rate. The outcome is as we found in Paper I. The solutions shown are magnetorotationally stable. Close to the edge of the solution manifold, the outflow is very vigorous, but if the field strength is increased by only a factor of two above the stability boundary, the outflow is suppressed by twenty orders of magnitude or so. As explained in Paper I, this happens because the disk becomes significantly sub-Keplerian (cf. Fig. 3) and the outflow experiences a large potential barrier (whose strength is roughly proportional to $`B_z^4`$). This means that, in the absence of additional heating, external irradiation, or other driving mechanisms, outflows are suppressed from strongly magnetized disks. For a fixed field strength, the outflow is maximized at an intermediate inclination angle of 40–$`50^{}`$. The corresponding results for the marginal stability hypothesis are shown in Figs 9 and 10. Note that this is a complementary region of parameter space corresponding to turbulent disks. The behavior of the mass loss rate is now completely different and perhaps more intuitive: it increases monotonically with increasing $`B_z`$ and with increasing $`i`$. However, the solution branch terminates before excessive mass loss rates are achieved. This suggests that a large potential barrier is not incurred for turbulent disks and that they are more promising for jet launching. The potential barrier is smaller because the magnetic field is weaker, resulting in a smaller Lorentz force and a smaller deviation from Keplerian rotation. These two sets of results were obtained under different physical assumptions. Nevertheless, they both suggest that efficient jet-launching solutions are found in a limited range of field strengths, and in a limited range of inclination angles in excess of $`30^{}`$. In both cases there are difficulties in interpreting the solutions. The more strongly magnetized solutions obtained under the fixed alpha hypothesis are magnetorotationally stable, and the origin of the dissipation in the disk remains unclear. The more weakly magnetized solutions obtained under the marginal stability hypothesis typically require values of $`\alpha `$ in excess of unity to bring the equilibrium to marginal magnetorotational stability. ## 6. Discussion We have developed a model of the local vertical structure of magnetized accretion disks that launch magnetocentrifugal outflows. Given certain assumptions concerning the dissipative processes (turbulent viscosity and magnetic diffusivity) in the disk, we have shown that it is possible to compute the mass loss rate in the outflow as a function of the surface density and the strength and inclination of the poloidal magnetic field. The net accretion rates of mass and magnetic flux are also determined. This information is precisely complementary to that obtained from numerical simulations of the subsequent acceleration and collimation of a ‘cold’ outflow (e.g. Krasnopolsky et al. 1999). The following result appears to be quite robust. For disks in which the mean poloidal magnetic field is sufficiently strong to stabilize the equilibrium against the magnetorotational instability, we find that the mass loss rate decreases extremely rapidly with increasing field strength, and is maximized at an inclination angle of 40–$`50^{}`$. For turbulent disks with weaker mean fields, we find that the mass loss rate increases monotonically with increasing strength and inclination of the field, but the solution branch terminates before excessive mass loss rates are achieved. This suggests that turbulent disks with moderate mean fields are more promising for jet launching, but there may be situations in which a steady solution is impossible. For each solution we have determined the net rate of flux migration, which depends on a competition between inward dragging by the accretion flow, and outward transport due to turbulent diffusion. The results depend on the effective magnetic Prandtl number of the turbulence, which has never been measured. However, we find that inward migration is more likely to occur in the case of a weak mean field, which bends at greater heights in the disk. In this case, inward dragging may be many times more effective than estimated in previous studies (Lubow et al. 1994a), with the result that a net inward migration of flux may occur even when the inclination angle is sufficient for jet launching, provided that the disk is not very thin. This issue requires further analysis and a better understanding of the kinematic behavior of the magnetic field. For stronger fields, we find fairly good agreement with previous studies, and the jet-launching configurations are much more difficult to maintain against dissipation. We speculate that this may lead naturally to a situation in which the flux is regulated to a value suitable for jet launching. However, we have not modeled the possible effect of an internal dynamo in the disk. In addition, it is quite probable that instabilities will arise when the global coupling between the outflow and the flux evolution is considered (Lubow, Papaloizou, & Pringle 1994b; Cao & Spruit, in preparation). Several other authors have considered the vertical structure of magnetized disks and the problem of the disk-jet connection. Our formulation is distinctive in that we have shown how to calculate the mass loss rate and the accretion rates of mass and magnetic flux at any radius from a knowledge of the local surface density and the strength and inclination of the magnetic field. When coupled with a numerical simulation of the outflow beyond the sonic point, this would provide a closed evolutionary scheme for a time-dependent magnetized accretion disk with an outflow. It is crucial that we allowed the external magnetic torque on the disk to be determined by the exterior outflow solution rather than the interior solution of the vertical disk structure. There are further differences in the detail of our approach. Wardle & Königl (1993) discussed many of the issues that we have addressed, but we differ from them in considering an optically thick disk with turbulent transport and energy dissipation, as opposed to an isothermal, inviscid disk with ambipolar diffusion. We are also able to calculate the net rate of flux accretion rather than specifying it as a parameter. Similar comparisons may be made with the analysis of Li (1995). The model of Campbell (1999) shares some features with the present paper, but his outflow solution treats the magnetic field lines as parabolae throughout the disc and the transonic region, while our solutions (e.g. Fig. 6) indicate that the field lines become straight in the atmosphere, provided that the plasma beta is less than unity. If $`\beta >1`$ in the atmosphere, the magnetocentrifugal mechanism does not work. Interesting comparisons may be made with the steady, self-similar model of Casse & Ferreira (2000). In common with those authors, we find that values of $`\alpha `$ exceeding unity, and/or anisotropic effective turbulent diffusion coefficients, may be helpful or necessary in obtaining plausible solutions for efficient jet-launching disks. Admittedly, our model has some limitations. Our simple modeling of the effects of the turbulence is consistent with the very limited information available from numerical simulations, although we have not included a dynamo $`\alpha `$-effect. We have assumed that the atmosphere is magnetically dominated so that the poloidal field lines act as rigid channels for the outflow. This constraint could be relaxed in future work at the expense of introducing considerable complications in solving for the atmospheric flow, and in matching to the exterior solution. Future work should aim at obtaining a better understanding of the interaction of the mean magnetic field with the turbulence. Numerical simulations might be used to measure the effective magnetic diffusivity tensor and how it depends on the strength of the field. A crucial issue is whether, when the mean field is a significant fraction of the value required for magnetorotational stability, the ‘channel solution’ of the instability persists or degenerates into turbulence. Ideally the whole problem that we have defined, including the correct boundary conditions, would be solved using a numerical simulation. We have included some non-local effects (Section 4.4) but most of the relevant terms are included in the shearing-box local accretion disk model (e.g. Stone et al. 1996). It may be difficult, however, to resolve the disk adequately from the mid-plane to the sonic point, although the techniques used by Miller & Stone (2000) should be helpful in achieving this. GIO acknowledges support from the STScI visitor program and from Clare College, Cambridge. ML acknowledges support from NASA Grants NAG5-6857 and GO-07378. We thank Jim Pringle and Henk Spruit for helpful discussions.
warning/0007/math0007009.html	ar5iv	text	# Multiple categories: the equivalence of a globular and a cubical approach ## Introduction An essential feature for the possibility of ‘higher dimensional group theory’ (see the expository article ) is the extension of the domain of discourse from groups to groupoids, that is from a set with a binary operation defined on all elements, to a set with an operation defined only on pairs satisfying a geometric condition. This fact itself leads to various equivalent candidates for ‘higher dimensional groups’, namely those based on different geometric structures, for example balls, globes, simplices, cubes and even polyhedra. The proofs of these equivalences are non trivial – the basic intuitions derive from the foundations of relative homotopy theory. Some of these equivalences have proved crucial for the applications: theorems may be easily proved in one context and then transferred into another, more computational context. Notable examples are the advantages of cubical methods for providing both a convenient ‘algebraic inverse to subdivision’, for use in local-to-global problems , and also a simple monoidal closed structure, which may then be translated into other situations . It has proved important to extend these ideas from groupoids to categories. The standard notion of (strict) higher dimensional category is that of globular $`\omega `$-category. Our main result is that there is an adjoint equivalence of categories $$\lambda :\text{globular }\omega \text{-categories}\text{}\text{cubical }\omega \text{-categories with connections }:\gamma .$$ Precise definitions are given below. The proof has interest because it is certainly much harder than the groupoid case, and because at one stage it uses braid relations among some key basic folding operations (Proposition 5.1, Theorem 5.2). The equivalence between the two forms should prove useful. In section 9 we use this equivalence to define the notion of ‘commutative $`n`$-cube’. In section 10 we follow methods of Brown and Higgins in to show that cubical $`\omega `$-categories with connection form a monoidal closed category. The equivalence of categories transfers this structure to the globular case – the resulting internal hom in the globular case gives various higher dimensional forms of ‘lax natural transformation’. Cubical $`\omega `$-categories with connection have been applied to concurrency theory by E. Goubault and by P. Gaucher , and again relations with the globular case are important for these studies. The origin of this equivalence is as follows. In developing the algebra of double groupoids as a framework for potential 2-dimensional Van Kampen Theorems, Brown and Spencer in were led to the notion of double groupoid with an extra structure of ‘connection’ – this was essential to obtain an equivalence of such a double groupoid with the classical notion of crossed module. This structure was also essential for the proof of the 2-dimensional Van Kampen Theorem given by Brown and Higgins in . The double groupoid case was generalised by Brown and Higgins to give an equivalence between crossed complexes and what were called there ‘$`\omega `$-groupoids’, and which we here call ‘cubical $`\omega `$-groupoids with connections’. It was also proved in that crossed complexes are equivalent to what were there called ‘$`\mathrm{}`$-groupoids’, and which we here call ‘globular $`\omega `$-groupoids’, following current fashions. Thus the globular and cubical cases of $`\omega `$-groupoids were known in 1981 to be equivalent, but the proof was via the category of crossed complexes. Other equivalences with crossed complexes were established, for example with: cubical $`T`$-complexes ; simplicial $`T`$-complexes by Nick Ashley ; and polyhedral $`T`$-complexes by David Jones . In $`T`$-complexes the basic concept is taken to be that of thin elements which determine a strengthening of the Kan extension condition. The notion of simplicial $`T`$-complex is due to Keith Dakin . Spencer observed in that the methods of allowed an equivalence between 2-categories and double categories with connections, using an ‘up-square’ construction of Bastiani and Ehresmann , but he gave no details. The full details of this have been recently given by Brown and Mosa in . The thesis of Mosa in 1987 attempted to give an equivalence between crossed complexes of algebroids and cubical $`\omega `$-algebroids, and while this was completed in dimension 2 even the case of dimension 3 proved hard, though some basic methods were established. This result raised the question of an equivalence between the globular $`\omega `$-categories defined in 1981 in and an appropriate form of cubical $`\omega `$-categories with connections, of which a definition was fairly easy to formulate as an extension of the previous definition of cubical $`\omega `$-groupoid. This problem was taken up in Al-Agl’s thesis of 1989 . The central idea, based on the groupoid methods of , was to define a ‘folding operation’ $`\mathrm{\Phi }`$ from a cubical $`\omega `$-category $`G`$ to the globular $`\omega `$-category $`\gamma G`$ it contained. This definition was successfully accomplished, but the problem of establishing some major properties of $`\mathrm{\Phi }`$, in particular the relation with the category structures, was solved only up to dimension 3. That is, the conjectured equivalence was proved in dimension 3. Steiner pursued the work of Al-Agl, and their joint paper does prove that globular $`\omega `$-categories are equivalent to cubical sets with extra structure, but, as stated in that paper, this extra structure is not described in finitary terms. Later, Steiner was stimulated by renewed interest in the cubical case coming from concurrency theory, in work of E. Goubault and P. Gaucher , and by the publication of the 2-dimensional case by Brown and Mosa in . He completed the programme given in and informed Brown, who announced the result at the Aalborg ‘Workshop on Geometric and Topological Methods in Concurrency’ in June, 1999. This paper is the result. It proves the conjecture implicit in , that a globular $`\omega `$-category is equivalent to a cubical set with extra structure directly analogous to the structure for cubical $`\omega `$-groupoids given in . There is considerable independent work on globular $`\omega `$-categories. The thesis of Sjoerd Crans already contains the adjoint pair $`(\lambda ,\gamma )`$ and also the closed monoidal structure on the category of globular $`\omega `$-categories. It also seems to be the first time that that the cube category (without connections) together with its $`\omega `$-category realisation is explicitly defined by generators and relations. The work in Australia by Ross Street has an initial aim to determine a simplicial nerve $`NX`$ of a globular $`\omega `$-category $`X`$. This developed into finding extra structure on $`NX`$ so that $`N`$ gave an equivalence between $`\omega `$-categories and certain structured simplicial sets, analogous to Ashley’s equivalence between $`\omega `$-groupoids and simplicial $`T`$-complexes. It is stated in that this programme has been completed by Dominic Verity, to verify the conjecture stated in . Street tells us that Verity also knew the equivalence proved in the present paper, but we have no further information. We also mention that Street’s paper implicitly contains our basic proposition (3.2), namely that the cells of the $`n`$-categorical $`n`$-cube compose in such a way that they give rise to the hemispherical (i.e. globular) decomposition $`_1^\pm \mathrm{\Phi }_n`$ of the $`n`$-cube. ## 1 $`\omega `$-categories An $`\omega `$-category arises when a sequence of categories $`C_0,C_1,\mathrm{}`$ all have the same set of morphisms $`X`$, the various category structures commute with one another, the identities for $`C_p`$ are also identities for $`C_q`$ when $`q>p`$, and every member of $`X`$ is an identity for some $`C_p`$. We write $`\#_p`$ for the composition in $`C_p`$. Given $`xX`$, we write $`d_p^{}x`$ and $`d_p^+x`$ for the identities of the source and target of $`x`$ in $`C_p`$, so that $`d_p^{}x\#_px=x\#_pd_p^+x=x`$. The structure can be expressed in terms of $`X`$, $`\#_p`$ and the $`d_p^\alpha `$ as follows. ###### Definition 1.1 An $`\omega `$-category is a set $`X`$ together with unary operations $`d_p^{}`$$`d_p^+`$ and partially defined binary operations $`\#_p`$ for $`p=0`$$`1`$, … such that the following conditions hold: 1. $`x\#_py`$ is defined if and only if $`d_p^+x=d_p^{}y`$; 2. $`d_q^\beta d_p^\alpha x=\{\begin{array}{cc}d_q^\beta x\hfill & \text{for }q<p,\hfill \\ d_p^\alpha x\hfill & \text{for }qp;\hfill \end{array}`$ 3. if $`x\#_py`$ is defined then $`d_p^{}(x\#_py)`$ $`=d_p^{}x,`$ $`d_p^+(x\#_py)`$ $`=d_p^+y,`$ $`d_q^\beta (x\#_py)`$ $`=d_q^\beta x\#_pd_q^\beta y\text{for }qp;`$ 4. $`d_p^{}x\#_px=x\#_pd_p^+x=x`$; 5. $`(x\#_py)\#_pz=x\#_p(y\#_pz)`$ if either side is defined; 6. if $`pq`$ then $$(x\#_py)\#_q(x^{}\#_py^{})=(x\#_qx^{})\#_p(y\#_qy^{})$$ whenever both sides are defined; 7. for each $`xX`$ there is a dimension $`dimx`$ such that $`d_p^\alpha x=x`$ if and only if $`pdimx`$. ###### Definition 1.2 An $`\omega `$-category of sets is an $`\omega `$-category $`X`$ whose members are sets such that $`x\#_py=xy`$ whenever $`x\#_py`$ is defined in $`X`$. The theory of pasting in $`\omega `$-categories associates $`\omega `$-categories of sets $`M(K)`$ with simple presentations to certain complexes $`K`$; the members of $`M(K)`$ are subcomplexes of $`K`$. Various types of complexes have been considered, but they certainly include the cartesian products of directed paths, and we will now describe the theory in that case. Let $`n`$ be a non-negative integer. We represent a directed path of length $`n`$ by the closed interval $`[0,n]`$; the vertices are the singleton subsets $`\{0\}`$, $`\{1\}`$, …, $`\{n\}`$ and the edges are the intervals $`[0,1]`$, $`[1,2]`$, …, $`[n1,n]`$, where $`[m1,m]`$ is directed from $`m1`$ to $`m`$. We write $$d^{}[m1,m]=\{m1\},d^+[m1,m]=\{m\}.$$ Now let $`K=K_1\times \mathrm{}\times K_p`$ be a cartesian product of directed paths. A product $`\sigma =\sigma _1\times \mathrm{}\times \sigma _p`$, where $`\sigma _i`$ is a vertex or edge in $`K_i`$, is called a cell in $`K`$. We can write a cell $`\sigma `$ in the form $$\sigma =P_0\times e_1\times P_1\times e_2\times P_2\times \mathrm{}\times P_{q1}\times e_q\times P_q,$$ where the $`P_j`$ are products of vertices and the $`e_j`$ are edges; the dimension of $`\sigma `$ is then $`q`$. The codimension $`1`$ faces of $`\sigma `$ are the subsets got by replacing one edge factor $`e_j`$ with $`d^{}e_j`$ or $`d^+e_j`$. The faces with $`d^{}e_1`$ or $`d^+e_2`$ or $`d^{}e_3`$ or … are called negative, and the faces with $`d^+e_1`$ or $`d^{}e_2`$ or $`d^+e_3`$ or … are called positive. The theory of pasting gives us the following result. ###### Theorem 1.3 Let $`K`$ be a cartesian product of directed paths. Then there is an $`\omega `$-category $`M(K)`$ of subsets of $`K`$ with the following presentation: the generators are the cells of $`K`$; if $`\sigma `$ is a cell of dimension $`q`$ then there are relations $`d_q^{}\sigma =d_q^+\sigma =\sigma `$; if $`\sigma `$ is a cell of dimension $`q`$ with $`q>0`$ then there are relations saying that $`d_{q1}^{}\sigma `$ and $`d_{q1}^+\sigma `$ are the unions of the negative and positive faces of $`\sigma `$ respectively. Every member of $`M(K)`$ is an iterated composite of cells. We will now describe the main examples. ###### Example 1.4 We write $`I=[0,1]`$ and $`I^n=[0,1]^n`$ for $`n1`$; for completeness we also write $`I^0=[0,0]`$. In this notation $`M(I^0)=\{I^0\}`$ and $`M(I)=\{I,d_0^{}I,d_0^+I\}`$; there are no members other than the generating cells. There are morphisms $$\stackrel{ˇ}{}^{},\stackrel{ˇ}{}^+:M(I^0)M(I),\stackrel{ˇ}{\epsilon }:M(I)M(I^0)$$ given by $$\stackrel{ˇ}{}^\alpha (I^0)=d_0^\alpha I,$$ $$\stackrel{ˇ}{\epsilon }(I)=\stackrel{ˇ}{\epsilon }(d_0^\alpha I)=I^0.$$ ###### Example 1.5 The members of $`M([0,2])`$ are the cells and the composite $$[0,2]=[0,1]\#_0[1,2].$$ There are morphisms $$\stackrel{ˇ}{\iota }^{},\stackrel{ˇ}{\iota }^+,\stackrel{ˇ}{\mu }:M(I)M([0,2])$$ given by $$\stackrel{ˇ}{\iota }^{}(d_0^{}I)=\{0\},\stackrel{ˇ}{\iota }^{}(I)=[0,1],\stackrel{ˇ}{\iota }^{}(d_0^+I)=\{1\},$$ $$\stackrel{ˇ}{\iota }^+(d_0^{}I)=\{1\},\stackrel{ˇ}{\iota }^+(I)=[1,2],\stackrel{ˇ}{\iota }^+(d_0^+I)=\{2\},$$ and $$\stackrel{ˇ}{\mu }(d_0^{}I)=\{0\},\stackrel{ˇ}{\mu }(I)=[0,2],\stackrel{ˇ}{\mu }(d_0^+I)=\{2\}.$$ ###### Example 1.6 The members of $`M(I^2)`$ are the cells and the composites $$d_1^{}I^2=(d_0^{}I\times I)\#_0(I\times d_0^+I),d_1^+I^2=(I\times d_0^{}I)\#_0(d_0^+I\times I).$$ There are morphisms $`\stackrel{ˇ}{\mathrm{\Gamma }}^+,\stackrel{ˇ}{\mathrm{\Gamma }}^{}:M(I^2)M(I)`$ given by $$\stackrel{ˇ}{\mathrm{\Gamma }}^\alpha (d_0^\alpha I\times d_0^\alpha I)=\stackrel{ˇ}{\mathrm{\Gamma }}^\alpha (d_0^\alpha I\times I)=\stackrel{ˇ}{\mathrm{\Gamma }}^\alpha (d_0^\alpha I\times d_0^\alpha I)=\stackrel{ˇ}{\mathrm{\Gamma }}^\alpha (I\times d_0^\alpha I)$$ $$=\stackrel{ˇ}{\mathrm{\Gamma }}^\alpha (d_0^\alpha I\times d_0^\alpha I)=d_0^\alpha I,$$ $$\stackrel{ˇ}{\mathrm{\Gamma }}^\alpha (I^2)=\stackrel{ˇ}{\mathrm{\Gamma }}^\alpha (I\times d_0^\alpha I)=\stackrel{ˇ}{\mathrm{\Gamma }}^\alpha (d_0^\alpha I\times I)=\stackrel{ˇ}{\mathrm{\Gamma }}^\alpha (d_1^{}I^2)=\stackrel{ˇ}{\mathrm{\Gamma }}^\alpha (d_1^+I^2)=I,$$ $$\stackrel{ˇ}{\mathrm{\Gamma }}^\alpha (d_0^\alpha I\times d_0^\alpha I)=d_0^\alpha I.$$ For cartesian products of members of the $`\omega `$-categories that we are considering, we have the following result. ###### Theorem 1.7 Let $`K`$ and $`L`$ be cartesian products of directed paths, let $`x`$ be a member of $`M(K)`$, and let $`y`$ be a member of $`M(L)`$. Then $`x\times y`$ is a member of $`M(K\times L)`$ and $$d_p^\alpha (x\times y)=\underset{i=0}{\overset{p}{}}\left(d_i^\alpha x\times d_{pi}^{()^i\alpha }y\right).$$ This has the following consequence. ###### Theorem 1.8 (i) Let $`K`$, $`K^{}`$, $`L`$, $`L^{}`$ be cartesian products of directed graphs, and let $`f:M(K)M(K^{})`$ and $`g:M(L)M(L^{})`$ be morphisms of $`\omega `$-categories. Then there is a unique morphism $$fg:M(K\times L)M(K^{}\times L^{})$$ such that $$(fg)(x\times y)=f(x)\times g(y)$$ for $`xM(K)`$ and $`yM(L)`$. (ii) The assignments $$(M(K),M(L))M(K\times L),(f,g)fg$$ form a bifunctor. Proof (i) From the presentation of $`M(K\times L)`$ and Theorem 1.7, there is a unique morphism $`fg`$ such that $`(fg)(x\times y)=f(x)\times g(y)`$ when $`x`$ and $`y`$ are cells. The formula then holds for a general product $`x\times y`$ because it is a composite of cells. (ii) One can check bifunctoriality by considering the values of the appropriate morphisms on generators. $`\mathrm{}`$ By applying the tensor product construction, we obtain further morphisms. ###### Example 1.9 Let $`id^r`$ denote the identity morphism from $`M(I^r)`$ to itself. There are morphisms $$\stackrel{ˇ}{}_i^{},\stackrel{ˇ}{}_i^+:M(I^{n1})M(I^n)(1in)$$ given by $$\stackrel{ˇ}{}_i^\alpha =id^{i1}\stackrel{ˇ}{}^\alpha id^{ni};$$ there are morphisms $$\stackrel{ˇ}{\epsilon }_i:M(I^n)M(I^{n1})(1in)$$ given by $$\stackrel{ˇ}{\epsilon }_i=id^{i1}\stackrel{ˇ}{\epsilon }id^{ni};$$ there are morphisms $$\stackrel{ˇ}{\iota }_i^{},\stackrel{ˇ}{\iota }_i^+,\stackrel{ˇ}{\mu }_i:M(I^n)M(I^{i1}\times [0,2]\times I^{ni})(1in)$$ given by $$\stackrel{ˇ}{\iota }_i^\alpha =id^{i1}\stackrel{ˇ}{\iota }^\alpha id^{ni},\stackrel{ˇ}{\mu }_i=id^{i1}\stackrel{ˇ}{\mu }id^{ni};$$ there are morphisms $$\stackrel{ˇ}{\mathrm{\Gamma }}_i^+,\stackrel{ˇ}{\mathrm{\Gamma }}_i^{}:M(I^n)M(I^{n1})(1in1)$$ given by $$\stackrel{ˇ}{\mathrm{\Gamma }}_i^\alpha =id^{i1}\stackrel{ˇ}{\mathrm{\Gamma }}^\alpha id^{ni1}.$$ Most of the morphisms in Example 1.9 map generators to generators, and one can verify their existence directly from Theorem 1.3. The exceptions are the $`\stackrel{ˇ}{\mu }_i`$, for which Theorem 1.8 is really necessary. ###### Remark 1.10 Suppose that $`K`$ is an $`n`$-dimensional product of directed paths. Then $`K`$ can be got from a family of $`n`$-cubes by gluing along $`(n1)`$-dimensional faces. From the presentation of $`M(K)`$, one sees that it is the colimit of a corresponding diagram in which the morphisms have the form $`\stackrel{ˇ}{}_i^\alpha :M(I^{n1})M(I^n)`$. In particular, $`\stackrel{ˇ}{\iota }^{}`$ and $`\stackrel{ˇ}{\iota }^+`$ exhibit $`M([0,2])`$ as the push-out of and $`\stackrel{ˇ}{\iota }_i^{}`$ and $`\stackrel{ˇ}{\iota }_i^+`$ exhibit $`M(I^{i1}\times [0,2]\times I^{ni})`$ as the push-out of ## 2 Cubical $`\omega `$-categories with connections Suppose that $`X`$ is an $`\omega `$-category. There is then a sequence of sets $$(\lambda X)_n=Hom[M(I^n),X](n=0,1,\mathrm{}),$$ and the morphisms of Example 1.9 induce functions between the $`(\lambda X)_n`$. It turns out that the $`(\lambda X)_n`$ form a cubical $`\omega `$-category with connections in the sense of the following definition. This definition is found in . The origin is in the definition of what was called ‘$`\omega `$-groupoid’ in , where the justification was the equivalence with crossed complexes (loc. cit.) and the use in the formulation and proof of a generalised Van Kampen Theorem . The corresponding definition for categories arose out of the work of Spencer and of Mosa . Let $`K`$ be a cubical set, that is, a family of sets $`\{K_n;n0\}`$ with for $`n1`$ face maps $`_i^\alpha :K_nK_{n1}(i=1,2,\mathrm{},n;\alpha =+,)`$ and degeneracy maps $`\epsilon _i:K_{n1}K_n(i=1,2,\mathrm{},n)`$ satisfying the usual cubical relations: $`_i^\alpha _j^\beta `$ $`=_{j1}^\beta _i^\alpha `$ $`(i<j),`$ (2.1)(i) $`\epsilon _i\epsilon _j`$ $`=\epsilon _{j+1}\epsilon _i`$ $`(ij),`$ (2.1)(ii) $`_i^\alpha \epsilon _j`$ $`=\{\begin{array}{cc}\epsilon _{j1}_i^\alpha \hfill & (i<j)\hfill \\ \epsilon _j_{i1}^\alpha \hfill & (i>j)\hfill \\ \mathrm{id}\hfill & (i=j)\hfill \end{array}`$ (2.1)(iii) We say that $`K`$ is a cubical set with connections if for $`n0`$ it has additional structure maps (called connections) $`\mathrm{\Gamma }_i^+,\mathrm{\Gamma }_i^{}:K_nK_{n+1}(i=1,2,\mathrm{},n)`$ satisfying the relations: $`\mathrm{\Gamma }_i^\alpha \mathrm{\Gamma }_j^\beta `$ $`=\mathrm{\Gamma }_{j+1}^\beta \mathrm{\Gamma }_i^\alpha `$ $`(i<j)`$ (2.2)(i) $`\mathrm{\Gamma }_i^\alpha \mathrm{\Gamma }_i^\alpha `$ $`=\mathrm{\Gamma }_{i+1}^\alpha \mathrm{\Gamma }_i^\alpha `$ (2.2)(ii) $`\mathrm{\Gamma }_i^\alpha \epsilon _j`$ $`=\{\begin{array}{cc}\epsilon _{j+1}\mathrm{\Gamma }_i^\alpha \hfill & (i<j)\hfill \\ \epsilon _j\mathrm{\Gamma }_{i1}^\alpha \hfill & (i>j)\hfill \end{array}`$ (2.2)(iii) $`\mathrm{\Gamma }_j^\alpha \epsilon _j`$ $`=\epsilon _j^2=\epsilon _{j+1}\epsilon _j,`$ (2.2)(iv) $`_i^\alpha \mathrm{\Gamma }_j^\beta `$ $`=\{\begin{array}{cc}\mathrm{\Gamma }_{j1}^\beta _i^\alpha \hfill & (i<j)\hfill \\ \mathrm{\Gamma }_j^\beta _{i1}^\alpha \hfill & (i>j+1),\hfill \end{array}`$ (2.2)(v) $`_j^\alpha \mathrm{\Gamma }_j^\alpha `$ $`=_{j+1}^\alpha \mathrm{\Gamma }_j^\alpha =id,`$ (2.2)(vi) $`_j^\alpha \mathrm{\Gamma }_j^\alpha `$ $`=_{j+1}^\alpha \mathrm{\Gamma }_j^\alpha =\epsilon _j_j^\alpha .`$ (2.2)(vii) The connections are to be thought of as extra ‘degeneracies’. (A degenerate cube of type $`\epsilon _jx`$ has a pair of opposite faces equal and all other faces degenerate. A cube of type $`\mathrm{\Gamma }_i^\alpha x`$ has a pair of adjacent faces equal and all other faces of type $`\mathrm{\Gamma }_j^\alpha y`$ or $`\epsilon _jy`$ .) Cubical complexes with this, and other, structures have also been considered by Evrard . The prime example of a cubical set with connections is the singular cubical complex $`KX`$ of a space $`X`$. Here for $`n0`$ $`K_n`$ is the set of singular $`n`$-cubes in $`X`$ (i.e. continuous maps $`I^nX`$) and the connection $`\mathrm{\Gamma }_i^\alpha :K_nK_{n+1}`$ is induced by the map $`\gamma _i^\alpha :I^{n+1}I^n`$ defined by $$\gamma _i^\alpha (t_1,t_2,\mathrm{},t_{n+1})=(t_1,t_2,\mathrm{},t_{i1},A(t_i,t_{i+1}),t_{i+2},\mathrm{},t_{n+1})$$ where $`A(s,t)=\mathrm{max}(s,t),\mathrm{min}(s,t)`$ as $`\alpha =,+`$ respectively. Here are pictures of $`\gamma _1^\alpha :I^2I^1`$ where the internal lines show lines of constancy of the map on $`I^2`$. The complex $`KX`$ has some further relevant structure, namely the composition of $`n`$-cubes in the $`n`$ different directions. Accordingly, we define a cubical complex with connections and compositions to be a cubical set $`K`$ with connections in which each $`K_n`$ has $`n`$ partial compositions $`_j(j=1,2,\mathrm{},n)`$ satisfying the following axioms. If $`a,bK_n`$, then $`a_jb`$ is defined if and only if $`_j^{}b=_j^+a`$ , and then $$\{\begin{array}{cc}_j^{}(a_jb)=_j^{}a\hfill & \\ _j^+(a_jb)=_j^+b\hfill & \end{array}_i^\alpha (a_jb)=\{\begin{array}{cc}_j^\alpha a_{j1}_i^\alpha b\hfill & (i<j)\hfill \\ _i^\alpha a_j_i^\alpha b\hfill & (i>j),\hfill \end{array}$$ (2.3) The interchange laws. If $`ij`$ then $$(a_ib)_j(c_id)=(a_jc)_i(b_jd)$$ (2.4) whenever both sides are defined. (The diagram $$\left[\begin{array}{cc}a& b\\ c& d\end{array}\right]\text{}$$ will be used to indicate that both sides of the above equation are defined and also to denote the unique composite of the four elements.) If $`ij`$ then $`\epsilon _i(a_jb)`$ $`=\{\begin{array}{cc}\epsilon _ia_{j+1}\epsilon _ib\hfill & (ij)\hfill \\ \epsilon _ia_j\epsilon _ib\hfill & (i>j)\hfill \end{array}`$ (2.5) $`\mathrm{\Gamma }_i^\alpha (a_jb)`$ $`=\{\begin{array}{cc}\mathrm{\Gamma }_i^\alpha a_{j+1}\mathrm{\Gamma }_i^\alpha b\hfill & (i<j)\hfill \\ \mathrm{\Gamma }_i^\alpha a_j\mathrm{\Gamma }_i^\alpha b\hfill & (i>j)\hfill \end{array}`$ (2.6)(i) $`\mathrm{\Gamma }_j^+(a_jb)`$ $`=\left[\begin{array}{cc}\mathrm{\Gamma }_j^+a& \epsilon _ja\\ \epsilon _{j+1}a& \mathrm{\Gamma }_j^+b\end{array}\right]\text{}`$ (2.6)(ii) $`\mathrm{\Gamma }_j^{}(a_jb)`$ $`=\left[\begin{array}{cc}\mathrm{\Gamma }_j^{}a& \epsilon _{j+1}b\\ \epsilon _jb& \mathrm{\Gamma }_j^{}b\end{array}\right]\text{}`$ (2.6)(iii) These last two equations are the transport laws<sup>2</sup><sup>2</sup>2Recall from that the term connection was chosen because of an analogy with path-connections in differential geometry. In particular, the transport law is a variation or special case of the transport law for a path-connection. . It is easily verified that the singular cubical complex $`KX`$ of a space $`X`$ satisfies these axioms if $`_j`$ is defined by $$(a_jb)(t_1,t_2,\mathrm{},t_n)=\{\begin{array}{cc}a(t_1,\mathrm{},t_{j1},2t_j,t_{j+1},\mathrm{},t_n)\hfill & (t_j\frac{1}{2})\hfill \\ b(t_1,\mathrm{},t_{j1},2t_j1,t_{j+1},\mathrm{},t_n)\hfill & (t_j\frac{1}{2})\hfill \end{array}$$ whenever $`_j^{}b=_j^+a`$. In this context the transport law for $`\mathrm{\Gamma }_1^{}(ab)`$ can be illustrated by the picture ###### Definition 2.1 A cubical $`\omega `$-category with connections $`G=\{G_n\}`$ is a cubical set with connections and compositions such that each $`_j`$ is a category structure on $`G_n`$ with identity elements $`\epsilon _jy(yG_{n1})`$, and in addition $$\mathrm{\Gamma }_i^+x_i\mathrm{\Gamma }_i^{}x=\epsilon _{i+1}x,\mathrm{\Gamma }_i^+x_{i+1}\mathrm{\Gamma }_i^{}x=\epsilon _ix.$$ (2.7) For simplicity, a cubical $`\omega `$-category with connections will be called a cubical $`\omega `$-category in the rest of this paper. ###### Remark 2.2 This list is a part of the list of structure and axioms which first appears in the thesis of Mosa \[23, Chapter V\], in the context of cubical algebroids with connection, and appears again in the thesis of Al-Agl . The rules for the connections are fairly clear extensions of the axioms given in , given the general notion of thin structure on a double category discussed by Spencer in . Note that a cubical $`\omega `$-category has an underlying cubical set under its face and degeneracy operations. It is now straightforward to construct a functor from $`\omega `$-categories to cubical $`\omega `$-categories. The following type of construction is well known. ###### Definition 2.3 The cubical nerve of an $`\omega `$-category $`X`$ is the cubical $`\omega `$-category $`\lambda X`$ defined as follows: $$(\lambda X)_n=Hom[M(I^n),X],$$ and the operations $`_i^\alpha `$, $`\epsilon _i`$, $`_i`$, $`\mathrm{\Gamma }_i^\alpha `$ are induced by $`\stackrel{ˇ}{}_i^\alpha `$, $`\stackrel{ˇ}{\epsilon }_i`$, $`\stackrel{ˇ}{\mu }_i`$, $`\stackrel{ˇ}{\mathrm{\Gamma }}_i^\alpha `$ according to the formulae $$_i^\alpha x=x\stackrel{ˇ}{}_i^\alpha :M(I^{n1})X$$ for $`x:M(I^n)X`$, etc. In particular, in Definition 2.3, note that the domain of $`_i`$ in $`(\lambda X)_n\times (\lambda X)_n`$ is precisely $$Hom[M(I^{i1}\times [0,2]\times I^{ni}),X]$$ according to Remark 1.10. To check that $`\lambda X`$ satisfies the conditions of Definition 2.1, one must check the corresponding identities for the $`\stackrel{ˇ}{}_i^\alpha `$, etc. Many relations essentially come from properties of the underlying morphisms $`\stackrel{ˇ}{}^\alpha `$, etc. The relation $`_i^\alpha \epsilon _i=id`$, for example, comes from the easily checked relation $`\stackrel{ˇ}{\epsilon }\stackrel{ˇ}{}^\alpha =id`$. For relations involving composition, one must use the morphisms $`\stackrel{ˇ}{\iota }^\alpha :M(I)M([0,2])`$ which present $`M([0,2])`$ as a push-out. Thus, to check the relation $`_i^{}(x_iy)=_i^{}x`$, which is a relation between binary operators, one must check that $$(\stackrel{ˇ}{\iota }^{})^1\stackrel{ˇ}{\mu }\stackrel{ˇ}{}^{}(\sigma )=\stackrel{ˇ}{}^{}(\sigma )$$ and $$(\stackrel{ˇ}{\iota }^+)^1\stackrel{ˇ}{\mu }\stackrel{ˇ}{}^{}(\sigma )=\mathrm{}$$ for every cell $`\sigma `$ in $`I^0`$. For the associative law, one must consider morphisms from $`M(I)`$ to $`M([0,3])`$. The functoriality of the tensor product is responsible for formulae looking like commutation rules, such as $`_i^\alpha \epsilon _j=\epsilon _{j1}_i^\alpha `$ for $`i<j`$. ###### Remark 2.4 Any natural operation $`\theta `$ on cubical $`\omega `$-categories determines an underlying homomorphism $`\stackrel{ˇ}{\theta }`$ between $`\omega `$-categories. For example, if $`\theta `$ maps $`G_n`$ to $`G_m`$, then in particular $`\theta `$ maps $$[\lambda M(I^n)]_n=Hom[M(I^n),M(I^n)]$$ to $`[\lambda M(I^n)]_m=Hom[M(I^m),M(I^n)]`$ and $`\stackrel{ˇ}{\theta }=\theta (id):M(I^m)M(I^n).`$ ## 3 The $`\omega `$-category associated to a cubical $`\omega `$-category In this section we construct a functor $`\gamma `$ associating an $`\omega `$-category to a cubical $`\omega `$-category. The idea is to recover an $`\omega `$-category from its nerve. We will use certain folding operations, which are defined as follows. ###### Definition 3.1 Let $`G`$ be a cubical $`\omega `$-category. The folding operations are the operations $$\psi _i,\mathrm{\Psi }_r,\mathrm{\Phi }_m:G_nG_n$$ defined for $`1in1`$, $`1rn`$ and $`0mn`$ by $$\psi _ix=\mathrm{\Gamma }_i^+_{i+1}^{}x_{i+1}x_{i+1}\mathrm{\Gamma }_i^{}_{i+1}^+x,$$ $$\mathrm{\Psi }_r=\psi _{r1}\psi _{r2}\mathrm{}\psi _1,$$ $$\mathrm{\Phi }_m=\mathrm{\Psi }_1\mathrm{\Psi }_2\mathrm{}\mathrm{\Psi }_m=\psi _1(\psi _2\psi _1)\mathrm{}(\psi _{m1}\mathrm{}\psi _1).$$ Note in particular that $`\mathrm{\Psi }_1`$, $`\mathrm{\Phi }_0`$ and $`\mathrm{\Phi }_1`$ are identity operations. Here is a picture of $`\psi _1:G_2G_1`$: The idea behind Definition 3.1 is best seen from the action of the underlying endomorphism $`\stackrel{ˇ}{\mathrm{\Phi }}_n`$ in the $`\omega `$-category of sets $`M(I^n)`$. ###### Proposition 3.2 The endomorphism $`\stackrel{ˇ}{\mathrm{\Phi }}_n:M(I^n)M(I^n)`$ underlying the folding operation $`\mathrm{\Phi }_n`$ is given by $`\stackrel{ˇ}{\mathrm{\Phi }}_n(I^n)=I^n`$ and $$\stackrel{ˇ}{\mathrm{\Phi }}_n(\sigma \times d_0^\alpha I\times I^p)=d_p^\alpha I^n$$ for any cell $`\sigma `$ in $`I^{np1}`$. Proof Let $`\stackrel{ˇ}{\psi }:M(I^2)M(I^2)`$ be the operation underlying $`\psi _1`$ in dimension $`2`$. The operations underlying $`\psi _i`$, $`\mathrm{\Psi }_r`$ and $`\mathrm{\Phi }_m`$ in dimension $`n`$ are then given by $$\stackrel{ˇ}{\psi }_i=id^{i1}\stackrel{ˇ}{\psi }id^{ni1},$$ $$\stackrel{ˇ}{\mathrm{\Psi }}_r=\stackrel{ˇ}{\psi }_1\stackrel{ˇ}{\psi }_2\mathrm{}\stackrel{ˇ}{\psi }_{r1},$$ $$\stackrel{ˇ}{\mathrm{\Phi }}_m=\stackrel{ˇ}{\mathrm{\Psi }}_m\stackrel{ˇ}{\mathrm{\Psi }}_{m1}\mathrm{}\stackrel{ˇ}{\mathrm{\Psi }}_1.$$ One finds that $`\stackrel{ˇ}{\psi }(I^2)=I^2`$, from which it follows that $`\stackrel{ˇ}{\psi }_i(I^n)=I^n`$ and then $`\stackrel{ˇ}{\mathrm{\Phi }}_n(I^n)=I^n`$. One also finds that $$\stackrel{ˇ}{\psi }(\tau \times d_0^\alpha I)=d_0^\alpha I\times d_0^\alpha I$$ for any cell $`\tau `$ in $`I`$. For a cell $`\sigma `$ in $`I^{np1}`$ it follows that $$(\stackrel{ˇ}{\mathrm{\Psi }}_{np1}\mathrm{}\stackrel{ˇ}{\mathrm{\Psi }}_1)(\sigma \times d_0^\alpha I\times I^p)I^{np1}\times d_0^\alpha I\times I^p$$ and $$\stackrel{ˇ}{\mathrm{\Psi }}_{np}(\stackrel{ˇ}{\mathrm{\Psi }}_{np1}\mathrm{}\stackrel{ˇ}{\mathrm{\Psi }}_1)(\sigma \times d_0^\alpha I\times I^p)=(d_0^\alpha I)^{np}\times I^p.$$ It then follows that $`\stackrel{ˇ}{\mathrm{\Phi }}_n(\sigma \times d_0^\alpha I\times I^p)`$ is independent of $`\sigma `$. It now suffices to show that $$\stackrel{ˇ}{\mathrm{\Phi }}_n[(d_0^\alpha I)^{np}\times I^p]=d_p^\alpha I^n.$$ Recall that $`d_{n1}^\alpha I^n`$ is the union of the $`(n1)`$-cells $$\tau _1=d_0^\alpha I\times I^{n1},\tau _2=I\times d_0^\alpha I\times I^{n2},\mathrm{}.$$ We see that $$\stackrel{ˇ}{\mathrm{\Phi }}_n(\tau _2)=\stackrel{ˇ}{\mathrm{\Phi }}_n(d_0^\alpha I\times d_0^\alpha I\times I^{n2})\stackrel{ˇ}{\mathrm{\Phi }}_n(\tau _1),$$ etc., so that $`\stackrel{ˇ}{\mathrm{\Phi }}(d_{n1}^\alpha I^n)=\stackrel{ˇ}{\mathrm{\Phi }}(\tau _1)`$. It follows that $$\stackrel{ˇ}{\mathrm{\Phi }}_n(d_0^\alpha I\times I^{n1})=\stackrel{ˇ}{\mathrm{\Phi }}_n(\tau _1)=\stackrel{ˇ}{\mathrm{\Phi }}_n(d_{n1}^\alpha I^n)=d_{n1}^\alpha \stackrel{ˇ}{\mathrm{\Phi }}_n(I^n)=d_{n1}^\alpha I^n.$$ By similar reasoning, $$\stackrel{ˇ}{\mathrm{\Phi }}_n[(d_0^\alpha I)^2\times I^{n2}]=d_{n2}^\alpha \stackrel{ˇ}{\mathrm{\Phi }}_n(d_0^\alpha I\times I^{n1})=d_{n2}^\alpha d_{n1}^\alpha I^n=d_{n2}^\alpha I^n,$$ and so on, eventually giving $$\stackrel{ˇ}{\mathrm{\Phi }}_n[(d_0^\alpha I)^{np}\times I^p]=d_p^\alpha I^n$$ as required. This completes the proof. $`\mathrm{}`$ It follows from Proposition 3.2 that $`\stackrel{ˇ}{\mathrm{\Phi }}_n:M(I^n)M(I^n)`$ is an idempotent endomorphism with image $$F_n=\{I^n,d_{n1}^{}I^n,d_{n1}^+I^n,\mathrm{},d_0^{}I^n,d_0^+I^n\}.$$ In fact $`F_n`$ is nothing else but the $`n`$-globe. For an $`\omega `$-category $`X`$, it follows that $$\mathrm{\Phi }_n[(\lambda X)_n]Hom(F_n,X).$$ Now, it is clear that $`F_n`$ has a presentation with generator $`I^n`$ and relations $`d_n^{}I^n=d_n^+I^n=I^n`$; therefore $$\mathrm{\Phi }_n(\lambda X)_n\{xX:d_n^{}x=d_n^+x=x\}.$$ It follows that $`X`$ can be recovered from $`\lambda X`$ as the colimit of a sequence $$\mathrm{\Phi }_0[(\lambda X)_0]\mathrm{\Phi }_1[(\lambda X)_1]\mathrm{}.$$ We will now explain how to perform this construction for cubical $`\omega `$-categories in general. We begin with some elementary relations. ###### Proposition 3.3 The folding operations satisfy the following relations: (i) $`\psi _j\epsilon _i`$ $`=\epsilon _i\psi _{j1}`$ $`\mathrm{𝑓𝑜𝑟}i<j,`$ $`\psi _j\epsilon _j`$ $`=\psi _j\epsilon _{j+1}=\psi _j\mathrm{\Gamma }_j^\alpha =\epsilon _j,`$ $`\psi _j\epsilon _i`$ $`=\epsilon _i\psi _j`$ $`\mathrm{𝑓𝑜𝑟}i>j+1,`$ $`_i^\alpha \psi _j`$ $`=\psi _{j1}_i^\alpha `$ $`\mathrm{𝑓𝑜𝑟}i<j,`$ $`_j^{}\psi _jx`$ $`=_j^{}x_j_{j+1}^+x,`$ $`_j^+\psi _jx`$ $`=_{j+1}^{}x_j_j^+x,`$ $`_{j+1}^\alpha \psi _j`$ $`=\epsilon _j_j^\alpha _{j+1}^\alpha ,`$ $`_i^\alpha \psi _j`$ $`=\psi _j_i^\alpha `$ $`\mathrm{𝑓𝑜𝑟}i>j+1;`$ (ii) $`\mathrm{\Psi }_1\epsilon _1`$ $`=\epsilon _1`$ $`\mathrm{\Psi }_r\epsilon _1`$ $`=\epsilon _1\mathrm{\Psi }_{r1}`$ $`\mathrm{𝑓𝑜𝑟}r>1,`$ $`\mathrm{\Psi }_r\epsilon _i`$ $`=\epsilon _{i1}\mathrm{\Psi }_r`$ $`\mathrm{𝑓𝑜𝑟}1<ir,`$ $`_i^\alpha \mathrm{\Psi }_r`$ $`=\mathrm{\Psi }_r_i^\alpha `$ $`\mathrm{𝑓𝑜𝑟}i>r,`$ $`_r^\alpha \mathrm{\Psi }_r`$ $`=\epsilon _1^{r1}(_1^\alpha )^r;`$ (iii) $`\mathrm{\Phi }_m\epsilon _i`$ $`=\epsilon _1\mathrm{\Phi }_{m1}`$ $`\mathrm{𝑓𝑜𝑟}1im,`$ $`_i^\alpha \mathrm{\Phi }_m`$ $`=\mathrm{\Phi }_m_i^\alpha `$ $`\mathrm{𝑓𝑜𝑟}i>m,`$ $`_m^\alpha \mathrm{\Phi }_m`$ $`=\epsilon _1^{m1}(_1^\alpha )^m.`$ Proof (i) These relations are straightforward consequences of the definitions. (ii) Since $`\mathrm{\Psi }_1=id`$, we have $`\mathrm{\Psi }_1\epsilon _1=\epsilon _1`$. From part (i), if $`r>1`$ then $$\mathrm{\Psi }_r\epsilon _1=(\psi _{r1}\mathrm{}\psi _2)\psi _1\epsilon _1=(\psi _{r1}\mathrm{}\psi _2)\epsilon _1=\epsilon _1(\psi _{r2}\mathrm{}\psi _1)=\epsilon _1\mathrm{\Psi }_{r1}.$$ Also from part (i), if $`1<ir`$ then $`\mathrm{\Psi }_r\epsilon _i`$ $`=(\psi _{r1}\mathrm{}\psi _i)\psi _{i1}(\psi _{i2}\mathrm{}\psi _1)\epsilon _i`$ $`=(\psi _{r1}\mathrm{}\psi _i)\psi _{i1}\epsilon _i(\psi _{i2}\mathrm{}\psi _1)`$ $`=(\psi _{r1}\mathrm{}\psi _i)\epsilon _{i1}(\psi _{i2}\mathrm{}\psi _1)`$ $`=\epsilon _{i1}(\psi _{r2}\mathrm{}\psi _{i1})(\psi _{i2}\mathrm{}\psi _1)`$ $`=\epsilon _{i1}\mathrm{\Psi }_{r1}.`$ From part (i), if $`i>r`$ then $$_i^\alpha \mathrm{\Psi }_r=_i^\alpha (\psi _{r1}\mathrm{}\psi _1)=(\psi _{r1}\mathrm{}\psi _1)_i^\alpha =\mathrm{\Psi }_r_i^\alpha .$$ It now follows that $`_r^\alpha \mathrm{\Psi }_r`$ $`=_r^\alpha \psi _{r1}\mathrm{\Psi }_{r1}`$ $`=\epsilon _{r1}_{r1}^\alpha _r^\alpha \mathrm{\Psi }_{r1}`$ $`=\epsilon _{r1}_{r1}^\alpha \mathrm{\Psi }_{r1}_r^\alpha `$ $`=\mathrm{}`$ $`=\epsilon _{r1}\mathrm{}\epsilon _2\epsilon _1_1^\alpha \mathrm{\Psi }_1_2^\alpha \mathrm{}_r^\alpha `$ $`=\epsilon _{r1}\mathrm{}\epsilon _2\epsilon _1_1^\alpha _2^\alpha \mathrm{}_r^\alpha `$ $`=\epsilon _1^{r1}(_1^\alpha )^r,`$ using (2.1). (iii) From part (ii), if $`1im`$ then $`\mathrm{\Phi }_m\epsilon _i`$ $`=\mathrm{\Psi }_1(\mathrm{\Psi }_2\mathrm{}\mathrm{\Psi }_{mi+1})(\mathrm{\Psi }_{mi+2}\mathrm{}\mathrm{\Psi }_m)\epsilon _i`$ $`=\mathrm{\Psi }_1(\mathrm{\Psi }_2\mathrm{}\mathrm{\Psi }_{mi+1})\epsilon _1(\mathrm{\Psi }_{mi+1}\mathrm{}\mathrm{\Psi }_{m1})`$ $`=\mathrm{\Psi }_1\epsilon _1(\mathrm{\Psi }_1\mathrm{}\mathrm{\Psi }_{mi})(\mathrm{\Psi }_{mi+1}\mathrm{}\mathrm{\Psi }_{m1})`$ $`=\epsilon _1(\mathrm{\Psi }_1\mathrm{}\mathrm{\Psi }_{mi})(\mathrm{\Psi }_{mi+1}\mathrm{}\mathrm{\Psi }_{m1})`$ $`=\epsilon _1\mathrm{\Phi }_{m1}.`$ Also from part (ii), if $`i>m`$ then $$_i^\alpha \mathrm{\Phi }_m=_i^\alpha (\mathrm{\Psi }_1\mathrm{}\mathrm{\Psi }_m)=(\mathrm{\Psi }_1\mathrm{}\mathrm{\Psi }_m)_i^\alpha =\mathrm{\Phi }_m_i^\alpha .$$ It now follows that $`_m^\alpha \mathrm{\Phi }_m`$ $`=_m^\alpha \mathrm{\Phi }_{m1}\mathrm{\Psi }_m`$ $`=\mathrm{\Phi }_{m1}_m^\alpha \mathrm{\Psi }_m`$ $`=\mathrm{\Phi }_{m1}\epsilon _1^{m1}(_1^\alpha )^m`$ $`=\epsilon _1^{m1}\mathrm{\Phi }_0(_1^\alpha )^m`$ $`=\epsilon _1^{m1}(_1^\alpha )^m.`$ $`\mathrm{}`$ We now observe that the operators $`\psi _i`$ are idempotent, and characterise their images. ###### Proposition 3.4 Let $`G`$ be a cubical $`\omega `$-category, and suppose that $`1in1`$. The operator $`\psi _i:G_nG_n`$ is idempotent. An element $`x`$ of $`G_n`$ is in $`\psi _i(G_n)`$ if and only if $`_{i+1}^{}x`$ and $`_{i+1}^+x`$ are in $`Im\epsilon _i`$. Proof From Proposition 3.3(i), if $`x\psi _i(G_n)`$ then $`_{i+1}^{}x`$ and $`_{i+1}^+x`$ are in $`Im\epsilon _i`$. To complete the proof, suppose that $`_{i+1}^{}x`$ and $`_{i+1}^+x`$ are in $`Im\epsilon _i`$; it suffices to show that $`\psi _ix=x`$. Now, $$\mathrm{\Gamma }_i^\alpha _{i+1}^\alpha xIm\mathrm{\Gamma }_i^\alpha \epsilon _i=Im\epsilon _i^2=Im\epsilon _{i+1}\epsilon _i,$$ so that the $`\mathrm{\Gamma }_i^\alpha _{i+1}^\alpha x`$ are identities for $`_{i+1}`$. It follows that $$\psi _ix=\mathrm{\Gamma }_i^+_{i+1}^{}x_{i+1}x_{i+1}\mathrm{\Gamma }_i^{}_{i+1}^+x=x$$ as required. This completes the proof. $`\mathrm{}`$ There is a similar result for $`\mathrm{\Phi }_n`$ as follows. ###### Proposition 3.5 Let $`G`$ be a cubical $`\omega `$-category. The operator $`\mathrm{\Phi }_n:G_nG_n`$ is idempotent. An element $`x`$ of $`G_n`$ is in $`\mathrm{\Phi }_n(G_n)`$ if and only if $`_m^\alpha xIm\epsilon _1^{m1}`$ for $`1mn`$ and $`\alpha =\pm `$. Proof Since $`\mathrm{\Phi }_n=\mathrm{\Phi }_m(\mathrm{\Psi }_{m+1}\mathrm{}\mathrm{\Psi }_n)`$, it follows from Proposition 3.3(iii) that $$Im_m^\alpha \mathrm{\Phi }_nIm_m^\alpha \mathrm{\Phi }_mIm\epsilon _1^{m1}.$$ Conversely, suppose that $`_m^\alpha xIm\epsilon _1^{m1}`$ for $`1mn`$ and $`\alpha =\pm `$; it suffices to show that $`\mathrm{\Phi }_nx=x`$, and for this it suffices to show that $`\psi _ix=x`$ for $`1in1`$. But $$_{i+1}^\alpha xIm\epsilon _1^i=Im\epsilon _i\epsilon _1^{i1}Im\epsilon _i$$ for $`\alpha =\pm `$, so that $`\psi _ix=x`$ by Proposition 3.4. This completes the proof. $`\mathrm{}`$ There is a useful result related to Proposition 3.5 as follows. ###### Proposition 3.6 If $`x\mathrm{\Phi }_n(G_n)`$ and $`1mn`$, then $$_m^\alpha x=\epsilon _1^{m1}(_1^\alpha )^mx\mathrm{𝑎𝑛𝑑}\epsilon _m_m^\alpha x=\epsilon _1^m(_1^\alpha )^mx.$$ Proof By Proposition 3.5, $`_m^\alpha x=\epsilon _1^{m1}x^{}`$ for some $`x^{}`$. It follows that $$x^{}=(_1^\alpha )^{m1}\epsilon _1^{m1}x^{}=(_1^\alpha )^{m1}_m^\alpha x=(_1^\alpha )^mx,$$ so that $`_m^\alpha x=\epsilon _1^{m1}x^{}=\epsilon _1^{m1}(_1^\alpha )^mx`$ and $`\epsilon _m_m^\alpha x=\epsilon _m\epsilon _1^{m1}(_1^\alpha )^mx=\epsilon _1^m(_1^\alpha )^mx`$ as required. $`\mathrm{}`$ We now deduce various closure properties for the family of sets $`\mathrm{\Phi }_n(G_n)`$. ###### Proposition 3.7 Let $`G`$ be a cubical $`\omega `$-category. The family of sets $`\mathrm{\Phi }_n(G_n)`$ ($`n0`$) is closed under the $`_i^\alpha `$ and under $`\epsilon _1`$. The individual sets $`\mathrm{\Phi }_n(G_n)`$ are closed under $`\epsilon _i_i^\alpha `$ and $`_i`$ for $`1in`$. Proof We use the characterisation in Proposition 3.5. We first show that the family is closed under $`_1^\alpha `$. Indeed, if $`x\mathrm{\Phi }_n(G_n)`$ then $$_m^\beta _1^\alpha x=_1^\alpha _{m+1}^\beta x_1^\alpha (Im\epsilon _1^m)=Im\epsilon _1^{m1},$$ since $`_1^\alpha \epsilon _1=id`$. Next we show that the family is closed under $`\epsilon _1`$. Indeed, if $`x\mathrm{\Phi }_n(G_n)`$, then $`_1^\beta \epsilon _1xIm\epsilon _1^0`$ trivially, and for $`m>1`$ we have $$_m^\beta \epsilon _1x=\epsilon _1_{m1}^\beta x\epsilon _1(Im\epsilon _1^{m2})=Im\epsilon _1^{m1}.$$ It now follows from Proposition 3.6 that the family is closed under $`_i^\alpha `$ for all $`i`$. Similarly, $`\mathrm{\Phi }_n(G_n)`$ is closed under $`\epsilon _i_i^\alpha `$. It remains to show that $`x_iy\mathrm{\Phi }_n(G_n)`$ when $`x`$ and $`y`$ are in $`\mathrm{\Phi }_n(G_n)`$ and the composite exists. Suppose that $`_m^\beta x=\epsilon _1^{m1}x^{}`$ and $`_m^\beta y=\epsilon _1^{m1}y^{}`$. If $`m<i`$ then $$_m^\beta (x_iy)=_m^\beta x_{i1}_m^\beta y=\epsilon _1^{m1}x^{}_{i1}\epsilon _1^{m1}y^{}$$ $$=\epsilon _1^{m1}(x^{}_{im}y^{})Im\epsilon _1^{m1};$$ if $`m=i`$ then $`_m^\beta (x_iy)`$ is $`_m^\beta x=\epsilon _1^{m1}x^{}`$ or $`_m^\beta y=\epsilon _1^{m1}y^{}`$, so $`_m^\beta (x_iy)`$ is certainly in $`Im\epsilon _1^{m1}`$; and if $`m>i`$ then $$_m^\beta (x_iy)=_m^\beta x_i_m^\beta y=\epsilon _1^{m1}x^{}_i\epsilon _1^{m1}y^{}=$$ $$=\epsilon _1^{i1}(\epsilon _1^{mi}x^{}_1\epsilon _1^{mi}y^{})=\epsilon _1^{i1}\epsilon _1^{mi}x^{}Im\epsilon _1^{m1}$$ (note that $`\epsilon _1^{mi}y^{}`$ is an identity for $`_1`$ because it lies in the image of $`\epsilon _1`$). This completes the proof. $`\mathrm{}`$ We now obtain the desired sequence of $`\omega `$-categories. ###### Theorem 3.8 Let $`G`$ be a cubical $`\omega `$-category. Then there is a sequence of $`\omega `$-categories and homomorphisms with the following structure on $`\mathrm{\Phi }_n(G_n)`$: if $`0p<n`$ then $$d_p^\alpha x=(\epsilon _1)^{np}(_1^\alpha )^{np}x$$ and $`x\#_py=x_{np}y`$ where defined; if $`pn`$ then $`d_p^\alpha x=x`$ and the only composites are given by $`x\#_px=x`$. Proof We first show that for a fixed value of $`n`$ the given structure maps $`d_p^\alpha `$ and $`\#_p`$ make $`\mathrm{\Phi }_n(G_n)`$ into an $`\omega `$-category. By Proposition 3.7, $`\mathrm{\Phi }_n(G_n)`$ is closed under the structure maps for $`0p<n`$, and the same result holds trivially for $`pn`$. From the identities in Section 2, if $`0p<n`$ then the triple $$(d_p^{},d_p^+,\#_p)=(\epsilon _{np}_{np}^{},\epsilon _{np}_{np}^+,_{np})$$ makes $`\mathrm{\Phi }_n(G_n)`$ into the morphism set of a category (with $`d_p^{}x`$ and $`d_p^+x`$ the left and right identities of $`x`$ and with $`\#_p`$ as composition), and these structures commute with one another. Trivially the triples $`(d_p^{},d_p^+,\#_p)`$ for $`np`$ provide further commuting category structures. To show that these structures make $`\mathrm{\Phi }_n(G_n)`$ into an $`\omega `$-category, it now suffices to show that an identity for $`\#_p`$ is also an identity for $`\#_q`$ if $`q>p`$; in other words, it suffices to show that $`d_q^\beta d_p^\alpha x=d_p^\alpha x`$ for $`x\mathrm{\Phi }_n(G_n)`$ and $`q>p`$. For $`qn`$, this is trivial; we may therefore assume that $`0p<q<n`$. But Proposition 3.6 gives us $`d_q^\beta d_p^\alpha x`$ $`=\epsilon _1^{nq}(_1^\beta )^{nq}\epsilon _1^{np}(_1^\alpha )^{np}x`$ $`=\epsilon _1^{nq}\epsilon _1^{qp}(_1^\alpha )^{np}x`$ $`=\epsilon _1^{np}(_1^\alpha )^{np}x`$ $`=d_p^\alpha x`$ as required. We have now shown that the $`\mathrm{\Phi }_n(G_n)`$ are $`\omega `$-categories. We know from Proposition 3.7 that $`\epsilon _1`$ maps $`\mathrm{\Phi }_n(G_n)`$ into $`\mathrm{\Phi }_{n+1}(G_{n+1})`$, and it remains to show that this function is a homomorphism. That is to say, we must show that $`\epsilon _1d_p^\alpha x=d_p^\alpha \epsilon _1x`$ for $`x\mathrm{\Phi }_n(G_n)`$, and we must show that $`\epsilon _1(x\#_py)=\epsilon _1x\#_p\epsilon _1y`$ for $`x\#_py`$ a composite in $`\mathrm{\Phi }_n(G_n)`$. But if $`0p<n`$ then $$\epsilon _1d_p^\alpha x=\epsilon _1\epsilon _{np}_{np}^\alpha x=\epsilon _{np+1}_{np+1}^\alpha \epsilon _1x=d_p^\alpha \epsilon _1x$$ and $$\epsilon _1(x\#_py)=\epsilon _1(x_{np}y)=\epsilon _1x_{np+1}\epsilon _1y=\epsilon _1x\#_p\epsilon _1y$$ by identities in Section 2; if $`p=n`$ we get $$\epsilon _1d_n^\alpha x=\epsilon _1x=\epsilon _1_1^\alpha \epsilon _1x=d_n^\alpha \epsilon _1x$$ and $$\epsilon _1(x\#_nx)=\epsilon _1x=\epsilon _1x_1\epsilon _1x=\epsilon _1x\#_n\epsilon _1x;$$ and if $`p>n`$ then $$\epsilon _1d_p^\alpha x=\epsilon _1x=d_p^\alpha \epsilon _1x$$ and $$\epsilon _1(x\#_px)=\epsilon _1x=\epsilon _1x\#_p\epsilon _1x$$ trivially. This completes the proof. $`\mathrm{}`$ We can now define a functor from cubical $`\omega `$-categories to $`\omega `$-categories. ###### Definition 3.9 Let $`G`$ be a cubical $`\omega `$-category. The $`\omega `$-category $`\gamma G`$ associated to $`G`$ is the colimit of the sequence ###### Remark 3.10 In Definition 3.9, one can identify $`\mathrm{\Phi }_n(G_n)`$ with the subset of $`\gamma G`$ consisting of elements $`x`$ such that $`d_n^{}x=d_n^+x=x`$. Indeed, the $`\epsilon _1`$ are injective, because $`_1^\alpha \epsilon _1=id`$, so that $`\mathrm{\Phi }_n(G_n)`$ can be identified with a subset of $`\gamma G`$; if $`x\mathrm{\Phi }_n(G_n)`$ then $`d_n^{}x=d_n^+x=x`$ by Theorem 3.8; if $`x\mathrm{\Phi }_m(G_m)`$ with $`m>n`$ and $`d_n^{}x=d_n^+x=x`$, then $$x=\epsilon _{mn}_{mn}^{}x=\epsilon _1^{mn}(_1^{})^{mn}x$$ (Proposition 3.6) with $`(_1^{})^{mn}x\mathrm{\Phi }_n(G_n)`$ (Proposition 3.7), and $`x`$ can be identified with $`(_1^{})^{mn}x`$. ###### Remark 3.11 It is convenient to describe the $`\omega `$-category $`\gamma G`$ in terms of the folding operations, but one can get a more direct description by using Proposition 3.5. The more direct description needs face maps, degeneracies and compositions, but not connections. ## 4 The natural isomorphism $`A:\gamma \lambda XX`$ Let $`X`$ be an $`\omega `$-category. From Definition 2.3 there is a cubical $`\omega `$-category $`\lambda X`$, and from Definition 3.9 there is an $`\omega `$-category $`\gamma \lambda X`$. We will now construct a natural isomorphism $`A:\gamma \lambda XX`$. Let $`F_n`$ be the $`\omega `$-category with one generator $`I^n`$ and with relations $`d_n^{}I^n=d_n^+I^n=I^n`$. By Proposition 3.2, $`F_n`$ can be realised as a sub-$`\omega `$-category of $`M(I^n)`$, and the morphism $`\stackrel{ˇ}{\mathrm{\Phi }}_n:M(I^n)M(I^n)`$ associated to the folding operation $`\mathrm{\Phi }_n`$ is an idempotent operation with image equal to $`F_n`$. Recalling that $`(\lambda X)_n=Hom[M(I^n),X]`$, we see that $$\mathrm{\Phi }_n[(\lambda X)_n]=\{xHom[M(I^n),X]:x\stackrel{ˇ}{\mathrm{\Phi }}_n=x\}.$$ Let $$A:\mathrm{\Phi }_n[(\lambda X)_n]X$$ be the function given by $$A(x)=x(I^n);$$ we see that $`A`$ is an injection with image equal to $$\{xX:d_n^{}x=d_n^+x=x\}.$$ These functions are compatible with the sequence indeed, if $`x\mathrm{\Phi }_n[(\lambda X)_n]`$ then $$A(\epsilon _1x)=(\epsilon _1x)(I^{n+1})=x\stackrel{ˇ}{\epsilon }_1(I^{n+1})=x(I^n)=A(x).$$ The functions $`A:\mathrm{\Phi }_n[(\lambda X)_n]X`$ therefore induce a bijection $`A:\gamma \lambda XX`$. We will now prove the following result. ###### Theorem 4.1 The functions $`A:\gamma \lambda XX`$ form a natural isomorphism of $`\omega `$-categories. Proof We have already shown that $`A:\gamma \lambda XX`$ is a bijection, and it is clearly natural. It remains to show that $`A`$ is a homomorphism. It suffices to show that $$A:\mathrm{\Phi }_n[(\lambda X)_n]X$$ is a homomorphism for each $`n`$; in other words, we must show that $`A(d_p^\alpha x)=d_p^\alpha A(x)`$ for $`x\mathrm{\Phi }_n[(\lambda X)_n]`$ and that $`A(x\#_py)=A(x)\#_pA(y)`$ for $`x\#_py`$ a composite in $`\mathrm{\Phi }_n[(\lambda X)_n]`$. Suppose that $`x\mathrm{\Phi }_n[(\lambda X)_n]`$ and $`0p<n`$. Noting that $`x=x\stackrel{ˇ}{\mathrm{\Phi }}_n`$ and using Proposition 3.2, we find that $`A(d_p^\alpha x)`$ $`=A(\epsilon _{np}_{np}^\alpha x)`$ $`=(\epsilon _{np}_{np}^\alpha x)(I^n)`$ $`=x\stackrel{ˇ}{}_{np}^\alpha \stackrel{ˇ}{\epsilon }_{np}(I^n)`$ $`=x(I^{np1}\times d_0^\alpha I\times I^p)`$ $`=x\stackrel{ˇ}{\mathrm{\Phi }}_n(I^{np1}\times d_0^\alpha I\times I^p)`$ $`=x(d_p^\alpha I^n)`$ $`=d_p^\alpha x(I^n)`$ $`=d_p^\alpha A(x).`$ Suppose that $`x\mathrm{\Phi }_n[(\lambda X)_n]`$ and $`pn`$. Then $$A(d_p^\alpha x)=A(x)=x(I^n)=x(d_p^\alpha I^n)=d_p^\alpha x(I^n)=d_p^\alpha x.$$ Suppose that $`x\mathrm{\#}_py`$ is a composite in $`\mathrm{\Phi }_n[(\lambda X)_n]`$ with $`0p<n`$. Let $$(x,y):M(I^{np1}\times [0,2]\times I^p)X$$ be the morphism such that $`(x,y)\stackrel{ˇ}{\iota }_{np}^{}=x`$ and $`(x,y)\stackrel{ˇ}{\iota }_{np}^+=y`$; then $$A(x\#_py)=A(x_{np}y)=(x,y)\stackrel{ˇ}{\mu }_{np}(I^n).$$ Let $`\eta :F_nM(I^n)`$ be the inclusion and let $`\pi :M(I^n)F_n`$ be $`\stackrel{ˇ}{\mathrm{\Phi }}_n`$ with its codomain restricted to $`F_n`$, so that $`\stackrel{ˇ}{\mathrm{\Phi }}_n=\eta \pi `$. Since $`x`$ and $`y`$ are in $`\mathrm{\Phi }_n[(\lambda X)_n]`$, we have $`x\stackrel{ˇ}{\mathrm{\Phi }}_n=x`$ and $`y=y\stackrel{ˇ}{\mathrm{\Phi }}_n`$; we therefore get $$(x,y)\stackrel{ˇ}{\mu }_{np}(I^n)=(x\eta \pi ,y\eta \pi )\stackrel{ˇ}{\mu }_{np}(I^n).$$ Now let $`F_p`$ be the $`\omega `$-category with one generator $`z`$ and with relations $`d_p^{}z=d_p^+z=z`$. We see that there is a factorisation $$(x\eta \pi ,y\eta \pi )=(x\eta ,y\eta )(\pi ,\pi )$$ through the obvious push-out of $$F_nF_pF_n.$$ We also see that $$(\pi ,\pi )\stackrel{ˇ}{\mu }_{np}(I^n)=\pi \stackrel{ˇ}{\iota }_{np}^{}(I^n)\#_p\pi \stackrel{ˇ}{\iota }_{np}^+(I^n).$$ It now follows that $`(x\eta \pi ,y\eta \pi )\stackrel{ˇ}{\mu }_{np}(I^n)`$ $`=(x\eta ,y\eta )(\pi ,\pi )\stackrel{ˇ}{\mu }_{np}(I^n)`$ $`=(x\eta ,y\eta )[\pi \stackrel{ˇ}{\iota }_{np}^{}(I^n)\#_p\pi \stackrel{ˇ}{\iota }_{np}^+(I^n)]`$ $`=(x\eta ,y\eta )\pi \stackrel{ˇ}{\iota }_{np}^{}(I^n)\#_p(x\eta ,y\eta )\pi \stackrel{ˇ}{\iota }_{np}^+(I^n)`$ $`=x\eta \pi (I^n)\#_py\eta \pi (I^n)`$ $`=x(I^n)\#_py(I^n)`$ $`=A(x)\#_pA(y);`$ therefore $$A(x\#_py)=A(x)\#_pA(y).$$ Finally, suppose that $`x\mathrm{\#}_py`$ is a composite in $`\mathrm{\Phi }_n[(\lambda X)_n]`$ with $`pn`$. We must have $`x=y`$, and we get $$A(x\#_px)=A(x)=A(x)\#_pA(x).$$ This completes the proof. $`\mathrm{}`$ ## 5 Foldings, degeneracies and connections According to Theorem 4.1, there are natural isomorphisms $`A:\gamma \lambda XX`$ for $`\omega `$-categories $`X`$. To prove that $`\omega `$-categories are equivalent to cubical $`\omega `$-categories, we will eventually construct natural isomorphisms $`B:G\lambda \gamma G`$ for cubical $`\omega `$-categories $`G`$. We will need properties of the folding operations, and we now begin to describe these. We first show that the operations $`\psi _i`$ behave like the standard generating transpositions of the symmetric groups (except of course that they are idempotent rather than involutory, by Proposition 3.4). There are two types of relation, the first of which is easy. ###### Proposition 5.1 If $`\|ij\|2`$, then $$\psi _i\psi _j=\psi _j\psi _i.$$ Proof This follows from the identities in Section 2. $`\mathrm{}`$ The next result is harder. ###### Theorem 5.2 If $`i>1`$ then $$\psi _i\psi _{i1}\psi _i=\psi _{i1}\psi _i\psi _{i1}.$$ Proof Recall the matrix notation used for certain composites: if $$(a_{11}_i\mathrm{}_ia_{1n})_{i+1}\mathrm{}_{i+1}(a_{m1}_i\mathrm{}_ia_{mn})$$ and $$(a_{11}_{i+1}\mathrm{}_{i+1}a_{m1})_i\mathrm{}_i(a_{1n}_{i+1}\mathrm{}_{i+1}a_{mn})$$ are equal by the interchange law, then we will write $$\left[\begin{array}{ccc}a_{11}& \mathrm{}& a_{1n}\\ \mathrm{}& & \mathrm{}\\ a_{m1}& \mathrm{}& a_{mn}\end{array}\right]\text{}$$ for the common value. In such a matrix, we write $``$ for elements in the image of $`\epsilon _i`$ (which are the identities for $`_i`$), and we write $`\|`$ for elements in the image of $`\epsilon _{i+1}`$ (which are the identities for $`_{i+1}`$). We first compute $`\psi _i\psi _{i1}\psi _ix`$. It is straightforward to check that $`\psi _{i1}\psi _ix`$ $`=\left[\begin{array}{ccc}\|& \mathrm{\Gamma }_i^+_{i+1}^{}x& \mathrm{\Gamma }_{i1}^{}_{i+1}^{}x\\ \mathrm{\Gamma }_{i1}^+_i^{}x& x& \mathrm{\Gamma }_{i1}^{}_i^+x\\ \mathrm{\Gamma }_{i1}^+_{i+1}^+x& \mathrm{\Gamma }_i^{}_{i+1}^+x& \|\end{array}\right]\text{}.`$ It follows that $`\mathrm{\Gamma }_i^+_{i+1}^{}\psi _{i1}\psi _ix`$ $`=\mathrm{\Gamma }_i^+(\mathrm{\Gamma }_{i1}^+_i^{}_i^{}x_i\epsilon _i_i^{}_i^{}x_i\mathrm{\Gamma }_{i1}^{}_i^{}_i^{}x)`$ $`=\mathrm{\Gamma }_i^+(\mathrm{\Gamma }_{i1}^+_i^{}_i^{}x_i\mathrm{\Gamma }_{i1}^{}_i^{}_i^{}x)`$ $`=\mathrm{\Gamma }_i^+\epsilon _{i1}_i^{}_i^{}x`$ $`=\epsilon _{i1}\mathrm{\Gamma }_{i1}^+_i^{}_i^{}x`$ and $`\mathrm{\Gamma }_i^{}_{i+1}^+\psi _{i1}\psi _ix`$ $`=\epsilon _{i1}\mathrm{\Gamma }_{i1}^{}_i^+_i^+x;`$ therefore $`\psi _i\psi _{i1}\psi _ix`$ $`=\epsilon _{i1}\mathrm{\Gamma }_{i1}^+_i^{}_i^{}x_{i+1}\psi _{i1}\psi _ix_{i+1}\epsilon _{i1}\mathrm{\Gamma }_{i1}^{}_i^+_i^+x.`$ Similarly, $`\psi _{i1}\psi _i\psi _{i1}x`$ is as a composite $$\left[\begin{array}{ccccc}& \mathrm{\Gamma }_{i1}^+\mathrm{\Gamma }_{i1}^+_i^{}_i^{}x& & & \mathrm{\Gamma }_{i1}^{}\mathrm{\Gamma }_{i1}^+_i^{}_i^{}x\\ & \|& \mathrm{\Gamma }_i^+_{i+1}^{}x& & \mathrm{\Gamma }_{i1}^{}_{i+1}^{}x\\ & \|& \|& \mathrm{\Gamma }_i^+\mathrm{\Gamma }_{i1}^{}_i^{}_i^+x& \mathrm{\Gamma }_i^{}\mathrm{\Gamma }_{i1}^{}_i^{}_i^+x\\ & \mathrm{\Gamma }_{i1}^+_i^{}x& x& \mathrm{\Gamma }_{i1}^{}_i^+x& \\ \mathrm{\Gamma }_i^+\mathrm{\Gamma }_{i1}^+_i^+_i^{}x& \mathrm{\Gamma }_i^{}\mathrm{\Gamma }_{i1}^+_i^+_i^{}x& \|& \|& \\ \mathrm{\Gamma }_{i1}^+_{i+1}^+x& & \mathrm{\Gamma }_i^{}_{i+1}^+x& \|& \\ \mathrm{\Gamma }_{i1}^+\mathrm{\Gamma }_{i1}^{}_i^+_i^+x& & & \mathrm{\Gamma }_{i1}^{}\mathrm{\Gamma }_{i1}^{}_i^+_i^+x& \end{array}\right]$$ We now evaluate the rows of the matrix for $`\psi _{i1}\psi _i\psi _{i1}x`$. The first and last rows yield $`\epsilon _{i1}\mathrm{\Gamma }_{i1}^+_i^{}_i^{}x`$ and $`\epsilon _{i1}\mathrm{\Gamma }_{i1}^{}_i^+_i^+x`$. The composite of the non-identity elements in the third row is $`\epsilon _{i+1}\mathrm{\Gamma }_{i1}^{}_i^{}_i^+x`$, which is an identity for $`_{i+1}`$, so the third row can be omitted. Similarly, the fifth row can be omitted. The second, fourth and sixth rows have the same values as the rows of the matrix for $`\psi _{i1}\psi _ix`$. It follows that $$\psi _{i1}\psi _i\psi _{i1}x=\epsilon _{i1}\mathrm{\Gamma }_{i1}^+_i^{}_i^{}x_{i+1}\psi _{i1}\psi _ix_{i+1}\epsilon _{i1}\mathrm{\Gamma }_{i1}^{}_i^+_i^+x$$ also. Therefore $`\psi _i\psi _{i1}\psi _ix=\psi _{i1}\psi _i\psi _{i1}x`$. This completes the proof. $`\mathrm{}`$ ###### Remark 5.3 Proposition 5.1 and Theorem 5.2 in some sense explain the formula for $`\mathrm{\Phi }_m`$ in Definition 3.1. The $`\psi _i`$ behave like the generating transpositions $`(i,i+1)`$ in the symmetric group of permutations of $`\{1,\mathrm{},m\}`$, and, as in the symmetric group, there are $`m!`$ distinct composites of $`\psi _1,\mathrm{},\psi _{m1}`$ given by $$\mathrm{\Psi }_{1,l(1)}\mathrm{}\mathrm{\Psi }_{m,l(m)}$$ for $`1l(r)r`$, where $`\mathrm{\Psi }_{r,l(r)}=\psi _{r1}\psi _{r2}\mathrm{}\psi _{l(r)}`$. Note that $`_1^\pm \mathrm{\Phi }_m`$ is the hemispherical decomposition of the $`m`$-categorical $`m`$-cube. This orders the top dimensional cells in the upper hemisphere in reverse order to the top dimensional cells of the lower hemisphere. Thus the composite $`\mathrm{\Phi }_m`$ corresponds to the order-reversing permutation $`pm+1p`$. In our context, we can characterise $`\mathrm{\Phi }_m`$ as the zero element in the semigroup generated by $`\psi _1,\mathrm{},\psi _{m1}`$ as follows. ###### Theorem 5.4 If $`1im1`$ then $$\mathrm{\Phi }_m\psi _i=\mathrm{\Phi }_m.$$ Proof If $`r>1`$ then $$\mathrm{\Psi }_r\psi _1=(\psi _{r1}\mathrm{}\psi _2)\psi _1\psi _1=(\psi _{r1}\mathrm{}\psi _2)\psi _1=\mathrm{\Psi }_r,$$ since $`\psi _1`$ is idempotent by Proposition 3.4. For $`1<i<r`$, it follows from Proposition 5.1 and Theorem 5.2 that $`\mathrm{\Psi }_r\psi _i`$ $`=(\psi _{r1}\mathrm{}\psi _{i+1})\psi _i\psi _{i1}(\psi _{i2}\mathrm{}\psi _1)\psi _i`$ $`=(\psi _{r1}\mathrm{}\psi _{i+1})\psi _i\psi _{i1}\psi _i(\psi _{i2}\mathrm{}\psi _1)`$ $`=(\psi _{r1}\mathrm{}\psi _{i+1})\psi _{i1}\psi _i\psi _{i1}(\psi _{i2}\mathrm{}\psi _1)`$ $`=\psi _{i1}(\psi _{r1}\mathrm{}\psi _{i+1})\psi _i\psi _{i1}(\psi _{i2}\mathrm{}\psi _1)`$ $`=\psi _{i1}\mathrm{\Psi }_r.`$ For $`1im1`$, it now follows that $`\mathrm{\Phi }_m\psi _i`$ $`=\mathrm{\Psi }_1(\mathrm{\Psi }_2\mathrm{}\mathrm{\Psi }_{mi})\mathrm{\Psi }_{mi+1}(\mathrm{\Psi }_{mi+2}\mathrm{}\mathrm{\Psi }_m)\psi _i`$ $`=\mathrm{\Psi }_1(\mathrm{\Psi }_2\mathrm{}\mathrm{\Psi }_{mi})\mathrm{\Psi }_{mi+1}\psi _1(\mathrm{\Psi }_{mi+2}\mathrm{}\mathrm{\Psi }_m)`$ $`=\mathrm{\Psi }_1(\mathrm{\Psi }_2\mathrm{}\mathrm{\Psi }_{mi})\mathrm{\Psi }_{mi+1}(\mathrm{\Psi }_{mi+2}\mathrm{}\mathrm{\Psi }_m)`$ $`=\mathrm{\Phi }_m,`$ as required. $`\mathrm{}`$ We can now give some interactions between $`\mathrm{\Phi }_m`$, degeneracies and connections. First we have the following result. ###### Proposition 5.5 For all $`i`$ there is a relation $$\psi _i\mathrm{\Gamma }_i^\alpha =\epsilon _i.$$ Proof From the definitions we get $`\psi _i\mathrm{\Gamma }_i^+x`$ $`=\mathrm{\Gamma }_i^+_{i+1}^{}\mathrm{\Gamma }_i^+x_{i+1}\mathrm{\Gamma }_i^+x_{i+1}\mathrm{\Gamma }_i^{}_{i+1}^+\mathrm{\Gamma }_i^+x`$ $`=\mathrm{\Gamma }_i^+\epsilon _i_i^{}x_{i+1}\mathrm{\Gamma }_i^+x_{i+1}\mathrm{\Gamma }_i^{}x`$ $`=\epsilon _i^2_i^{}x_{i+1}\epsilon _ix`$ $`=\epsilon _{i+1}\epsilon _i_i^{}x_{i+1}\epsilon _ix`$ $`=\epsilon _ix,`$ and we similarly get $`\psi _i\mathrm{\Gamma }_i^{}x=\epsilon _ix`$. $`\mathrm{}`$ We draw the following conclusions. ###### Theorem 5.6 If $`1im`$ then $$\mathrm{\Phi }_m\epsilon _i=\epsilon _1\mathrm{\Phi }_{m1}.$$ If $`1im1`$ then $$\mathrm{\Phi }_m\mathrm{\Gamma }_i^\alpha =\epsilon _1\mathrm{\Phi }_{m1}.$$ Proof The first of these results was given in Proposition 3.3(iii). The second result then follows from Theorem 5.4 and Proposition 5.5: indeed, we get $$\mathrm{\Phi }_m\mathrm{\Gamma }_i^\alpha =\mathrm{\Phi }_m\psi _i\mathrm{\Gamma }_i^\alpha =\mathrm{\Phi }_m\epsilon _i=\epsilon _1\mathrm{\Phi }_{m1}$$ as required. $`\mathrm{}`$ ## 6 Foldings, face maps and compositions In this section we describe interactions between the $`\mathrm{\Phi }_n`$, face maps and compositions. For face maps, the basic results are given in Proposition 3.3. For compositions, the basic results are as follows, of which the first two cases correspond to the 2-dimensional case in \[14, Proposition 5.1\]. ###### Proposition 6.1 In a cubical $`\omega `$-category $$\psi _i(x_jy)=\{\begin{array}{cc}(\psi _ix_{i+1}\epsilon _i_{i+1}^+y)_i(\epsilon _i_{i+1}^{}x_{i+1}\psi _iy)\hfill & \text{if }j=i,\hfill \\ (\epsilon _i_i^{}x_{i+1}\psi _iy)_i(\psi _ix_{i+1}\epsilon _i_i^+y)\hfill & \text{if }j=i+1,\hfill \\ \psi _ix_j\psi _iy\hfill & \text{otherwise.}\hfill \end{array}$$ Proof Note that we have $`_j^+x=_j^{}y`$ for $`x_jy`$ to be defined. The proof for the cases $`j=i`$ and $`j=i+1`$ consists in evaluating in two ways each of the matrices $$\left[\begin{array}{cc}\mathrm{\Gamma }_i^+_{i+1}^{}x& \epsilon _i_{i+1}^{}x\\ \epsilon _{i+1}_{i+1}^{}x& \mathrm{\Gamma }_i^+_{i+1}^{}y\\ x& y\\ \mathrm{\Gamma }_i^{}_{i+1}^+x& \epsilon _{i+1}_{i+1}^+y\\ \epsilon _i_{i+1}^+y& \mathrm{\Gamma }_i^{}_{i+1}^+y\end{array}\right]\left[\begin{array}{cc}\epsilon _i\epsilon _i_i^{}_i^{}x& \mathrm{\Gamma }_i^+_{i+1}^{}x\\ \epsilon _i_i^{}x& x\\ \mathrm{\Gamma }_i^+_{i+1}^{}y& \mathrm{\Gamma }_i^{}_{i+1}^+x\\ y& \epsilon _i_i^+y\\ \mathrm{\Gamma }_i^{}_{i+1}^+y& \epsilon _i\epsilon _i_i^+_i^+y\end{array}\right]\text{}$$ (Note that $`\epsilon _i\epsilon _i_i^{}_i^{}x`$ and $`\epsilon _i\epsilon _i_i^+_i^+y`$ are identities for $`_{i+1}`$ because $`\epsilon _i\epsilon _i=\epsilon _{i+1}\epsilon _i)`$. The other case follows from the identities in Section 2. $`\mathrm{}`$ Because of Proposition 6.1, it is convenient to regard $`\epsilon _i_i^\alpha `$ and $`\epsilon _i_{i+1}^\alpha `$ as generalisations of $`\psi _i`$. We extend this idea to $`\mathrm{\Psi }_r`$ and $`\mathrm{\Phi }_m`$, and arrive at the following definition. ###### Definition 6.2 A generalised $`\psi _i`$ is an operator of the form $`\psi _i`$ or $`\epsilon _i_i^\alpha `$ or $`\epsilon _i_{i+1}^\alpha `$. A generalised $`\mathrm{\Psi }_r`$ is an operator of the form $`\psi _{r1}^{}\psi _{r2}^{}\mathrm{}\psi _1^{}`$, where $`\psi _i^{}`$ is a generalised $`\psi _i`$. A generalised $`\mathrm{\Phi }_m`$ is an operator of the form $`\mathrm{\Psi }_1^{}\mathrm{\Psi }_2^{}\mathrm{}\mathrm{\Psi }_m^{}`$, where $`\mathrm{\Psi }_r^{}`$ is a generalised $`\mathrm{\Psi }_r`$. From (2.3) and (2.5), there are results for $`\epsilon _i`$ and $`_i^\alpha `$ analogous to Proposition 6.1: $`\epsilon _i(x_jy)`$ is a composite of $`\epsilon _ix`$ and $`\epsilon _iy`$; if $`j=i`$ then $`_i^\alpha (x_jy)`$ is $`_i^\alpha x`$ or $`_i^\alpha y`$; if $`ji`$ then $`_i^\alpha (x_jy)`$ is a composite of $`_i^\alpha x`$ and $`_i^\alpha y`$. From these observations and from Proposition 6.1 we immediately get the following result. ###### Proposition 6.3 Let $`\psi _i^{}`$ be a generalised $`\psi _i`$. Then $`\psi _i^{}(x^{}_jx^+)`$ is naturally equal to a composite of factors $`\psi _i^{\prime \prime }x^\alpha `$ with $`\psi _i^{\prime \prime }`$ a generalised $`\psi _i`$. Let $`\mathrm{\Psi }_r^{}`$ be a generalised $`\mathrm{\Psi }_r`$. Then $`\mathrm{\Psi }_r^{}(x^{}_jx^+)`$ is naturally equal to a composite of factors $`\mathrm{\Psi }_r^{\prime \prime }x^\alpha `$ with $`\mathrm{\Psi }_r^{\prime \prime }`$ a generalised $`\mathrm{\Psi }_r`$. Let $`\mathrm{\Phi }_m^{}`$ be a generalised $`\mathrm{\Phi }_m`$. Then $`\mathrm{\Phi }_m^{}(x^{}_jx^+)`$ is naturally equal to a composite of factors $`\mathrm{\Phi }_m^{\prime \prime }x^\alpha `$ with $`\mathrm{\Phi }_m^{\prime \prime }`$ a generalised $`\mathrm{\Phi }_m`$. We will eventually express a generalised $`\mathrm{\Phi }_n`$ in terms of the genuine folding operators $`\mathrm{\Phi }_m`$. In order to do this, we now investigate the faces of generalised foldings. ###### Proposition 6.4 Let $`\psi _j^{}`$ be a generalised $`\psi _j`$. If $`i<j`$, then $`_i^\alpha \psi _j^{}=\psi _{j1}^{\prime \prime }_i^\alpha `$ with $`\psi _{j1}^{\prime \prime }`$ a generalised $`\psi _{j1}`$. If $`i=j`$, then $`_i^\alpha \psi _j^{}x`$ is naturally equal to $`_j^\beta x`$ or $`_{j+1}^\beta x`$ for some $`\beta `$, or to a composite of two such factors. If $`i>j`$ then $`_i^\alpha \psi _j^{}=\psi _j^{\prime \prime }_i^\beta `$ for some $`\beta `$, with $`\psi _j^{\prime \prime }`$ a generalised $`\psi _j`$. Proof We use relations from Section 2 and Proposition 3.3. For $`i<j`$ we have $`_i^\alpha \psi _j=\psi _{j1}_i^\alpha `$ or $`_i^\alpha \epsilon _j_j^\gamma =\epsilon _{j1}_i^\alpha _j^\gamma =\epsilon _{j1}_{j1}^\gamma _i^\alpha `$ or $`_i^\alpha \epsilon _j_{j+1}^\gamma =\epsilon _{j1}_i^\alpha _{j+1}^\gamma =\epsilon _{j1}_j^\gamma _i^\alpha `$. For $`i=j`$ we have $`_j^{}\psi _jx=_j^{}x_j_{j+1}^+x`$ or $`_j^+\psi _jx=_{j+1}^{}x_j_j^+x`$ or $`_j^\alpha \epsilon _j_j^\gamma x=_j^\gamma x`$ or $`_j^\alpha \epsilon _j_{j+1}^\gamma x=_{j+1}^\gamma x`$. For $`i=j+1`$ we have $`_{j+1}^\alpha \psi _j=\epsilon _j_j^\alpha _{j+1}^\alpha `$ or $`_{j+1}^\alpha \epsilon _j_j^\gamma =\epsilon _j_j^\alpha _j^\gamma =\epsilon _j_j^\gamma _{j+1}^\alpha `$ or $`_{j+1}^\alpha \epsilon _j_{j+1}^\gamma =\epsilon _j_j^\alpha _{j+1}^\gamma `$. For $`i>j+1`$ we have $`_i^\alpha \psi _j=\psi _j_i^\alpha `$ or $`_i^\alpha \epsilon _j_j^\gamma =\epsilon _j_{i1}^\alpha _j^\gamma =\epsilon _j_j^\gamma _i^\alpha `$ or $`_i^\alpha \epsilon _j_{j+1}^\gamma =\epsilon _j_{i1}^\alpha _{j+1}^\gamma =\epsilon _j_{j+1}^\gamma _i^\alpha `$. $`\mathrm{}`$ For a generalised $`\mathrm{\Psi }_r`$ we get the following results. ###### Proposition 6.5 Let $`\mathrm{\Psi }_r^{}`$ be a generalised $`\mathrm{\Psi }_r`$. If $`ir`$, then $`_i^\alpha \mathrm{\Psi }_r^{}=\mathrm{\Psi }_r^{\prime \prime }_i^\beta `$ for some $`\beta `$, with $`\mathrm{\Psi }_r^{\prime \prime }`$ a generalised $`\mathrm{\Psi }_r`$. If $`i<r`$, then $`_i^\alpha \mathrm{\Psi }_r^{}x`$ is naturally equal to a composite of factors $`\mathrm{\Psi }_{r1}^{\prime \prime }_h^\beta x`$ with $`hr`$ and with $`\mathrm{\Psi }_{r1}^{\prime \prime }`$ a generalised $`\mathrm{\Psi }_{r1}`$. Proof If $`ir`$ then $`_i^\alpha \mathrm{\Psi }_r^{}=_i^\alpha (\psi _{r1}^{}\mathrm{}\psi _1^{})`$ with $`\psi _j^{}`$ a generalised $`\psi _j`$, and the result is immediate from Proposition 6.4. Now suppose that $`i<r`$. Then $$_i^\alpha \mathrm{\Psi }_r^{}x=_i^\alpha (\psi _{r1}^{}\mathrm{}\psi _{i+1}^{})\psi _i^{}\mathrm{\Psi }_i^{}x,$$ with $`\psi _j^{}`$ a generalised $`\psi _j`$ and with $`\mathrm{\Psi }_i^{}`$ a generalised $`\mathrm{\Psi }_i`$. By Proposition 6.4 $$_i^\alpha \mathrm{\Psi }_r^{}x=(\psi _{r2}^{\prime \prime }\mathrm{}\psi _i^{\prime \prime })_i^\alpha \psi _i^{}\mathrm{\Psi }_i^{}x$$ with $`\psi _j^{\prime \prime }`$ a generalised $`\psi _j`$. By Propositions 6.4 and 6.3, this is a composite of factors of the form $$(\psi _{r2}^{\prime \prime \prime }\mathrm{}\psi _i^{\prime \prime \prime })_h^\gamma \mathrm{\Psi }_i^{}x,$$ with $`ihi+1r`$ and with $`\psi _j^{\prime \prime \prime }`$ a generalised $`\psi _j`$. Since $`hi`$, it follows from the case already covered that the factors can be written as $$(\psi _{r2}^{\prime \prime \prime }\mathrm{}\psi _i^{\prime \prime \prime })\mathrm{\Psi }_i^{\prime \prime }_h^\beta x,$$ with $`\mathrm{\Psi }_i^{\prime \prime }`$ a generalised $`\mathrm{\Psi }_i`$. The factors now have the form $`\mathrm{\Psi }_{r1}^{\prime \prime }_h^\beta x`$ with $`\mathrm{\Psi }_r^{\prime \prime }`$ a generalised $`\mathrm{\Psi }_r`$, as required. $`\mathrm{}`$ By iterating Proposition 6.5, we get the following result. ###### Proposition 6.6 If $`imn`$ and $`\mathrm{\Psi }_r^{}`$ is a generalised $`\mathrm{\Psi }_r`$ for $`m<rn`$, then $`_i^\alpha (\mathrm{\Psi }_{m+1}^{}\mathrm{}\mathrm{\Psi }_n^{})x`$ is naturally equal to a composite of factors $`(\mathrm{\Psi }_m^{\prime \prime }\mathrm{}\mathrm{\Psi }_{n1}^{\prime \prime })_h^\beta x`$ with $`hn`$ and with $`\mathrm{\Psi }_r^{\prime \prime }`$ a generalised $`\mathrm{\Psi }_r`$. Proof This follows from Propositions 6.5 and 6.3. $`\mathrm{}`$ Now let $`\mathrm{\Phi }_n^{}`$ be a generalised $`\mathrm{\Phi }_n`$; we aim to express $`\mathrm{\Phi }_n^{}`$ in terms of the genuine folding operators $`\mathrm{\Phi }_m`$. If $`n=0`$ or $`n=1`$, then necessarily $`\mathrm{\Phi }_n^{}=\mathrm{\Phi }_n`$ already. In general, we use an inductive process; the inductive step is as follows. ###### Proposition 6.7 Let $`\mathrm{\Phi }_n^{}`$ be a generalised $`\mathrm{\Phi }_n`$ which is distinct from $`\mathrm{\Phi }_n`$. Then $`\mathrm{\Phi }_n^{}x`$ is naturally a composite of factors $`\epsilon _1\mathrm{\Phi }_{n1}^{}_h^\beta x`$ with $`\mathrm{\Phi }_{n1}^{}`$ a generalised $`\mathrm{\Phi }_{n1}`$. Proof By considering the first place where $`\mathrm{\Phi }_n^{}`$ and $`\mathrm{\Phi }_n`$ differ, we see that $$\mathrm{\Phi }_n^{}x=[\mathrm{\Phi }_{m1}(\psi _{m1}\mathrm{}\psi _{j+1})\epsilon _j][_i^\alpha \mathrm{\Psi }_j^{}(\mathrm{\Psi }_{m+1}^{}\mathrm{}\mathrm{\Psi }_n^{})x]$$ for some $`m`$ and $`j`$ such that $`1j<mn`$, with $`i=j`$ or $`i=j+1`$ and with $`\mathrm{\Psi }_r^{}`$ a generalised $`\mathrm{\Psi }_r`$. Since $`jm1`$, it follows from Proposition 3.3 that $$\mathrm{\Phi }_{m1}(\psi _{m1}\mathrm{}\psi _{j+1})\epsilon _j=\mathrm{\Phi }_{m1}\epsilon _j(\psi _{m2}\mathrm{}\psi _j)=\epsilon _1\mathrm{\Phi }_{m2}(\psi _{m2}\mathrm{}\psi _j);$$ since $`jim`$, it follows from Propositions 6.5, 6.6 and 6.3 that $$_i^\alpha \mathrm{\Psi }_j^{}(\mathrm{\Psi }_{m+1}^{}\mathrm{}\mathrm{\Psi }_n^{})x$$ is a composite of factors $`\mathrm{\Psi }_j^{\prime \prime }(\mathrm{\Psi }_m^{\prime \prime }\mathrm{}\mathrm{\Psi }_{n1}^{\prime \prime })_h^\beta x`$ with $`\mathrm{\Psi }_r^{\prime \prime }`$ a generalised $`\mathrm{\Psi }_r`$. By Proposition 6.3, $`\mathrm{\Phi }_n^{}x`$ is then a composite of factors of the form $$\epsilon _1\mathrm{\Phi }_{m2}^{}(\psi _{m2}^{}\mathrm{}\psi _j^{})\mathrm{\Psi }_j^{\prime \prime }(\mathrm{\Psi }_m^{\prime \prime }\mathrm{}\mathrm{\Psi }_{n1}^{\prime \prime })_h^\beta x$$ with $`\mathrm{\Phi }_{m2}^{}`$ a generalised $`\mathrm{\Phi }_{m2}`$ and with $`\psi _k^{}`$ a generalised $`\psi _k`$. These factors have the form $`\epsilon _1\mathrm{\Phi }_{n1}^{}_h^\beta x`$ with $`\mathrm{\Phi }_{n1}^{}`$ a generalised $`\mathrm{\Phi }_{n1}`$, as required. $`\mathrm{}`$ We can now describe the interaction of $`\mathrm{\Phi }_n`$ with compositions and face maps in general terms as follows. ###### Proposition 6.8 If a composite $`x^{}_ix^+`$ exists, then $`\mathrm{\Phi }_n(x^{}_ix^+)`$ is naturally equal to a composite of factors $`\epsilon _1^{nm}\mathrm{\Phi }_mDx^\alpha `$ with $`D`$ an $`(nm)`$-fold product of face operators. If $`in`$ then $`_i^\alpha \mathrm{\Phi }_nx`$ is naturally equal to a composite of factors $`\epsilon _1^{nm1}\mathrm{\Phi }_mDx`$ with $`D`$ an $`(nm)`$-fold product of face operators. Proof The result for $`\mathrm{\Phi }_n(x^{}_ix^+)`$ comes from Proposition 6.3 by iterated application of Proposition 6.7; recall from (2.5) that $`\epsilon _1(y^{}_jy^+)`$ is a composite of $`\epsilon _1y^{}`$ and $`\epsilon _1y^+`$. Now suppose that $`in`$. By Proposition 3.3(iii) $$_i^\alpha \mathrm{\Phi }_nx=_i^\alpha \mathrm{\Phi }_i(\mathrm{\Phi }_{i+1}\mathrm{}\mathrm{\Phi }_n)x=\epsilon _1^{i1}(_1^\alpha )^i(\mathrm{\Phi }_{i+1}\mathrm{}\mathrm{\Phi }_n)x.$$ From Proposition 6.6, this is a composite of factors $`\epsilon _1^{i1}\mathrm{\Phi }_{ni}^{}D^{}x`$ with $`\mathrm{\Phi }_{ni}^{}`$ a generalised $`\mathrm{\Phi }_{ni}`$ and with $`D^{}`$ an $`i`$-fold product of face operators. By repeated application of Proposition 6.7, there is a further decomposition into factors $`\epsilon _1^{nm1}\mathrm{\Phi }_mDx`$ with $`D`$ an $`(nm)`$-fold product of face operators. This completes the proof. $`\mathrm{}`$ We will now specify the composites in Proposition 6.8 more precisely. Let $`G`$ be a cubical $`\omega `$-category, and consider $`\mathrm{\Phi }_n(x^{}_ix^+)`$, where $`x^{}x^+`$ is a composite in $`G_n`$. The factors $`\epsilon _1^{nm}\mathrm{\Phi }_mDx^\alpha `$ lie in $`\mathrm{\Phi }_n(G_n)`$ (see Proposition 3.7), and their composite can be regarded as a composite in the $`\omega `$-category $`\mathrm{\Phi }_n(G_n)`$ (see Theorem 3.8). To identify the composite, we take the universal case $$G=\lambda M(I^{i1}\times [0,2]\times I^{ni});$$ we may then identify $`\mathrm{\Phi }_n(G_n)`$ with $`M(I^{i1}\times [0,2]\times I^{ni})`$ by Theorem 4.1. The universal elements $$x^\alpha [\lambda M(I^{i1}\times [0,2]\times I^{ni})]_n=Hom[M(I^n),M(I^{i1}\times [0,2]\times I^{ni})]$$ are the inclusions $`\stackrel{ˇ}{\iota }_i^\alpha `$ representing $`M(I^{i1}\times [0,2]\times I^{ni})`$ as a push-out. Evaluating $`\mathrm{\Phi }_n(x^{}_ix^+)`$ and the corresponding composite on $`I^n`$, and using Proposition 3.2, we see that $`\mathrm{\Phi }_n(x^{}x^+)`$ gives us $`I^{i1}\times [0,2]\times I^{ni}`$ and the factors give us cells in $`I^{i1}\times [0,2]\times I^{ni}`$. The composite for $`\mathrm{\Phi }_n(x^{}_ix^+)`$ in Proposition 6.8 is an $`\omega `$-category formula expressing $`I^{i1}\times [0,2]\times I^{ni}`$ as a composite of cells. All such formulae are equivalent in all $`\omega `$-categories because of the presentation of $`M(I^{i1}\times [0,2]\times I^{ni})`$ in Theorem 1.3. The formula uses $`\#_p`$ only for $`0p<n`$ (see Theorem 3.8). Similarly, the formula for $`_i^\alpha \mathrm{\Phi }_nx`$ is an $`\omega `$-category formula expressing $`d_{ni}^\alpha I^n`$ as a composite of cells (see Propositions 3.6 and 3.2). In order to state these results more clearly, we introduce the following notation. ###### Definition 6.9 Let $`\sigma `$ be a cell in $`I^n`$, and let the dimension of $`\sigma `$ be $`m`$. Then $`_\sigma :G_mG_n`$ is the cubical $`\omega `$-category operation of the form $`_{i(1)}^{\alpha (1)}\mathrm{}_{i(nm)}^{\alpha (nm)}`$ such that the underlying homomorphism $$\stackrel{ˇ}{}_\sigma :M(I^m)M(I^n)$$ sends $`I^m`$ to $`\sigma `$. Note that $`_\sigma `$ is uniquely determined by $`\sigma `$, because of relation 2.1(i). In this notation, we can state the following theorem. ###### Theorem 6.10 (i) Let $`f`$ be a formula expressing $`I^{i1}\times [0,2]\times I^{ni}`$ as a $`(\#_0,\mathrm{},\#_{n1})`$-composite of cells $`\stackrel{ˇ}{\iota }_i^{}(\sigma )`$ and $`\stackrel{ˇ}{\iota }_i^+(\sigma )`$, where the $$\stackrel{ˇ}{\iota }_i^\alpha :M(I^n)M(I^{i1}\times [0,2]\times I^{ni})$$ are the inclusions expressing $`M(I^{i1}\times [0,2]\times I^{ni})`$ as a push-out. Let $`x^{}_ix^+`$ be a composite in a cubical $`\omega `$-category. Then $`\mathrm{\Phi }_n(x_iy)`$ can be got from $`f`$ by replacing $`\stackrel{ˇ}{\iota }_i^\alpha (\sigma )`$ with $`\epsilon _1^{nm}\mathrm{\Phi }_m_\sigma x^\alpha `$, where $`m=dim\sigma `$, and by replacing $`\#_p`$ with $`_{np}`$. (ii) Let $`g`$ be a formula expressing $`d_{ni}^\alpha I^n`$ as a $`(\#_0,\mathrm{},\#_{n2})`$-composite of cells, where $`1in`$. In a cubical $`\omega `$-category, $`_i^\alpha \mathrm{\Phi }_nx`$ can be got from $`g`$ by replacing $`\sigma `$ with $`\epsilon _1^{nm1}\mathrm{\Phi }_m_\sigma x`$, where $`m=dim\sigma `$, and by replacing $`\#_p`$ with $`_{np}`$. ## 7 The natural homomorphism $`B:G\lambda \gamma G`$ Let $`G`$ be a cubical $`\omega `$-category. We will now use Theorem 6.10 to construct a natural homomorphism $`B:G\lambda \gamma G`$. Let $`x`$ be a member of $`G_n`$. We must define $$B(x)(\lambda \gamma G)_n=Hom[M(I^n),\gamma G].$$ Now, $`M(I^n)`$ is generated by the cells in $`I^n`$ (see Theorem 1.3), and $`\gamma G`$ is the colimit of the sequence (see Definition 3.9). We can therefore define $`B(x)`$ by giving a suitable value to $`[B(x)](\sigma )`$ for $`\sigma `$ a cell in $`I^n`$; these values must lie in the $`\mathrm{\Phi }_m(G_m)`$, and a value $`\epsilon _1^sy`$ can be identified with $`y`$. The precise result is as follows. ###### Theorem 7.1 There is a natural homomorphism $`B:G\lambda \gamma G`$ for $`G`$ a cubical $`\omega `$-category given by $$[B(x)](\sigma )=\mathrm{\Phi }_m_\sigma x$$ for $`\sigma `$ a cell in $`I^n`$, where $`m=dim\sigma `$. Proof We first show that the values prescribed for the $`[B(x)](\sigma )`$ really define a homomorphism on $`M(I^n)`$; in other words, we must show that they respect the relations given in Theorem 1.3. Let $`\sigma `$ be an $`m`$-dimensional cell in $`I^n`$. We must show that $`d_m^\alpha (\mathrm{\Phi }_m_\sigma x)=\mathrm{\Phi }_m_\sigma x`$; if $`m>0`$ we must also show that $`d_{m1}^\alpha (\mathrm{\Phi }_m_\sigma x)`$ is the appropriate composite of the $`\mathrm{\Phi }_l_\tau x`$, where $`\tau \sigma `$. The first of these equations, $`d_m^\alpha (\mathrm{\Phi }_m_\sigma x)=\mathrm{\Phi }_m_\sigma x`$, is an immediate consequence of Theorem 3.8. For the second equation, let $`\sigma `$ be a cell of positive dimension $`m`$. By Theorem 3.8, $$d_{m1}^\alpha (\mathrm{\Phi }_m_\sigma x)=\epsilon _1_1^\alpha \mathrm{\Phi }_m_\sigma x,$$ which may be identified with $`_1^\alpha \mathrm{\Phi }_m_\sigma x`$. By Theorem 6.10, this is the appropriate composite of the $`\mathrm{\Phi }_l_\tau x`$, as required. We have now constructed functions $`B:G_n(\lambda \gamma G)_n`$, and we must show that these functions form a homomorphism of cubical $`\omega `$-categories. We must therefore show that $`B(_i^\alpha x)=_i^\alpha B(x)`$, that $`B(\epsilon _ix)=\epsilon _iB(x)`$, that $`B(x^{}_ix^+)=B(x^{})_iB(x^+)`$, and that $`B(\mathrm{\Gamma }_i^\alpha x)=\mathrm{\Gamma }_i^\alpha B(x)`$. First we consider $`B(_i^\alpha x)`$, where $`xG_n`$. Let $`\sigma `$ be a cell in $`I^{n1}`$ of dimension $`m`$, and let $`\tau =\stackrel{ˇ}{}_i^\alpha (\sigma )`$. We then have $`\tau =\stackrel{ˇ}{}_i^\alpha \stackrel{ˇ}{}_\sigma (I^m)`$, so $`\stackrel{ˇ}{}_\tau =\stackrel{ˇ}{}_i^\alpha \stackrel{ˇ}{}_\sigma `$ and $`_\tau =_\sigma _i^\alpha `$. It follows that $$[B(_i^\alpha x)](\sigma )=\mathrm{\Phi }_m_\sigma _i^\alpha x=\mathrm{\Phi }_m_\tau x$$ and $$[_i^\alpha B(x)](\sigma )=[B(x)][\stackrel{ˇ}{}_i^\alpha (\sigma )]=[B(x)](\tau )=\mathrm{\Phi }_m_\tau x;$$ Therefore $`[B(_i^\alpha x)](\sigma )=[_i^\alpha B(x)](\sigma )`$ as required. Next we consider $`B(\epsilon _ix)`$, where $`xG_n`$. Let $`\sigma `$ be a cell in $`I^{n+1}`$ of dimension $`m`$. From Definition 2.1, we see that $`_\sigma \epsilon _i`$ has the form $`id_\tau `$ or $`\epsilon _j_\tau `$. Let $`l=dim\tau `$, so that $`l=m`$ in the first case and $`l=m1`$ in the second case. Let $`\theta :G_lG_m`$ be $`id`$ or $`\epsilon _j:G_lG_m`$ as the case may be, and let $`\stackrel{ˇ}{\theta }:M(I^m)M(I^l)`$ be the underlying $`\omega `$-category homomorphism. We now see that $`_\sigma \epsilon _i=\theta _\tau `$ and $`\stackrel{ˇ}{\epsilon }_i\stackrel{ˇ}{}_\sigma =\stackrel{ˇ}{}_\tau \stackrel{ˇ}{\theta }`$ with $`\stackrel{ˇ}{\theta }(I^m)=I^l`$. It follows that $$\stackrel{ˇ}{\epsilon }_i(\sigma )=\stackrel{ˇ}{\epsilon }_i\stackrel{ˇ}{}_\sigma (I^m)=\stackrel{ˇ}{}_\tau \stackrel{ˇ}{\theta }(I^m)=\stackrel{ˇ}{}_\tau (I^l)=\tau .$$ Using Theorem 5.6, we also see that $`\mathrm{\Phi }_m\theta =\epsilon _1^{ml}\mathrm{\Phi }_l`$. We now get $$[B(\epsilon _ix)](\sigma )=\mathrm{\Phi }_m_\sigma \epsilon _ix=\mathrm{\Phi }_m\theta _\tau x=\epsilon _1^{ml}\mathrm{\Phi }_l_\tau x=\mathrm{\Phi }_l_\tau x$$ (recall that $`\epsilon _1^sy`$ is to be identified with $`y`$) and $$[\epsilon _iB(x)](\sigma )=[B(x)][\stackrel{ˇ}{\epsilon }_i(\sigma )]=[B(x)](\tau )=\mathrm{\Phi }_i_\tau x,$$ so that $`[B(\epsilon _ix)](\sigma )=[\epsilon _iB(x)](\sigma )`$ as required. Next we consider $`B(x^{}_ix^+)`$, where $`x^{}_ix^+`$ is a composite in $`G_n`$. Let $`\sigma `$ be a cell in $`I^n`$ of dimension $`m`$. From Definition 2.1, $$[B(x^{}_ix^+)](\sigma )=\mathrm{\Phi }_m_\sigma (x^{}_ix^+)$$ is equal to $`\mathrm{\Phi }_m_\sigma x^{}`$ or $`\mathrm{\Phi }_m_\sigma x^+`$ or to $`\mathrm{\Phi }_m(_\sigma x^{}_j_\sigma x^+)`$ for some $`j`$. In any case, using Theorem 6.10 if necessary, we see that $`[B(x^{}_ix^+)](\sigma )`$ is a composite of factors $`\mathrm{\Phi }_l_\tau x^\alpha `$ such that $`\stackrel{ˇ}{\mu }_i(\sigma )`$ is the corresponding composite of the $`\stackrel{ˇ}{\iota }_i^\alpha (\tau )`$, where $$\stackrel{ˇ}{\iota }_i^{},\stackrel{ˇ}{\iota }_i^+:M(I^n)M(I^{i1}\times [0,2]\times I^{ni})$$ are the functions expressing $`M(I^{i1}\times [0,2]\times I^{ni})`$ as a push-out. Let $$(B(x^{}),B(x^+)):M(I^{i1}\times [0,2]\times I^{ni})\gamma G$$ be the function such that $$(B(x^{}),B(x^+))\stackrel{ˇ}{\iota }_i^\alpha =B(x^\alpha );$$ we see that $$[B(x^{}_ix^+)](\sigma )=(B(x^{}),B(x^+))\stackrel{ˇ}{\mu }_i(\sigma )=[B(x^{})_iB(x^+)](\sigma )$$ as required. Finally we consider $`B(\mathrm{\Gamma }_i^\alpha x)`$, where $`xG_n`$. Let $`\sigma `$ be a cell in $`I^{n+1}`$ of dimension $`m`$. From Definition 2.1, $`_\sigma \mathrm{\Gamma }_i^\alpha `$ has the form $`_\tau `$ or $`\epsilon _i_\tau `$ or $`\mathrm{\Gamma }_i^\alpha _\tau `$. We can now use the same argument as for $`B(\epsilon _ix)`$, noting that $`\stackrel{ˇ}{\mathrm{\Gamma }}_i^\alpha (I^m)=I^{m1}`$ and that $`\mathrm{\Phi }_m\mathrm{\Gamma }_i^\alpha =\epsilon _1\mathrm{\Phi }_{m1}`$ by Theorem 5.6. This completes the proof. $`\mathrm{}`$ ## 8 The natural isomorphism $`B:G\lambda \gamma G`$ In Theorem 4.1, we have constructed a natural isomorphism $`A:\gamma \lambda XX`$ for $`X`$ an $`\omega `$-category. In Theorem 7.1 we have constructed a natural homomorphism $`B:G\lambda \gamma G`$ for $`G`$ a cubical $`\omega `$-category. We will now show that $`\omega `$-categories and cubical $`\omega `$-categories are equivalent by showing that $`B`$ is an isomorphism. We begin with the following observation. ###### Proposition 8.1 Let $`G`$ be a cubical $`\omega `$-category. Then $`\gamma B:\gamma G\gamma \lambda \gamma G`$ is an isomorphism. Proof Consider the composite $$A(\gamma B):\gamma G\gamma G.$$ By Theorem 4.1, $`A`$ is an isomorphism; it therefore suffices to show that the composite $`A(\gamma B)`$ is the identity. This amounts to showing that $`AB(x)=x`$ for $`x\mathrm{\Phi }_n(G_n)`$. Now, from the definitions of $`A`$ and $`B`$, we find that $$AB(x)=[B(x)](I^n)=\mathrm{\Phi }_nx;$$ since $`x\mathrm{\Phi }_n(G_n)`$ and $`\mathrm{\Phi }_n`$ is idempotent (Proposition 3.5), it follows that $`AB(x)=x`$ as required. This completes the proof. $`\mathrm{}`$ Because of Proposition 8.1, to show that $`B`$ is an isomorphism it suffices to show that a cubical $`\omega `$-category $`G`$ is determined by the $`\omega `$-category $`\gamma G`$. Because of Remark 3.10, this is the same as showing that $`G`$ is determined by the $`\mathrm{\Phi }_n(G_n)`$. We will work inductively, showing that an element $`x`$ of $`G_n`$ is determined by $`\mathrm{\Phi }_nx`$ and by its faces. To handle the family of faces of $`x`$, we will use the following terminology. ###### Definition 8.2 Let $`G`$ be a cubical $`\omega `$-category and let $`n`$ be a positive integer. An $`n`$-shell in $`G`$ an ordered $`(2n)`$-tuple $$z=(z_1^{},z_1^+,\mathrm{},z_n^{},z_n^+)$$ of members of $`G_{n1}`$ such that $`_i^\alpha z_j^\beta =_{j1}^\beta z_i^\alpha `$ whenever $`i<j`$. The set of $`n`$-shells is denoted $`\mathrm{}G_{n1}`$. ###### Remark 8.3 This construction is used in \[9, Section 5\] to construct a coskeleton functor from $`(n1)`$-truncated cubical $`\omega `$-groupoids to $`n`$-truncated $`\omega `$-groupoids determined by $$(G_0,G_1,\mathrm{},G_{n1})(G_0,G_1,\mathrm{},G_{n1},\mathrm{}G_{n1}),$$ and the same construction clearly works for the category case. It follows that the folding operations are also defined on $`\mathrm{}G_{n1}`$. In the following we take a slightly more direct route. First, by Definition 2.1, it is easy to check the following result. ###### Proposition 8.4 Let $`G`$ be a cubical $`\omega `$-category and let $`n`$ be a positive integer. There is a boundary map $`:G_n\mathrm{}G_{n1}`$ given by $$x=(_1^{}x,_1^+x,\mathrm{},_n^{}x,_n^+x).$$ Now we define folding operations on shells directly. ###### Proposition 8.5 Let $`G`$ be a cubical $`\omega `$-category. For $`1jn1`$ the cubical structure of $`(G_0,\mathrm{},G_{n1})`$ yields a natural function $`\psi _j:\mathrm{}G_{n1}\mathrm{}G_{n1}`$ such that $$\psi _j=\psi _j:G_n\mathrm{}G_{n1}.$$ Proof Let $`z=(z_i^\alpha )`$ be an $`n`$-shell. Guided by Proposition 3.3(i), we let $`\psi _iz`$ be the $`(2n)`$-tuple $`w=(w_i^\alpha )`$ such that $$w_i^\alpha =\{\begin{array}{cc}\psi _{j1}z_i^\alpha \hfill & \text{for }i<j,\hfill \\ z_j^{}_jz_{j+1}^+\hfill & \text{for }(\alpha ,i)=(,j),\hfill \\ z_{j+1}^{}_jz_j^+\hfill & \text{for }(\alpha ,i)=(+,j),\hfill \\ \epsilon _j_j^\alpha z_{j+1}^\alpha \hfill & \text{for }i=j+1,\hfill \\ \psi _jz_i^\alpha \hfill & \text{for }i>j+1.\hfill \end{array}$$ From Proposition 3.1(i) and the identities in Section 2, it is straightforward to check that $`\psi _j`$ is a well-defined function from $`\mathrm{}G_{n1}`$ to itself, and it is easy to see that $`\psi _j=\psi _j`$. $`\mathrm{}`$ We will now show that the $`n`$-dimensional elements $`(n>0)`$ in a cubical $`\omega `$-category are determined by the lower-dimensional elements and by the image of $`\mathrm{\Phi }_n`$. ###### Theorem 8.6 Let $`G`$ be a cubical $`\omega `$-category, let $`n`$ be a positive integer, and let $`\mathrm{\Phi }_n:\mathrm{}G_{n1}\mathrm{}G_{n1}`$ be the function given by $$\mathrm{\Phi }_n=\psi _1(\psi _2\psi _1)(\psi _3\psi _2\psi _1)\mathrm{}(\psi _{n1}\mathrm{}\psi _1).$$ Then there is a bijection $`x(x,\mathrm{\Phi }_nx)`$ from $`G_n`$ to the pull-back $$\mathrm{}G_{n1}\times _{G_n}\mathrm{\Phi }_n(G_n)=\{(z,y)\mathrm{}G_{n1}\times \mathrm{\Phi }_n(G_n):\mathrm{\Phi }_nz=y\}.$$ Proof This amounts to showing that $$\mathrm{\Phi }_n:^1(z)^1(\mathrm{\Phi }_nz)$$ is a bijection for each $`z`$ in $`\mathrm{}G_{n1}`$. Since $`\mathrm{\Phi }_n`$ is a composite of operators $`\psi _j`$, it suffices to show that $$\psi _j:^1(z)^1(\psi _jz)$$ is a bijection for each $`z`$ in $`\mathrm{}G_{n1}`$. Given $`y^1(\psi _jz)`$, it is straightforward to check that there is a composite $$\theta y=(\epsilon _jz_j^{}_{j+1}\mathrm{\Gamma }_j^+z_{j+1}^+)_jy_j(\mathrm{\Gamma }_j^{}z_{j+1}^{}_{j+1}\epsilon _jz_j^+)$$ and that $`\theta y^1(z)`$. We will carry out the proof by showing that $`\theta \psi _jx=x`$ for $`x^1(z)`$ and that $`\psi _j\theta y=y`$ for $`y^1(\psi _jz)`$. Let $`x`$ be a member of $`^1(z)`$. Then $$\theta \psi _jx=\left[\begin{array}{ccc}\|& \mathrm{\Gamma }_j^+z_{j+1}^{}& \mathrm{\Gamma }_j^{}z_{j+1}^{}\\ \epsilon _jz_j^{}& x& \epsilon _jz_j^+\\ \mathrm{\Gamma }_j^+z_{j+1}^+& \mathrm{\Gamma }_j^{}z_{j+1}^+& \|\end{array}\right]\text{}.$$ The first and third rows are in the image of $`\epsilon _{j+1}`$ by (2.5) and (2.7), so they are identities for $`_{j+1}`$ and can therefore be omitted. This leaves the second row in which $`\epsilon _jz_j^{}`$ and $`\epsilon _jz_j^+`$ are identities for $`_j`$. It follows that $`\theta \psi _jx=x`$. Now let $`y`$ be a member of $`^1(\psi _jz)`$. By (2.2)(vi) and (2.1)(ii), $`\epsilon _j_j^{}\mathrm{\Gamma }_j^+=\epsilon _j\epsilon _j_j^{}=\epsilon _{j+1}\epsilon _j_j^{}`$, so $`\mathrm{\Gamma }_j^+_{j+1}^{}\theta y`$ $`=\mathrm{\Gamma }_j^+z_{j+1}^{}`$ $`=\epsilon _j_j^{}\mathrm{\Gamma }_j^+z_{j+1}^{}_j\epsilon _j_j^{}\mathrm{\Gamma }_j^+z_{j+1}^{}_j\mathrm{\Gamma }_j^+z_{j+1}^{}`$ $`=\epsilon _{j+1}\epsilon _j_j^{}z_{j+1}^{}_j\epsilon _{j+1}\epsilon _j_j^{}z_{j+1}^{}_j\mathrm{\Gamma }_j^+z_{j+1}^{}.`$ Similarly $`\mathrm{\Gamma }_j^{}_{j+1}^+\theta y`$ $`=\mathrm{\Gamma }_j^{}z_{j+1}^+_j\epsilon _{j+1}\epsilon _j_j^+z_{j+1}^+_j\epsilon _{j+1}\epsilon _j_j^+z_{j+1}^+.`$ It follows that $`\psi _j\theta y`$ $`=\mathrm{\Gamma }_j^+_{j+1}^{}\theta y_{j+1}\theta y_{j+1}\mathrm{\Gamma }_j^{}_{j+1}^+\theta y`$ $`=\left[\begin{array}{ccc}\|& \|& \mathrm{\Gamma }_j^+z_{j+1}^{}\\ \epsilon _jz_j^{}_{j+1}\mathrm{\Gamma }+_jz_{j+1}^+& y& \mathrm{\Gamma }_j^{}z_{j+1}^{}_{j+1}\epsilon _jz_j^+\\ \mathrm{\Gamma }_j^{}z_{j+1}^+& \|& \|\end{array}\right]\text{}.`$ By (2.7) and (2.5) the first and third columns are in the image of $`\epsilon _j`$, so they are identities for $`_j`$ and can be omitted. This leaves the second column so that $`\psi _j\theta y=y`$. This completes the proof. $`\mathrm{}`$ From Theorem 8.6 we deduce the following result. ###### Theorem 8.7 Let $`f:GH`$ be a morphism of cubical $`\omega `$-categories such that $`\gamma f:\gamma G\gamma H`$ is an isomorphism. Then $`f`$ is an isomorphism. Proof By Remark 3.10, $`f`$ induces isomorphisms from $`\mathrm{\Phi }_n(G_n)`$ to $`\mathrm{\Phi }_n(H_n)`$. Since $`\mathrm{\Phi }_0`$ is the identity operation, $`f`$ induces a bijection from $`G_0`$ to $`H_0`$. By an inductive argument using Theorem 8.6, $`f`$ induces a bijection from $`G_n`$ to $`H_n`$ for all $`n`$. Therefore $`f`$ is an isomorphism. $`\mathrm{}`$ It follows from Proposition 8.1 and Theorem 8.7 that $`B:G\lambda \gamma G`$ is a natural isomorphism for cubical $`\omega `$-categories $`G`$. From Theorem 4.1, $`A:\gamma \lambda XX`$ is a natural isomorphism for $`\omega `$-categories $`X`$. We draw the following conclusion. ###### Theorem 8.8 The categories of $`\omega `$-categories and of cubical $`\omega `$-categories are equivalent under the functors $`\lambda `$ and $`\gamma `$. ## 9 Thin elements and commutative shells in a cubical $`\omega `$-category In this section we use the equivalence of Theorem 8.8 to clarify two concepts in the theory of cubical $`\omega `$-categories: thin elements and commutative shells. Thin elements (sometimes called hollow elements) were introduced in the thesis of K. Dakin , and were developed in the cubical $`\omega `$-groupoid context by Brown and Higgins in . They are used by Ashley in and by Street in . In the cubical nerve of an $`\omega `$-category they arise as follows. Throughout this section, let $`G`$ be a cubical $`\omega `$-category. Whenever convenient, we will identify $`G`$ with the nerve of $`\gamma G`$; in other words, each element $`x`$ of $`G_n`$ is identified with a homomorphism $`x:M(I^n)\gamma G`$. First we deal with thin elements. Intuitively, an element is thin if its real dimension is less than its apparent dimension. In the nerve of an $`\omega `$-category we can make this precise as follows. ###### Definition 9.1 Let $`x`$ be a member of $`G_n`$. Then $`x`$ is thin if $$\mathrm{dim}x(I^n)<n.$$ Given an element $`x`$ of $`G_n`$, we can identify $`x(I^n)`$ with $`\mathrm{\Phi }_nx`$ by Theorem 7.1. By Remark 3.10, dim $`\mathrm{\Phi }_nx<n`$ if and only if $`\mathrm{\Phi }_nx`$ is in the image of $`\epsilon _1`$. We therefore have the following characterisation. ###### Proposition 9.2 Let $`x`$ be a member of $`G_n`$. Then $`x`$ is thin if and only if $`\mathrm{\Phi }_nx`$ is in the image of $`\epsilon _1`$. There is also a less obvious characterisation in more elementary cubical terms: the thin elements of $`G_n`$ are those generated by the $`G_m`$ with $`m<n`$. The precise statement is as follows. ###### Theorem 9.3 Let $`x`$ be a member of $`G_n`$. Then $`x`$ is thin if and only if it is a composite of elements of the forms $`\epsilon _iy`$ and $`\mathrm{\Gamma }_j^\alpha z`$ for various values of $`i,j,\alpha ,y,z`$. Proof Suppose that $`x`$ is a composite of elements of the forms $`\epsilon _iy`$ and $`\mathrm{\Gamma }_j^\alpha z`$. Then $`\mathrm{\Phi }_nx`$ is in the image of $`\epsilon _1`$ by Theorems 5.6 and 6.10, so x is thin by Proposition 9.2. Conversely, suppose that $`x`$ is thin. It follows from the proof of Theorem 8.5 that $`x`$ is a composite of $`\mathrm{\Phi }_nx`$ with elements of the forms $`\epsilon _iy`$ and $`\mathrm{\Gamma }_j^\alpha z`$. By Proposition 9.2, $`\mathrm{\Phi }_nx`$ is in the image of $`\epsilon _1`$, so $`x`$ is itself a composite of elements of the forms $`\epsilon _iy`$ and $`\mathrm{\Gamma }_j^\alpha z`$. $`\mathrm{}`$ Next we deal with commutative shells. There is an obvious concept of commutative square, or commutative 2-shell; we want commutative $`n`$-shells for arbitrary positive $`n`$. Now an $`n`$-shell $`z`$ in $`G`$ can be identified with a homomorphism $`z:M(d_{n1}^{}I^nd_{n1}^+I^n)\gamma G`$, and we must obviously define a commutative $`n`$-shell as follows. ###### Definition 9.4 For $`n>0`$ an $`n`$-shell $`z`$ in $`G`$ is commutative if $$z(d_{n1}^{}I^n)=z(d_{n1}^+I^n).$$ By Theorem 6.10(ii), if $`z`$ is an $`n`$-shell with $`n>0`$ then $`z(d_{n1}^\alpha )`$ can be identified with $`(\mathrm{\Phi }_nz)_1^\alpha `$, the $`(\alpha ,1)`$ face of the $`n`$-shell $`\mathrm{\Phi }_nz`$ as in theorem 8.6. We can therefore describe commutative $`n`$-shells in cubical terms as follows. ###### Proposition 9.5 For $`n>0`$ an $`n`$-shell $`z`$ in $`G`$ is commutative if and only if $$(\mathrm{\Phi }_nz)_1^{}=(\mathrm{\Phi }_nz)_1^+.$$ ## 10 Monoidal closed structures In Al-Agl and Steiner constructed a monoidal closed structure on the category $`\omega \text{-}\mathrm{𝖢𝖺𝗍}^{}`$ of (globular) $`\omega `$-categories by using a cubical description of that category. Now that we have a more explicit cubical description we can give a more explicit description of the monoidal closed structure; we modify the construction which is given by Brown and Higgins in for the case of a single connection and for groupoids rather than categories. Following the method there, we first define the closed structure on the category $`\omega \text{-}\mathrm{𝖢𝖺𝗍}^{\mathrm{}}`$ of cubical $`\omega `$-categories using a notion of $`n`$-fold left homotopy which we outline below, and then obtain the tensor product as the adjoint to the closed structure. This gives: ###### Theorem 10.1 The category $`\omega \text{-}\mathrm{𝖢𝖺𝗍}^{\mathrm{}}`$ admits a monoidal closed structure with an adjoint relationship $$\omega \text{-}\mathrm{𝖢𝖺𝗍}^{\mathrm{}}(GH,K)\omega \text{-}\mathrm{𝖢𝖺𝗍}^{\mathrm{}}(G,\omega \text{-}\mathrm{𝖢𝖠𝖳}^{\mathrm{}}(H,K))$$ in which $`\omega \text{-}\mathrm{𝖢𝖠𝖳}^{\mathrm{}}(H,K)_0`$ is the set of morphisms $`HK`$, while for $`n1`$ $`\omega \text{-}\mathrm{𝖢𝖠𝖳}^{\mathrm{}}(H,K)_n`$ is the set of $`n`$-fold left homotopies $`HK`$. The proof is given below. Because of the equivalence between $`\omega \text{-}\mathrm{𝖢𝖺𝗍}^{\mathrm{}}`$ and the category $`\omega \text{-}\mathrm{𝖢𝖺𝗍}^{}`$ of $`\omega `$-categories we then have: ###### Corollary 10.2 The category $`\omega \text{-}\mathrm{𝖢𝖺𝗍}^{}`$ admits a monoidal closed structure with an adjoint relationship $$\omega \text{-}\mathrm{𝖢𝖺𝗍}^{}(XY,Z)\omega \text{-}\mathrm{𝖢𝖺𝗍}^{}(X,\omega \text{-}\mathrm{𝖢𝖠𝖳}^{}(Y,Z))$$ in which $`\omega ^{}\text{-}\mathrm{𝖢𝖠𝖳}(Y,Z)_0`$ is the set of morphisms $`YZ`$, while for $`n1`$ $`\omega \text{-}\mathrm{𝖢𝖠𝖳}^{}(Y,Z)_n`$ is the set of $`n`$-fold left homotopies $`YZ`$ corresponding to the cubical homotopies. The tensor product in Corollary 10.2 is an extension of the tensor product in Theorem 1.8. Note that by Remark 3.11, the set $`\omega \text{-}\mathrm{𝖢𝖠𝖳}^{}(Y,Z)_n`$ of globular $`n`$-fold left homotopies may be thought of as an explicitly described subset of the set of cubical $`n`$-fold left homotopies $`\lambda Y\lambda Z`$. Because of the complications of the folding operations, explicit descriptions of the globular monoidal closed structure are not so easy, but have been partly accomplished by Steiner in . See also Crans . We now give details of these cubical constructions, following directly the methods of . Let $`H`$ be a cubical $`\omega `$-category and $`n`$ be a non-negative integer. We can construct a cubical $`\omega `$-category $`P^nH`$ called the $`n`$-fold (left) path cubical $`\omega `$-category of $`H`$ as follows: $`(P^nH)_r=H_{n+r}`$; the operations $`_i^\alpha `$, $`\epsilon _i`$, $`\mathrm{\Gamma }_i^\alpha `$ and $`_i`$ of $`P^nH`$ are the operations $`_{n+i}^\alpha `$, $`\epsilon _{n+i}`$, $`\mathrm{\Gamma }_{n+i}^\alpha `$ and $`_{n+i}`$ of $`H`$. The operations $`_1^\alpha ,\mathrm{},_n^\alpha `$ not used in $`P^nH`$ give us morphisms of cubical $`\omega `$-categories from $`P^nH`$ to $`P^{n1}H`$, etc., and we get an internal cubical $`\omega `$-category $$𝖯H=(H,P^1H,P^2H,\mathrm{})$$ in the category $`\omega `$-$`\mathrm{𝖢𝖺𝗍}^{\mathrm{}}`$. For any cubical $`\omega `$-categories $`G,H`$ we now define $$\omega \text{-}\mathrm{𝖢𝖠𝖳}^{\mathrm{}}(G,H)=\omega \text{-}\mathrm{𝖢𝖺𝗍}^{\mathrm{}}(G,𝖯H);$$ that is, $`\omega \text{-}\mathrm{𝖢𝖠𝖳}_m^{\mathrm{}}(G,H)=\omega \text{-}\mathrm{𝖢𝖺𝗍}^{\mathrm{}}(G,P^mH)`$, and the cubical $`\omega `$-category structure on $`\omega \text{-}\mathrm{𝖢𝖠𝖳}_m^{\mathrm{}}(G,H)`$ is induced by the internal cubical $`\omega `$-category structure on $`𝖯H`$. Ultimately, this means that the operations $`_i^\alpha `$, etc. on $`\omega \text{-}\mathrm{𝖢𝖠𝖳}_m^{\mathrm{}}(G,H)`$ are induced by the similarly numbered operations on $`H`$. In dimension $`0`$, $`\omega \text{-}\mathrm{𝖢𝖠𝖳}^{\mathrm{}}(G,H)`$ consists of all morphisms $`GH`$, while in dimension $`n`$ it consists of $`n`$-fold (left) homotopies $`GH`$. We make $`\omega `$-$`\mathrm{𝖢𝖠𝖳}^{\mathrm{}}`$$`(G,H)`$ a functor in $`G`$ and $`H`$ (contravariant in $`G`$) in the obvious way. The definition of tensor product of cubical $`\omega `$-categories is harder. We require that $`G`$ be left adjoint to $`\omega `$-$`\mathrm{𝖢𝖠𝖳}^{\mathrm{}}`$$`(G,)`$ as a functor from $`\omega `$-$`\mathrm{𝖢𝖺𝗍}^{\mathrm{}}`$ to $`\omega `$-$`\mathrm{𝖢𝖺𝗍}^{\mathrm{}}`$, and this determines $``$ up to natural isomorphism. Its existence, that is, the representability of the functor $`\omega \text{-}\mathrm{𝖢𝖺𝗍}^{\mathrm{}}(F,\omega \text{-}\mathrm{𝖢𝖠𝖳}^{\mathrm{}}(G,))`$, can be asserted on general grounds. Indeed, $`\omega `$-$`\mathrm{𝖢𝖺𝗍}^{\mathrm{}}`$ is an equationally defined category of many sorted algebras in which the domains of the operations are defined by finite limit diagrams, and general theorems on such algebraic categories imply that $`\omega `$-$`\mathrm{𝖢𝖺𝗍}^{\mathrm{}}`$ is complete and cocomplete. We can also specify the tensor product cubical $`\omega `$-category by a presentation; that is, we give a set of generators in each dimension and a set of relations of the form $`u=v`$, where $`u,v`$ are well formed formulae of the same dimension made from generators and the operators $`_i^\alpha ,\epsilon _i,\mathrm{\Gamma }_i^\alpha ,_i`$. This is analogous to the standard tensor product of modules over a ring, and the universal property of the presentation gives the required adjointness. The details are as follows. ###### Definition 10.3 Let $`F,G`$ be cubical $`\omega `$-categories. Then $`FG`$ is the cubical $`\omega `$-category generated by elements in dimension $`n0`$ of the form $`xy`$ where $`xF_p,yG_q`$ and $`p+q=n`$, subject to the following defining relations (plus, of course, the laws for cubical $`\omega `$-categories): 1. $`_i^\alpha (xy)=\{\begin{array}{cc}(_i^\alpha x)y\hfill & \text{if }1ip,\hfill \\ x(_{ip}^\alpha y)\hfill & \text{if }p+1in;\hfill \end{array}`$ 2. $`\epsilon _i(xy)=\{\begin{array}{cc}(\epsilon _ix)y\hfill & \text{if }1ip+1,\hfill \\ x(\epsilon _{ip}y)\hfill & \text{if }p+1in+1;\hfill \end{array}`$ 3. $`\mathrm{\Gamma }_i^\alpha (xy)=\{\begin{array}{cc}(\mathrm{\Gamma }_i^\alpha x)y\hfill & \text{if }1ip,\hfill \\ x(\mathrm{\Gamma }_{ip}^\alpha y)\hfill & \text{if }p+1in;\hfill \end{array}`$ 4. $`(x_ix^{})y=(xy)_i(x^{}y)`$ if $`1ip,`$ and $`x_ix^{}`$ is defined in $`F;`$ 5. $`x(y_jy^{})=(xy)_{p+j}(xy^{})`$ if $`1jq,`$ and $`y_jy^{}`$ is defined in $`G;`$ We note that the relation 1. $`(\epsilon _{p+1}x)y=x(\epsilon _1y)`$ follows from (ii). An alternative way of stating this definition is to define a bimorphism $`(F,G)A`$, where $`F,G,A`$ are cubical $`\omega `$-categories, to be a family of maps $`F_p\times G_qA_{p+q}(p,q0)`$, denoted by $`(x,y)\chi (x,y)`$ such that 1. for each $`xF_p`$, the map $`y\chi (x,y)`$ is a morphism of cubical $`\omega `$-categories $`GP^pA`$; 2. for each $`gG_q`$ the map $`x\chi (x,y)`$ is a morphism of cubical $`\omega `$-categories $`FTP^qTA`$, where the cubical $`\omega `$-category $`TX`$ has the same elements as $`X`$ but its cubical operations, connections and compositions are numbered in reverse order. The cubical $`\omega `$-category $`FG`$ is now defined up to natural isomorphisms by the two properties 1. the map $`(x,y)xy`$ is a bimorphism $`(F,G)FG`$; 2. every bimorphism $`(F,G)A`$ is uniquely of the form $`(x,y)\sigma (xy)`$ where $`\sigma :FGA`$ is a morphism of cubical $`\omega `$-categories. In the definition of a bimorphism $`(F,G)A`$, condition (a) gives maps $`F_p\omega \text{-}\mathrm{𝖢𝖠𝖳}_p^{\mathrm{}}(G,A)`$ for each $`p`$, and condition (b) states that these combine to give a morphism of cubical $`\omega `$-categories $`F\omega \text{-}\mathrm{𝖢𝖠𝖳}^{\mathrm{}}(G,A)`$. This observation yields a natural bijection between bimorphisms $`(F,G)A`$ and morphisms $`F\omega \text{-}\mathrm{𝖢𝖠𝖳}^{\mathrm{}}(G,A)`$. Since we also have a natural bijection between bimorphisms $`(F,G)A`$ and morphisms $`FGA`$, we have ###### Proposition 10.4 The functor $`G`$ is left adjoint to the functor $`\omega \text{-}\mathrm{𝖢𝖠𝖳}^{\mathrm{}}(G,)`$ from $`\omega \text{-}\mathrm{𝖢𝖺𝗍}^{\mathrm{}}`$ to $`\omega \text{-}\mathrm{𝖢𝖺𝗍}^{\mathrm{}}`$. $`\mathrm{}`$ ###### Proposition 10.5 For cubical $`\omega `$-categories $`F,G,H`$, there are natural isomorphisms of cubical $`\omega `$-categories 1. $`(FG)HF(GH)`$, and 2. $`\omega \text{-}\mathrm{𝖢𝖠𝖳}^{\mathrm{}}(FG,H)\omega \text{-}\mathrm{𝖢𝖠𝖳}^{\mathrm{}}(F,\omega \text{-}\mathrm{𝖢𝖠𝖳}^{\mathrm{}}(G,H))`$ giving $`\omega \text{-}\mathrm{𝖢𝖺𝗍}^{\mathrm{}}`$ the structure of a monoidal closed category. Proof (ii) In dimension $`r`$ there is by adjointness a natural bijection $`\omega \text{-}\mathrm{𝖢𝖠𝖳}_r^{\mathrm{}}(FG,H)`$ $`=\omega \text{-}\mathrm{𝖢𝖺𝗍}^{\mathrm{}}(FG,P^rH)`$ $`\omega \text{-}\mathrm{𝖢𝖺𝗍}^{\mathrm{}}(F,\omega \text{-}\mathrm{𝖢𝖠𝖳}^{\mathrm{}}(G,P^rH))`$ $`=\omega \text{-}\mathrm{𝖢𝖺𝗍}^{\mathrm{}}(F,P^r(\omega \text{-}\mathrm{𝖢𝖠𝖳}^{\mathrm{}}(G,H)))`$ $`=\omega \text{-}\mathrm{𝖢𝖠𝖳}_r^{\mathrm{}}(F,\omega \text{-}\mathrm{𝖢𝖠𝖳}^{\mathrm{}}(G,H)).`$ These bijections combine to form the natural isomorphism (ii) of cubical $`\omega `$-categories because, on both sides, the cubical $`\omega `$-category structures are induced by the corresponding operators $`_i^\alpha ,\epsilon _j`$, etc. in $`H`$. (i) This isomorphism may be proved directly, or, as is well known, be deduced from the axioms for a monoidal closed category. $`\mathrm{}`$ We can also relate the construction to the category of cubical sets, which we denote $`\mathrm{𝖢𝗎𝖻}`$. The underlying cubical set functor $`U:\omega \text{-}\mathrm{𝖢𝖺𝗍}^{\mathrm{}}\mathrm{𝖢𝗎𝖻}`$ has a left adjoint $`\sigma :\mathrm{𝖢𝗎𝖻}\omega \text{-}\mathrm{𝖢𝖺𝗍}^{\mathrm{}}`$, and we call $`\sigma (K)`$ the free cubical $`\omega `$-category on the cubical set $`K`$. The category $`\mathrm{𝖢𝗎𝖻}`$ has a monoidal closed structure in the same way as $`\omega \text{-}\mathrm{𝖢𝖺𝗍}^{\mathrm{}}`$ (see ); the internal hom $`\mathrm{𝖢𝖴𝖡}`$ is given by $`\mathrm{𝖢𝖴𝖡}(L,M)_r=\mathrm{𝖢𝗎𝖻}(L,P^rM)`$ where $`P^r`$ is now the $`n`$-fold path functor on cubical sets. We have the following results. ###### Proposition 10.6 For a cubical set $`L`$ and cubical $`\omega `$-category $`G`$, there is a natural isomorphism of cubical sets $$U(\omega \text{-}\mathrm{𝖢𝖠𝖳}^{\mathrm{}}(\sigma (L),G))\mathrm{𝖢𝖴𝖡}(L,UG).$$ Proof The functor $`\sigma :\mathrm{𝖢𝗎𝖻}\omega \text{-}\mathrm{𝖢𝖺𝗍}^{\mathrm{}}`$ is left adjoint to $`U:\omega \text{-}\mathrm{𝖢𝖺𝗍}^{\mathrm{}}\mathrm{𝖢𝗎𝖻}`$, and this is what the proposition says in dimension 0. In dimension $`r`$ we have a natural bijection $`\omega \text{-}\mathrm{𝖢𝖠𝖳}_r^{\mathrm{}}(\sigma (L),G)`$ $`=\omega \text{-}\mathrm{𝖢𝖺𝗍}^{\mathrm{}}(\sigma (L),P^rG)`$ $`\mathrm{𝖢𝗎𝖻}(L,UP^rG)`$ $`=\mathrm{𝖢𝖴𝖡}_r(L,UG)`$ and these bijections are compatible with the cubical operators. $`\mathrm{}`$ ###### Proposition 10.7 If $`K,L`$ are cubical sets, there is a natural isomorphism of cubical $`\omega `$-categories $$\sigma (K)\sigma (L)\sigma (KL).$$ Proof For any cubical $`\omega `$-category $`G`$, there are natural isomorphisms of cubical sets $`U(\omega \text{-}\mathrm{𝖢𝖠𝖳}^{\mathrm{}}(\sigma (K)\sigma (L),G))`$ $`U(\omega \text{-}\mathrm{𝖢𝖠𝖳}^{\mathrm{}}(\sigma (K),\omega \text{-}\mathrm{𝖢𝖠𝖳}^{\mathrm{}}(\sigma (L),G)))`$ $`\mathrm{𝖢𝖴𝖡}(K,U(\omega \text{-}\mathrm{𝖢𝖠𝖳}^{\mathrm{}}(\sigma (L),G)))`$ $`\mathrm{𝖢𝖴𝖡}(K,\mathrm{𝖢𝖴𝖡}(L,UG))`$ $`\mathrm{𝖢𝖴𝖡}(KL,UG)`$ $`U(\omega \text{-}\mathrm{𝖢𝖠𝖳}^{\mathrm{}}(\sigma (KL),G).`$ The proposition follows from the information in dimension $`0`$, namely $`\omega \text{-}\mathrm{𝖢𝖺𝗍}^{\mathrm{}}(\sigma (KL),G)`$ $`\omega \text{-}\mathrm{𝖢𝖺𝗍}^{\mathrm{}}(\sigma (K)\sigma (L),G)).`$ $`\mathrm{}`$ The $`\omega `$-categories $`M(I^n)`$ of Section 1 can be fitted into this framework if one regards them as cubical $`\omega `$-categories. Indeed, as a cubical $`\omega `$-category, $`M(I^n)`$ is freely generated by one element in dimension $`n`$; therefore $`M(I^n)=\sigma (𝕀^n)`$ where $`𝕀^n`$ is the cubical set freely generated by one element in dimension $`n`$. Calculations with cubical sets show that $`𝕀^m𝕀^n𝕀^{m+n}`$, and we get the following result. ###### Corollary 10.8 The are natural isomorphisms of cubical $`\omega `$-categories $$M(I^m)M(I^n)M(I^{m+n}).$$ $`\mathrm{}`$ ###### Proposition 10.9 1. $`M(I^n)`$ is left adjoint to $`P^n:\omega \text{-}\mathrm{𝖢𝖺𝗍}^{\mathrm{}}\omega \text{-}\mathrm{𝖢𝖺𝗍}^{\mathrm{}}`$. 2. $`M(I^n)`$ is left adjoint to $`\omega \text{-}\mathrm{𝖢𝖠𝖳}^{\mathrm{}}(M(I^n),)`$. 3. $`\omega \text{-}\mathrm{𝖢𝖠𝖳}^{\mathrm{}}(M(I^n),)`$ is naturally isomorphic to $`TP^nT`$. Proof (i) There are natural bijections $`\omega \text{-}\mathrm{𝖢𝖺𝗍}^{\mathrm{}}(M(I^n)H,K)`$ $`\omega \text{-}\mathrm{𝖢𝖺𝗍}^{\mathrm{}}(M(I^n),\omega \text{-}\mathrm{𝖢𝖠𝖳}^{\mathrm{}}(H,K))`$ $`\omega \text{-}\mathrm{𝖢𝖠𝖳}_n^{\mathrm{}}(H,K)`$ $`=\omega \text{-}\mathrm{𝖢𝖺𝗍}^{\mathrm{}}(H,P^nK).`$ (ii) This is a special case of Proposition 10.5. (iii) It follows from (i) that $`TP^nT:\omega \text{-}\mathrm{𝖢𝖺𝗍}^{\mathrm{}}\omega \text{-}\mathrm{𝖢𝖺𝗍}^{\mathrm{}}`$ has left adjoint $`T(M(I^n)T())TM(I^n)`$. But the obvious isomorphism $`T𝕀𝕀`$ induces an isomorphism $`TM(I^n)M(I^n)`$, so $`TM(I^n)`$ is naturally isomorphic to $`M(I^n)`$. The result now follows from (ii). $`\mathrm{}`$ The free cubical $`\omega `$-category on a cubical set is important in applications to concurrency theory. The data for a concurrent process can be given as a cubical set $`K`$, and the evolution of the data can be reasonably described by the free cubical $`\omega `$-category $`\sigma (K)`$; indeed, $`\sigma (K)`$ is the higher-dimensional analogue of the path category on a directed graph. The idea is pursued by Gaucher in . ## Acknowledgements We would like to thank the Department of Mathematical Sciences at Aalborg University for supporting Brown at the June, 1999, GETCO Workshop; the London Mathematical Society for support of a small workshop on multiple categories and concurrency in September, 1999, at Bangor; and a referee for helpful comments. This paper is dedicated to Philip Higgins in recognition of the power of his insights into ‘the algebra of cubes’ and in thanks for a long and happy collaboration with the second author.
warning/0007/astro-ph0007134.html	ar5iv	text	# Untitled Document WANDERING AMONG NEWTON WONDERS By D. Lynden-Bell Institute of Astronomy, Cambridge In the electromagnetic fields $`𝐁=(\omega t)^2𝐛(𝐫/\omega t),𝐄=𝐫\times 𝐁/(ct)`$ the trajectories of non-relativistic charged particles conserve $`(𝐫𝐯t)^2`$. The transformation $`\stackrel{~}{𝐫}=𝐫/(\omega t)`$, $`\tau =\omega ^2t^1`$ maps such trajectories into orbits in the constant magnetic field $`\stackrel{~}{𝐁}=𝐛(\stackrel{~}{𝐫})`$ all of which conserve $`\stackrel{~}{𝐯}^2`$. $`\omega `$ is a constant. The transformation may also be used to transform any fields obeying $`\mathrm{curl}𝐄=c^1𝐁/t,\mathrm{div}𝐁=0`$ into others and relates the particle trajectories in them. Introduction Newton’s equal areas in equal times theorem for non-coplanar motions<sup>1,2</sup> stimulated me to ask under what circumstances could the magnitudes of other vector constants of the motion remain constant even when extra forces were introduced which change their direction<sup>3</sup>. In particular $`(𝐫𝐯t)^2`$ is conserved when the force is of the form $$𝐅=(q/c)(𝐯\times 𝐁𝐫\times 𝐁/t)$$ $`(1)`$ for any $`𝐁(𝐫,t)`$, which may or may not have zero divergence as in electricity. The orbits in such a field were derived analytically for the special case of the electromagnetic field<sup>3</sup> $$𝐁=(\omega t)^2𝐛,𝐄=𝐫\times 𝐁/(ct)$$ with $`𝐛`$ constant and $`\omega `$ a constant of dimension $`t^1`$ put in to make $`𝐁`$ and $`𝐛`$ have the same dimensions. Here we demonstrate an unexpected correspondence between orbits under any force law of the form (1) and those under the conservative force law $$𝐅=(\stackrel{~}{𝐯}/c)\times \stackrel{~}{𝐁}$$ $`(2)`$ where $$\stackrel{~}{𝐫}=𝐫/\theta ,\theta =\omega t,\tau =\omega ^1/\theta ,\stackrel{~}{𝐯}=d\stackrel{~}{𝐫}/d\tau ,$$ $`(3)`$ and $$\stackrel{~}{𝐁}(\stackrel{~}{𝐫},\tau )=(\omega \tau )^2𝐁(\stackrel{~}{𝐫}/\omega \tau ,\omega ^2\tau ^1)=\theta ^2𝐁(𝐫,t).$$ $`(4)`$ The conservation of $`(𝐫𝐯t)^2`$ under the force (1) translates into the conservation of $`\stackrel{~}{𝐯}^2`$ under force (2), because $$\stackrel{~}{𝐯}=d\stackrel{~}{𝐫}/d\tau =\omega \theta ^2d(𝐫/\theta )/d\theta =𝐫𝐯t.$$ The equation of motion of a unit mass particle of charge $`q`$ under the force law (1) is, $$d^2𝐫/dt^2=(q/c)\left[d𝐫/dt\times 𝐁𝐫\times 𝐁/t\right].$$ $`(5)`$ Now under the transformation (3) equation (5) can be written $$\omega ^2\theta ^3d^2(\stackrel{~}{𝐫}\theta )/d\theta ^2=(q/c)\omega ^2\theta ^2d\stackrel{~}{𝐫}/d\theta \times \theta ^2𝐁.$$ $`(6)`$ But $`\theta ^3d^2(\stackrel{~}{𝐫}\theta )/d\theta ^2=\theta ^2d/d\theta (\theta ^2d\stackrel{~}{𝐫}/d\theta )=\omega ^2d^2\stackrel{~}{𝐫}/d\tau ^2`$ so we may rewrite (6) in the form, $$d^2\stackrel{~}{𝐫}/𝐝\tau ^2=(q/c)d\stackrel{~}{𝐫}/d\tau \times \stackrel{~}{𝐁},$$ $`(7)`$ where $`\stackrel{~}{𝐁}`$ is given by (4). Equation (7) is the equation of motion of a unit mass particle of charge $`q`$ in a field $`\stackrel{~}{𝐁}`$ of magnetic type with $`\tau `$ as the time. So if $`𝐫=𝐑(t)`$ is a solution of (1), $`\stackrel{~}{𝐫}=\omega \tau 𝐑(\omega ^2\tau ^1)`$ solves (7). Application to Electromagnetic orbits Under what circumstances can forces of the form (1) be delivered on a charge $`q`$ by electromagnetic fields? Clearly the first term in the bracket of (1) is of magnetic type (provided $`𝐁`$ has zero divergence) and for the other to be electric we need $$𝐄=𝐫\times 𝐁/(ct)$$ Thus $`\mathrm{Curl}𝐄=(2𝐁+𝐫.\mathbf{}𝐁)/(ct).`$ Putting this equal to $`(1/c)𝐁/t`$ leads to $`t/t(𝐁t^2)+𝐫\mathbf{}(𝐁t^2)=0`$ of which the general solution is $$𝐁=\theta ^2𝐛(𝐫/\theta )$$ $`(8)`$ where $`𝐛`$ is any (vector) function of its argument subject to Maxwell’s condition $`\mathrm{div}𝐛=0`$. Under the transformation (4) we find that (8) reduces to $`\stackrel{~}{𝐁}=𝐛(\stackrel{~}{𝐫})`$, so $`\stackrel{~}{𝐁}`$ can be any stationary magnetic field. Thus, if $`\stackrel{~}{𝐫}=\stackrel{~}{𝐫}(\tau )`$ is the trajectory of a charged test particle that moves non-relativistically in any stationary magnetic field $`\stackrel{~}{𝐁}=𝐛(\stackrel{~}{𝐫})`$, then the electromagnetic fields $`𝐁=\theta ^2𝐛(𝐫/\theta ),𝐄=𝐫\times 𝐁/(ct)`$ satisfy Maxwell’s $`\mathrm{curl}𝐄=(1/c)𝐁/t,\mathrm{div}𝐁=0`$ and in them there is a corresponding trajectory $`𝐫=\theta \stackrel{~}{𝐫}(\omega ^1/\theta )`$ which conserves $`(𝐫𝐯t)^2.`$ Neither Maxwell’s equations, nor the equations of motion involve a particular zero point for time, so the same arguments hold if we write to $`tt_0`$ wherever we have written $`t`$ above and the trajectory then preserves $`\left[𝐫𝐯(tt_0)\right]^2`$. A specific example is given by particles trapped by a dipolar field to form Van Allen belts; the field is $`\stackrel{~}{𝐁}=(𝐌+3𝐌\widehat{𝐫}\widehat{𝐫})/\stackrel{~}{r}^3`$ where $`\widehat{𝐫}`$ is the unit vector $`\stackrel{~}{𝐫}/\stackrel{~}{r}=𝐫/r`$. Our theorem relates the orbits $`\stackrel{~}{𝐫}=\stackrel{~}{𝐫}(\tau )`$ in this constant dipole to the orbits $`𝐫=\omega (tt_0)\stackrel{~}{𝐫}\left(\omega ^2(tt_0)^1\right)`$ in the dipolar field $`𝐁`$ with moment $`\omega (tt_0)𝐌`$ $$𝐁=\omega (tt_0)(𝐌3𝐌\widehat{𝐫}\widehat{𝐫})/r^3$$ $`(9)`$ with $`𝐄=𝐫\times 𝐌/(cr^3).`$ Now in a uniform field the orbits expand as the field weakens but in this dipolar field we see the orbits shrink as $`t`$ approaches $`t_0`$ and the dipole weakens. This is because $`\stackrel{~}{𝐫}=\stackrel{~}{𝐫}(\tau )`$ is confined to the ‘radiation belts’ for all $`\tau `$ so the whole of that motion is now shrunk by the initial factor $`(tt_0)`$ in $`𝐫`$. The physics behind this shrinking of the orbit with the dipole strength lies in the $`c𝐄\times 𝐁/B^2`$ drift of the orbit. The $`𝐄`$, which is an inevitable consequence of the changing $`𝐁`$, is directed toroidally. For a decreasing dipole the drift is directed inwards. It is then of interest to ask whether the field strength at the guiding centre increases or decreases as the dipole strength diminishes. For an equatorial gyration the drift velocity is $`𝐯_d=𝐫/(tt_0)`$ and $`DB^2/Dt=2𝐁\left(𝐁/t+𝐯_d\mathbf{}B\right)=4B^2/(tt_0)`$ which is clearly positive for $`tt_0`$ when the field is decreasing. Thus the drift pushes the orbits into regions of higher field strength even though the field strength at a given position decreases. This result emphasises the importance of considering the electric fields inevitably associated with changing magnetic fields. These electric fields have been included in our theorem on changing orbits. In this particular example the electric field is steady and toroidal but its lack of change depends crucially on the $`r^3`$ dependence of the field strength of a dipole and in general the $`E`$ field is time dependent. Of course (9) is only the correct magnetic field of a changing dipole in the ‘near zone’ in which its radiation may be neglected. Notice that $`\mathrm{div}𝐄=0`$ so no charge density is associated with this ‘pure induction’ $`𝐄`$ field. Transformations of more general electromagnetic fields Although we found the above results by looking at forces that preserve $`(𝐫𝐯t)^2)`$, we have arrived at a strange transformation that preserves the two Maxwell equations that do not involve sources, so the transformation can be applied to any electromagnetic field. Under the transformation $`\begin{array}{cc}\stackrel{~}{𝐫}=𝐫/(\omega t)\hfill & 𝐫=\stackrel{~}{𝐫}/(\omega \tau )\hfill \\ \omega \tau =(\omega t)^1\hfill & \omega t=(\omega \tau )^1\hfill \\ \stackrel{~}{𝐁}(\stackrel{~}{𝐫},\tau )=(\omega t)^2𝐁(𝐫,t)=(\omega \tau )^2𝐁(\stackrel{~}{𝐫}/\omega \tau ,\omega ^2\tau ^1)\hfill & 𝐁(𝐫,t)=(\omega \tau )^2\stackrel{~}{𝐁}(\stackrel{~}{𝐫},\tau )\hfill \\ \stackrel{~}{𝐄}(\stackrel{~}{𝐫},\tau )=(\omega t)^3\left[𝐄(𝐫,t)+𝐫\times 𝐁(𝐫,t)/ct\right]\hfill & 𝐄(𝐫,t)=(\omega \tau )^3[\stackrel{~}{𝐄}(\stackrel{~}{𝐫},\tau )+\hfill \\ =(\omega \tau )^3\left[𝐄(\stackrel{~}{𝐫}\omega ^1\tau ^1,\omega ^2\tau ^1)+\stackrel{~}{𝐫}\times 𝐁(\stackrel{~}{𝐫}\omega ^1\tau ^1,\omega ^2\tau ^1)/c\right]\hfill & +\stackrel{~}{𝐫}\times \stackrel{~}{𝐁}(\stackrel{~}{𝐫},\tau )/(c\tau )]\hfill \end{array}`$ we write $`/𝐫`$ in place of $`\mathbf{}`$ and deduce $`\begin{array}{cccc}& /\stackrel{~}{𝐫}\stackrel{~}{𝐁}=0& & /𝐫𝐁=0\\ & /\stackrel{~}{𝐫}\times \stackrel{~}{𝐄}+(1/c)\stackrel{~}{𝐁}/\tau =0& & /𝐫\times 𝐄+(1/c)𝐁/t=0\\ \mathrm{and}\hfill & & & \\ & d^2\stackrel{~}{𝐫}/d\tau ^2=q\left[\stackrel{~}{𝐯}/c\times \stackrel{~}{𝐁}+\stackrel{~}{𝐄}\right]& & d^2𝐫/dt^2=q(𝐯/c\times 𝐁+𝐄)\end{array}`$ Thus if $`𝐄`$ and $`𝐁`$ obey Maxwell’s equations with suitable sources $`\rho ,𝐣`$ then $`\stackrel{~}{𝐄}`$ and $`\stackrel{~}{𝐁}`$ will also (with other sources) and the trajectories of classical non-relativistic particles $`𝐫=𝐑(t)`$ map via the transformation into the trajectories $`\stackrel{~}{𝐫}=\omega \tau 𝐑(\omega ^2\tau ^1)`$ of particles of the same charge/mass ratio under the fields $`\stackrel{~}{𝐄}`$ and $`\stackrel{~}{𝐁}`$. We notice that the field of an electric dipole is invariant while that of a magnetic monopole reverses under the transformation. However if we added a time reversal to the transformation then both the electric dipole and the magnetic monopole would be invariant. For a somewhat related theorem on the gravitational $`N`$ body problem see<sup>4</sup>, but note there should be a dot over $`f`$ the fourth time it appears there in the mathematics. References (1) I. Newton, Principia, Scholium to Proposition 2 (R.Soc., London), 1687. (2) D. Lynden-Bell & M. Nouri-Zonoz, Revs of Mod Phys, 70, 427, 1999. (3) D.Lynden-Bell, The Observatory, (accepted), 2000. (4) D.Lynden-Bell, The Observatory, 102, No 1048, 86, 1982.
warning/0007/hep-ph0007338.html	ar5iv	text	# Tests of Large–𝑁_𝑐 QCD from Hadronic 𝜏 Decay ## Abstract We use the Aleph data on vector and axial-vector spectral functions to test simple duality properties of QCD in the large $`N_c`$ limit, which emerge in the approximation of a minimal hadronic ansatz of a spectrum of narrow states. These duality properties relate the short– and long–distance behaviours of specific correlation functions, which are order parameters of spontaneous chiral symmetry breaking, in a way that we find well supported by the data. preprint: CPT-2000/P.4047preprint: UAB-FT-494 1. At first sight, the hadronic world predicted by QCD in the limit of a large number of colours $`N_c`$ tH74 may seem rather different from the real world. The hadronic spectrum of vector and axial–vector states, observed e.g. in $`e^+e^{}`$ annihilations and in $`\tau `$ decays, has certainly much more structure than the infinite set of narrow states predicted by large $`N_c`$ QCD Wi79 (QCD$`_{\mathrm{}}`$). There are, however, many instances in Particle Physics where one is only interested in certain weighted integrals of hadronic spectral functions. In these cases, it may be enough to know a few global properties of the hadronic spectrum; one does not expect the integrals to depend crucially on the details of the spectrum at all energies. Typical examples of that are the coupling constants of the effective chiral Lagrangian of QCD at low energies, as well as the coupling constants of the effective chiral Lagrangian of the electroweak interactions of pseudoscalar particles in the Standard Model, which are needed to understand $`K`$–Physics in particular, (see e.g. the review article in ref. Pi99 and references therein.) It is in these examples that the hadronic world predicted by QCD$`_{\mathrm{}}`$may provide a good approximation to the real hadronic spectrum. If so, QCD$`_{\mathrm{}}`$could then become a useful phenomenological approach for understanding non–perturbative QCD physics at low energies. There are indeed a number of successful calculations which have already been made within the framework of QCD$`_{\mathrm{}}`$, (see ref. KPdeR00 and references therein.) The picture which emerges from these applications is one of a remarkable simplicity. It is found that, when dealing with Green’s functions that are order parameters of spontaneous chiral symmetry breaking, the restriction of the infinite set of large $`N_c`$ narrow states to a minimal hadronic ansatz which is needed to satisfy the leading short– and long–distance behaviours of the relevant Green’s functions, provides already a very good approximation to the observables one computes. The purpose of this note is to investigate this minimal hadronic ansatz approximation in a case where one can compare, in detail, the theoretical predictions to the phenomenological results evaluated with experimental data. 2. Of particular interest for our purposes is the correlation function ($`Q^2q^20`$ for $`q^2`$ space–like) $$\mathrm{\Pi }_{LR}^{\mu \nu }(q)=2id^4xe^{iqx}0\|\text{T}\left(L^\mu (x)R^\nu (0)^{}\right)\|0,$$ (1) with colour singlet currents $$R^\mu \left(L^\mu \right)=\overline{d}(x)\gamma ^\mu \frac{1}{2}(1\pm \gamma _5)u(x).$$ (2) In the chiral limit, $`m_{u,d,s}0`$ , this correlation function has only a transverse component, $$\mathrm{\Pi }_{LR}^{\mu \nu }(Q^2)=(q^\mu q^\nu g^{\mu \nu }q^2)\mathrm{\Pi }_{LR}(Q^2).$$ (3) The self-energy–like function $`\mathrm{\Pi }_{LR}(Q^2)`$ vanishes order by order in perturbative QCD (pQCD) and is an order parameter of S$`\chi `$SB for all values of $`Q^2`$; therefore it obeys an unsubtracted dispersion relation, $$\mathrm{\Pi }_{LR}(Q^2)=_0^{\mathrm{}}𝑑t\frac{1}{t+Q^2}\frac{1}{\pi }\text{Im}\mathrm{\Pi }_{LR}(t).$$ (4) In QCD$`_{\mathrm{}}`$the spectral function $`\frac{1}{\pi }\text{Im}\mathrm{\Pi }_{LR}(t)`$ consists of the difference of an infinite number of narrow vector and axial–vector states, together with the Goldstone pole of the pion: $`{\displaystyle \frac{1}{\pi }}\text{Im}\mathrm{\Pi }_{LR}(t)={\displaystyle \underset{V}{}}f_V^2M_V^2\delta (tM_V^2)`$ (5) $`F_0^2\delta (t){\displaystyle \underset{A}{}}f_A^2M_A^2\delta (tM_A^2).`$ The low $`Q^2`$ behaviour of $`\mathrm{\Pi }_{LR}(Q^2)`$, i.e. the long–distance behaviour of the correlation function in Eq. (1), is governed by chiral perturbation theory: $$Q^2\mathrm{\Pi }_{LR}(Q^2)\|_{Q^20}=F_0^2+4L_{10}Q^2+𝒪(Q^4),$$ (6) where $`F_0`$ is the pion coupling constant in the chiral limit, and $`L_{10}`$ is one of the coupling constants of the $`𝒪(p^4)`$ effective chiral Lagrangian. The high $`Q^2`$ behaviour of $`\mathrm{\Pi }_{LR}(Q^2)`$, i.e. the short–distance behaviour of the correlation function in Eq. (1), is governed by the operator product expansion (OPE) of the two local currents in Eq. (1SVZ79 , $$\underset{Q^2\mathrm{}}{lim}Q^6\mathrm{\Pi }_{LR}(Q^2)=\left[4\pi ^2\frac{\alpha _s}{\pi }+𝒪(\alpha _s^2)\right]\overline{\psi }\psi ^2,$$ (7) which implies the two Weinberg sum rules: $$_0^{\mathrm{}}𝑑t\text{Im}\mathrm{\Pi }_{LR}(t)=\underset{V}{}f_V^2M_V^2\underset{A}{}f_A^2M_A^2F_0^2=0,$$ (8) and $$_0^{\mathrm{}}𝑑tt\text{Im}\mathrm{\Pi }_{LR}(t)=\underset{V}{}f_V^2M_V^4\underset{A}{}f_A^2M_A^4=0,$$ (9) as well as the sum rule KdeR98 $$\underset{V}{}f_V^2M_V^6\underset{A}{}f_A^2M_A^6=\left[4\pi \alpha _s+𝒪(\alpha _s^2)\right]\overline{\psi }\psi ^2.$$ (10) In fact, as pointed out in ref. KdeR98 , in QCD$`_{\mathrm{}}`$there exist an infinite number of Weinberg–like sum rules. In full generality, the moments of the $`\mathrm{\Pi }_{LR}`$ spectral function with $`n=3,4,\mathrm{}`$, $`{\displaystyle _0^{\mathrm{}}}𝑑tt^{n1}\left[{\displaystyle \frac{1}{\pi }}\text{Im}\mathrm{\Pi }_V(t){\displaystyle \frac{1}{\pi }}\text{Im}\mathrm{\Pi }_A(t)\right]=`$ (11) $`{\displaystyle \underset{V}{}}f_V^2M_V^{2n}{\displaystyle \underset{A}{}}f_A^2M_A^{2n},`$ govern the short–distance expansion of the $`\mathrm{\Pi }_{LR}(Q^2)`$ function $`\mathrm{\Pi }_{LR}(Q^2)\|_{Q^2\mathrm{}}=\left({\displaystyle \underset{V}{}}f_V^2M_V^6{\displaystyle \underset{A}{}}f_A^2M_A^6\right){\displaystyle \frac{1}{Q^6}}`$ (12) $`+\left({\displaystyle \underset{V}{}}f_V^2M_V^8{\displaystyle \underset{A}{}}f_A^2M_A^8\right){\displaystyle \frac{1}{Q^8}}+\mathrm{}.`$ On the other hand inverse moments of the $`\mathrm{\Pi }_{LR}`$ spectral function with the pion pole removed (which we denote by $`\text{Im}\stackrel{~}{\mathrm{\Pi }}_A(t)`$) determine a class of coupling constants of the low–energy effective chiral Lagrangian. For example, $`{\displaystyle _0^{\mathrm{}}}𝑑t{\displaystyle \frac{1}{t}}\left[{\displaystyle \frac{1}{\pi }}\text{Im}\mathrm{\Pi }_V(t){\displaystyle \frac{1}{\pi }}\text{Im}\stackrel{~}{\mathrm{\Pi }}_A(t)\right]=`$ (13) $`{\displaystyle \underset{V}{}}f_V^2{\displaystyle \underset{A}{}}f_A^2=4L_{10}.`$ Moments with higher inverse powers of $`t`$ are associated with couplings of composite operators of higher dimension in the chiral Lagrangian. Tests of the two Weinberg sum rules in Eqs. (8) and (9) and of the $`L_{10}`$ sum rule in Eq. (13), in a different context from the one we are interested in here, have often appeared in the literature (see e.g. refs. DHGS98 and DS99 for recent discussions where earlier references can also be found). 3. The minimal hadronic ansatz which satisfies the two Weinberg sum rules in Eqs. (8) and (9) is a spectrum of one vector state $`V`$, one axial–vector state $`A`$ and the Goldstone pion, with the ordering KdeR98 $`M_V<M_A`$. In this approximation, $`\mathrm{\Pi }_{LR}(Q^2)`$ has a very simple form $`Q^2\mathrm{\Pi }_{LR}(Q^2)`$ $`=`$ $`{\displaystyle \frac{F_0^2}{\left(1+\frac{Q^2}{M_V^2}\right)\left(1+\frac{Q^2}{M_A^2}\right)}}`$ (14) $`=`$ $`{\displaystyle \frac{M_A^2M_V^2}{Q^4}}{\displaystyle \frac{F_0^2}{\left(1+\frac{M_V^2}{Q^2}\right)\left(1+\frac{M_A^2}{Q^2}\right)}}.`$ (15) This equation shows, explicitly, a remarkable short–distance $``$ long–distance duality deR99 . Indeed, with $`g_A`$ defined so that $`M_V^2=g_AM_A^2`$ and $`z\frac{Q^2}{M_V^2}`$, the non–local order parameters corresponding to the long–distance expansion for $`z0`$, which are couplings of the effective chiral Lagrangian i.e., $`Q^2\mathrm{\Pi }_{LR}(Q^2)\|_{z0}=F_0^2\{1(1+g_A)z`$ (16) $`+(1+g_A+g_A^2)z^2+\mathrm{}\},`$ are correlated to the local order parameters of the short–distance OPE for $`z\mathrm{}`$ in a very simple way: $`Q^2\mathrm{\Pi }_{LR}(Q^2)\|_z\mathrm{}=F_0^2{\displaystyle \frac{1}{g_A}}{\displaystyle \frac{1}{z^2}}\{1(1+{\displaystyle \frac{1}{g_A}}){\displaystyle \frac{1}{z}}`$ (17) $`+(1+{\displaystyle \frac{1}{g_A}}+{\displaystyle \frac{1}{g_A^2}}){\displaystyle \frac{1}{z^2}}+\mathrm{}\};`$ in other words, there is a one-to-one correspondance between the two expansions by changing $$g_A\frac{1}{g_A}\text{and}z^n\frac{1}{g_A}\frac{1}{z^{n+2}}.$$ (18) The moments of the $`\mathrm{\Pi }_{LR}`$ spectral function, when evaluated in the minimal hadronic ansatz approximation, can be converted into a very simple set of finite energy sum rules (FESR’s), corresponding to the OPE in Eq. (17) $`{\displaystyle _0^{s_0}}𝑑tt^2{\displaystyle \frac{1}{\pi }}\text{Im}\mathrm{\Pi }_{LR}(t)`$ $`=`$ $`F_0^2M_V^4{\displaystyle \frac{1}{g_A}},`$ (19) $`{\displaystyle _0^{s_0}}𝑑tt^3{\displaystyle \frac{1}{\pi }}\text{Im}\mathrm{\Pi }_{LR}(t)`$ $`=`$ $`F_0^2M_V^6{\displaystyle \frac{1+\frac{1}{g_A}}{g_A}},`$ (20) $`{\displaystyle _0^{s_0}}𝑑tt^4{\displaystyle \frac{1}{\pi }}\text{Im}\mathrm{\Pi }_{LR}(t)`$ $`=`$ $`F_0^2M_V^8{\displaystyle \frac{1+\frac{1}{g_A}+\frac{1}{g_A^2}}{g_A}},`$ (21) $`\mathrm{}`$ $`\mathrm{}.`$ where the upper limit of integration $`s_0`$ denotes the onset of the pQCD continuum which, in the chiral limit, is common to the vector and axial–vector spectral functions. It is important to realize that $`s_0`$ is not a free parameter. Its value is fixed by the requirement that the OPE of the correlation function of two vector currents, (or two axial–vector currents,) in the chiral limit, have no $`1/Q^2`$ term, which results in an implicit equation for $`s_0`$ BLdeR85 ; PPdeR98 . In the minimal hadronic ansatz approximation the onset of the pQCD continuum, which we shall call $`s_0^{}`$, is then fixed by the equation $$\frac{N_c}{16\pi ^2}\frac{2}{3}s_0^{}\left(1+𝒪(\alpha _\text{s})\right)=F_0^2\frac{1}{1g_A}.$$ (22) Also, the moments which correspond to the chiral expansion in Eq. (16) are given by another simple set of FESR’s: $`{\displaystyle _0^{s_0}}𝑑t{\displaystyle \frac{1}{\pi }}\text{Im}\stackrel{~}{\mathrm{\Pi }}_{LR}(t)`$ $`=`$ $`F_0^2,`$ (23) $`{\displaystyle _0^{s_0}}{\displaystyle \frac{dt}{t}}{\displaystyle \frac{1}{\pi }}\text{Im}\stackrel{~}{\mathrm{\Pi }}_{LR}(t)`$ $`=`$ $`{\displaystyle \frac{F_0^2}{M_V^2}}(1+g_A),`$ (24) $`{\displaystyle _0^{s_0}}{\displaystyle \frac{dt}{t^2}}{\displaystyle \frac{1}{\pi }}\text{Im}\stackrel{~}{\mathrm{\Pi }}_{LR}(t)`$ $`=`$ $`{\displaystyle \frac{F_0^2}{M_V^4}}(1+g_A+g_A^2),`$ (25) $`\mathrm{}`$ $`\mathrm{}.`$ We propose to test these duality relations by comparing moments of the physical spectral function $`\frac{1}{\pi }\text{Im}\mathrm{\Pi }_{LR}(t)`$ to the predictions of the minimal hadronic ansatz. 4. The ALEPH collaboration at LEP has measured the inclusive invariant mass–squared distribution of hadronic $`\tau `$ decays ALEPH into non–strange particles. They have been able to extract from their data both the vector current spectral function $`\frac{1}{\pi }\text{Im}\mathrm{\Pi }_V^{\text{exp.}}(t)`$ and the axial–vector current spectral function $`\frac{1}{\pi }\text{Im}\mathrm{\Pi }_A^{\text{exp.}}(t)`$ up to $`t3\text{GeV}^2`$. In fact, in the real world, the correlation function in Eq. (3) has a non–transverse term as well, which is dominated by the pion pole contribution to the axial–vector component. The vector contribution to the non–transverse term vanishes in the limit of isospin invariance. In order to compare the moments of the experimental spectral function $`\frac{1}{\pi }\text{Im}\mathrm{\Pi }_{LR}^{\text{exp.}}(t)`$ to those in Eqs. (19)-(21) and (23)-(25) we still have to correct for the fact that the FESR’s in these equations apply in the chiral limit where $`m_{u,d}0`$. This we do by exploiting the analyticity properties of the two–point function $`\mathrm{\Pi }_{LR}`$ in the complex $`q^2`$–plane. Integration over a standard contour relates weighted integrals of the spectral function $`\frac{1}{\pi }\text{Im}\mathrm{\Pi }_{LR}^{\text{exp.}}(t)`$ in a finite interval on the real axis to integrals of $`\mathrm{\Pi }_{LR}(q^2)`$ over a small circle $`\|q^2\|=s_{\text{th}}`$ and a large circle $`\|q^2\|=s_0`$: $`{\displaystyle _{s_{th}}^{s_0}}𝑑tf(t)\text{Im}\mathrm{\Pi }_{LR}(t)=`$ $`{\displaystyle \frac{1}{2i}}{\displaystyle \underset{\|q^2\|=s_{th}}{}}𝑑q^2f(q^2)\mathrm{\Pi }_{LR}(q^2){\displaystyle \frac{1}{2i}}{\displaystyle \underset{\|q^2\|=s_0}{}}𝑑q^2f(q^2)\mathrm{\Pi }_{LR}(q^2),`$ (26) where the weight function $`f(q^2)`$ is a conveniently chosen analytic function inside the contour; in our case simple powers and inverse powers of $`q^2`$. The chiral corrections in the small circle are particularly important in the evaluation of the inverse moments. We have evaluated them by taking into account the one loop expression of $`\mathrm{\Pi }_{LR}(z)`$ in chiral perturbation theory GL84 . The chiral corrections in the large circle are rather small. They appear as leading $`1/Q^2`$ and next–to–leading $`1/Q^4`$ power corrections in the OPE of $`\mathrm{\Pi }_{LR}(Q^2)`$ at large $`Q^2`$ but their coefficients, proportional to quark masses, are small FNdeR79 . With these corrections incorporated, we proceed now to the comparison we are looking for. This is shown in Figs. 1 and 2 below where we show the various moments as a function of $`s_0`$. Fig. 1 Plot of the experimental moments in Eqs. (19), (20) and (21) normalized to the minimal hadronic ansatz predictions on the r.h.s. The six plots in Figs. 1 and 2 show the experimental moments on the l.h.s. of Eqs. (19)-(21) and Eqs. (23)-(25), respectively, as a function of $`s_0`$, extrapolated to the chiral limit as discussed above and normalized to the corresponding minimal hadronic ansatz predictions on the r.h.s. The horizontal bands on these plots correspond to the induced error of the minimal hadronic ansatz predictions from the input values: $`F_0=87\pm 3.5\text{MeV}`$, $`M_V=748\pm 29\text{MeV}`$ and $`g_A=0.50\pm 0.06`$. These are the values favored by a global fit of the minimal hadronic ansatz to low–energy observables PPdeR98 . The moments $`_n`$, with the experimental error propagation included, are the curved bands in the figures. Fig. 2 Plot of the experimental moments in Eqs. (23), (24) and (25) normalized to the minimal hadronic ansatz predictions on the r.h.s. The remarkable feature which the curves in Figs. 1 and 2 show is that, within errors, the first crossing of all the experimental moments with the minimal hadronic ansatz band takes place in the same $`s_0`$ region, around $`s_01.4\text{GeV}^2`$ rather close indeed to the $`s_0^{}`$ value which follows from Eq. (22): $`s_0^{}=(1.2\pm 0.2)\text{GeV}^2`$. We have also checked that for the 2nd Weinberg sum rule in Eq. (9), not shown in the figures. In fact, the agreement for the inverse moments is excellent. This is due to the fact that inverse moments put more and more weight on the low energy tail of the spectral function, which is known to be dominated by the $`\rho `$–resonance. By contrast, the positive moments are very sensitive to the cancellations between opposite parity hadronic states; this is why the experimental curves show larger and larger oscillations as one increases the power of the moment. In spite of that, it is quite impressive that, when restricted to the $`s_0`$ region of duality, the experimental moments agree well with the minimal hadronic ansatz prediction, even for rather large powers which correspond to vacuum expectation values of operators of higher and higher dimension in the OPE. We conclude that the experimental data from ALEPH is consistent with the simple pattern of duality properties between short and long–distances which follow from the minimal hadronic ansatz of a narrow vector and axial-vector states plus the Goldstone pion in large–$`N_c`$ QCD. Work supported in part by TMR, EC-Contract No. ERBFMRX-CT980169 (EuroDa$`\varphi `$ne). The work of S.P. is also supported by CICYT-AEN99-0766.
warning/0007/quant-ph0007053.html	ar5iv	text	# Remarks on 2–q-bit states ## 1 Introduction Entangled q-bits (binary quantum alternatives) are exploited in most schemes proposed for quantum communication purposes, for quantum information processing, or for the secure key distribution procedures known as quantum cryptography. The basic units are entangled q-bit pairs. Obviously then, a thorough understanding of the properties of 2–q-bit states is desirable. Although there has been considerable progress in this matter recently, the situation is still quite unsatisfactory. The characterization of the 2–q-bit states produced by some source requires the experimental determination of 15 real parameters. Ideally, this is done by measuring a suitably chosen set of five observables that constitute “a complete set of five pairs of complementary propositions” . In an optical model , which makes use of single-photon 2–q-bit states, these measurements can be realized, and other experimental studies of 2–q-bit states can be performed as well. Then, based on the knowledge of the 15 state-specifying parameters, one can classify the 2–q-bit state. We distinguish, in Sec. 2, six classes of families of locally equivalent states. Roughly speaking, local equivalence means that the difference is of a geometrical, not a physical nature. In a certain sense, the 15 parameters can be regarded as consisting of 6 geometrical ones and 9 physical ones. The classification of Sec. 2 is straightforward but not sufficient. One also needs to know if the 2–q-bit state in question is useful for quantum communication purposes. In the technical terms of Sec. 4, we ask for its degree of separability as a numerical measure for this usefulness. The degree of separability is part and parcel of the so-called optimal Lewenstein-Sanpera decomposition of a 2–q-bit state. This decomposition is known for a number of relevant types of states but, despite a good understanding of its properties, presently we do not have a method for finding it for any arbitrary 2–q-bit state. In Sec. 5 we remark briefly on the so-called concurrence of a 2–q-bit state and surmise that the sum of the degree of separability and the concurrence does not exceed unity. The Appendix reports some technical details. ## 2 Families of 2–q-bit states We employ the terminology and the notational conventions of . As usual, we describe the individual q-bits with the aid of analogs of Pauli’s spin vector operator: $`\stackrel{}{\sigma }`$ for the first q-bit, $`\stackrel{}{\tau }`$ for the second. These row vectors refer to two three-dimensional vector spaces that are unrelated, which is to say that in $`\stackrel{}{\sigma }`$ $`=`$ $`{\displaystyle \underset{\alpha =x,y,z}{}}\sigma _\alpha \underset{\alpha }{\overset{}{e}}=(\sigma _x,\sigma _y,\sigma _z)\left(\begin{array}{c}\underset{x}{\overset{}{e}}\\ \underset{y}{\overset{}{e}}\\ \underset{z}{\overset{}{e}}\end{array}\right),`$ (4) $`\stackrel{}{\tau }`$ $`=`$ $`{\displaystyle \underset{\beta =x,y,z}{}}\tau _\beta \underset{\beta }{\overset{}{n}}=(\tau _x,\tau _y,\tau _z)\left(\begin{array}{c}\underset{x}{\overset{}{n}}\\ \underset{y}{\overset{}{n}}\\ \underset{z}{\overset{}{n}}\end{array}\right)`$ (8) the orthonormal right-handed vector sets $`\underset{x}{\overset{}{e}},\underset{y}{\overset{}{e}},\underset{z}{\overset{}{e}}`$ and $`\underset{x}{\overset{}{n}},\underset{y}{\overset{}{n}},\underset{z}{\overset{}{n}}`$ have nothing to do with each other. In addition to these pre-chosen $`xyz`$ coordinate systems, we’ll also consider 123 coordinate systems that are adapted to the 2–q-bit state of interest. Then $$\stackrel{}{\sigma }=(\sigma _1,\sigma _2,\sigma _3)\left(\begin{array}{c}\underset{1}{\overset{}{e}}\\ \underset{2}{\overset{}{e}}\\ \underset{3}{\overset{}{e}}\end{array}\right),\stackrel{}{\tau }=(\tau _1,\tau _2,\tau _3)\left(\begin{array}{c}\underset{1}{\overset{}{n}}\\ \underset{2}{\overset{}{n}}\\ \underset{3}{\overset{}{n}}\end{array}\right)$$ (9) are the respective parameterizations of $`\stackrel{}{\sigma }`$ and $`\stackrel{}{\tau }`$. As an elementary illustration think of the statistical operator of the first q-bit, specified by the Pauli vector $`\stackrel{}{s}=\stackrel{}{\sigma }`$, $`\rho _1={\displaystyle \frac{1}{2}}\left(1+\stackrel{}{s}\sigma ^{}\right)`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(1+s_x\sigma _x+s_y\sigma _y+s_z\sigma _z\right)`$ (10) $`=`$ $`{\displaystyle \frac{1}{2}}(1+s\sigma _1),`$ where $`s0`$ is the length of $`\stackrel{}{s}`$ and the 123 system has $`\underset{1}{\overset{}{e}}`$ in the direction of $`\stackrel{}{s}`$ by definition. Note that we are carefully distinguishing row vectors from column vectors, as in the scalar product $`\stackrel{}{s}\sigma ^{}`$ of row $`\stackrel{}{s}`$ and column $`\sigma ^{}`$; of course, columns and rows are transposes of each other, $`s^{}=\stackrel{}{s}^\text{T}`$. Admittedly, this distinction is somewhat pedantic, but it makes book keeping much easier. Unitary transformations that affect only one of the q-bits or both q-bits independently are local transformations. Geometrically speaking, local transformations rotate $`\stackrel{}{\sigma }`$ and $`\stackrel{}{\tau }`$. Two states that can be turned into each other by a local transformation are locally equivalent. For instance, two first–q-bit states (10) are equivalent if their Pauli vectors have the same length; at most the two states can differ by the direction of $`\underset{1}{\overset{}{e}}`$. In other words: The difference of two equivalent states is only in the 123 bases that go with the generic form $`\frac{1}{2}(1+s\sigma _1)`$. Likewise, there are families of locally equivalent 2–q-bit states. To decide whether two given 2–q-bit states belong to the same family, one may put them into a generic form that is uniquely fixed by convenient conventions. The following set of conventions seems to be quite natural. It all begins with recalling that the general form of a 2–q-bit state is given by $$\mathrm{P}=\frac{1}{4}\left(1+\stackrel{}{\sigma }s^{}+\stackrel{}{t}\tau ^{}+\stackrel{}{\sigma }\text{}\tau ^{}\right).$$ (11) It involves the cross dyadic $`\text{}`$, $$\text{}=\sigma ^{}\stackrel{}{\tau }=(e_x^{},e_y^{},e_z^{})\left(\begin{array}{ccc}C_{xx}& C_{xy}& C_{xz}\\ C_{yx}& C_{yy}& C_{yz}\\ C_{zx}& C_{zy}& C_{zz}\end{array}\right)\left(\begin{array}{c}\underset{x}{\overset{}{n}}\\ \underset{y}{\overset{}{n}}\\ \underset{z}{\overset{}{n}}\end{array}\right),$$ (12) in addition to the Pauli vectors $`s^{}`$ and $`\stackrel{}{t}`$, $`s^{}`$ $`=`$ $`\sigma ^{}=(e_x^{},e_y^{},e_z^{})\left(\begin{array}{c}s_x\\ s_y\\ s_z\end{array}\right),`$ (16) $`\stackrel{}{t}`$ $`=`$ $`\stackrel{}{\tau }=(t_x,t_y,t_z)\left(\begin{array}{c}\underset{x}{\overset{}{n}}\\ \underset{y}{\overset{}{n}}\\ \underset{z}{\overset{}{n}}\end{array}\right).`$ (20) In a first step we bring $`\text{}`$ into the diagonal form $`\text{}=\pm \left(e_1^{}c_1\underset{1}{\overset{}{n}}+e_2^{}c_2\underset{2}{\overset{}{n}}+e_3^{}c_3\underset{3}{\overset{}{n}}\right)`$ $`\text{for}\{\begin{array}{c}\mathrm{det}\left\{\text{}\right\}0,\hfill \\ \mathrm{det}\left\{\text{}\right\}<0,\hfill \end{array}`$ (23) with its characteristic values ordered in accordance with $$c_1c_2c_30.$$ (24) Their squares are the eigenvalues of $`\text{}\text{}^\text{T}`$ or $`\text{}^\text{T}\text{}`$; the eigencolumns of $`\text{}\text{}^\text{T}`$ constitute the orthonormal right-handed set $`e_1^{},e_2^{},e_3^{}`$, and the corresponding $`\underset{1}{\overset{}{n}},\underset{2}{\overset{}{n}},\underset{3}{\overset{}{n}}`$ are eigenrows of $`\text{}^\text{T}\text{}`$. Whereas the sign in (23) and the values of the $`c_k`$s are determined by the three local invariants<sup>1</sup><sup>1</sup>1We write $`\mathrm{Sp}\left\{\right\}`$ for the trace of a dyadic in order to avoid confusion with quantum mechanical traces such as $`C_{xy}=\sigma _x\tau _y=\mathrm{Tr}\left\{\sigma _x\tau _y\mathrm{P}\right\}`$. $$\mathrm{Sp}\left\{\text{}^\text{T}\text{}\right\},\mathrm{det}\left\{\text{}\right\},\mathrm{Sp}\left\{\left(\text{}^\text{T}\text{}\right)^2\right\},$$ (25) the 123 bases are not uniquely specified by (23), however, because the simultaneous replacements $$(e_1^{},e_2^{},e_3^{})(e_1^{},e_2^{},e_3^{}),\left(\begin{array}{c}\underset{1}{\overset{}{n}}\\ \underset{2}{\overset{}{n}}\\ \underset{3}{\overset{}{n}}\end{array}\right)\left(\begin{array}{c}\underset{1}{\overset{}{n}}\\ \underset{2}{\overset{}{n}}\\ \underset{3}{\overset{}{n}}\end{array}\right),$$ (26) for example, do not change the right-hand side of (23). The resulting freedom in choosing $`e_1^{}`$, $`e_2^{}`$, $`e_3^{}`$ (which then fixes $`\underset{1}{\overset{}{n}}`$, $`\underset{2}{\overset{}{n}}`$, $`\underset{3}{\overset{}{n}}`$ unless $`c_2=0`$) is then used to enforce conventions imposed on the coefficients in $`s^{}`$ $`=`$ $`\sigma ^{}=(e_1^{},e_2^{},e_3^{})\left(\begin{array}{c}s_1\\ s_2\\ s_3\end{array}\right),`$ (30) $`\stackrel{}{t}`$ $`=`$ $`\stackrel{}{\tau }=(t_1,t_2,t_3)\left(\begin{array}{c}\underset{1}{\overset{}{n}}\\ \underset{2}{\overset{}{n}}\\ \underset{3}{\overset{}{n}}\end{array}\right).`$ (34) In brief terms, these conventions aim at making as many of the $`s_k`$s and $`t_k`$s vanish as possible and to give definite signs to as many as possible of the remaining ones. Eventually, each family is characterized by nine numbers: the values of the three local invariants in (25), the three $`s_k`$ ($`k=1,2,3`$) coefficients, and the three $`t_k`$s, some of them equal to zero and others with a known sign. Roughly speaking, of the 15 coefficients appearing in (12) and (20), six are thus used up in defining the two 123 coordinate systems, and nine to identify the family of locally equivalent states to which the given $`\mathrm{P}`$ belongs. Clearly, the nine family-defining parameters are invariant under local transformations. Degeneracy among the characteristic values of the cross dyadic distinguishes six classes of families, $$\begin{array}{cc}\hfill c_1=c_2=c_3=0:\text{class }& \mathrm{A},\hfill \\ \hfill c_1=c_2=c_3>0:\text{class }& \mathrm{B},\hfill \\ \hfill c_1>c_2=c_3=0:\text{class }& \mathrm{C},\hfill \\ \hfill c_1>c_2=c_3>0:\text{class }& \mathrm{D},\hfill \\ \hfill c_1=c_2>c_3:\text{class }& \mathrm{E},\hfill \\ \hfill c_1>c_2>c_3:\text{class }& \mathrm{F}.\hfill \end{array}$$ (35) In classes A and C the $`+`$ sign in (23) applies; both signs can occur in classes B, D, E, and F which, therefore, consist of two subclasses each. Given the local invariants of (25), we find the respective class as follows. First calculate the auxiliary quantities $`a`$ $`=`$ $`{\displaystyle \frac{9}{4}}\mathrm{Sp}\left\{\text{}^\text{T}\text{}\right\}\mathrm{Sp}\left\{\left(\text{}^\text{T}\text{}\right)^2\right\}`$ $`{\displaystyle \frac{5}{4}}\left[\mathrm{Sp}\left\{\text{}^\text{T}\text{}\right\}\right]^3+{\displaystyle \frac{27}{2}}\left[\mathrm{det}\left\{\text{}\right\}\right]^2,`$ $`b`$ $`=`$ $`{\displaystyle \frac{3}{2}}\mathrm{Sp}\left\{\left(\text{}^\text{T}\text{}\right)^2\right\}{\displaystyle \frac{1}{2}}\left[\mathrm{Sp}\left\{\text{}^\text{T}\text{}\right\}\right]^2,`$ which are subject to $`a^2b^3`$. Then we have the classification $$\begin{array}{cc}\text{class }\mathrm{A}\text{if}\hfill & a^2=b^3=0\text{and}\mathrm{det}\left\{\text{}\right\}=0,\hfill \\ \text{class }\mathrm{B}\text{if}\hfill & a^2=b^3=0\text{and}\mathrm{det}\left\{\text{}\right\}0,\hfill \\ \text{class }\mathrm{C}\text{if}\hfill & a^2=b^3>0\begin{array}{c}\text{and}a>0\hfill \\ \text{and}\mathrm{det}\left\{\text{}\right\}=0,\hfill \end{array}\hfill \\ \text{class }\mathrm{D}\text{if}\hfill & a^2=b^3>0\begin{array}{c}\text{and}a>0\hfill \\ \text{and}\mathrm{det}\left\{\text{}\right\}0,\hfill \end{array}\hfill \\ \text{class }\mathrm{E}\text{if}\hfill & a^2=b^3>0\text{and}a<0,\hfill \\ \text{class }\mathrm{F}\text{if}\hfill & a^2<b^3.\hfill \end{array}$$ (37) The generic forms for the various classes are as follows. Class A: Since $`\text{}=0`$ here, we can choose the two sets of 123 coordinates independently, and $`s_10`$, $`s_2=s_3=0`$ as well as $`t_10`$, $`t_2=t_3=0`$ specify the conventions. This class consists of a two-parametric set of families of the generic form $$\mathrm{P}=\frac{1}{4}\left(1+s\sigma _1+t\tau _1\right)\text{with}s0,t0.$$ (38) For $`s=0`$, $`t=0`$ we have the chaotic state $`\mathrm{P}_{\mathrm{chaos}}=\frac{1}{4}`$ which forms a single-state family all by itself. Class B: Here we can choose $`e_1^{},e_2^{},e_3^{}`$ freely and the conventional choice is specified by $`s^{}=e_1^{}s,\stackrel{}{t}=t_1\underset{1}{\overset{}{n}}+t_3\underset{3}{\overset{}{n}}`$ (42) with $`\left\{\begin{array}{c}s>0\text{and}t_30\\ \text{or}\\ s=0\text{and}t_1=t0,t_3=0\end{array}\right\}.`$ Each subclass \[$`\pm `$ in (23)\] consists of four-parametric sets of families. In passing we note that the so-called Werner states constitute the two class-B families with $`s=0`$ and $`t=0`$. Class C: Here the replacement (26) is used to enforce $`s_10`$ or $`t_10`$ if $`s_1=0`$. Then $`s_2=0,s_30`$ and $`t_2=0,t_30`$ are achieved by suitable rotations of $`e_2^{},e_3^{}`$ and, independently, of $`\underset{2}{\overset{}{n}},\underset{3}{\overset{}{n}}`$. In summary, this establishes $`s^{}=e_1^{}s_1+e_3^{}s_3,\stackrel{}{t}=t_1\underset{1}{\overset{}{n}}+t_3\underset{3}{\overset{}{n}},`$ (46) $`\text{}=e_1^{}c_1\underset{1}{\overset{}{n}}`$ with $`\left\{\begin{array}{c}s_10\\ \text{or}\\ s_1=0\text{and}t_10\end{array}\right\}`$ and $`s_30,t_30`$ (47) for the five-parametric sets of families. Class D: In distinction from class C, the rotations in the 23 sectors are not independent here. Thus we get $`s^{}=e_1^{}s_1+e_3^{}s_3,\stackrel{}{t}=t_1\underset{1}{\overset{}{n}}+t_2\underset{2}{\overset{}{n}}+t_3\underset{3}{\overset{}{n}},`$ (51) $`\text{}=\pm \left(e_1^{}c_1\underset{1}{\overset{}{n}}+e_2^{}c_2\underset{2}{\overset{}{n}}+e_3^{}c_2\underset{3}{\overset{}{n}}\right)`$ with $`\left\{\begin{array}{c}s_10\\ \text{or}\\ s_1=0\text{and}t_10\end{array}\right\}`$ and $`\left\{\begin{array}{c}s_30\text{and}t_20\\ \text{or}\\ s_3=0\text{and}t_2=0,t_30\end{array}\right\}.`$ (55) Each subclass contains seven-parametric sets of families. Class E: This class is very similar to class D, but now the degeneracy is in the 12 sector, and so we have $`s^{}=e_1^{}s_1+e_3^{}s_3,\stackrel{}{t}=t_1\underset{1}{\overset{}{n}}+t_2\underset{2}{\overset{}{n}}+t_3\underset{3}{\overset{}{n}},`$ (59) $`\text{}=\pm \left(e_1^{}c_1\underset{1}{\overset{}{n}}+e_2^{}c_1\underset{2}{\overset{}{n}}+e_3^{}c_3\underset{3}{\overset{}{n}}\right)`$ with $`\left\{\begin{array}{c}s_10\text{and}t_20\\ \text{or}\\ s_1=0\text{and}t_10,t_2=0\end{array}\right\}`$ and $`\left\{\begin{array}{c}s_30\\ \text{or}\\ s_3=0\text{and}t_30\end{array}\right\}.`$ (63) Here, too, each subclass is made up of seven-parametric sets of families. Class F: The lack of degeneracy limits changes of the 123 bases to discrete $`180^{}`$ rotations as in (26) where the rotation is around the 3rd axes. The generic form is then defined by that choice of 123 coordinates for which as many as possible of the coefficients $`s_1`$, $`t_1`$, $`s_2`$, $`t_2`$, $`s_3`$, $`t_3`$, are non-negative (in this order, say). Here we get, in each subclass, sets of families specified by the full number of nine parameters, of which five or more are non-negative. Arbitrary local unitary transformations turn members of a family into other members of the same family — this, we recall, is the defining property of a family of locally equivalent states. It is possible that some local transformations have no effect at all, as exemplified by $`U_{\mathrm{loc}}=\mathrm{exp}(i\phi \sigma _1+i\varphi \tau _1)`$ acting on the class-A state (38). Therefore, some families are larger than others, and determining a family’s size is a problem of considerable interest. Recent progress on this front is reported by Kuś and Życzkowski . We close this section with a single example. Pure states are of the generic form $`\mathrm{P}_{\mathrm{pure}}={\displaystyle \frac{1}{4}}\left(1+p\sigma _1p\tau _1\sigma _1\tau _1q\sigma _2\tau _2q\sigma _3\tau _3\right)`$ $`\text{with}0p1,q=\sqrt{1p^2}0.`$ (64) One verifies easily the purity condition $$\mathrm{P}_{\mathrm{pure}}\left(1\mathrm{P}_{\mathrm{pure}}\right)=0.$$ (65) For $`p=0,q=1`$ we have the family of Bell states, $$\mathrm{P}_{\mathrm{Bell}}=\frac{1}{4}\left(1\sigma _1\tau _1\sigma _2\tau _2\sigma _3\tau _3\right),$$ (66) which is in class B; the $`p=1,q=0`$ family consists of the product states $`\frac{1}{2}(1+\sigma _1)\frac{1}{2}(1\tau _1)`$ and is in class C; and the $`0<p<1`$ families belong to class D. These families are of different sizes : three-dimensional, four-dimensional, and five-dimensional, respectively. ## 3 Positivity and separability An arbitrary choice for the (real) coefficients in (8) and (12) or, equivalently, of the nine family-defining parameters plus the 123 coordinate systems specifies a hermitian $`\mathrm{P}`$ of unit trace, but its positivity must be ensured by imposing restrictions on the Pauli vectors $`s^{}`$, $`\stackrel{}{t}`$, and the cross dyadic $`\text{}`$. It is expedient to switch the emphasis from $`\mathrm{P}`$ to the traceless operator $`K`$ introduced by $$\mathrm{P}=\frac{1}{4}(1K),K=14\mathrm{P},$$ (67) so that $`\mathrm{P}0`$ requires $$K1.$$ (68) Convex sums of two states are weighted sums of their $`K`$s. Admixing $`\mathrm{P}_{\mathrm{chaos}}`$ to a given $`\mathrm{P}`$ amounts to multiplying its $`K`$ by a factor. One could, of course, check the positivity criterion (68) by calculating the eigenvalues of $`K`$. For such purposes, it is often very convenient to use a $`4\times 4`$-matrix representation in which $`\sigma _{1,2,3}`$ and $`\tau _{1,2,3}`$ have imaginary antisymmetric matrices, $`\stackrel{}{\sigma }s^{}+\stackrel{}{t}\tau ^{}`$ (73) $`\widehat{=}\left(\begin{array}{cccc}0& i\left(s_1+t_1\right)& +i\left(s_2+t_2\right)& i\left(s_3t_3\right)\\ +i\left(s_1+t_1\right)& 0& +i\left(s_3+t_3\right)& +i\left(s_2t_2\right)\\ i\left(s_2+t_2\right)& i\left(s_3+t_3\right)& 0& +i\left(s_1t_1\right)\\ +i\left(s_3t_3\right)& i\left(s_2t_2\right)& i\left(s_1t_1\right)& 0\end{array}\right),`$ and products $`\sigma _j\tau _k`$ ($`j,k=1,2,3`$) have real symmetric matrices, in particular $`\stackrel{}{\sigma }\text{}t^{}=\pm {\displaystyle \underset{k=1}{\overset{3}{}}}\sigma _kc_k\tau _k`$ (79) $`\widehat{=}\pm \left(\begin{array}{cccc}c_1+c_2c_3& 0& 0& 0\\ 0& c_1c_2+c_3& 0& 0\\ 0& 0& c_1+c_2+c_3& 0\\ 0& 0& 0& c_1c_2c_3\end{array}\right)`$ is diagonal. But precise knowledge of the actual eigenvalues of $`K`$ is not needed if we only want to verify (68). Since $`K`$ is traceless, its eigenvalues $`\kappa _j`$ ($`j=1,2,3,4`$) have a vanishing sum and solve a quartic equation without a cubic term, $$\kappa ^4A_2\kappa ^2+A_1\kappa A_0=0,$$ (81) where $`A_2`$ $`=`$ $`{\displaystyle \frac{1}{2}}\mathrm{Tr}\left\{K^2\right\},`$ $`A_1`$ $`=`$ $`{\displaystyle \frac{1}{3}}\mathrm{Tr}\left\{K^3\right\},`$ $`A_0`$ $`=`$ $`{\displaystyle \frac{1}{4}}\mathrm{Tr}\left\{K^4\right\}{\displaystyle \frac{1}{8}}\left[\mathrm{Tr}\left\{K^2\right\}\right]^2.`$ (82) These three numbers are invariant under arbitrary (local or not) unitary transformations; they are three independent global invariants of the given $`\mathrm{P}`$. Expressed in terms of $`s^{}`$, $`\stackrel{}{t}`$, and $`\text{}`$ they read $`A_2`$ $`=`$ $`2\mathrm{Sp}\left\{\text{}^\text{T}\text{}\right\}+2(s^2+t^2),`$ $`A_1`$ $`=`$ $`8\mathrm{det}\left\{\text{}\right\}8\stackrel{}{s}\text{}t^{},`$ $`A_0`$ $`=`$ $`\left({\displaystyle \frac{1}{2}}A_2\right)^2+2\left[\mathrm{Sp}\left\{\text{}^\text{T}\text{}\right\}\right]^2+4s^2t^2`$ (83) $`+4\stackrel{}{s}\text{}\text{}^\text{T}s^{}+4\stackrel{}{t}\text{}^\text{T}\text{}t^{}`$ $`+8\mathrm{det}\left\{\text{}\right\}8\mathrm{det}\left\{\text{}\right\}`$ where $$\text{}=\text{}s^{}\stackrel{}{t}$$ (84) is the entanglement dyadic. As we see, the traces of (82) involve nine different local polynomial invariants of $`s^{}`$, $`\stackrel{}{t}`$, and $`\text{}`$, and it is clear that their values are determined by the nine family-specifying parameters of classes A,…, F. Suggestive as it is, the converse is not true,<sup>2</sup><sup>2</sup>2Therefore, the assertion “All other local invariants …” shortly after (19) in is false. as can be demonstrated by a counter example. Consider, for instance, the two states $`\mathrm{P}_1`$ $`=`$ $`{\displaystyle \frac{1}{4}}\left(1+{\displaystyle \frac{1}{4}}\sigma _1+{\displaystyle \frac{1}{2}}\sigma _3+{\displaystyle \frac{1}{2}}\sigma _1\tau _1+{\displaystyle \frac{1}{4}}\sigma _2\tau _2\right),`$ $`\mathrm{P}_2`$ $`=`$ $`{\displaystyle \frac{1}{4}}\left(1+{\displaystyle \frac{1}{2}}\sigma _2+{\displaystyle \frac{1}{4}}\sigma _3+{\displaystyle \frac{1}{2}}\sigma _1\tau _1+{\displaystyle \frac{1}{4}}\sigma _2\tau _2\right),`$ (85) which belong to two different class-F families, but all terms in (83) are the same for $`\mathrm{P}_1`$ and $`\mathrm{P}_2`$. Whereas the nine polynomial invariants of (83) do not always suffice to determine the values of all local invariants, the nine parameters that specify the family certainly do. They are, however, not given by (traces of) polynomials of $`s^{}`$, $`\stackrel{}{t}`$, and $`\text{}`$. According to Makhlin , there are 18 polynomial invariants whose values uniquely characterize the family in question (actually, of nine of them only the sign matters). In addition to the nine invariants in (83), which exhaust the polynomials of degree 4 or lower, there are nine invariants of higher degree in Makhlin’s set, which do not enter the three global invariants $`A_0`$, $`A_1`$, $`A_2`$. All solutions of the quartic equation (81) are real by construction — it is, after all, the characteristic polynomial of a hermitian operator. Then, if all solutions are in the range $`\kappa 1`$, this polynomial and its derivatives must be non-negative for $`\kappa 1`$. Consequently, the positivity requirement (68) implies $$A_2A_1+A_01,2A_2A_14,A_26.$$ (86) The converse is also true: If these three inequalities are obeyed, the four real solutions of (81) are in the range $`\kappa 1`$, so that $`K1`$ and $`\mathrm{P}0`$. In other words, the restrictions on $`s^{}`$, $`\stackrel{}{t}`$, and $`\text{}`$ alluded to at the beginning of this section are just the inequalities (86). Although the equivalence of (68) and (86) is rather obvious, a clear-cut demonstration of the case could be of interest to some readers. We give one in the Appendix. If the entanglement dyadic $`\text{}`$ vanishes, the state in question is of product form, $$\mathrm{P}=\frac{1}{2}\left(1+\stackrel{}{\sigma }s^{}\right)\frac{1}{2}\left(1+\stackrel{}{t}\tau ^{}\right),$$ (87) so that results of measurements on the first q-bit show no correlations whatsoever with measurement results concerning the second q-bit. Under these circumstances the 2–q-bit system is not entangled. Entangled q-bit pairs, $`\text{}0`$, may be in a mixed state blended from disentangled ingredients, $`\mathrm{P}={\displaystyle \underset{n}{}}w_n{\displaystyle \frac{1}{2}}\left(1+\stackrel{}{\sigma }s_n^{}\right){\displaystyle \frac{1}{2}}\left(1+\underset{n}{\overset{}{t}}\tau ^{}\right)`$ $`\text{with}w_n>0,{\displaystyle \underset{n}{}}w_n=1;`$ (88) then all correlations found in the measurement data can be understood classically. States of this kind are called separable. As an elementary example, consider the pure states (64): For $`p=1`$ they are not entangled and therefore separable, for $`p<1`$ they are entangled and not separable. Correlations of a genuine quantum character require a non-separable state $`\mathrm{P}`$. According to a criterion that we owe to Peres as well as M., P., and R. Horodecki , a given state $`\mathrm{P}`$ is separable if $`\stackrel{~}{\mathrm{P}}`$ $`=`$ $`{\displaystyle \frac{1}{4}}\left(1\stackrel{}{\sigma }s^{}+\stackrel{}{t}\tau ^{}\stackrel{}{\sigma }\text{}\tau ^{}\right)`$ (89) $`=`$ $`{\displaystyle \frac{1}{2}}\left(\stackrel{}{\sigma }\mathrm{P}\sigma ^{}\mathrm{P}\right)`$ is positive and only then. Let’s call $`\stackrel{~}{\mathrm{P}}`$ the PH<sub>3</sub> transform of $`\mathrm{P}`$; it is unitarily equivalent to the partial transpose originally considered by Peres and the Horodeckis. The positivity of $`\stackrel{~}{\mathrm{P}}`$, or $$14\stackrel{~}{\mathrm{P}}=\stackrel{~}{K}=\frac{1}{2}\left(\stackrel{}{\sigma }K\sigma ^{}K\right)1,$$ (90) can be checked analogously to the positivity of $`\mathrm{P}`$. Now, the quartic equation solved by the eigenvalues $`\stackrel{~}{\kappa }_1,\mathrm{},\stackrel{~}{\kappa }_4`$ of $`\stackrel{~}{K}`$ is $`\stackrel{~}{\kappa }^4A_2\stackrel{~}{\kappa }^2+\left(A_1+16\mathrm{det}\left\{\text{}\right\}\right)\stackrel{~}{\kappa }`$ $`\left(A_016\mathrm{det}\left\{\text{}\right\}+16\mathrm{det}\left\{\text{}\right\}\right)=0,`$ (91) so that $`A_2A_1+A_0`$ $``$ $`1+16\mathrm{det}\left\{\text{}\right\},`$ $`2A_2A_1`$ $``$ $`4+16\mathrm{det}\left\{\text{}\right\},`$ $`A_2`$ $``$ $`6.`$ (92) are equivalent to (90); the third is always obeyed by a positive $`\mathrm{P}`$. So, a non-separable state must violate either the first or the second inequality, or both. The equal sign holds in the first inequality, if the PH<sub>3</sub> transform of the given $`\mathrm{P}`$ has a zero eigenvalue; the first and the second are equalities, if the PH<sub>3</sub> transform has two zero eigenvalues. Accordingly, the $`\stackrel{~}{\mathrm{P}}`$ of a non-separable $`\mathrm{P}`$ can at most have one zero eigenvalue and thus must be of rank 3 or 4. While we are at it, let us also mention that the PH<sub>3</sub> transform of any state $`\mathrm{P}`$ can have at most a single negative eigenvalue \[see below at Eq. (98)\]. Thus the separability of a given $`\mathrm{P}`$ is checked as easily as its positivity. Neither test requires actual knowledge of the solutions of (81) or (91). They could, of course, be stated analytically but these explicit expressions are not very transparent unless special relations exist among the coefficients of the quartic equations. As an immediate implication of the PH<sub>3</sub> criterion, in the form of the inequalities (92), we note that a state $`\mathrm{P}`$ with $`\mathrm{det}\left\{\text{}\right\}0`$ and $`\mathrm{det}\left\{\text{}\right\}0`$ is surely separable. Therefore, for example, all states in classes A and C are separable. ## 4 Lewenstein-Sanpera decompositions As Lewenstein and Sanpera observed , any 2–q-bit state $`\mathrm{P}`$ can be written as a convex sum of a separable state and a pure state, $$\mathrm{P}=\lambda \mathrm{P}_{\mathrm{sep}}+(1\lambda )\mathrm{P}_{\mathrm{pure}}\text{with}0\lambda 1.$$ (93) Rare exceptions aside, the Lewenstein-Sanpera decomposition (LSD) of a given (non-separable) $`\mathrm{P}`$ is not unique, there is usually a continuum of LSDs to choose from. Among them is the optimal LSD, the one with the largest value of $`\lambda `$, $$\mathrm{P}=𝒮\mathrm{P}_{\mathrm{sep}}^{(\mathrm{opt})}+(1𝒮)\mathrm{P}_{\mathrm{pure}}^{(\mathrm{opt})}\text{with}𝒮=\mathrm{max}\left\{\lambda \right\},$$ (94) and we call $`𝒮`$, the maximal $`\lambda `$ value, the degree of separability of $`\mathrm{P}`$. Without presently attempting to be precise about this matter, we repeat the remark in that “a state $`\mathrm{P}`$ is the more useful for quantum communication purposes, the smaller its degree of separability.” The spectral decomposition of the PH<sub>3</sub> transform of a pure state (64) is of the generic form $`\stackrel{~}{\mathrm{P}}_{\mathrm{pure}}`$ $`=`$ $`{\displaystyle \frac{1}{4}}\left(1p\sigma _1p\tau _1+\sigma _1\tau _1+q\sigma _2\tau _2+q\sigma _3\tau _3\right)`$ $`=`$ $`{\displaystyle \frac{1+p}{2}}\mathrm{P}_{\mathrm{pure}}^{(1)}+{\displaystyle \frac{1p}{2}}\mathrm{P}_{\mathrm{pure}}^{(2)}+{\displaystyle \frac{q}{2}}\mathrm{P}_{\mathrm{pure}}^{(3)}{\displaystyle \frac{q}{2}}\mathrm{P}_{\mathrm{pure}}^{(4)}`$ with $$\begin{array}{c}\mathrm{P}_{\mathrm{pure}}^{(1)}\hfill \\ \mathrm{P}_{\mathrm{pure}}^{(2)}\hfill \end{array}\}=\frac{1}{4}\left(1\sigma _1\tau _1+\sigma _1\tau _1\right),$$ (96) which are pure states of the separable class-C kind, and $$\begin{array}{c}\mathrm{P}_{\mathrm{pure}}^{(3)}\hfill \\ \mathrm{P}_{\mathrm{pure}}^{(4)}\hfill \end{array}\}=\frac{1}{4}\left(1\sigma _1\tau _1\pm \sigma _2\tau _2\pm \sigma _3\tau _3\right),$$ (97) which are Bell states (non-separable, class B). Therefore, the PH<sub>3</sub> transform $`\stackrel{~}{\mathrm{P}}`$ of any 2–q-bit state $`\mathrm{P}`$ can be written as $$\stackrel{~}{\mathrm{P}}=(1+x)\mathrm{P}^{}x\mathrm{P}_{\mathrm{Bell}},0x\frac{1}{2}(1𝒮)$$ (98) with some state $`\mathrm{P}^{}`$ and a Bell state $`\mathrm{P}_{\mathrm{Bell}}`$. As a consequence, $`\stackrel{~}{\mathrm{P}}`$ can have at most one negative eigenvalue, so that only one solution of (91) can be in the range $`\stackrel{~}{\kappa }>1`$, as noted above. Since $`\mathrm{P}^{}`$ is a mixture of four or fewer pure states, (98) shows that the PH<sub>3</sub> transform of a non-separable state is a pseudo-mixture of up to five pure states with one negative weight only, carried by a Bell state. There is a very similar observation by Sanpera, Tarrach, and Vidal about $`\mathrm{P}`$ itself: It can always be presented as a pseudo-mixture of four or five separable pure states; as an immediate consequence its PH<sub>3</sub> transform is also such a pseudo-mixture. The optimal LSD (94) has a number of properties that help in decomposing given states in the optimal way. Let’s briefly consider some particularly important ones. Existence: The degree of separability $`𝒮`$ is really the maximum of all possible $`\lambda `$ values in (93), not just their supremum, because the subset of separable states is compact. Therefore, a LSD with $`\lambda =𝒮`$ does exist. Uniqueness: If we have two different LSDs with the same non-zero value of $`\lambda `$, their symmetric convex sum also equals the given $`\mathrm{P}`$. It contains the convex sum of the two different $`\mathrm{P}_{\mathrm{sep}}`$s, which is separable, and the convex sum of the two $`\mathrm{P}_{\mathrm{pure}}`$s, which has LSDs of its own. Either one of them contains a separable part, so that we get a new LSD of the $`\mathrm{P}`$ in question with a larger $`\lambda `$ value. Consequently, the common $`\lambda `$ of the original two LSDs is not maximal, and it follows that the optimal LSD is unique. This does not imply that one can always find another LSD with the same $`\lambda `$ value if $`\lambda <𝒮`$. There are $`\mathrm{P}`$s with a continuum of LSDs in which each value of $`\lambda `$ occurs only once.<sup>3</sup><sup>3</sup>3Therefore, the ‘only’ is too strong in the assertion “Only $`\mathrm{P}_{\mathrm{sep}}`$ and $`\mathrm{P}_{\mathrm{pure}}`$ …” after (15) in . Examples are the rank-2 states of (56) in that obey inequality (61) in . $`\mathrm{P}_{\mathrm{sep}}^{(\mathrm{opt})}`$ is barely separable: Consider the optimal LSD of some non-separable $`\mathrm{P}`$ and a parameter $`ϵ`$ in the range $`0<ϵ1𝒮`$. In $`\mathrm{P}`$ $`=`$ $`(𝒮+ϵ)\left[{\displaystyle \frac{𝒮}{𝒮+ϵ}}\mathrm{P}_{\mathrm{sep}}^{(\mathrm{opt})}+{\displaystyle \frac{ϵ}{𝒮+ϵ}}\mathrm{P}_{\mathrm{pure}}^{(\mathrm{opt})}\right]`$ (99) $`+(1𝒮ϵ)\mathrm{P}_{\mathrm{pure}}^{(\mathrm{opt})}`$ the convex sum in square brackets is surely non-negative, but cannot be separable. Because, if it were, we would have found a LSD with $`\lambda >𝒮`$. Thus the PH<sub>3</sub> transform of $`\left[\mathrm{}\right]`$ has a negative eigenvalue for $`ϵ>0`$, but none for $`ϵ=0`$. Since the eigenvalues are continuous functions of $`ϵ`$, the PH<sub>3</sub> transform of $`\mathrm{P}_{\mathrm{sep}}^{(\mathrm{opt})}`$ must have at least one zero eigenvalue. Formally, $$\stackrel{~}{\mathrm{P}}_{\mathrm{sep}}^{(\mathrm{opt})}0\text{but not}\stackrel{~}{\mathrm{P}}_{\mathrm{sep}}^{(\mathrm{opt})}>0;$$ (100) for $`\mathrm{P}_{\mathrm{sep}}^{(\mathrm{opt})}`$, the equal sign holds in the first inequality of (92). A useful terminology calls $`\mathrm{P}_{\mathrm{sep}}^{(\mathrm{opt})}`$ barely separable with respect to $`\mathrm{P}_{\mathrm{pure}}^{(\mathrm{opt})}`$. When searching for the optimal LSD of a given $`\mathrm{P}`$ it is, therefore, sufficient to consider LSDs with $`\mathrm{P}_{\mathrm{sep}}`$s that are barely separable with respect to the $`\mathrm{P}_{\mathrm{pure}}`$ with which they are paired in (93). If the $`\mathrm{P}_{\mathrm{sep}}`$ of some LSD does not have this property, one adds the appropriate amount of the respective $`\mathrm{P}_{\mathrm{pure}}`$ to it (in the sense of a convex sum, of course) and gets a barely separable $`\mathrm{P}_{\mathrm{sep}}`$. Local invariance is passed on: Suppose that the $`\mathrm{P}`$ considered is invariant under some local unitary transformation, $$U_{\mathrm{loc}}^{}\mathrm{P}U_{\mathrm{loc}}^{}=\mathrm{P}.$$ (101) Then its $`\mathrm{P}_{\mathrm{sep}}^{(\mathrm{opt})}`$ and $`\mathrm{P}_{\mathrm{pure}}^{(\mathrm{opt})}`$ must be invariant under this local transformation as well. Otherwise we could apply it to the optimal LSD and get another LSD with the same $`\lambda `$ value, in conflict with the uniqueness of the optimal LSD. This argument builds on the elementary observation that local transformations do not affect the purity and separability of a state. The limitations resulting from this “inheritance of local invariance” can facilitate the search for the optimal LSD substantially. The optimal decompositions of the (generalized) Werner states reported in were found this way. Swapping invariance is passed on: Similarly one finds that the $`\mathrm{P}_{\mathrm{sep}}^{(\mathrm{opt})}`$ and $`\mathrm{P}_{\mathrm{pure}}^{(\mathrm{opt})}`$ of a $`\mathrm{P}`$ that is invariant under the swapping transformation $$\sigma _k\tau _k\text{for}k=1,2,3$$ (102) must be invariant themselves because swapping does not affect the separability or the purity of a state. Clearly, this swapping invariance is only possible if the Pauli vectors $`s^{}`$ and $`\stackrel{}{t}`$ are of equal length. Orthogonality is passed on: If the $`\mathrm{P}`$ in question is orthogonal to a certain other state $`\mathrm{P}_{}`$, $$\mathrm{Tr}\left\{\mathrm{PP}_{}\right\}=0,$$ (103) then the $`\mathrm{P}_{\mathrm{sep}}`$s and $`\mathrm{P}_{\mathrm{pure}}`$s of all LSDs of $`\mathrm{P}`$ are also orthogonal to $`\mathrm{P}_{}`$ because both traces in $$0=\lambda \mathrm{Tr}\left\{\mathrm{P}_{\mathrm{sep}}\mathrm{P}_{}\right\}+(1\lambda )\mathrm{Tr}\left\{\mathrm{P}_{\mathrm{pure}}\mathrm{P}_{}\right\}$$ (104) must be non-negative, so both must vanish. In particular, the $`\mathrm{P}_{\mathrm{sep}}^{(\mathrm{opt})}`$ and $`\mathrm{P}_{\mathrm{pure}}^{(\mathrm{opt})}`$ of $`\mathrm{P}`$ must have this orthogonality property. The optimal LSDs of rank-2 states reported in were found by exploiting this “inheritance of orthogonality.” ## 5 Degree of separability and concurrence In their studies of what they call “entanglement of formation,” Hill and Wootters consider $`\overline{\mathrm{P}}`$ $`=`$ $`{\displaystyle \frac{1}{4}}\left(1\stackrel{}{\sigma }s^{}\stackrel{}{t}\tau ^{}+\stackrel{}{\sigma }\text{}\tau ^{}\right)`$ (105) $`=`$ $`{\displaystyle \frac{1}{2}}\left(\stackrel{}{\tau }\stackrel{~}{\mathrm{P}}\tau ^{}\stackrel{~}{\mathrm{P}}\right);`$ let’s call it the HW transform of $`\mathrm{P}`$. Since the replacement $`(s^{},\stackrel{}{t},\text{})(s^{},\stackrel{}{t},\text{})`$ changes none of the local invariants in (83), $`\overline{\mathrm{P}}`$ has the same eigenvalues as $`\mathrm{P}`$ and, therefore, $`\overline{\mathrm{P}}`$ is unitarily equivalent to $`\mathrm{P}`$. Equally well we could argue that the matrix representations of $`\mathrm{P}`$ and $`\overline{\mathrm{P}}`$, composed of the ingredients in (LABEL:eq:sigtau-matr) and (LABEL:eq:sCt-matr), are complex conjugates or transposes of each other, and so they must have the same real eigenvalues. Note that a 2–q-bit state and its HW transform are always in the same class of states but they may or may not belong to the same family; their unitary equivalence may be local or global. Hill and Wootters use the HW transform to introduce the concurrence $`𝒞`$ of $`\mathrm{P}`$. It is given by $$𝒞=\mathrm{max}\{0,2\underset{k}{\mathrm{max}}\left\{r_k\right\}\underset{k}{}r_k\},$$ (106) where $`r_1,r_2,r_3,r_4`$ are the four non-negative eigenvalues of $$\text{ }\sqrt{\mathrm{P}}\sqrt{\overline{\mathrm{P}}}\text{ }=\sqrt{\sqrt{\overline{\mathrm{P}}}\mathrm{P}\sqrt{\overline{\mathrm{P}}}}.$$ (107) The roles of $`\mathrm{P}`$ and $`\overline{\mathrm{P}}`$ can be interchanged in this definition of $`𝒞`$; thus the concurrence of $`\overline{\mathrm{P}}`$ is equal to the concurrence of $`\mathrm{P}`$. For instance, the concurrence of a pure state (64) is $`q`$. Separable states ($`𝒮=1`$) have $`𝒞=0`$ and non-separable states ($`𝒮<1`$) have $`𝒞>0`$. This suggests that there might be a close relation between the degree of separability and the concurrence. Indeed, $`𝒮+𝒞=1`$ holds if $`s^{}=0`$ and $`\stackrel{}{t}=0`$ — such $`\mathrm{P}`$s are generalized Werner states of the first kind in the terminology of — but more generally we find $$0<𝒮+𝒞1.$$ (108) Pure states (64) have $`𝒮+𝒞=1`$ if $`q=0`$ and $`𝒮+𝒞=q`$ if $`q>0`$. A set of non-pure states exploring the whole range of (108) is given by the rank-2 states $$\mathrm{P}=\frac{1}{4}\left(1+(\sigma _3+x\tau _3)\mathrm{sin}\theta +(\sigma _1\tau _1x\sigma _2\tau _2)\mathrm{cos}\theta +x\sigma _3\tau _3\right)$$ (109) with $`1<x<1`$, for which $`𝒮=\left\{\begin{array}{cc}\hfill 1\text{if}& \mathrm{cos}\theta =0\hfill \\ \hfill 1\text{ }x\text{ }\text{if}& \mathrm{cos}\theta 0\hfill \end{array}\right\},𝒞=\text{ }x\mathrm{cos}\theta \text{ }.`$ (112) We surmise that (108) is obeyed by all 2–q-bit states. ## Appendix We write the four solutions of (81) in terms of three parameters, $`\begin{array}{c}\kappa _1\hfill \\ \kappa _2\hfill \end{array}\}`$ $`=`$ $`\pm (\lambda _1\lambda _2)\lambda _3,`$ (115) $`\begin{array}{c}\kappa _3\hfill \\ \kappa _4\hfill \end{array}\}`$ $`=`$ $`(\lambda _1+\lambda _2)+\lambda _3,`$ (118) thereby taking care of $`\mathrm{Tr}\left\{K\right\}=\kappa _1+\kappa _2+\kappa _3+\kappa _4=0`$. The coefficients in (81) are then given by $`A_2`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(\kappa _1^2+\kappa _2^2+\kappa _3^2+\kappa _4^2\right)`$ $`=`$ $`2\left(\lambda _1^2+\lambda _2^2+\lambda _3^2\right),`$ $`A_1`$ $`=`$ $`\kappa _1\kappa _2(\kappa _1+\kappa _2)+\kappa _3\kappa _4(\kappa _3+\kappa _4)`$ $`=`$ $`8\lambda _1\lambda _2\lambda _3,`$ $`A_0`$ $`=`$ $`\kappa _1\kappa _2\kappa _3\kappa _4=2\left(\lambda _1^2\lambda _2^2+\lambda _2^2\lambda _3^2+\lambda _3^2\lambda _1^2\right)`$ (119) $`\left(\lambda _1^4+\lambda _2^4+\lambda _3^4\right).`$ Now look at the second inequality in (86). It says $$\lambda _1^2+\lambda _2^2+\lambda _3^2+2\lambda _1\lambda _2\lambda _31$$ (120) or, singling out $`\lambda _3`$, $$\left(\lambda _1\lambda _2+\lambda _3\right)^2\left(1\lambda _1^2\right)\left(1\lambda _2^2\right),$$ (121) and cyclic permutations produce two analogous statements in which $`\lambda _1`$ and $`\lambda _2`$ are privileged. Since the left-hand sides of these three equations are non-negative, it follows that all $`\lambda _j^2`$ exceed unity if one of them does. Combined with the third inequality in (86), here reading $$\lambda _1^2+\lambda _2^2+\lambda _3^23,$$ (122) this implies $$\lambda _1^21,\lambda _2^21,\lambda _3^21.$$ (123) In conjunction with $`\lambda _1`$ $`=`$ $`{\displaystyle \frac{1}{2}}(\kappa _1+\kappa _4)={\displaystyle \frac{1}{2}}(\kappa _2+\kappa _3),`$ $`\lambda _2`$ $`=`$ $`{\displaystyle \frac{1}{2}}(\kappa _2+\kappa _4)={\displaystyle \frac{1}{2}}(\kappa _1+\kappa _3),`$ $`\lambda _3`$ $`=`$ $`{\displaystyle \frac{1}{2}}(\kappa _3+\kappa _4)={\displaystyle \frac{1}{2}}(\kappa _1+\kappa _2)`$ (124) this means that at most one of the four $`\kappa `$s can be larger than $`1`$, and that then the other three must be less than $`1`$. But the first inequality in (86), $$1A_2+A_1A_0=(1\kappa _1)(1\kappa _2)(1\kappa _3)(1\kappa _4)0,$$ (125) excludes this possibility because one negative factor and three positive factors would yield a negative product. Therefore, all three inequalities (86) can only be obeyed if all four $`\kappa `$s are less than $`1`$. In other words: (86) implies (68) indeed. Note that this reasoning is only valid if one knows, as we do, that all solutions of the quartic equation (81) are real. This property itself is not guaranteed by the inequalities (86), as shown by $`A_2=2`$, $`A_1=0`$, $`A_0=2`$ when $`(\kappa ^21)^2+1=0`$. ## Acknowledgments BGE thanks Marek Kuś, Maciej Lewenstein, and Karol Życzkowski for particularly helpful discussions. NM would like to thank the Egyptian government for granting a fellowship. This paper has been submitted to Applied Physics B as a contribution to the Proceedings of the DPG Spring Meeting, held in Bonn, April 2000.
warning/0007/cond-mat0007109.html	ar5iv	text	# Dynamics and Critical Behaviour of the q-model ## I Introduction It is fortunate that in physics the same equations sometimes arise in contexts that are apparently very different. Feynman illustrates this through elementary examples in his introductory lectures on physics to impart the lesson that the “same equations have the same solutions” . Our purpose is to study a model, recently dubbed the $`q`$model, that provides another such instance. The $`q`$model has been used to describe the merging of tributaries to form rivers ; the aggregation of diffusing charges ; the flow of stress in a granular medium ; and can be mapped onto the abelian sandpile, a model studied in context of self-organised criticality . It is also closely related to models that describe the surface of a quantum Hall multilayer and passive scalar turbulence . Here we focus on the application to granular matter, river networks and the quantum Hall multilayer. Granular matter exhibits fascinating behaviour that is little understood . Examples of granular matter include sand, powders and agricultural grains stored in silos. An important problem is the propagation of stress through a granular medium at rest. This has been studied by ingenious experiments, in which a vertical load is applied to an amorphous pack of beads, and the loads on the beads in the top and bottom layers are recorded using carbon paper . Such experiments yield the distribution of load on the beads and reveal that there are no horizontal correlations in load even amongst neighbouring beads. The $`q`$-model was introduced by Coppersmith and coworkers to account for the distribution of load . As we shall see, it also correctly predicts the lack of horizontal correlation. For simplicity we describe the $`q`$-model in a plane. Since the vertical and horizontal directions are treated asymmetrically we call this the 1+1 dimensional $`q`$-model. The extension to 2+1 dimensions (relevant to experiments on bead packs) and higher, is straightforward and is discussed in section V. In the $`q`$-model it is assumed that the beads sit on a regular lattice shown in fig 1. The location of the beads is specified by the co-ordinates $`t`$ (the depth of the layer) and $`n`$ (the location of the bead within the layer). Note that $`n`$ takes only even values for $`t`$ even; only odd, for $`t`$ odd. Each bead is assumed to be supported by its two nearest neighbours in the layer directly below. More precisely, it is assumed that a random fraction $`f_n(t)`$ of the load of bead $`(n,t)`$ is supported by the neighbour to the left, bead $`(n1,t+1)`$; the remainder, $`1f_n(t)`$, by the neighbour to the right, bead $`(n+1,t+1)`$. Denoting the load on a bead $`w`$ and its weight $`I`$ we may write $`w_n(t)`$ $`=`$ $`w_{n1}(t1)[1f_{n1}(t1)]`$ (2) $`+w_{n+1}(t1)f_{n+1}(t1)+I_n(t).`$ The content of eq (1) is that the load on a bead is the sum of the loads transmitted to it by its neighbours in the layer above plus its own weight. The last term in eq (1) is called the injection term. Once the fractions are specified, a given load on the top layer can be propagated downward by use of eq (1). In the $`q`$-model it is assumed that the fractions are independent, identically distributed random variables. The distribution is assumed to be symmetric about $`f=1/2`$ to avoid introducing a horizontal drift to the flow of stress; in other words it is assumed $`P(f)=P(1f)`$. There is no other restriction. Thus the $`q`$-models really constitute an enormous family of models corresponding to different symmetric distributions $`P(f)`$. To fully specify a particular model it is necessary to choose the distribution $`P(f)`$. One obvious possibility is to take $`P(f)`$ to be uniformly distributed over the unit interval; another is to assume that the fractions must be 0 or 1 with equal probability. The latter is called the singular distribution. Mathematically, the $`q`$-model is a problem of random walkers that coalesce upon contact and fission spontaneously. The singular distribution corresponds to the case that the walkers coalesce but do not fission. Coppersmith et al. argued that, neglecting injection, at sufficient depth the distribution of load would attain a steady state . They studied $`\mathrm{\Pi }(w,t\mathrm{})`$, the probability distribution of load on beads in a sufficiently deep layer. For almost all distributions $`P(f)`$, except the singular distribution, they concluded that $`\mathrm{\Pi }(w,t\mathrm{})`$ decays exponentially for large $`w`$. This agrees with experiment and constitutes an important success of the $`q`$-model. For the singular distribution, Coppersmith et al. argued that $`\mathrm{\Pi }(w,\mathrm{})`$ follows a power law. Hence they conjectured that the singular distribution constitutes a critical point in the family of $`q`$-models. A major goal of this paper is to make this analogy to thermodynamic critical phenomena precise by detailed analysis of the critical point. In spite of the success mentioned above the $`q`$-model cannot be considered a complete theory of stress propagation in granular matter. This is clear both empirically and on grounds of internal consistency. Since the publication of the $`q`$-model, interesting new ideas on the subject of stress flow have appeared , but in this paper we restrict attention to the $`q`$-model. This seems justified because the $`q`$-model does capture some elements of the physics correctly and because it exhibits non-trivial critical behaviour that is interesting in its own right. Further motivation to study the $`q`$-model and particularly its critical point comes from hydrology. To make contact with that subject consider a singular $`q`$-model with zero injection and imagine that only a few beads in the top layer are loaded. The load then zig-zags downwards, perhaps along the lines shown in fig 2. If we interpret these lines as tributaries merging to form a river we arrive at Scheidegger’s model which appeared in the hydrology literature more than thirty years ago <sup></sup><sup></sup>*Parenthetically we note that Scheidegger’s model is purely descriptive in the sense that it is a recipe to draw statistically realistic networks. Somewhat different in spirit are models that seek to represent physical processes, sometimes very crudely, by which the network forms. Two examples of such models in the recent Physics literature are refs . The model of Leheny and Nagel for example describes an apocalyptic lattice world with discrete time. Each time step brings precipitation, and in its wake, erosion and avalanches. Realistic networks result.. Networks of tributaries in river basins are known empirically to be scale invariant structures that obey a variety of power laws. Scheidegger networks too obey these laws and are in this statistical sense extremely realistic representations of river basins. An excellent discussion of river basin power laws is given in refs . Ref presents some discussion of data; ref provides a detailed comparison between real and Scheidegger networks. Here we wish to point out that non-singular $`q`$-models too can be interpreted as models of river networks. For example, consider a model in which the fractions can take only the values 0, 1/2 and 1 with probability $`(1\delta )/2`$, $`\delta `$ and $`(1\delta )/2`$ respectively. This model reduces to Scheidegger’s as $`\delta 0`$. It produces networks similar to Scheidegger’s except that occasionally streams split to form distributaries. Thus this network is topologically distinct from Scheidegger networks. More significantly, as we show below, a network with non-zero $`\delta `$ is not scale-invariant. This is reminiscent of a river network model studied by Narayan and Fisher . In their “rocky-river” model too the network is not scale invariant except if a model parameter is tuned to a special (critical) value. Effectively this tuning parameter also controls river splitting. Taken together, these results suggest that river splitting is a relevant perturbation that spoils the scale invariant structure of networks. In this paper we concentrate on showing that $`q`$-model networks with river splitting are not scale invariant. We do not explore whether such non-scale invariant networks are realised in nature (for further discussion and speculation in this regard, however, see section VII). A quantum Hall multilayer consists of layers of two-dimensional electron gases stacked vertically. Multilayers can be realised by fabricating an appropriate GaAs heterostructure . They are also realised naturally in some organic salts. In a quantum Hall multilayer a sufficiently large magnetic field is applied perpendicular to the layers so that the lowest Landau level in each layer is fully occupied. Under this circumstance the only important electronic states in each layer are the chiral edge states that propagate in one direction only as shown in Fig 3(a). These edge states are coupled by tunneling between layers. Thus the surface of a multilayer is covered by a chiral sheath of coupled edge states. These surface states control the electrical transport properties of the multilayer. A central question from a quantum transport point of view is whether these surface states are localized or extended in the direction of the field . Fig 3(b) shows a network model of the multilayer surface introduced by Saul, Kardar and Read and studied by many authors subsequently. In this model it is assumed that tunneling between edges takes place only at discrete nodes (dashed vertical lines in Fig 3b) that appear at regular intervals along an edge. The edges are separated by nodes into horizontal segments called links. The wavefunction has a definite value on each link. Each node is visited by two incoming links and by two outgoing links. Each node is characterised by a $`2\times 2`$ S-matrix that relates the wavefunction on the outgoing links to the incoming amplitudes. Once the S-matrices are specified, given the wavefunction through a vertical slice, we can propagate it to the right. The S-matrices are chosen at random from some ensemble to incorporate the effect of disorder. To fully specify the model it is necessary to choose a distribution for the S-matrices. Periodic boundary conditions are imposed in the horizontal direction . The directed network model above is quantum mechanical but in the limit of infinite circumference and for a special choice of disorder, Saul, Kardar and Read have shown that it reduces to a classical model, the $`q`$-model with uniform distribution of fractions and zero injection . In section VI we discuss some respects in which more generic models of the multilayer surface, that do not reduce to classical models, still do show behaviour similar to the $`q`$-model . At the same time we show that in case of finite circumference quantum interference effects become important and there is little to be learnt from the study of the classical $`q`$-model. Instead a mapping to a ferromagnetic supersymmetric spin-chain has proved fruitful in this case . In section VI we discuss aspects of this mapping. A detailed summary of our results is given in section VII. The reader interested in first obtaining an overview of the paper or interested only in the results should proceed directly to section VII. ## II Critical Behaviour in 1+1 Dimensions Coppersmith et al. analysed the distribution of load $`\mathrm{\Pi }(w,t\mathrm{})`$ at very large depth where presumably a steady state is achieved . Here we study how the distribution evolves as a function of depth to this asymptotic steady state. We assume that a uniform load is applied to the top layer, $$w_n(t)=1\mathrm{for}\mathrm{all}n.$$ (3) In this section we neglect the weight of each bead (the injection term). In partial support of this neglect we note that in the experiment of ref typically a total load of 7600 N was applied to the bead pack. In comparison we estimate that the weight of a single bead was less than a mN; of the entire pack, less than 100 N. However, right at the critical point injection is a relevant perturbation, and at sufficiently large depth must be taken into account even if the weight of a single bead is small. We return to the effects of injection in section IV. To make the problem tractable we study not the entire distribution $`\mathrm{\Pi }(w,t)`$ but only its lowest non-trivial moment. With the neglect of injection it follows that the total load on every layer is the same; the $`q`$-model dynamics (eq 1) just shuffles this load. Hence the average load in layer $`t`$ $`w(t)`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}𝑑ww\mathrm{\Pi }(w,t)`$ (4) $`=`$ $`1.`$ (5) The lowest non-trivial moment is therefore the variance $$\delta w^2(t)=_0^{\mathrm{}}𝑑ww^2\mathrm{\Pi }(w,t)1.$$ (6) Since a uniform load is applied to the top layer the variance in that layer vanishes. As the load propagates downward, the fluctuations must grow and saturate. Our purpose is to analyse this evolution for different distributions $`P(f)`$, particularly those that are close to the singular distribution. Right at the critical point the asymptotic distribution $`\mathrm{\Pi }(w,\mathrm{})`$ is believed to be a power law. If we assume that it does not have a well defined variance, then by analogy to critical phenomena we surmise that close to the critical point the variance must diverge as $$\delta w^2(t\mathrm{})\frac{1}{\delta ^\theta }.$$ (7) Here $`\delta `$ measures the distance of a distribution $`P(f)`$ from the singular distribution; $`\delta `$ will be defined precisely below. We also expect that the depth-scale $`\xi _{\mathrm{corr}}`$ at which the steady state is attained will diverge as the critical point is approached. Thus $$\xi _{\mathrm{corr}}\frac{1}{\delta ^\phi }.$$ (8) $`\xi _{\mathrm{corr}}`$ is a vertical correlation length that diverges as the critical point is approached. Combining eqs (5) and (6) we expect that close to the critical point the fluctuations must have a scaling form $$\delta w^2(t)=\frac{1}{\delta ^\theta }(t\delta ^\phi ).$$ (9) To be consistent with eq (5) we expect that the scaling function $`(u)`$ const as $`u\mathrm{}`$. For short times we expect that the system should behave as it would at the critical point. The $`\delta `$ dependence should cancel and so we expect $`(u)u^{\theta /\phi }`$ for $`u1`$ so that $`\delta w^2(t)t^{\theta /\phi }`$ at the critical point. In the remainder of this section we will confirm that eq (7) and these inferences are valid. We will determine the exponents $`\theta `$ and $`\phi `$ and the scaling function $`(u)`$. As an aside to experts we note that it may have been more natural to name the exponents $`\theta (3\tau )/\sigma `$ and $`\phi \nu z`$. These names follow from a more general scaling hypothesis for the entire distribution (eq 174). However in this section we have elected to make the more restricted hypothesis eq (7) and to give the exponents single letter names taking care to avoid common exponent names such as $`\alpha ,\beta `$ and $`\nu `$. ### A Disorder Average Consider the correlation function $$c_m(t)=\frac{1}{N}\underset{n}{}w_n(t)w_{n+m}(t).$$ (10) We assume there are $`N`$ beads in each layer and we impose periodic boundary conditions in the horizontal direction. Ultimately we are interested in taking $`N\mathrm{}`$. Note that $`m`$ is even for both $`t=`$ even and $`t=`$ odd. In terms of the correlation function the variance is given by $$\delta w^2(t)=c_0(t)1.$$ (11) The correlation function obeys a remarkably simple evolution equation. This equation can be solved by straightforward classical analysis to yield the evolution of the variance. It is not difficult to obtain the entire correlation function by this method, and thereby obtain information on the horizontal correlation length, but we do not attempt this here. To analyse the evolution of the correlation function we write $`c_m(t+1)`$ $`=`$ $`{\displaystyle \frac{1}{N}}{\displaystyle \underset{n}{}}w_n(t+1)w_{n+m}(t+1)`$ (12) $`=`$ $`{\displaystyle \frac{1}{N}}{\displaystyle \underset{n}{}}w_{n+1}(t)w_{n+m+1}(t)f_{n+1}(t)f_{n+m+1}(t)`$ (14) $`+\mathrm{others}.`$ $`=`$ $`{\displaystyle \frac{1}{N}}{\displaystyle \underset{n}{}}w_{n+1}(t)w_{n+m+1}(t)f_{n+1}(t)f_{n+m+1}(t)`$ (16) $`+\mathrm{others}.`$ To obtain the second line of eq (10) we have used eq (1). Four terms result; we have written only one for illustration. To obtain the third line it is crucial to observe that $`w_n(t)`$ depends only on fractions in the layers above. It is not correlated with the fractions in layer $`t`$, allowing us to factorise the average as shown. To perform the average we need information about the distribution $`P(f)`$. By symmetry for any choice of distribution $$f=_0^1𝑑ffP(f)=\frac{1}{2}.$$ (17) For the variance we write $$\left(f\frac{1}{2}\right)^2=\frac{ϵ}{4}.$$ (18) $`ϵ`$ is a parameter that characterises the distribution $`P(f)`$. For example, $`ϵ=1/3`$ for the uniform distribution. For the singular distribution the parameter takes its maximum possible value $`ϵ=1`$. Since the fractions for different beads are assumed to be independently distributed we conclude $$f_n(t_1)f_m(t_2)=\frac{1}{4}+\frac{ϵ}{4}\delta _{n,m}\delta _{t_1,t_2}.$$ (19) Substituting eq (13) in eq (10) we obtain $`c_m(t+1)`$ $`=`$ $`\left({\displaystyle \frac{1}{4}}+{\displaystyle \frac{ϵ}{4}}\delta _{m,0}\right)c_m(t)`$ (21) $`+\mathrm{others}`$ $`=`$ $`\left({\displaystyle \frac{1}{2}}+{\displaystyle \frac{ϵ}{2}}\delta _{m,0}\right)c_m(t)`$ (24) $`+\left({\displaystyle \frac{1}{4}}{\displaystyle \frac{ϵ}{4}}\delta _{m,2}\right)c_{m2}(t)`$ $`+\left({\displaystyle \frac{1}{4}}{\displaystyle \frac{ϵ}{4}}\delta _{m,2}\right)c_{m+2}(t).`$ In the second line of eq (14) the other terms have been unveiled. Recall that $`m`$ takes even integer values. It is convenient to replace $`mm/2`$ to obtain $`c_m(t+1)`$ $`=`$ $`\left({\displaystyle \frac{1}{2}}+{\displaystyle \frac{ϵ}{2}}\delta _{m,0}\right)c_m(t)`$ (27) $`+\left({\displaystyle \frac{1}{4}}{\displaystyle \frac{ϵ}{4}}\delta _{m,1}\right)c_{m1}(t)`$ $`+\left({\displaystyle \frac{1}{4}}{\displaystyle \frac{ϵ}{4}}\delta _{m,1}\right)c_{m+1}(t).`$ Eq (15) is the main result of this subsection. It governs the evolution of the correlation function. We wish to solve it subject to the initial condition $$c_m(t0)=1\mathrm{for}\mathrm{all}m.$$ (28) The initial condition follows from the definition of $`c_m`$ (eq 8) and the assumed uniform load on the top layer. Note that the distribution $`P(f)`$ enters the evolution equation only through the parameter $`ϵ`$. Since the parameter takes its maximum value $`ϵ=1`$ for the singular distribution we may define $$\delta =1ϵ$$ (29) as the distance of a distribution $`P(f)`$ from the critical point. ### B Scattering solution It is easy to verify that a steady state solution to eq (15) is $`c_0`$ $`=`$ $`{\displaystyle \frac{1}{1ϵ}};`$ (30) $`c_n`$ $`=`$ $`1\mathrm{for}n0.`$ (31) Assuming this is the unique steady state towards which our initial condition evolves, eq (18) reveals that the variance does diverge as the singular distribution is approached. Using eq (9) $`\delta ^2w(t\mathrm{})`$ $`=`$ $`{\displaystyle \frac{ϵ}{1ϵ}}`$ (32) $``$ $`{\displaystyle \frac{1}{\delta }}\mathrm{as}ϵ1.`$ (33) Comparing eq (5) we see that the exponent $`\theta =1`$. Eq (18) also reveals that in steady state the fluctuations in load are uncorrelated for all pairs of beads including neighbours. This is in agreement with experiment . A full solution of evolution dynamics needs more work. Schematically eq (15) states $$c(t+1)=Hc(t).$$ (34) The strategy we adopt here is to seek the eigenvectors of $`H`$, $$H\varphi ^\lambda =\lambda \varphi ^\lambda ,$$ (35) and to expand the initial correlation vector $`c(0)`$ in terms of the eigenvectors, $$c(0)=\underset{\lambda }{}a_\lambda \varphi ^\lambda .$$ (36) The correlation vector at depth $`t`$ is then $$c(t)=\underset{\lambda }{}\lambda ^ta_\lambda \varphi ^\lambda .$$ (37) A complication we must negotiate is that $`H`$ is non-Hermitian. According to the standard theory of biorthogonal expansion (briefly recounted in Appendix A) to execute the plan above we must prove that the eigenvectors of $`H`$ span the vector space. Then we must find the eigenvectors of $`H^{}`$, called the left eigenvectors of $`H`$ in this context. The eigenvalues of $`H^{}`$ are the complex conjugate of the eigenvalues of $`H`$. Thus $$H^{}\psi ^\lambda =\lambda ^{}\psi ^\lambda .$$ (38) $`\psi ^\lambda `$ denotes the left eigenvector with eigenvalue $`\lambda ^{}`$. Having completed these tasks we may write the completeness relation $$\underset{\lambda }{}(\psi _m^\lambda )^{}\varphi _n^\lambda =\delta _{mn}.$$ (39) Using eq (25) we conclude that the expansion coefficients in eq (22) are determined by the left eigenvectors: $$a_\lambda =\underset{m}{}(\psi _m^\lambda )^{}c_m(0).$$ (40) Implementing the plan we first write the eigenvalue equation for $`H`$ $`{\displaystyle \frac{1}{2}}\varphi _r^\lambda +{\displaystyle \frac{1}{4}}\varphi _{r+1}^\lambda +{\displaystyle \frac{1}{4}}\varphi _{r1}^\lambda `$ $`=`$ $`\lambda \varphi _r^\lambda \mathrm{for}\|r\|2;`$ (41) $`{\displaystyle \frac{1}{2}}\varphi _1^\lambda +{\displaystyle \frac{1}{4}}\varphi _2^\lambda +{\displaystyle \frac{1ϵ}{4}}\varphi _0^\lambda `$ $`=`$ $`\lambda \varphi _1^\lambda ;`$ (42) $`{\displaystyle \frac{1+ϵ}{2}}\varphi _0^\lambda +{\displaystyle \frac{1}{4}}\varphi _1^\lambda +{\displaystyle \frac{1}{4}}\varphi _1^\lambda `$ $`=`$ $`\lambda \varphi _0^\lambda ;`$ (43) $`{\displaystyle \frac{1}{2}}\varphi _1^\lambda +{\displaystyle \frac{1ϵ}{4}}\varphi _0^\lambda +{\displaystyle \frac{1}{4}}\varphi _2^\lambda `$ $`=`$ $`\lambda \varphi _1^\lambda .`$ (44) Note that for $`ϵ=0`$ eq (27) may be interpreted as the Schrödinger equation for a free particle on a tightbinding lattice, familiar from elementary solid state physics. For non-zero $`ϵ`$ the particle may be viewed as scattering off a (non-Hermitian) barrier at the origin. Thus we seek a solution of the scattering form $`\varphi _n^{(+)k}`$ $`=`$ $`T(k)e^{ikn}\mathrm{for}n1;`$ (45) $`=`$ $`A(k)\mathrm{for}n=0;`$ (46) $`=`$ $`e^{ikn}+R(k)e^{ikn}\mathrm{for}n1.`$ (47) Here $`0<k<\pi `$. The first line of eq (27) then yields the eigenvalue $$\lambda (k)=\frac{1}{2}+\frac{1}{2}\mathrm{cos}k.$$ (48) The next three lines yield the scattering coefficients $`A(k)`$ $`=`$ $`{\displaystyle \frac{i\mathrm{sin}k}{(1ϵ)e^{ik}+ϵ\mathrm{cos}k}};`$ (49) $`T(k)`$ $`=`$ $`(1ϵ)A(k);`$ (50) $`R(k)`$ $`=`$ $`(1ϵ)A(k)1.`$ (51) There are also scattering solutions to eq (27) corresponding to the fictitious particle coming in from the right $`\varphi _n^{()k}`$ $`=`$ $`e^{ikn}+R(k)e^{ikn}\mathrm{for}n1;`$ (52) $`=`$ $`A(k)\mathrm{for}n=0;`$ (53) $`=`$ $`T(k)e^{ikn}\mathrm{for}n1.`$ (54) By symmetry the scattering coefficients for this state are also given by eq (30). There are no bound state solutions to eq (27). The scattering solutions we have found all have real eigenvalues. In principle, since $`H`$ is non-Hermitian, complex eigenvalues are also possible. However it turns out there are no solutions with complex eigenvalue that are biorthonormalisable. It will be seen that the scattering solutions we have found constitute a complete set. The next step is to find the left eigenvectors that obey $`{\displaystyle \frac{1}{2}}\psi _r^\lambda +{\displaystyle \frac{1}{4}}\psi _{r+1}^\lambda +{\displaystyle \frac{1}{4}}\psi _{r1}^\lambda `$ $`=`$ $`\lambda \psi _r^\lambda \mathrm{for}\|r\|2;`$ (55) $`{\displaystyle \frac{1}{2}}\psi _1^\lambda +{\displaystyle \frac{1}{4}}\psi _2^\lambda +{\displaystyle \frac{1}{4}}\psi _0^\lambda `$ $`=`$ $`\lambda \psi _1^\lambda ;`$ (56) $`{\displaystyle \frac{1+ϵ}{2}}\psi _0^\lambda +{\displaystyle \frac{1ϵ}{4}}\psi _1^\lambda +{\displaystyle \frac{1ϵ}{4}}\psi _1^\lambda `$ $`=`$ $`\lambda \psi _0^\lambda ;`$ (57) $`{\displaystyle \frac{1}{2}}\psi _1^\lambda +{\displaystyle \frac{1}{4}}\psi _0^\lambda +{\displaystyle \frac{1}{4}}\psi _2^\lambda `$ $`=`$ $`\lambda \psi _1^\lambda .`$ (58) Eq (32) is the transpose of eq (27). The left eigenvectors are $`\psi _n^{(+)k}`$ $`=`$ $`𝒯(k)e^{ikn}\mathrm{for}n1;`$ (59) $`=`$ $`𝒜(k)\mathrm{for}n=0;`$ (60) $`=`$ $`e^{ikn}+(k)e^{ikn}\mathrm{for}n1.`$ (61) and $`\psi _n^{()k}`$ $`=`$ $`e^{ikn}+(k)e^{ikn}\mathrm{for}n1;`$ (62) $`=`$ $`𝒜(k)\mathrm{for}n=0;`$ (63) $`=`$ $`𝒯(k)e^{ikn}n1.`$ (64) The scattering coefficients are given by $`𝒜(k)`$ $`=`$ $`{\displaystyle \frac{i(1ϵ)\mathrm{sin}k}{(ϵ\mathrm{cos}k)+(1ϵ)e^{ik}}};`$ (65) $`𝒯(k)`$ $`=`$ $`𝒜(k);`$ (66) $`(k)`$ $`=`$ $`𝒜(k)1.`$ (67) Having found the left and right eigenvectors, by analogy with eq (25), we now posit the completeness relation $$_0^\pi \frac{dk}{2\pi }\left(\psi _m^{(+)k}\varphi _n^{(+)k}+\psi _m^{()k}\varphi _n^{()k}\right)=\delta _{mn}.$$ (68) The proof of this completeness relation, an important element of the analysis, is carried out in Appendix A. The expansion of the initial correlation vector indicated schematically in eq (22) may now be written $$c_m(0)=_0^\pi \frac{dk}{2\pi }\left[a^{(+)}(k)\varphi _m^{(+)k}+a^{()}(k)\varphi _m^{()k}\right].$$ (69) The correlation vector at depth $`t`$ is now $$c_m(t)=_0^\pi \frac{dk}{2\pi }\lambda (k)^t\left[a^{(+)}(k)\varphi _m^{(+)k}+a^{()}(k)\varphi _m^{()k}\right]$$ (70) as previously shown schematically in eq (23). The expansion coefficients $`a(k)`$, obtained using the completeness relation (eq 36), are $`a^{(+)}(k)={\displaystyle \underset{n=\mathrm{}}{\overset{+\mathrm{}}{}}}c_n(0)\psi _n^{(+)k};`$ (71) $`a^{()}(k)={\displaystyle \underset{n=\mathrm{}}{\overset{+\mathrm{}}{}}}c_n(0)\psi _n^{()k};`$ (72) as previously indicated schematically in eq (26). To ensure convergence of the sums in eq (39) we set $`c_m(0)e^{\eta \|m\|}`$ and take $`\eta 0`$ at the end. Using eqs (33), (34) and (35) we perform the sums exactly to obtain $`a^{(+)}(k)`$ $`=`$ $`a^{()}(k)`$ (73) $`=`$ $`𝒜(k)^{}+2𝒜(k)^{}{\displaystyle \frac{e^{ik\eta }}{1e^{ik\eta }}}+\left({\displaystyle \frac{e^{ik\eta }}{1e^{ik\eta }}}\mathrm{cc}\right)`$ (74) $`=`$ $`2\pi 𝒜(k)^{}\delta (k)+[1𝒜(k)^{}]i\mathrm{cot}{\displaystyle \frac{k}{2}}.`$ (75) The last line of eq (40) is obtained by taking the limit $`\eta 0`$. Substituting eq (40) in eq (38) and making use of eqs (28), (29), (30), (31) and (35) we finally obtain $$c_0(t)=\frac{1}{1ϵ}\frac{ϵ}{\pi }_0^\pi 𝑑k\frac{\mathrm{cos}^{2(t+1)}(k/2)}{ϵ^2(2ϵ1)\mathrm{cos}^2k}.$$ (76) Eq (41) is the exact expression for the evolution of $`c_0(t)`$ that we sought. Finally we would like to re-express eq (41) in terms of standard special functions. Some of the manipulations will prove useful later in the analysis of injection. Introduce the $`z`$-transform $`c_0(z)`$ $`=`$ $`{\displaystyle \underset{t=0}{\overset{\mathrm{}}{}}}z^tc_0(t)={\displaystyle \frac{1}{1ϵ}}{\displaystyle \frac{1}{1z}}`$ (78) $`{\displaystyle \frac{ϵ}{\pi }}{\displaystyle _0^\pi }𝑑k{\displaystyle \frac{\mathrm{cos}^2(k/2)}{ϵ^2(2ϵ1)\mathrm{cos}^2(k/2)}}{\displaystyle \frac{1}{1z\mathrm{cos}^2(k/2)}}.`$ The $`k`$-integral may now be performedExtend the range from 0 to $`2\pi `$. Then use the standard trick for turning an angular integral into a contour integral around the unit circle in the $`\zeta `$ plane via the substitution $`\zeta e^{ik}`$. See for example , p 409. to yield $`c_0(z)`$ $`=`$ $`{\displaystyle \frac{1}{(1ϵ)(1z)}}{\displaystyle \frac{ϵ^2}{(1ϵ)(2ϵ1)(1\alpha (ϵ)z)}}`$ (80) $`{\displaystyle \frac{ϵ}{(12ϵ)\sqrt{1z}(1\alpha (ϵ)z)}}.`$ For brevity $`\alpha (ϵ)=ϵ^2/(2ϵ1)`$. Upon inversion of the $`z`$-transform (details relegated to Appendix B) we obtain $$c_0(t)=\frac{1}{1ϵ}\frac{1}{\pi ϵ}_1^{\mathrm{}}𝑑x(x1)^{1/2}\frac{1}{x^{t+1}}\left(x\frac{2ϵ1}{ϵ^2}\right)^1.$$ (81) Comparing an integral representation for the hypergeometric function $$F(a,b,c;s)=\frac{\mathrm{\Gamma }(c)}{\mathrm{\Gamma }(cb)\mathrm{\Gamma }(b)}_1^{\mathrm{}}𝑑x(x1)^{cb1}x^{ac}(xs)^a,$$ (82) valid for Re $`c>`$ Re $`b>0`$ and $`\|s\|<1`$, we conclude $`\delta w^2(t)`$ $`=`$ $`{\displaystyle \frac{ϵ}{1ϵ}}{\displaystyle \frac{1}{\pi ϵ}}{\displaystyle \frac{\mathrm{\Gamma }(1/2)\mathrm{\Gamma }(t+3/2)}{\mathrm{\Gamma }(t+2)}}`$ (84) $`\times F(1,t+3/2,t+2;{\displaystyle \frac{2ϵ1}{ϵ^2}})`$ Eq (46) is the final result of this section. It is an exact formula for the evolution of fluctuations with depth, in terms of known special functions. As a practical matter eqs (41) and (44) are equivalent to eq (46) and will prove more useful. ### C Scaling Limit Eq (46) gives the exact evolution of load fluctuations for the $`q`$-model without injection. It is valid for all $`t`$ and all distributions, $`P(f)`$. From our point of view however it is more interesting to examine the scaling limit of large depth behaviour near the critical point. To derive the scaling limit we start with eq (44)—eq (41) would have served just as well—and consider the limit $`t1`$ and $`\delta =1ϵ0`$. We do not make any assumption about the relationship between $`t`$ and $`1/\delta `$. We obtain $`\delta ^2w(t)`$ $``$ $`{\displaystyle \frac{1}{\delta }}{\displaystyle \frac{1}{\pi }}{\displaystyle _0^{\mathrm{}}}𝑑ss^{1/2}e^{t\mathrm{ln}(1+s)}(s+\delta ^2)^1`$ (85) $``$ $`{\displaystyle \frac{1}{\delta }}{\displaystyle \frac{1}{\pi }}{\displaystyle _0^{\mathrm{}}}𝑑ss^{1/2}e^{ts}(s+\delta ^2)^1.`$ (86) In the first line of eq (47) we have changed the integration variable from $`x`$ to $`s=x1`$. Again changing the integration variable from $`s`$ to $`\overline{s}=s/\delta ^2`$ we obtain $$\delta ^2w(t)=\frac{1}{\delta }\left[1\frac{2}{\pi }_0^{\mathrm{}}𝑑\overline{s}\frac{e^{\overline{s}^2t\delta ^2}}{1+\overline{s}^2}\right].$$ (87) Comparing eq (7) we conclude that close to the critical point and in the large depth limit $`\delta ^2w(t)`$ does indeed have a scaling form with exponents $$\theta =1,\phi =2$$ (88) and scaling function $$(u)=1\frac{2}{\pi }_0^{\mathrm{}}𝑑s\frac{e^{us^2}}{1+s^2}.$$ (89) Fig 4 shows a plot of $`(u)`$. As anticipated the asymptotic behaviour of the scaling function is $`(u)`$ $``$ $`1{\displaystyle \frac{1}{\sqrt{\pi u}}}\mathrm{for}u\mathrm{}`$ (90) $``$ $`{\displaystyle \frac{2}{\sqrt{\pi }}}\sqrt{u}\mathrm{for}u0.`$ (91) We conclude that the saturation depth scale $`\xi _{\mathrm{corr}}1/\delta ^2`$. For very great depths $`t\xi _{\mathrm{corr}}`$, the fluctuations saturate to the value $`1/\delta `$ as found earlier by analysis of the steady state (eq 19). For small depths, $`1t\xi _{\mathrm{corr}}`$ they grow as $$\delta ^2w(t)=\frac{2}{\sqrt{\pi }}\sqrt{t}.$$ (92) This behaviour must persist at all depths right at the critical point as will be explicitly confirmed in section III. In summary, we have shown that the singular distribution is an isolated critical point in the space of $`q`$-models. There is a (vertical) correlation length that diverges as the critical point is approached. We have determined the exponents $`\theta `$ and $`\phi `$ and the scaling function $`(x)`$ introduced in eqs (5), (6) and (7). In context of river networks we have found that any $`q`$-model with stream splitting (hence non-zero $`\delta `$) has a (possibly very long) correlation length in the direction of flow. Such a network is therefore not scale invariant on sufficiently long scalesStrictly, to analyse a river network the appropriate initial condition is to load a fraction of randomly chosen sites in the $`t=0`$ layer, rather than the uniform load analysed here. However we do not expect our conclusion regarding correlation lengths is sensitive to initial conditions.. ## III Critical Point Distribution Right at the critical point in $`1+1`$ dimensions it is possible to analyse the dynamics of the entire distribution $`\mathrm{\Pi }(w,t)`$. Since there is no vertical length scale at the critical point we expect that in the large depth, scaling limit $$\mathrm{\Pi }(w,t)=\frac{1}{t^\omega }(wt^\mathrm{{\rm Y}}).$$ (93) Eq (53) implies that at the critical point the variance should grow as $`t^{3\mathrm{{\rm Y}}\omega }`$ in the scaling limit $`t1`$. From eq (7) we had surmised that the variance would grow as $`t^{\theta /\phi }`$ for $`\delta =0`$. Hence the exponents $`\theta ,\phi `$ of the preceding section and $`\omega ,\mathrm{{\rm Y}}`$ of this section are not independent; they satisfy the relation $`3\mathrm{{\rm Y}}+\omega +\theta /\phi =0`$. Below we calculate the exponents $`\omega `$ and $`\mathrm{{\rm Y}}`$, explicitly verify the exponent relationship and obtain the scaling function $`(s)`$. Again as an aside to experts we note that the exponents $`\omega `$ and $`\mathrm{{\rm Y}}`$ might more naturally have been written $`\omega \tau /\nu z\sigma `$, $`\mathrm{{\rm Y}}1/\nu z\sigma `$. These expressions follow from the $`\delta 0`$ limit of the more general scaling hypothesis for the entire distribution close to the critical point (eq 174). However for this section we have elected to make the more restricted hypothesis, eq (53), and to give the exponents single character names. In this section too we neglect injection. At the critical point injection is a relevant perturbation. The form we derive is therefore a transient that will break down at sufficient depth. Provided the injection is weak however that depth could be very great. Majumdar and Sire have analysed the scaling limit of $`\mathrm{\Pi }(w,t)`$ when injection is present; however it does not appear straightforward to take the injection $`0`$ limit in their expression. It would also be desirable for the case of non-zero injection to have a simple explicit formula for the crossover of $`\mathrm{\Pi }(w,t)`$ from the transient we derive (eq 53) to the injection dominated, large depth limit. Presumably this can be accomplished by extracting the suitable limit of the results of ref , or by direct calculation, but we do not attempt it here. ### A Disorder Average As in section II we assume a uniform load is applied to the top layer (eq 2). To obtain the distribution $`\mathrm{\Pi }(w,t)`$ following ref we consider the quantities $$Z_r(\rho ,t)=\mathrm{exp}i\rho \underset{n=1}{\overset{r}{}}w_n(t)$$ (94) where $`r=1,2,3,\mathrm{}`$ By translational invariance $`Z_1(\rho ,t)`$ $`=`$ $`\mathrm{exp}i\rho w_1(t)`$ (95) $`=`$ $`{\displaystyle \underset{w=0}{\overset{\mathrm{}}{}}}e^{i\rho w}\mathrm{\Pi }(w,t).`$ (96) Note that for the critical $`q`$-model without injection the load on a site is an integer. Thus $`Z_1(\rho ,t)`$ is the discrete Fourier or $`z`$-transform of the distribution $`\mathrm{\Pi }(w,t)`$; $`\rho `$ is the transform domain variable conjugate to $`w`$. $`Z_2(\rho ,t)`$ similarly encodes the joint probability distribution of load on neighbouring sites and so on. For the business at hand the imaginary part of $`Z_r(\rho ,t)`$, $$𝒵_r(\rho ,t)=\mathrm{Im}Z_r(\rho ,t),$$ (97) is especially valuable. It is evident from eq (55) that $$𝒵_1(\rho ,t)=\underset{w=0}{\overset{\mathrm{}}{}}\mathrm{sin}(\rho w)\mathrm{\Pi }(w,t).$$ (98) By using Fourier’s identity $$\frac{2}{\pi }_0^\pi 𝑑k\mathrm{sin}kn\mathrm{sin}km=\delta _{mn}$$ (99) and eq (57) we can extract the distribution $`\mathrm{\Pi }(w,t)`$ from $`𝒵_1(\rho ,t)`$ via $$\mathrm{\Pi }(w,t)=\frac{2}{\pi }_0^\pi 𝑑\rho \mathrm{sin}(\rho w)𝒵_1(\rho ,t)$$ (100) for $`w=1,2,3,\mathrm{}`$ We cannot obtain $`\mathrm{\Pi }(w=0,t)`$ in this way from $`𝒵_1(\rho ,t)`$, but we can obtain it from the normalisation of $`\mathrm{\Pi }(w,t)`$; $$\mathrm{\Pi }(w0,t)=1\underset{w=1}{\overset{\mathrm{}}{}}\mathrm{\Pi }(w,t).$$ (101) The benefit of considering the quantities $`Z_r(\rho ,t)`$ is that they obey a simple linear evolution equation. Following ref write $`Z_r(\rho ,t+1)`$ $`=`$ $`\mathrm{exp}i\rho {\displaystyle \underset{n=1}{\overset{r}{}}}w_n(t+1)`$ (102) $`=`$ $`\mathrm{exp}i\rho \{w_1(t)f_1(t)+{\displaystyle \underset{n=2}{\overset{r}{}}}w_n(t)`$ (104) $`+w_{r+1}(t)[1f_{r+1}(t)]\}.`$ To obtain the second line we have used the $`q`$-model evolution eq (1). Since the weights in layer $`t`$ depend only on fractions in the preceding layers we can perform the average over $`f_1(t)`$ and $`f_{r+1}(t)`$ separately in eq (61): $`\mathrm{exp}i\rho w_1(t)f_1(t)_{f_1}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left[1+\mathrm{exp}i\rho w_1(t)\right];`$ (105) $`\mathrm{exp}i\rho w_{r+1}(t)[1f_{r+1}]_{f_{r+1}}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left[1+\mathrm{exp}i\rho w_{r+1}(t)\right].`$ (106) Substituting eq (62) in eq (61) we obtain the evolution equation $$Z_r(\rho ,t+1)=\frac{1}{4}Z_{r1}(\rho ,t)+\frac{1}{2}Z_r(\rho ,t)+\frac{1}{4}Z_{r+1}(\rho ,t)$$ (107) where we have again made use of horizontal translational invariance. Note that eq (63) is linear. Hence it is obeyed separately by the real and imaginary parts of $`Z`$. $`𝒵`$ therefore evolves according to $$𝒵_r(\rho ,t+1)=\frac{1}{4}𝒵_{r1}(\rho ,t)+\frac{1}{2}𝒵_r(\rho ,t)+\frac{1}{4}𝒵_{r+1}(\rho ,t)$$ (108) Eq (64) is reminiscent of a tight-binding lattice Schrödinger equation for a free particle on a half-line (since the site index $`r1`$). The main results of this subsection are eqs (57) and (59) that define the relationship between $`\mathrm{\Pi }(w,t)`$ and $`𝒵_1(\rho ,t)`$ and eq (64) that controls the evolution of $`𝒵_r(\rho ,t)`$ with depth. ### B Solution and Scaling Limit We wish to solve eq (64) subject to the initial condition $$𝒵_r(\rho ,t0)=\mathrm{sin}\rho r.$$ (109) This follows from the assumed uniform load applied to the top layer and eqs (54) and (56). Schematically, eq (64) has the form $$𝒵_r(\rho ,t+1)=\underset{s}{}H_{rs}𝒵_s(\rho ,t).$$ (110) It is easy to verify that our initial condition is an eigenfunction of $`H`$; $$\underset{s}{}H_{rs}\mathrm{sin}\rho s=\left(\frac{1}{2}+\frac{1}{2}\mathrm{cos}\rho \right)\mathrm{sin}\rho r.$$ (111) Hence eq (64) has the remarkably simple solution $$𝒵_r(\rho ,t)=\left(\frac{1}{2}+\frac{1}{2}\mathrm{cos}\rho \right)^t\mathrm{sin}\rho r.$$ (112) Substituting eq (68) in eq (59) we obtain the desired expression for $$\mathrm{\Pi }(w,t)=\frac{2}{\pi }_0^\pi 𝑑\rho \mathrm{sin}(\rho w)\left(\frac{1}{2}+\frac{1}{2}\mathrm{cos}\rho \right)^t\mathrm{sin}\rho $$ (113) for $`w=1,2,3,\mathrm{}`$ The integral over $`\rho `$ can be performed exactly by a standard contour integration trick (see footnote 2) to yield $`\mathrm{\Pi }(w,t)`$ $`=`$ $`{\displaystyle \frac{1}{4^t}}{\displaystyle \frac{(2t)!}{(t+1w)!(t1+w)!}}`$ (116) $`{\displaystyle \frac{1}{4^t}}{\displaystyle \frac{(2t)!}{(t1w)!(t+1+w)!}}`$ $`\mathrm{for}w=1,2,\mathrm{},t1`$ $`=`$ $`{\displaystyle \frac{1}{4^t}}{\displaystyle \frac{(2t)!}{(t+1w)!(t1+w)!}}\mathrm{for}w=t,t+1`$ (117) $`=`$ $`0\mathrm{for}w>t+1.`$ (118) We now use eq (60) and (70) to obtain $`\mathrm{\Pi }(w0,t)`$. The sum proves tractable and yields $$\mathrm{\Pi }(w0,t)=1\frac{1}{4^t}\frac{(2t+1)!}{(t+1)!t!}.$$ (119) Eq (70) and (71) are the exact expressions for $`\mathrm{\Pi }(w,t)`$ for the critical $`q`$-model without injection. Much more interesting than the exact formula is the scaling limit of large depth. We now assume $`t1`$ but we will make no assumptions about the relative size of $`w`$ and $`t`$. To derive this limit we return to eq (69) and write $$\left(\frac{1}{2}+\frac{1}{2}\mathrm{cos}\rho \right)^te^{t\rho ^2/4}$$ (120) justified (inside the integral) for large $`t`$. Hence we obtain a Gaussian integral $`\mathrm{\Pi }(w,t)`$ $`=`$ $`{\displaystyle \frac{1}{\pi }}{\displaystyle _\pi ^\pi }𝑑\rho \mathrm{sin}(\rho w)\rho e^{t\rho ^2/4}`$ (121) $`=`$ $`{\displaystyle \frac{4}{\sqrt{\pi }}}{\displaystyle \frac{w}{t^{3/2}}}e^{w^2/t}.`$ (122) Comparing eq (53) and (73) we see that at large depth $`\mathrm{\Pi }`$ has the anticipated scaling form with exponents $$\omega =1,\mathrm{{\rm Y}}=\frac{1}{2}$$ (123) and scaling function $$(s)=\frac{4}{\sqrt{\pi }}se^{s^2}.$$ (124) Eq (73) holds for $`w1`$. In the same large depth limit $$\mathrm{\Pi }(w0,t)1\frac{2}{\sqrt{\pi }}\frac{1}{\sqrt{t}}.$$ (125) The distribution of load thus consists of a spike at zero load followed by smooth behaviour for non-zero load given by eq (73). At great depths it is extremely probable that the load on a given bead is zero; most of the weight of the distribution is in the spike. From the distribution of load, eq (73), it is easy to confirm that its variance (eq 4) grows without bound as the square root of depth, as we had earlier inferred from the scaling function $``$ (cf eq 52). It is instructive that the exact formula, eqs (70) and (71), is so cumbersome; the scaling limit, eqs (53), (74) and (75), emerges only when we plumb the depths. ## IV Effect of Injection In this section we consider the $`q`$-model in 1+1 dimensions taking into account injection. We will assume that the weights of the beads are independent and identically distributed with mean $`I`$ and variance $`\delta I^2`$. To probe the behaviour of the model we will assume that a uniform load is applied to the top layer (eq 2). We will study how the mean square load $`w^2(t)`$ evolves with depth since the mean load has the trivial variation $$w(t)=1+It.$$ (126) Near the critical point we expect that the mean square load should have a scaling form $$w^2(t)=\frac{1}{\delta ^\theta }𝒞(t\delta ^\phi ,\delta I^2\delta ^\mu ,I\delta ^\kappa ).$$ (127) We can guess all the exponents and obtain some features of the scaling function from simple arguments. The load on a particular bead at depth $`t`$ is a random linear combination of the weights of the beads in the layer above plus a term, due to the applied load, that does not depend on the weights, $`I_n`$. Hence the scaling function has to be of the form $`w^2(t)`$ $`=`$ $`{\displaystyle \frac{1}{\delta ^\theta }}(t\delta ^\phi )+{\displaystyle \frac{\delta I^2}{\delta ^{\theta +\mu }}}(t\delta ^\phi )+{\displaystyle \frac{I}{\delta ^{\theta +\kappa }}}𝒦(t\delta ^\phi )`$ (129) $`+{\displaystyle \frac{I^2}{\delta ^{\theta +2\kappa }}}(t\delta ^\phi ).`$ In the limit of zero injection eq (79) should reduce to our result in section II. Thus $$\theta =1,\phi =2$$ (130) and $``$ has the same form (eq 50) as in section II justifying the recycling of these particular symbols. By rewriting the average weight at depth $`t`$ (eq 77) as $`1+t\delta ^2I/\delta ^2`$ we conjecture $$\kappa =2.$$ (131) To obtain $`\mu `$ we imagine that the system is very close to the critical point. Then for times that are not too long, effectively, it will behave as it would right at the critical point. At that point the weight of each bead zig-zags down lines that merge but do not split. If we add the squares of the loads on all the beads on layer $`t`$ we will obtain the sum, over all the beads above layer $`t`$, of their squared deviation from the average weight $`I`$ plus other terms. Hence $`_nw_n(t)^2=\delta I^2Nt+`$ other pieces that do not depend on $`\delta I^2`$. Here $`N`$ is the number of beads in a layer. By translational invariance we conclude $$w^2(t)\delta I^2t+\mathrm{others}.$$ (132) In eq (82) “others” represents contributions to $`w^2(t)`$ that do not depend on $`\delta I^2`$. Comparing eq (82) and (79) we see that for small values of its argument $$(u)u$$ (133) and the exponent $$\mu =1,$$ (134) needed to cancel the $`\delta `$ dependence at small depths With the exponents in hand we can analyse the behaviour of $`w^2(t)`$ at small depths (compared to $`1/\delta ^2`$). This behaviour would persist out to all depths right at the critical point. For the term independent of injection we have already obtained the exact result, eq (52). For the term that depends on $`\delta I^2`$ we have just worked out the behaviour in this limit, including the precise coefficient (eq 82). For the term that is proportional to $`I`$ we argue that for small $`u`$, $`𝒦(u)u^{3/2}`$ to cancel the $`\delta `$ dependence, leading to $$w^2(t)It^{3/2}+\mathrm{others}.$$ (135) Similarly the contribution of the term that is proportional to $`I^2`$ is $$w^2(t)I^2t^{5/2}+\mathrm{others}.$$ (136) The last result has a simple interpretation. We have seen in section II that without injection at the critical point the mean weight at depth $`t`$ is 1; the mean square weight $`\sqrt{t}`$. With injection the average weight at sufficient depth is $`It`$. If we assume that uniform injection does not change the distribution, only its scale, then since the mean is inflated by a factor $`It`$, the mean square should be inflated by a factor $`I^2t^2`$, leading to eq (86). The same interpretation can be used to derive the behaviour of the last term in eq (79) in the limit $`t1/\delta ^2`$, the opposite of the limit we have so far considered. In that limit, in the absence of injection, the fluctuations saturate at the value $`1/\delta `$. Hence we expect this term to behave as $$w^2(t)\frac{1}{\delta }I^2t^2.$$ (137) We can check some of these deductions by making contact with Majumdar and Sire, who have analysed the entire distribution of load at the critical point . Following these authors let us imagine that the injection term is very small, with the squared mean $`I^2`$, significantly smaller than the variance $`\delta I^2`$. According to our analysis ultimately the fluctuations at the critical point should grow as $`t^{5/2}`$, but the depth at which the $`t^{5/2}`$ term (eq 86) overtakes the term linear in $`t`$ (eq 82) could be very great; it diverges as $`1/I^{4/3}`$. Majumdar and Sire arrived at the same value $`4/3`$ for this crossover exponent. Moreover, since they argued that right at the critical point (the only case they considered) there is only one independent exponent, we have made contact with their entire analysis as regards exponents. In summary we anticipate that near the critical point the mean square load will follow the scaling form (eq 78). Using simple arguments we have conjectured values for all the exponents \[eqs (80), (81) and (84)\] and guessed some features of the scaling function. As a check we have made contact with the critical point analysis of Majumdar and Sire and recovered the known value of the crossover exponent, 4/3 . In the remainder of this section we will fully confirm the deductions we have made above. We will obtain an exact formula for the evolution of the mean square load; the exponents, $`\theta ,\phi ,\mu `$ and $`\kappa `$; and the scaling function, $`𝒞`$. ### A Disorder Average and Exact Solution As in section II our strategy is to analyse the evolution of the correlation function, $`c_n(t)`$; the mean-squared weight $`w^2(t)=c_0(t)`$. The analysis is given in outline since most of the needed technical elements have already been described in section II. Here we shall focus on the new complications introduced by consideration of injection. Following the method of section IIB we first obtain the evolution equation for the correlation function, now including injection. Schematically this equation has the form $$c_m(t+1)=\underset{n}{}H_{mn}c_n(t)+\xi _m(t).$$ (138) $`H_{mn}`$ is the same matrix as in eq (15). The effect of injection appears in the inhomogeneous term $`\xi _m`$. Explicitly $$\xi _m=2I+(2t+1)I^2+\delta I^2\delta _{m=0}.$$ (139) Our strategy to solve eq (88) is to first expand $`c(t)`$ and $`\xi (t)`$ in terms of the right eigenvectors of $`H`$: $$c_m(t)=\underset{\lambda }{}a_\lambda (t)\varphi _m^\lambda ;\xi _m(t)=\underset{\lambda }{}\xi _\lambda (t)\varphi _m^\lambda .$$ (140) As discussed before, the expansion amplitudes $`a_\lambda `$ and $`\xi _\lambda `$ are calculated by use of the left eigenvectors $$a_\lambda (t)=\underset{m}{}(\psi _m^\lambda )^{}c_m(t);\xi _\lambda (t)=\underset{m}{}(\psi _m^\lambda )^{}\xi _m(t).$$ (141) In section IIB we have calculated the amplitudes for $`c_m(t0)`$. We found $`a^{(+)}(k,t0)=a^{()}(k,t0)a(k,t0)`$ with $$a(k,t0)=2\pi 𝒜(k)^{}\delta (k)+[1𝒜(k)^{}]i\mathrm{cot}\frac{k}{2}.$$ (142) Here $`𝒜(k)`$ is given by eq (35). Similarly $`\xi ^{(+)}(k,t0)=\xi ^{()}(k,t0)\xi (k,t0)`$ with $`\xi (k,t)`$ $`=`$ $`\delta I^2𝒜(k)^{}+\{2I+(2t+1)I^2\}`$ (144) $`\{2\pi 𝒜(k)^{}\delta (k)+[1𝒜(k)^{}]i\mathrm{cot}{\displaystyle \frac{k}{2}}\}.`$ Substituting the expansions eq (91) into the evolution eq (88) shows that the dynamics of the amplitudes for different right eigenvectors is decoupled and is given by $$a_\lambda (t+1)=\lambda a_\lambda (t)+\xi _\lambda (t).$$ (145) To solve this dynamics we introduce the $`z`$-transforms $`a_\lambda (z)`$ $`=`$ $`{\displaystyle \underset{t=0}{\overset{\mathrm{}}{}}}a_\lambda (t)z^t,`$ (146) $`\xi _\lambda (z)`$ $`=`$ $`{\displaystyle \underset{t=0}{\overset{\mathrm{}}{}}}\xi _\lambda (t)z^t,`$ (147) to obtain $$a_\lambda (z)=\frac{a_\lambda (t0)}{1z\lambda }+\frac{z\xi _\lambda (z)}{1z\lambda }.$$ (148) Combining eq (90) and (96) we conclude $$c_m(z)=\underset{\lambda }{}\left[\frac{a_\lambda (t0)}{1z\lambda }+\frac{z\xi _\lambda (z)}{1z\lambda }\right]\varphi _m^\lambda .$$ (149) More explicitly $`c_0(z)`$ $`=`$ $`2{\displaystyle _0^\pi }𝑑k{\displaystyle \frac{a(k,t0)A(k)}{1z\lambda (k)}}`$ (151) $`+2z{\displaystyle _0^\pi }𝑑k{\displaystyle \frac{\xi (k,z)A(k)}{1z\lambda (k)}}.`$ Here $`c_0(z)`$ is the $`z`$-transform of $`c_0(t)`$; $`A(k)`$ is given by eq (30); $`\lambda (k)`$, by eq (29); and $`a(k,t0)`$, by eq (92). $`\xi (k,z)`$ is to be obtained by $`z`$-transforming eq (93). Now all the pieces have been assembled. It remains to perform the $`k`$ integral and invert the $`z`$-transform. The $`k`$-integrals may be performed by the standard contour integration method mentioned in footnote 2. The $`z`$-transforms can all be inverted as illustrated in Appendix B. After much calculation we find $$w^2(t)=\overline{F}(t,ϵ)+\delta I^2M(t,ϵ)+IK(t,ϵ)+I^2L(t,ϵ)$$ (152) with $`\overline{F}(t,ϵ)`$ $`=`$ $`{\displaystyle \frac{ϵ}{1ϵ}}{\displaystyle \frac{1}{\pi ϵ}}\mathrm{\Gamma }F_1;`$ (153) $`M(t,ϵ)`$ $`=`$ $`{\displaystyle \frac{ϵ}{(1ϵ)^2}}+{\displaystyle \frac{2}{\pi }}{\displaystyle \frac{(1ϵ)}{ϵ^2}}\mathrm{\Gamma }[tF_1+F_2];`$ (154) $`K(t,ϵ)`$ $`=`$ $`{\displaystyle \frac{2}{1ϵ}}t+{\displaystyle \frac{2ϵ^2}{(1ϵ)^3}}{\displaystyle \frac{4}{\pi ϵ}}\mathrm{\Gamma }[tF_1+F_2];`$ (155) $`L(t,ϵ)`$ $`=`$ $`{\displaystyle \frac{ϵ^4+2ϵ^3ϵ^2}{(1ϵ)^5}}+{\displaystyle \frac{2ϵ^2}{(1ϵ)^3}}t`$ (157) $`+{\displaystyle \frac{2}{3\pi ϵ}}\mathrm{\Gamma }[(4t^2t)F_1+(8t5)F_2+8F_3].`$ We have put an overline on $`\overline{F}(t,ϵ)`$ to avoid confusion with a hypergeometric function. For brevity we have written $`\mathrm{\Gamma }`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Gamma }(1/2)\mathrm{\Gamma }(t+3/2)}{\mathrm{\Gamma }(t+2)}};`$ (158) $`F_n`$ $`=`$ $`F(n,t+{\displaystyle \frac{3}{2}},t+2;{\displaystyle \frac{2ϵ1}{ϵ^2}})`$ (159) in eq (100). Eq (100) is the final result of this subsection. It gives the evolution of load fluctuations for the $`q`$-model with injection in 1+1 dimensions. It holds for any distribution of fractions $`P(f)`$ and at any depth. ### B Scaling Limit More interesting than the exact results is the scaling behaviour that emerges for $`t1`$ and $`\delta =1ϵ0`$. To derive this behaviour it is useful to express the hypergeometric functions in eq (100) via the integral representation, eq (45). The asymptotic behaviour of $`\mathrm{\Gamma }F_1`$ has been analysed in section IIC \[cf. eqs (47) and (48)\]. The corresponding analysis of $`\mathrm{\Gamma }F_2`$ and $`\mathrm{\Gamma }F_3`$ is very similar and finally leads to $`\overline{F}(t,ϵ)`$ $``$ $`{\displaystyle \frac{1}{\delta }}\left\{1{\displaystyle \frac{2}{\pi }}\mathrm{\Phi }_1(u)\right\};`$ (160) $`M(t,ϵ)`$ $``$ $`{\displaystyle \frac{1}{\delta ^2}}\left\{1+{\displaystyle \frac{4}{\pi }}u\mathrm{\Phi }_1(u)+{\displaystyle \frac{4}{\pi }}\mathrm{\Phi }_2(u)\right\};`$ (161) $`K(t,ϵ)`$ $``$ $`{\displaystyle \frac{1}{\delta ^3}}\left\{2u+2{\displaystyle \frac{8u}{\pi }}\mathrm{\Phi }_1(u){\displaystyle \frac{8}{\pi }}\mathrm{\Phi }_2(u)\right\};`$ (162) $`L(t,ϵ)`$ $``$ $`{\displaystyle \frac{1}{\delta ^5}}\{2+2u+u^2`$ (164) $`{\displaystyle \frac{16}{3\pi }}[u^2\mathrm{\Phi }_1(u)+2u\mathrm{\Phi }_2(u)+2\mathrm{\Phi }_3(u)]\}.`$ Here $`u=t\delta ^2`$. For brevity we have written $$\mathrm{\Phi }_n(u)=_0^{\mathrm{}}𝑑s\frac{e^{us^2}}{(1+s^2)^n}.$$ (165) Comparing eq (79) to (102) we conclude that the exponents are $`\theta =1,\phi =2,\kappa =1`$ and $`\mu =1`$ as conjectured. It is also straightforward to extract the scaling functions $`(u),(u),𝒦(u)`$ and $`(u)`$ from eq (102). The scaling functions are plotted in figs 4, 5, 6 and 7 respectively. The asymptotics of the integrals $`\mathrm{\Phi }_n(u)`$ are analysed in Appendix C. Using those results we conclude that for small $`u`$ $`(u)`$ $``$ $`{\displaystyle \frac{2}{\sqrt{\pi }}}u^{1/2};`$ (166) $`(u)`$ $``$ $`u;`$ (167) $`𝒦(u)`$ $``$ $`{\displaystyle \frac{8}{3\sqrt{\pi }}}u^{3/2};`$ (168) $`(u)`$ $``$ $`{\displaystyle \frac{16}{15\sqrt{\pi }}}u^{5/2}.`$ (169) For large $`u`$ $`(u)`$ $``$ $`1;`$ (170) $`(u)`$ $``$ $`{\displaystyle \frac{2}{\sqrt{\pi }}}\sqrt{u};`$ (171) $`𝒦(u)`$ $``$ $`2u;`$ (172) $`(u)`$ $``$ $`u^2.`$ (173) Substituting the small $`u`$ asymptotics in eq (79) we obtain the behaviour for depths small compared to $`1/\delta ^2`$. $$w^2(t)\frac{2}{\sqrt{\pi }}t^{1/2}+\delta I^2t+I\frac{8}{3\sqrt{\pi }}t^{3/2}+I^2\frac{16}{15\sqrt{\pi }}t^{5/2}.$$ (174) This behaviour would persist for all depths right at the critical point. Note that eq (106) agrees with the forms conjectured in eq (82), (85) and (86) (including the numerical coefficient in the first case). It is hardly necessary to add that eq (106) is consistent with the critical point analysis of Majumdar and Sire since it leads, by the arguments given earlier in this section, to their crossover exponent 4/3 . The large $`u`$ asymptotics give the behaviour at depths large compared to $`1/\delta ^2`$. We find $$w^2(t)=\frac{1}{\delta }+\delta I^2\frac{1}{\delta }\frac{2}{\sqrt{\pi }}t^{1/2}+I\frac{1}{\delta }2t+I^2\frac{1}{\delta }t^2.$$ (175) The term proportional to $`I^2`$ has the form anticipated in eq (87); at the greatest depths this term is dominant. In summary we have shown that the singular distribution is an isolated critical point. Near the critical point the fluctuations in load have the scaling form eq (79). We have derived this scaling form and all the exponents. The results are in agreement with expectations based on simpler (non-rigorous) arguments. ## V Higher Dimensions We now turn to the $`q`$-model in $`D+1`$ dimensions. The quantum Hall multilayer and river networks are both 1+1 dimensional systems; bead-packs however are described by the 2+1 dimensional $`q`$-model. The behaviour of the model as a function of $`D`$ is of intrinsic interest moreover. We will find that right at the critical point the growth exponents vary smoothly with dimension for $`D<2`$. Above $`D=2`$ they become fixed, revealing $`D=2`$ as the upper critical dimension for the critical case. Off the critical point we expect the fluctuations to grow according to a scaling function $`(x)`$ (eq 7). We will study how the function and exponents vary with dimensionality below $`D=2`$. For simplicity in this section we neglect injection. ### A Model and Disorder Average First we must generalise the description of the $`q`$-model, so far confined to 1+1 dimensions. The case of 2+1 dimensions is easy to visualise. Fig 8 illustrates a square lattice composed of two interpenetrating square sublattices. The co-ordinates of sites $`\stackrel{}{n}=(n_1,n_2)`$ are both even for the black sites; both odd for the grey. The displacements from a site on either sublattice to its four nearest neighbours on the other sublattice are $`(\pm 1,\pm 1)`$. We will denote these displacements $`\stackrel{}{u}`$. In the $`q`$-model planes of such square lattices are stacked vertically. The beads alternately occupy only even or odd sublattices. Denoting the depth of a layer $`t`$, for $`t`$ even only the even sublattice is occupied; for $`t`$ odd, only the odd sublattice. Viewed in three dimensions the beads occupy a body-centered cubic stucture. In the same sense, Fig 1 can be viewed as a body-centered square structure. Now consider a $`D`$ dimensional simple cubic lattice. The co-ordinates of a site are specified by $`\stackrel{}{n}=(n_1,n_2,\mathrm{},n_D)`$ where $`n_i`$ are integers. For the even sublattice the $`n_i`$ are even; for the odd sublattice, they are odd. Each site has $`2^D`$ nearest neighbours on the other sublattice. We denote the displacements $`(\pm 1,\pm 1,\mathrm{},\pm )`$ to these neighbours $`\stackrel{}{u}`$. The $`D+1`$ dimensional $`q`$-model consists of $`D`$ dimensional cubic lattices stacked in the “vertical” $`t`$ direction. In alternate $`t`$ slices only the even or odd sublattices are occupied by beads. It is assumed that a random fraction of the load on each bead is supported by its $`2^D`$ neighbours in the layer below. The fractions must sum to one; $$\underset{\stackrel{}{u}}{}f_\stackrel{}{u}=1.$$ (176) Here $`f_\stackrel{}{u}`$ is the fraction of load transmitted by the bead to the neighbour separated by a horizontal displacement of $`\stackrel{}{u}`$. Hence the dynamics of the $`q`$-model is governed by $$w(\stackrel{}{n},t+1)=\underset{\stackrel{}{u}}{}f_\stackrel{}{u}(\stackrel{}{n}\stackrel{}{u},t)w(\stackrel{}{n}\stackrel{}{u},t).$$ (177) Eq (109) is the $`D+1`$ dimensional generalisation of eq (1). The fractions for a particular bead are assumed to be drawn from a distribution that is symmetric with respect to direction and respects the constraint eq (108). It follows $$f_\stackrel{}{u}=\frac{1}{2^D}.$$ (178) We write $$f_\stackrel{}{u}^2=\frac{1}{2^{2D}}+\frac{ϵ}{2^{2D}}$$ (179) where $`ϵ`$ is a parameter that characterises the distribution of fractions. From the sum constraint eq (108) it follows $$f_{\stackrel{}{u}_1}f_{\stackrel{}{u}_2}=\frac{1}{2^{2D}}\frac{ϵ}{(2^D1)}\frac{1}{2^{2D}}.$$ (180) for $`\stackrel{}{u}_1\stackrel{}{u}_2`$. The fractions are assumed to be independently and identically distributed for different beads. For the singular distribution all the fractions are zero except one. The probability for each fraction to be one is $`1/2^D`$. It is easy to calculate $`ϵ=2^D1`$ for the singular distribution using eq (111) and to verify eq (110) and (112) are satisfied. Since $`ϵ=2^D1`$ for the singular distribution we shall use $`\delta `$, defined by $$\delta =1\frac{ϵ}{(2^D1)},$$ (181) as our measure of the distance of a distribution from the critical point. As before it is useful to consider the correlation function $$c(\stackrel{}{m},t)=\underset{\stackrel{}{n}}{}w(\stackrel{}{n},t)w(\stackrel{}{n}+\stackrel{}{m},t).$$ (182) Note that $`\stackrel{}{m}`$ is a $`D`$ dimensional vector with even integer entries for both $`t`$ even and $`t`$ odd. The correlation function therefore lives on a simple cubic lattice in $`D`$ dimensions. By rescaling, as in section IIA, we reduce the lattice constant of this lattice to one so that the components of the vector $`\stackrel{}{m}`$ are now integers. The variance in load is related to the on-site correlation by $$\delta w^2(t)=c(\stackrel{}{m}0,t)1.$$ (183) Following the discussion of section IIA and using eqs (109), (110), (111) and (112) it is easy to show that the correlation function evolves with depth according to $`c(\stackrel{}{m},t+1)`$ $`=`$ $`{\displaystyle \frac{1}{2^D}}c(\stackrel{}{m},t)+{\displaystyle \frac{1}{2^{D+1}}}{\displaystyle \underset{\stackrel{}{b}=\mathrm{nn}}{}}c(\stackrel{}{m}+\stackrel{}{b},t)`$ (189) $`+{\displaystyle \frac{1}{2^{D+2}}}{\displaystyle \underset{\stackrel{}{b}=\mathrm{nnn}}{}}c(\stackrel{}{m}+\stackrel{}{b},t)+\mathrm{}`$ $`+{\displaystyle \frac{1}{2^{2D}}}{\displaystyle \underset{\stackrel{}{b}=\mathrm{n}\mathrm{}\mathrm{n}}{}}c(\stackrel{}{m}+\stackrel{}{b},t)+{\displaystyle \frac{ϵ}{2^D}}c(\stackrel{}{m},t)\delta _{\stackrel{}{m}=0}`$ $`{\displaystyle \frac{ϵ}{(2^D1)}}{\displaystyle \frac{1}{2^{D+1}}}{\displaystyle \underset{\stackrel{}{b}=\mathrm{nn}}{}}c(\stackrel{}{m}+\stackrel{}{b},t)\delta _{\stackrel{}{m}+\stackrel{}{b}=0}`$ $`{\displaystyle \frac{ϵ}{(2^D1)}}{\displaystyle \frac{1}{2^{D+2}}}{\displaystyle \underset{\stackrel{}{b}=\mathrm{nnn}}{}}c(\stackrel{}{m}+\stackrel{}{b},t)\delta _{\stackrel{}{m}+\stackrel{}{b}=0}`$ $`\mathrm{}{\displaystyle \frac{ϵ}{(2^D1)}}{\displaystyle \frac{1}{2^{2D}}}{\displaystyle \underset{\stackrel{}{b}=\mathrm{n}\mathrm{}\mathrm{n}}{}}c(\stackrel{}{m}+\stackrel{}{b},t)\delta _{\stackrel{}{m}+\stackrel{}{b}=0}`$ While reading eq (116) it is useful to recall that the correlation function lives on a $`D`$-dimensional cubic lattice. For $`D=2`$ each site has four nearest neighbours and four next-nearest neighbours. For general $`D`$, each site has $`2D`$ nearest neighbours; $`2^2C(D,2)`$ next nearest neighbours; $`2^3C(D,3)`$ third nearest neighbours; and $`2^DC(D,D)`$ $`D^{\mathrm{th}}`$ nearest neighbours. In eq (116) $`\stackrel{}{b}`$ denotes the displacement from a site to any of these neighbours; nn denotes nearest neighbour; nnn, next nearest; and so forth. In the next subsection we will solve eq (116) for $`c(\stackrel{}{m}0,t)`$ subject to the initial condition that a uniform load has been applied to the top layer. Thus $`c(\stackrel{}{m},t0)=1`$ for all $`\stackrel{}{m}`$. ### B Solution It is easy to verify that $`c(\stackrel{}{m},t\mathrm{})`$ $`=`$ $`{\displaystyle \frac{1}{\delta }}\mathrm{for}\stackrel{}{m}=0,`$ (191) $`=1\mathrm{otherwise},`$ is a steady state solution to eq (116). Eq (117) shows that the variance $`\delta w^2`$ saturates at sufficient depth in all dimensions for all distributions except the singular. We now calculate the evolution of the variance with depth using a method different from that of section II . First we $`z`$-transform the (discrete) $`t`$ dependence of the correlation function, $$c(\stackrel{}{m},z)=\underset{t=0}{\overset{\mathrm{}}{}}c(\stackrel{}{m},t)z^t,$$ (192) and Fourier transform the space dependence, $$c(\stackrel{}{p},z)=\underset{m}{}e^{i\stackrel{}{p}.\stackrel{}{m}}c(\stackrel{}{m},z).$$ (193) The use of the same symbol for the correlation and its transforms, although customary, is potentially confusing. For example, $`c(\stackrel{}{p},t0)`$ denotes the Fourier transform of $`c(\stackrel{}{m},t)`$ at $`t=0`$; no $`z`$-transform is implied. Performing both transforms on eq (116) we obtain $`c(\stackrel{}{p},z)`$ $`=`$ $`c(\stackrel{}{p},t0)+zc(\stackrel{}{p},z)S(\stackrel{}{p})`$ (195) $`+{\displaystyle \frac{ϵ}{2^D1}}zc(\stackrel{}{m}0,z).`$ Here $`S(\stackrel{}{p})`$ $`=`$ $`{\displaystyle \frac{1}{2^D}}\{1+{\displaystyle \frac{1}{2}}{\displaystyle \underset{b=\mathrm{nn}}{}}e^{i\stackrel{}{p}.\stackrel{}{b}}+{\displaystyle \frac{1}{2^2}}{\displaystyle \underset{b=\mathrm{nnn}}{}}e^{i\stackrel{}{p}.\stackrel{}{b}}`$ (197) $`+\mathrm{}+{\displaystyle \frac{1}{2^D}}{\displaystyle \underset{b=\mathrm{n}\mathrm{}\mathrm{n}}{}}e^{i\stackrel{}{p}.\stackrel{}{b}}\}`$ $`=`$ $`{\displaystyle \frac{(1+\mathrm{cos}p_1)}{2}}{\displaystyle \frac{(1+\mathrm{cos}p_2)}{2}}\mathrm{}{\displaystyle \frac{(1+\mathrm{cos}p_D)}{2}}`$ (198) is a “structure factor” for the cubic lattice. It will also prove convenient to define $$G(\stackrel{}{p},z)=\frac{1}{1zS(\stackrel{}{p})}.$$ (199) Both $`S(\stackrel{}{p})`$ and $`G(\stackrel{}{p},z)`$ have helpful physical interpretations that we shall make use of below. For the moment we rearrange eq (120) to obtain $`c(\stackrel{}{p},z)`$ $`=`$ $`c(\stackrel{}{p},t0)G(\stackrel{}{p},z)`$ (201) $`+(1\delta )zc(\stackrel{}{m}0,z)[1S(\stackrel{}{p})]G(\stackrel{}{p},z).`$ By inverting the Fourier transform we can turn eq (123) into an expression for $`c(\stackrel{}{m}0,z)`$. After further re-arrangement $`c(\stackrel{}{m}0,z)`$ $`=`$ $`{\displaystyle \frac{\frac{d\stackrel{}{p}}{(2\pi )^D}c(\stackrel{}{p},t0)G(\stackrel{}{p},z)}{1(1\delta )z\frac{d\stackrel{}{p}}{(2\pi )^D}[1S(\stackrel{}{p})]G(\stackrel{}{p},z)}}.`$ (202) Eq (124) is a general expression for $`c(\stackrel{}{m}0,z)`$ for an arbitrary initial condition. For uniform loading of the top layer $$c(\stackrel{}{p},t0)=(2\pi )^D\delta (\stackrel{}{p}).$$ (203) It follows from eq (121) and (122) that $`G(\stackrel{}{p}0,z)=1/(1z)`$; hence eq (124) simplifies to $`c(\stackrel{}{m}0,z)=(1z)^1`$ (204) $`\times \left\{1(1\delta )z{\displaystyle \frac{d\stackrel{}{p}}{(2\pi )^D}[1S(\stackrel{}{p})]G(\stackrel{}{p},z)}\right\}^1.`$ (205) Eq (126), together with the definitions of the structure factor (eq 121) and $`G(\stackrel{}{p},z)`$ (eq 122), constitutes an exact formal evaluation of the variance with depth. To obtain $`\delta w^2(t)`$ explicitly it only remains to peform the integral over $`\stackrel{}{p}`$ and to invert the $`z`$-transform. We return to this task in the next subsection. We conclude this subsection with a useful interpretation of $`S(\stackrel{}{p})`$ and $`G(\stackrel{}{p},z)`$. Eq (116) with $`ϵ0`$ resembles the Schrödinger equation for a particle on a $`D`$-dimensional cubic lattice with hopping to the nearest neighbours, the next nearest neighbours, and so on to the $`D^{\mathrm{th}}`$ nearest neighbours. It is not difficult to see that the eigenstates of this Schrödinger equation are plane waves. $`S(\stackrel{}{p})`$ is the dispersion relation, the eigenvalue at wave vector $`\stackrel{}{p}`$. From eq (121) we see that the energy level spectrum is a continuous band between zero and one. The momentum space Green’s function for this tight-binding lattice would normally be written $$𝒢(\stackrel{}{p},E)=\frac{1}{ES(\stackrel{}{p})}.$$ (206) Comparing eq (127) to eq (122) we see that $`G(\stackrel{}{p},z)`$ is essentially the Green’s function with $`E1/z`$. It is familiar from quantum mechanics that the real space Green’s function at the origin, $$𝒢(\stackrel{}{m}0,E)=\frac{d\stackrel{}{p}}{(2\pi )^D}\frac{1}{ES(\stackrel{}{p})},$$ (207) regarded as a function of (complex) $`E`$, has a branch cut running from $`E=0`$ to $`E=1`$, the interval that supports the eigenvalue band. It is not difficult to use the familiar arguments to conclude that, regarded as a function of complex $`z`$, $`c(\stackrel{}{m}0,z)`$ has a branch cut along the line $`z=1`$ to $`\mathrm{}`$ (onto which the segment maps under the transformation $`E1/z`$). The analytic properties of $`c(\stackrel{}{m}0,z)`$ will prove useful in the next sub-section. ### C Scaling Limit In this subsection we study the evolution of the variance in the large depth scaling limit. Thus $`t1`$ and $`\delta `$ is zero or very close to it throughout. An advantage of studying the large depth limit is that we do not have to calculate $`c(\stackrel{}{m}0,z)`$ exactly; it is only necessary to calculate the leading behaviour as $`z1`$. One way to understand this is to consider the critical case $`\delta =0`$. In this case we expect that at great depth $$c(\stackrel{}{m}0,t)t^x.$$ (208) It is easy to show that for $`f(t)=t^x`$, the $`z`$-transform is $`\mathrm{\Gamma }(x+1)/(1z)^{x+1}`$ plus less singular terms. Thus for a function that behaves as $`t^x`$ for large $`t`$ also the $`z`$-transform is $$t^x\frac{\mathrm{\Gamma }(x+1)}{(1z)^{x+1}}+\mathrm{less}\mathrm{singular}.$$ (209) If we know the leading singularity of $`c(\stackrel{}{m}0,z)`$ as $`z1`$ we can use eq (130) to read off the large depth behaviour. Another way to see that we only need the behaviour of $`c(\stackrel{}{m}0,z)`$ as $`z1`$ is to consider inverting the $`z`$-transform by the contour integral method of Appendix B. This is accomplished by folding the contour over the branch point of $`c(\stackrel{}{m}0,z)`$ at $`z=1`$ and integrating along the cut. In that integral $`c(\stackrel{}{m}0,z)`$ is weighted by a factor that decays extremely rapidly away from $`z=1`$ at large depths. Our goal therefore is to analyse the $`z1`$ behaviour of $$G(z)=\frac{d\stackrel{}{p}}{(2\pi )^D}\frac{1}{1zS(\stackrel{}{p})}$$ (210) since by a straightforward re-arrangement the integral in eq (126) simplifies to $$\frac{d\stackrel{}{p}}{(2\pi )^D}[1S(\stackrel{}{p})]G(\stackrel{}{p},z)=\left(1\frac{1}{z}\right)G(z)+\frac{1}{z}.$$ (211) Insight into the behaviour of $`G(z)`$ can be gained by expanding $`S(\stackrel{}{p})`$ around $`\stackrel{}{p}=0`$ to obtain $$G(z)\frac{d\stackrel{}{p}}{(2\pi )^D}\frac{1}{(1z)+\stackrel{}{p}^2/4}.$$ (212) If we set $`z=1`$ in eq (133) the integrand diverges as $`\stackrel{}{p}0`$ for $`D2`$; it is regular in more than two dimensions. Thus in more than two dimensions $`G(z)`$ has a branch point at $`z=1`$ but there is no actual divergence. In two dimensions or less there is an actual divergence. The leading behaviour of $`G(z)`$ above two dimensions is thus simply obtained by setting $`z=1`$ in eq (131): $$G(z)G(1)\mathrm{for}D>2.$$ (213) In two dimensions we can obtain the singularity by recognising $`G(z)`$ to be a Jacobi elliptic integral. Square lattice Green’s functions are known to be related to Jacobi’s elliptic functions; but since our lattice features next-nearest neighbour hopping, in addition to the customary nearest neighbour hopping, we outline the analysis in Appendix D. The result is that for $`z1`$ $$G(z)\frac{1}{\pi }\mathrm{ln}(1z)\mathrm{for}D=2.$$ (214) For $`D<2`$ we obtain the singular behaviour of $`G(z)`$ in Appendix D. The result is $$G(z)=\frac{\mathrm{\Gamma }(1D/2)}{\sqrt{\pi }^D}(1z)^{D/21}\mathrm{for}D<2.$$ (215) An important feature revealed by this calculation is that the singular behaviour of $`G(z)`$ is controlled by the long wavelength behaviour of $`G(p,z)`$ for all $`D<2`$; it breaks down as $`D2`$. Although it is instructive to do the calculation for continuous $`D`$ to examine the $`D2`$ limit, the only case that is physically relevant is of course the integer dimension $`D=1`$. Equipped with the leading behaviour of $`G(z)`$ in all dimensions we now obtain the long time behaviour of $`\delta w^2(t)`$. At the critical point we set $`\delta =0`$ and substitute eqs (132), (134), (135) and (136) in eq (126). Except in two dimensions the $`z`$-transforms may be inverted by inspection of eq (130). For two dimensions we must resort to the method of Appendix B and finally obtain $`\delta w^2(t)`$ $`=`$ $`\pi ^{D/21}\mathrm{sin}\left({\displaystyle \frac{\pi D}{2}}\right)t^{D/2}\mathrm{for}D<2`$ (216) $`=`$ $`{\displaystyle \frac{t}{\mathrm{ln}t}}\mathrm{for}D=2`$ (217) $`=`$ $`{\displaystyle \frac{1}{G(1)}}t\mathrm{for}D>2.`$ (218) As indicated by the simple steady state solution, at the critical point the fluctuations grow without bound as a power of $`t`$ for all dimensions. The exponent becomes independent of $`D`$ for $`D>2`$ revealing $`D=2`$ as the upper critical dimension. By substituting eq (132) and (136) in eq (126) we can also obtain the behaviour of $`\delta w^2(t)`$ away from the critical point for less than two dimensions. Inverting the $`z`$-transform by the method of Appendix B we find $$\delta w^2(t)=\frac{1}{\delta ^\theta }(t\delta ^\phi ).$$ (219) Here the exponents $$\theta =1,\phi =\frac{2}{D}$$ (220) and the scaling function $$(u)=\frac{1}{D}\frac{2}{\pi D}_0^{\mathrm{}}𝑑s\frac{1}{1+s^2}\mathrm{exp}\frac{us^{2/D}}{q_D^{2/D}}$$ (221) with $`q_D=\mathrm{\Gamma }(1D/2)\mathrm{sin}(\pi D/2)/\sqrt{\pi }^D`$ a dimension dependent constant. Again, only the result for $`D=1`$ is physically meaningful; in this case eq (140) coincides with the result of section II. In summary the main results of this section are that for all distributions, except the singular, at sufficient depth the load fluctuations saturate and (in agreement with experiment) there are no horizontal correlations in load (eq 117). The saturation value of the load variance diverges as the critical point is approached. At the critical point the load fluctuations grow without bound as a power of depth (eq 137). Below two dimensions this exponent depends on dimensionality; above two dimensions it is constant, revealing $`D=2`$ as the critical dimension. At the critical dimension the growth of fluctuations is tempered by a logarithmic factor as might be expected at a critical dimension. We have also evaluated the scaling function that describes the growth and saturation of load fluctuations near the critical point for $`D<2`$. ## VI Quantum Hall Multilayer ### A Models In this section we turn to the chiral wave models that are believed to adequately describe the surface electronic states of a quantum Hall multilayer. We begin by examining the circumstances under which the quantum network model of Saul, Kardar and Read discussed in section I is equivalent to a $`q`$-model. Following Saul, Kardar and Read, the first step is to identify pairs of links (joined by vertical grey bars in Fig 9) as “beads”. The “load” on a bead is the total probability that the electron is on either of its two constituent links. Load propagates from left to right now rather than top to bottom as it did in our earlier depictions of the $`q`$-model. For this reason we will label the vertical co-ordinate $`n`$ and the horizontal co-ordinate $`t`$ here (see fig 3). To analyse how load propagates consider an elementary vertex of the Saul, Kardar and Read model shown in Fig 9. The wave function amplitudes are related via $$\left(\begin{array}{c}\varphi _2\\ \varphi _3\end{array}\right)=S\left(\begin{array}{c}\psi _1\\ \psi _2\end{array}\right);$$ (222) here $`S`$ is a random $`2\times 2`$ su(2) rotation matrix. Saul, Kardar and Read assumed the $`S`$-matrices were drawn from the invariant distribution for the su(2) group . The loads on beads A, B and C are respectively $`\|\psi _1\|^2+\|\psi _2\|^2`$, $`\|\varphi _1\|^2+\|\varphi _2\|^2`$ and $`\|\varphi _3\|^2+\|\varphi _4\|^2`$ . By unitarity $`\|\psi _1\|^2+\|\psi _2\|^2=\|\varphi _3\|^2+\|\varphi _4\|^2`$. Thus bead A sends a fraction $`f`$ of its load to neighbour B and the remainder $`1f`$ to neighbour C. A key feature of the Saul, Kardar and Read model is that the distribution of the fractions, $`P(f)`$ is independent of the input amplitudes $`\psi _1`$ and $`\psi _2`$. This follows from the assumed group invariant distribution for the $`S`$-matrices. It is this feature that allows the Saul, Kardar and Read model to be mapped onto a $`q`$-model. To derive the distribution of the fractions recall that an su(2) matrix may be parametrized $`S=x_0+i\stackrel{}{x}.\stackrel{}{\sigma }`$ with $`(x_0,\stackrel{}{x})`$ real and subject to $`x_0^2+\stackrel{}{x}^2=1`$. If we take $`\psi _1=1,\psi _2=0`$ then $`f=x_0^2+x_1^2`$. From the invariant distribution for su(2) matrices, $$P(x_0,\stackrel{}{x})=\frac{1}{\pi }\delta (x_0^2+\stackrel{}{x}^21),$$ (223) it is not difficult to show that the fraction $`f`$ follows the uniform distribution, $`P(f)=1`$ for $`0<f<1`$. Now suppose the wave function is known through the vertical slice $`t=0`$. We could propagate the wave function $`t`$ slices to the right using the quantum Saul, Kardar and Read model. Alternatively we could calculate the load in the initial layer and propagate it to the right using the $`q`$-model with uniform distribution. Either way the load we obtain in layer $`t`$ would be the same statistically. This is the sense in which the Saul, Kardar and Read model is equivalent to the $`q`$-model. Note that the $`q`$-model does not keep track of phase information. The mapping is useful only under circumstances that the phase information is unimportant. Below we will discuss some problems of wave packet dynamics for which the mapping is useful. The mapping can also be used to study vertical transport in the quantum Hall multilayer in the limit of large circumference but we do not discuss that application here. An obvious circumstance when the phase information is important and the mapping cannot be used is if periodic boundary conditions are imposed in the horizontal $`t`$-direction, as would be appropriate for a multilayer in the fully phase-coherent, mesoscopic regime. Phase information is needed to match the wavefunction after it is propagated around the circumference. We will develop this point in a more technical way in subsection VI C. Another case in which a quantum network model will map onto a $`q`$-model is if the wavefunctions and $`S`$-matrices are chosen to be real and the $`S`$-matrices are further assumed to be distributed over the subgroup of rotations about the $`y`$-axis with appropriate invariant measure. The fraction distribution $`P(f)=(1/\pi )f^{1/2}(1f)^{1/2}`$ for the $`q`$model that results. For most distributions of the $`S`$-matrix however it is not possible to obtain even the limited mapping between the quantum network model and the classical $`q`$-model obtainable in this and in the Saul, Kardar and Read case. Finally we present a convenient continuum model of the multilayer surface governed by the Schrödinger equation $$i\frac{}{t}\psi _n(t)=m_n(t)\psi _{n+1}(t)+m_{n1}^{}(t)\psi _{n1}(t).$$ (224) Since the equation is first order in $`t`$, given the wavefunction at a fixed $`t`$ slice we can use it to propagate the wavefunction to the right, just as in the discrete network model. In the transverse direction the model is discrete and second -order. Disorder is incorporated by taking the hopping elements $`m_n(t)`$ to be random. For a discussion of the relationship between onsite and hopping disorder see ref . Evidently this model cannot be reduced to a classical $`q`$-model. ### B Wave-packet dynamics In this section we briefly discuss wave packet dynamics for the models of the previous section. Mathematically this problem is identical to the motion of a wave-packet in a crystal with noise (temporal randomness). It also bears formal resemblance to the directed polymer model, an important problem in statistical mechanics. Hence it is a problem of general interest and has been studied since at least the 1970s from various points of view (see ref and refs therein). A considerable amount is now known. For the Saul, Kardar and Read model wave-packet dynamics can be studied using the mapping to the $`q`$-model; indeed the mapping was introduced for this purpose. In this section we will formulate the problem and summarize known results. These results reveal that the $`q`$-model and the continuum wave model introduced in the last section behave in qualitatively similar ways. Consider an electron localized at $`n=0`$ at $`t=0`$. This wavepacket can be propagated to the right using eq (143). As it propagates it will broaden and its mean position will deflect. It is interesting to know how the breadth and deflection grow with displacement and to analyse the distribution of “load” at sufficiently great displacement that a steady state is reached. The root mean square width of the wave-packet grows as the square root of the displacement. This was derived for the continuum model in the 1970s and it is easy to show that the same form is obtained in the Saul, Kardar and Read model. The root mean square deflection grows as the fourth root of the displacement. This result has been obtained numerically and analytically for both the Saul, Kardar and Read and continuum models . To compare the distribution of load, for the continuum model we define the load on an edge as $`w_n(t)=\|\psi _n(t)\|^2`$. The asymptotic distribution of load, $`\mathrm{\Pi }(w,t\mathrm{})`$ was obtained by Coppersmith et al. for the $`q`$-model . For various distributions of the fractions, $`P(f)`$, they found that $`\mathrm{\Pi }(w,t\mathrm{})`$ decayed exponentially with $`w`$ with a power law prefactor that depended on the distribution $`P(f)`$. For the uniform distribution the prefactor was a constant. The corresponding result for the continuum wave model was obtained by ref by mapping the problem onto an su(1,1) quantum ferromagnet. Here too the result for the load distribution is an exponential with a prefactor linear in $`w`$. ### C Field Theory Formulation We have emphasized above that the equivalence between the Saul, Kardar and Read model and the $`q`$-model is useful only when open boundary conditions are imposed in the horizontal $`t`$-direction; it breaks down for periodic boundary conditions needed to describe transport in phase-coherent multilayers. The importance of boundary conditions is also reflected in field theory formulations of these models. In ref the continuum model with open boundary conditions was mapped onto a Heisenberg ferromagnet. In contrast, with periodic boundary conditions a mapping to a supersymmetric analogue of the Heisenberg ferromagnet was obtained in refs In this section we derive the supersymmetric spin representation following the operator methods of ref . This derivation highlights the role of boundary conditions, the feature we wish to emphasize here. It only makes use of operator methods and is in this sense more elementary than the functional methods of ref . Moreover mappings to supersymmetric spin models have recently been used fruitfully not only to study the multilayer but also to provide non-perturbative insights into various other problems of electron localization . It is hoped that the present derivation, with its emphasis on boundary conditions<sup>§</sup><sup>§</sup>§For the effect of boundary conditions on the supersymmetry mapping for models such as the Chalker model of the quantum Hall transition see refs and where the random hopping model in one-dimension is analysed with periodic and open (scattering) boundary conditions respectively. and its use of operator methods will prove of interest in this broader context also. #### 1 Fermion Representation We wish to evaluate $`G(n,t;n^{},t^{})`$, the Green’s function for the continuum model governed by the Schrödinger equation, $`i{\displaystyle \frac{}{t}}G(n,t;n^{},t^{})`$ $`=`$ $`m_n(t)G(n+1,t;n^{},t^{})`$ (227) $`+m_{n1}^{}(t)G(n1,t;n^{},t^{})`$ $`i\delta (tt^{})\delta _{nn^{}},`$ and subject to the periodic boundary condition $$G(n,t+T;n^{},t^{})=G(n,t;n^{},t^{}).$$ (228) Here $`T`$ is the period in the $`t`$-direction. In ref , the Green’s function was calculated subject to the chiral boundary condition, $`G(n,t;n^{},t^{})=0`$ for $`t<t^{}`$, leading to a simpler field theory formulation. The key idea is to reinterpret the co-ordinate $`t`$ as time. Eq (144) then describes a particle on a one-dimensional lattice with noise. In second quantised notation the (time-dependent) Hamiltonian that governs the motion of this fictitious particle is $$H_F^R(t)=\underset{n}{}\left[m_n(t)c_n^Rc_{n+1}^R+m_{n1}^{}(t)c_n^Rc_{n1}^R\right].$$ (229) Here $`c_n^R`$ creates a Fermion at site n; $`c_n^R`$ annhilates it. The reasons for the superscript on the Fermion Hamiltonian and on the creation and annhilation operators will become apparent shortly. The $`S`$-matrix for this model obeys $$i\frac{}{t}S_F^R(t)=H_F^R(t)S_F^R(t)$$ (230) subject to $`S_F^R(t0)=1`$. From $`S_F^{R1}S_F^R=1`$ it is easy to verify the useful result $$i\frac{}{t}S_F^{R1}=S_F^{R1}(t)H_F^R(t).$$ (231) We define $$c_n^R(t)=S_F^{R1}(t)c_nS_F^R(t)$$ (232) and similarly for $`c_n^R(t)`$. Now by analogy with finite temperature field theory we write the Green’s function $`G(n,t;n^{},t^{})`$ $`=`$ $`\mathrm{Tr}[S_F^R(T)c_n^R(t)c_n^{}^R(t^{})]/Z_F^R(T)`$ (234) $`\mathrm{for}t>t^{}`$ $`=`$ $`\mathrm{Tr}[S_F^R(T)c_n^{}^R(t^{})c_n^R(t)]/Z_F^R(T)`$ (236) $`\mathrm{for}t<t^{};`$ $`Z_F^R(T)`$ $`=`$ $`\mathrm{Tr}[S_F^R(T)].`$ (237) $`Z_F^R(T)`$ is analogous to the partition function in finite temperature field theory. It is easy to verify that $`G`$ obeys the differential eq (144) by making use of eqs (147) and (148) and the commutation relation $$[H_F^R(t),c_n^R]=m_n(t)c_{n+1}^Rm_{n1}^{}(t)c_{n1}^R.$$ (238) However $`G(n,T;n^{},t^{})`$ $`=`$ $`\mathrm{Tr}[S_F^R(T)S_F^{R1}(T)c_n^RS_F^R(T)c_n^{}^R(t^{})]/Z_F^R(T)`$ (239) $`=`$ $`\mathrm{Tr}[S_F^R(T)c_n^{}^R(t^{})c_n^R]/Z_F^R(T)`$ (240) $`=`$ $`G(n,0;n^{},t^{}).`$ (241) Thus $`G`$ obeys antiperiodic rather than periodic boundary conditions. This problem is fixed by adding a term to the Hamiltonian $$H_F^R(t)H_F^R(t)+\frac{\pi }{T}\underset{n}{}c_n^Rc_n^R.$$ (242) Alternatively we may replace $`\mathrm{Tr}\mathrm{STr}`$ in eq (150). By STr we mean the trace of an operator over all states with an even number of fermions minus the trace over states with an odd number of fermions. We also need an expression for the complex conjugate of the Green’s function since our ultimate purpose is to calculate the disorder average of $`\|G(n,t;n^{},t^{})\|^2`$, the diffuson propagator. To this end we complex conjugate eq (144) to obtain the differential equation obeyed by $`G^{}`$. Comparison to eq (144) reveals that we should consider $`A`$ fermions governed by the Hamiltonian $$H_F^A(t)=\underset{n}{}\left[m_n(t)c_{n+1}^Ac_n^A+m_n^{}(t)c_n^Ac_n^A\right].$$ (243) $`G^{}(n,t;n^{},t^{})`$ is then given by the right hand side of eq (150) if we replace $`RA`$ and $`\mathrm{Tr}\mathrm{STr}`$. As might be expected the Hamiltonian for the $`A`$ fermions is related to that for the $`R`$ fermions via a particle hole transformation. This symmetry between the $`R`$ fermions and the $`A`$ holes leads to an su(2) symmetry in the fermion sector of the complete field theory formulation that we obtain below (eq 171). It is also at the root of the supersymmetry of the field theory formulation. In summary the Green’s function with periodic boundary conditions may be generated from the second quantized Hamiltonian, $`H_F^R(t)`$ \[eq (153) and (146)\] using the definition eq (150). The complex conjugate of the Green’s function may be obtained similarly using the Hamiltonian $`H_F^A(t)`$ (eq 154). Eq (150) and its A fermion analogue provide exact formal expressions for the Green’s function for a particular realization of the random tunneling $`m_n(t)`$. These expressions are not particularly convenient to average since $`m_n(t)`$ appears in both numerator and denominator. #### 2 Boson Representation Alternatively we could interpret eq (144) as a time dependent Schrödinger equation for bosonic particles on a one dimensional lattice. The corresponding “time”-dependent bosonic Hamiltonian in second quantized notation is $$H_B^R(t)=\underset{n}{}\left[m_n(t)b_n^Rb_{n+1}^R+m_{n1}^{}(t)b_n^Rb_{n1}^R\right].$$ (244) Here $`b_n^R`$ creates an $`R`$ boson at site $`n`$; $`b_n^R`$ annhilates it. The Green’s function is now defined as $`G(n,t;n^{},t^{})`$ $`=`$ $`\mathrm{Tr}[S_B^R(T)b_n^R(t)b_n^{}^R(t^{})]/Z_B^R(T)`$ (246) $`\mathrm{for}t>t^{}`$ $`=`$ $`\mathrm{Tr}[S_B^R(T)b_n^{}^R(t^{})b_n^R(t)]/Z_B^R(T)`$ (248) $`\mathrm{for}t<t^{};`$ $`Z_B^R(T)`$ $`=`$ $`\mathrm{Tr}S_B^R(T).`$ (249) Here $`S_B^R`$ is the bosonic S-matrix and $`Z_B^R(T)`$ is the bosonic analogue of the partition function. For greater rigour we must regulate the traces to ensure convergence but for brevity we do not discuss this explicitly here. The complex conjugate of the Green’s function is generated similarly if instead of the R bosons we consider A bosons governed by $$H_B^A(t)=\underset{n}{}[m_n(t)b_{n+1}^Ab_n^A+m_n^{}(t)b_n^Ab_{n+1}^A].$$ (250) The main result of this subsubsection is eq (156). It provides a formal bosonic expression for the exact Green’s function for a particular realization of random tunneling, $`m_n(t)`$. A similar expression for $`G^{}`$ may be obtained by working with the Hamiltonian eq (157). Like their fermionic counterparts these bosonic expressions are not particularly well suited for averaging over disorder. #### 3 Supersymmetry We now develop an expression for the diffuson suitable for averaging over disorder. In Appendix E it is shown that $$Z_F^R(T)Z_B^R(T)=1.$$ (251) Thus we consider a model that includes both A and R fermions and bosons governed by the Hamiltonian $`H_{\mathrm{SUSY}}(t)`$ $`=`$ $`H_F^R(t)+H_F^A(t)+H_B^R(t)+H_B^A(t)`$ (252) $`=`$ $`{\displaystyle \underset{n}{}}[m_n(t)A_n+m_n^{}(t)A_n^{}].`$ (253) Here $$A_n=c_n^Rc_{n+1}^Rc_{n+1}^Ac_n^A+b_n^Rb_{n+1}^Rb_{n+1}^Ab_n^A.$$ (254) The corresponding S-matrix obeys $$i\frac{}{t}S_{\mathrm{SUSY}}(t)=H_{\mathrm{SUSY}}(t)S_{\mathrm{SUSY}}(t)$$ (255) subject to $`S_{\mathrm{SUSY}}(t0)=1`$. A formal solution to eq (161) is given by $$S_{\mathrm{SUSY}}(t)=P\mathrm{exp}\left(i_0^t𝑑t_1H_{\mathrm{SUSY}}(t_1)\right).$$ (256) Here $`P`$ is the chronological ordering operator. Hence the diffuson is given by $`\|G(n,t;n^{},t^{})\|^2`$ $`=`$ $`\mathrm{STr}[S_{\mathrm{SUSY}}(T)c_n^R(t)c_n^A(t)c_n^{}^A(t^{})c_n^{}^R(t^{})]`$ (258) $`\mathrm{for}t>t^{}`$ $`=`$ $`\mathrm{STr}[S_{\mathrm{SUSY}}(T)c_n^{}^A(t^{})c_n^{}^R(t^{})c_n^R(t)c_n^A(t)]`$ (260) $`\mathrm{for}t<t^{}`$ The content of eq (163) is that to calculate the diffuson we must create or annhilate a pair of R and A fermions (depending on the time order). Then we must propagate this state in accordance with $`H_{\mathrm{SUSY}}`$ and perform an S-matrix weighted trace. The Hamiltonian $`H_{\mathrm{SUSY}}`$ is non-interacting but it is random and time dependent. Eq (163) is an exact formal expression for the diffuson. Note the lack of a denominator, eliminated by virtue of eq (158). This feature allows us to perform the average over disorder easily. For example, $$S_{\mathrm{SUSY}}(t)=\mathrm{exp}[_{\mathrm{SUSY}}t]$$ (261) with $$_{\mathrm{SUSY}}=\frac{D}{2}\underset{n}{}(A_n^{}A_n+A_nA_n^{}).$$ (262) Here we have assumed that the tunneling $`m_n(t)`$ is a Gaussian white noise process with zero mean and variance $$m_n^{(\alpha )}(t)m_n^{}^{(\beta )}(t^{})=D\delta (tt^{})\delta _{nn^{}}\delta _{\alpha \beta }.$$ (263) Here $`m_n^{(1)}(t)=`$ real part of $`m_n(t)`$; $`m_n^{(2)}(t)=`$ imaginary part of $`m_n(t)`$. Recall that for a single Gaussian random variable $`y`$, the phase average $`e^{iy}=e^{y^2/2}`$. Eq (164) is analogous to this result but with the added complications that $`S_{\mathrm{SUSY}}`$ is an ordered exponential, not a simple exponential, and the average is over a random process rather than a single random variable. To derive eq (165) it is simplest to expand the time ordered exponential (eq 162) and average term by term. Proceeding in this manner we obtain an expression for the average diffuson $`\|G(n,t;n^{},t^{})\|^2`$ $`=`$ $`\mathrm{STr}\{\mathrm{exp}[_{\mathrm{SUSY}}(Tt+t^{})]c_n^Rc_n^A`$ (266) $`\times \mathrm{exp}[_{\mathrm{SUSY}}(tt^{})c_n^Ac_n^R\}`$ $`\mathrm{for}t>t^{}.`$ A similar expression may be written for the case $`t<t^{}`$. The content of eq (167) is that to calculate the average diffuson we must create (or for the other time order, annhilate) a pair of R and A fermions and propagate the resulting state according to the effective Hamiltonian $`_{\mathrm{SUSY}}`$. In contrast to $`H_{\mathrm{SUSY}}`$ the effective Hamiltonian is not time dependent or random but it is interacting. This completes our formulation of the continuum directed wave model of section VI A as a superspin field theory. The main results are the superspin Hamiltonian (eq 165) and eq (167) which shows how interesting correlation functions are calculated in the superspin formulation. The usefulness of this formulation depends on the extent to which the superspin model can be analysed. In the remainder of this section we discuss the form and symmetry of the superspin Hamiltonian (eq 165). To this end it is helpful to introduce special notation for the boson and fermion bilinears of which $`_{\mathrm{SUSY}}`$ is composed. We denote the fermion bilinears $`J_+=c^Rc^A=J_x+iJ_y,`$ (267) $`J_{}=c^Ac^R=J_xiJ_y,`$ (268) $`J_z={\displaystyle \frac{1}{2}}(c^Rc^R+c^Ac^A1),`$ (269) $`J={\displaystyle \frac{1}{2}}(c^Rc^Rc^Ac^A+1);`$ (270) the boson bilinears, $`K_+=b^Rb^A=K_x+iK_y,`$ (271) $`K_{}=b^Ab^R=K_xiK_y,`$ (272) $`K_z={\displaystyle \frac{1}{2}}(b^Rb^R+b^Ab^A+1),`$ (273) $`K={\displaystyle \frac{1}{2}}(b^Rb^Rb^Ab^A1);`$ (274) and the mixed bilinears, $`M_1=b^Rc^R,M_2=b^Ac^A,`$ (275) $`L_1=b^Ac^R,L_2=b^Rc^A.`$ (276) In eqs (168), (169) and (170) the site indices have been suppressed for brevity. In terms of these bilinears we may write $`_{\mathrm{SUSY}}`$ $`=`$ $`2D{\displaystyle \underset{n}{}}(\stackrel{}{J}_{n+1}.\stackrel{}{J}_n+J_{n+1}J_nJ_n)`$ (280) $`+2D{\displaystyle \underset{n}{}}(\stackrel{}{K}_{n+1}.\stackrel{}{K}_n+K_{n+1}K_n+K_n)`$ $`+D{\displaystyle \underset{n}{}}\left(M_{n+1}^{(1)}M_n^{(1)}+M_{n+1}^{(2)}M_n^{(2)}+\mathrm{hc}\right)`$ $`+D{\displaystyle \underset{n}{}}\left(L_{n+1}^{(1)}L_n^{(1)}+L_{n+1}^{(2)}L_n^{(2)}+\mathrm{hc}\right).`$ Here hc denotes Hermitian conjugate and $`\stackrel{}{K}_{n+1}.\stackrel{}{K}_n=K_{n+1}^zK_n^zK_{n+1}^xK_n^xK_{n+1}^yK_n^y`$ . It is instructive to study the commutation relations for bilinears at the same site $`n`$ (bilinears at different sites simply commute or anticommute). It is easy to verify that $`J_+,J_{}`$ and $`J_z`$ satisfy angular momentum or su(2) commutation relations and $`J`$ commutes with the other three. Similarly $`K_+,K_{}`$ and $`K_z`$ satisfy the su(1,1) or hyperbolic angular momentum algebra—essentially the angular momentum algebra but with a sign change for the $`K_+,K_{}`$ commutator . $`K`$ commutes with the other three. The anticommutators of $`L_i,L_i^{},M_i`$ and $`M_i^{}`$ are linear combinations of the $`K`$’s and $`J`$’s. The commutators of the $`J`$’s or $`K`$’s with the $`L`$’s or $`M`$’s are linear combinations of the $`L`$’s and $`M`$’s. Hence these bilinears constitute a superalgebra. The $`J`$’s and $`K`$’s are commuting elements of the superalgebra; the $`L`$’s and $`M`$’s, anticommuting elements. The superalgebra is called u(1,1—2). It includes the Lie algebras su(2) and su(1,1) as subalgebras. Further insight into the superalgebra is obtained by considering the Hilbert space at a single site. This is a direct product of the four dimensional fermion space and the infinite dimensional two-boson space. The fermion space may be decomposed into irreducible representations of the su(2) algebra. The fermion vacuum and the state with both R and A fermions present constitute a doublet or spin 1/2 representation; the two states with one fermion present are singlets. The boson space similarly decomposes into an infinity of infinite dimensional irreducible representations of the su(1,1) algebraLet $`\|n+m,n`$ denote a state with $`(n+m)`$ R-bosons and $`n`$ A-bosons on the site. The infinite dimensional subspace with $`m`$ a fixed integer and $`n=0,1,2,\mathrm{}`$, for $`m0`$, or $`n=m,m+1,=m+2,\mathrm{}`$, for $`m<0`$, is invariant under the four $`K`$ operators. These subspaces corresponding to different values of $`m`$ constitute the irreducible representations of the su(1,1) algebra.. The single site Hilbert space thus decomposes rather simply into irreducible representations of the direct sum of the su(2) and su(1,1) algebra. These subspaces do not constitute a representation of the whole superalgebra. The anticommuting elements mix different irreducible representations of su(2) and su(1,1). In particular they mix representations with different spins—a celebrated feature of supersymmetry. It is not difficult to decompose the single site Hilbert space into blocks irreducible under the superalgebra; however this would carry us too far afield. More details on the superalgebra are given in ref and refs therein. Finally we define $$𝒥_{\mathrm{tot}}=\underset{n}{}𝒥_n.$$ (281) Here $`𝒥`$ denotes any element of the superalgebra such as $`J_+,K_z,L_1`$ etc. After some algebra we find $$[_{\mathrm{SUSY}},𝒥_{\mathrm{tot}}]=0$$ (282) revealing the supersymmetry of the field theory formulation. ## VII Summary and Conclusion Much of this paper is concerned with the behaviour of the $`q`$-model close to the critical point. To probe this behaviour we imagine that a uniform load is applied to the top layer. As the load propagates downward fluctuations develop in the distribution of load. Coppersmith et al. studied the entire distribution of load at very great depth where it was presumed that a steady state had been reached. In contrast we study only the variance of the distribution of load but we analyse its evolution with depth. Our purpose is to study this evolution for all distributions of the fractions, $`P(f)`$, particularly those close to the singular distribution (the critical point). In section II we consider the $`q`$-model in 1+1 dimensions without injection (the weight of the beads is neglected). In this case the average load does not vary with depth since the total load is the same in every layer; it is merely redistributed by the q-model dynamics. For the growth of the variance, by analogy to critical phenomena, we make the following hypotheses: For all distributions $`P(f)`$ except the singular distribution we posit that the variance will saturate at sufficient depth. Both the saturation depth and the saturated variance are expected to diverge as the distribution approaches the singular distribution. We introduce $`\delta `$, a measure of the distance of a distribution $`P(f)`$ from the singular distribution, and conjecture that the saturation depth $`\xi _{\mathrm{corr}}`$ will diverge as $`1/\delta ^\phi `$; the saturated variance, as $`1/\delta ^\theta `$. More specifically, we expect that close to the critical point the variance will have a scaling form, eq (7). For the singular distribution we expect that the variance will grow indefinitely as a power of the depth. Close to the critical point and at depths shallow compared to the saturation depth the variance should grow as it would right at the critical point. From this and from eq (7) we deduce a relationship between the critical exponents $`\theta `$ and $`\phi `$ and the exponent that describes the growth of the variance right at the critical point; namely we expect that at the critical point the variance will grow as $`t^{\theta /\phi }`$. In sections IIB and IIC we derive an exact formula for the variance as a function of depth (eq 46) and study its scaling limit ($`t1`$, $`\delta 0`$ but with $`t\delta ^\phi `$ arbitrary). These calculations bear out all the expectations enumerated above, provide the precise form of the scaling function \[eq (50) and Fig 4\] and yield the exact exponents (eq 49). In section III we characterise the critical point more fully by analysing the evolution with depth of the entire distribution of load right at the critical point in 1+1 dimensions. In the absence of injection the critical point is a simple model of random walkers that coalesce upon contact; hence it is quite straightforward to derive these results. We present them because they illuminate the results of the previous section. At large depth it is found that the distribution of load consists of a large spike at zero load together with a smooth part \[eq (73) and (76)\]. It is overwhelmingly probable that the load on a bead is zero; most of the weight of the distribution is in the spike. The smooth part follows the anticipated scaling form (eq 53). Its width grows as the square root of the depth, consistent with the exponent found in section II to describe the growth of the variance of load at the critical point. In section IV the effect of injection is included. For simplicity we consider only 1+1 dimensions. We assume that the weights of the beads are independent and identically distributed random variables. The behaviour of the mean load is still not very interesting. It grows linearly with depth (eq 77). Close to the critical point we conjecture that the variance will have the form eq (78). We are able to deduce all the exponents in eq (78) and to obtain some limiting behaviours of the scaling function through simple (non-rigourous) arguments. These conjectures are all verified by the exact calculation of sections IV A and IV B which provides the precise form of the scaling function \[eqs (99), (102) and (103)\] and yields all the exponents \[eqs (80), (81) and (84)\]. We find that beyond a crossover depth the variance (normalized by the squared mean) saturates. The saturation value and the crossover depth both diverge as the critical point is approached. At depths less than the crossover depth the variance grows as it would right at the critical point (eq 106). The behaviour at the critical point has many crossovers if the weight of the beads is small compared to the applied load. In this case at first the variance grows as the square root of the depth as it was found to do in section II in the absence of injection. At greater depths there are crossovers to growth as $`t`$ and $`t^{5/2}`$, as first the effects of large rare fluctuations in the weight of a bead and then mean injection assert themselves. Ultimately at the critical point the variance grows with the 5/2 exponent but the depth at which this behaviour sets in can be very great if the mean injection is small. This depth diverges as $`I^{4/3}`$. The crossover exponent 4/3, deduced by simple arguments and then via exact calculation in section IV, agrees with the value previously obtained by a different method by Majumdar and Sire . In their work Majumdar and Sire only study the behaviour right at the critical point. However at this point they calculate the dynamics of the entire distribution of load whereas we study only the variance. In section V we turn to the $`q`$-model in D+1 dimensions. For simplicity we neglect injection in this section. We find that right at the critical point the variance grows as a power of depth in all dimensions except two (eq 137). The power is given by $`D/2`$ for $`D<2`$. For all dimensions above two the growth is linear. This shows that $`D=2`$ is the upper critical dimension for this problem. For $`D=2`$ we find a linear growth of the variance tempered by a log factor as might be expected at the critical dimension. An intriguing feature of the critical behaviour we obtain is that it is exhibited at all. For ordinary continuous phase transitions the renormalization group provides a framework to understand the critical behaviour. We are not aware of any such framework for the $`q`$-model. Random critical points are notoriously difficult to analyse in general. The feature that allows us to analyse the $`q`$-model is that the two point load correlation function \[defined by eqs (8) and (114)\] evolves with depth according to a simple linear equation. In section II we analyse the evolution by expanding in the eigenvectors of an appropriate linear operator. There are some subtleties posed by the non-Hermiticity of the linear operator, making it necessary to prove that its eigenvectors are complete (further complicated by the infinite dimensionality of the vector space). Nonetheless we like this approach because it parallels transfer matrix methods used for equilibrium critical phenomena. We find that the large depth scaling behaviour is controlled by the low energy long wavelength eigenfunctions of the non-Hermitian “Hamiltonian”. Another virtue of this approach is that with about the same effort it yields both the variance and the correlation functions. However we have left analysis of the correlation functions open for later work. Here we focus entirely on the variance of the load. In section V we analyse the variance using another technique based on transform methods. Our analysis, neglecting injection, confirms that the q-model has essentially no horizontal correlations in the steady state for any distribution except the singular. This agrees with experiments on bead packs. The bulk of our results however are concerned with the q-model close to the critical point. Bead pack experiments such as those of ref appear to be far from the critical point. We estimate $`\delta 0.5`$ for this experiment. It is not obvious how to tune the parameter for bead packs to access the critical behaviour we analyse here. Claudin et al. have also studied the horizontal steady state correlations of the $`q`$-model without injection away from the critical point . They employ a continuum limit and arrive at conclusions similar to eq (19) in section IIB. The main focus of their work however is to explore a tensor model of stress propagation in granular matter, intended to supplant the $`q`$-model. Interpreted in terms of river networks our results show that allowing a small amount of river splitting in a Scheidegger network introduces a length scale in the vertical direction. On sufficiently long length scales such a network is not scale invariant. This resembles the finding of Narayan and Fisher . In their model too there was a parameter that controlled river splitting. Their networks were not scale invariant unless river splitting was tuned to zero. However their model appears to be in a different universality class as its vertical correlation exponent is different from the value $`\phi =2`$ we obtain here in section II. Presumably the difference is because their rule for stream splitting was non-local and depended on the entire history of the network upstream from the split. Taken together with the model of Narayan and Fisher it appears that river splitting is a perturbation that spoils the scale invariance of river networks. It is therefore interesting to ask whether such networks exist in Nature. River deltas are one possibility. Traced backwards they may constitute networks of merging streams that occasionally split. Even for river basins it might be interesting to examine the extent to which streams split. In this context it is worth noting that some of the data against which river scaling laws are tested are based not on actual maps of the river network but on networks that are indirectly inferred according to certain rules from digital elevation data obtained from satellite images. The rules by which the network is inferred from the elevation maps exclude the possibility of splitting . In summary the $`q`$-model is rich in applications and behaviour and yet analytically tractable by elementary means a combination of circumstances that invites further exploration. Among the many problems that remain open we conclude by mentioning two: For the $`q`$-model beyond the saturation depth there is no correlation in the horizontal direction but in the vertical direction there are very strong and long ranged correlations . We have not obtained the precise form of these vertical correlations for the $`q`$-model either in steady state or at the critical point. It would be very interesting (and straightforward) to obtain these forms and the crossover between them. Second it would be interesting to obtain the dynamics of the entire distribution of load near the critical point. We have not attempted to do this except right at the critical point. A natural scaling hypothesis is that the full distribution of load, neglecting injection, will be of the form $$\mathrm{\Pi }(w,t,\delta )=w^\tau 𝒬(w\delta ^{1/\sigma },t\delta ^{\nu z}).$$ (283) The exponents in eq (174) are $`\tau =2,\nu z=2`$ and $`\sigma =1`$. Their values are fixed by our result for the variance away from the critical point derived in section II and the result for the entire distribution at the critical point $`\delta =0`$ derived in section III. We also know that for $`x0`$ and $`y0`$, the presently unknown function $`𝒬`$ has the asymptotic behaviour $$𝒬(x,y)\frac{4}{\sqrt{\pi }}\frac{x}{y^{3/2}}e^{x^2/y}$$ (284) to be consistent with the critical point distribution (eq 73) derived in section III. In the second part of this paper we turn to chiral wave models that are believed to describe the surface electronic states of a quantum Hall multilayer. In section VI A we discuss circumstances under which the quantum network model of Saul, Kardar and Read (described in the introduction) is equivalent to the $`q`$-model. In section VI B we compare known results about the behaviour of the $`q`$-model to a continuum chiral wave model that cannot be mapped onto a $`q`$-model under any circumstance. The two are found to behave in qualitatively similar ways. A circumstance under which the mapping to the $`q`$-model is not useful is when periodic boundary conditions must be imposed in the chiral direction. Physically this is because of the interference of electron paths that wind around the quantum Hall multilayer. Such long range interference cannot be captured by the classical $`q`$-model. In this phase-coherent or mesoscopic regime, the chiral wave model has been studied via a mapping to a supersymmetric spin model . In section VI C we derive this mapping in a way that emphasizes boundary conditions. Our derivation makes use of operator methods and is hence more elementary than the derivation of ref that makes use of mixed functional integrals over Grassman and bosonic variables. We do not attempt further analysis of the superspin model here; the interested reader should consult papers on multilayer transport, particularly refs and that provide a nice overview of the early work on this problem. Mappings to superspin models have been useful not only in the study of the quantum Hall multilayer but have also recently lead to new non-perturbative results and insights into other important problems of electron localization . Hence it is hoped that our derivation, with its emphasis on boundary conditions and use of elementary operator methods, will be of interest in this general context. Note added: While writing this paper we learnt of an e-print by Rajesh and Majumdar on spatio-temporal correlations in the Takayasu model and the $`q`$-model . These authors derive many interesting results complementary to ours. In this paper we concentrate on the behaviour close to the critical point. For the $`q`$-model Rajesh and Majumdar concentrate on length scales long compared to our vertical correlation length, $`\xi _{\mathrm{corr}}`$; the crossovers and scaling functions that we study are transients that are invisible in their asymptotic formulae. On the other hand they have derived both vertical and horizontal load correlation functions; this paper is limited (in practice but not in principle) to the study of the variance of load. Among their interesting findings are (i) They find power law correlations in the vertical direction both at the critical point and away from it addressing in part a question raised above. (ii) They emphasize the interesting structure of the horizontal correlation function, including injection, at great depth. Although the goals are a bit different, there are points of intersection between the two papers with regard to technique. Rajesh and Majumdar too exploit the linearity of the relation that describes the evolution of the correlations with depth and solve it using the method of section V. An overlapping result is a formula for the variance at the critical point in 1+1 dimensions including injection. At the large depths studied by Rajesh and Majumdar the last term in our eq (106) should dominate. Rajesh and Majumdar obtain the same exponent 5/2 and the same numerical prefactor $`16/15\sqrt{\pi }`$ providing a nice check on both calculations. It is a pleasure to acknowledge stimulating discussions with Sue Coppersmith and Onuttom Narayan. We thank Onuttom Narayan in particular for patient explanation of refs , and , for encouraging us to study the critical point and for explaining to us the significance of river splitting. This work was supported in part by NSF Grant DMR 98-04983 and by the Alfred P. Sloan Foundation. HM acknowledges the hospitality of the Aspen Center for Physics where this work was completed. ## A Proof of Completeness First let us recall the principles of biorthogonal expansion (see, for example, ref , p884). We discuss the simplest case of a finite $`N\times N`$ dimensional non-Hermitian matrix $`H_{mn}`$. Consider its eigenvectors $$\underset{n}{}H_{mn}\varphi _n^\lambda =\lambda \varphi _m^\lambda .$$ (A1) In context of biorthogonal expansion these eigenvectors are called the right eigenvectors. For simplicity we will assume that the right eigenvalues are non-degenerate in this case. The bad news regarding the right eigenvectors is: (i) $`\lambda `$ may be complex. (ii) There is no guarantee that there are $`N`$ eigenvectors (needed to span the vector space). (iii) Eigenvectors corresponding to different eigenvalues are not necessarily orthogonal. Now consider the left eigenvectors, defined as the eigenvectors of $`H^{}`$. (i) If $`\lambda `$ is a right eigenvalue then $`\lambda ^{}`$ is a left eigenvalue (Proof: The coefficients for the characteristic polynomials of $`H`$ and $`H^{}`$ are complex conjugates of one another). (ii) There are as many left eigenvectors as right. (iii) Left eigenvectors are orthogonal to right eigenvectors. The last point merits elaboration. Let $`\psi _n^\lambda `$ denote the left eigenvector with left eigenvalue $`\lambda ^{}`$. Thus $$\underset{n}{}H_{mn}^{}\psi _n^\lambda =\lambda ^{}\psi _m^\lambda .$$ (A2) According to (iii) above $$\underset{n}{}(\psi _n^\lambda )^{}\varphi _n^\lambda ^{}=\delta _{\lambda \lambda ^{}}.$$ (A3) Eq (A3) is the biorthogonality relation. It may be proved by noting $`{\displaystyle \underset{mn}{}}(\psi _n^\lambda )^{}H_{nm}\varphi _m^\lambda ^{}`$ $`=`$ $`\lambda {\displaystyle \underset{n}{}}(\psi _n^\lambda )^{}\varphi _n^\lambda ^{}`$ (A4) $`=`$ $`\lambda ^{}{\displaystyle \underset{n}{}}(\psi _n^\lambda )^{}\varphi _n^\lambda ^{}`$ (A5) whence $`_n(\psi _n^\lambda )^{}\varphi _n^\lambda ^{}=0`$ for $`\lambda \lambda ^{}`$. In general there is no guarantee of completeness, but in this case assume that $`N`$ eigenvectors have been found. Then we can prove the completeness relation $$\underset{\lambda }{}(\psi _m^\lambda )^{}\varphi _n^\lambda =\delta _{mn}.$$ (A6) The proof follows from the observation that if there are $`N`$ eigenvectors, any vector $`a_n`$ may be expanded as $$a_n=\underset{\lambda }{}a_\lambda \varphi _n^\lambda .$$ (A7) Completeness then follows from biorthogonality, eq (A3). The problem in section IIB presents some complications not present in the pedagogical discussion above. Among them are degeneracy, an infinite dimensional vector space and a continuous spectrum. Nonetheless the broad strategy is the same. In section IIB we found left and right eigenvectors and we conjectured biorthogonality and completeness relations. To justify the analysis of section IIB we must prove the completeness relation. That is the purpose of this appendix. Note that we cannot simply assume completeness is true—because the matrix $`H`$ is non-Hermitian there are no theorems to guarantee it. Nor can we prove completeness by counting eigenvectors as in the finite dimensional discussion above. The proof of completeness is remarkably simple and direct. We substitute the exact expressions for $`\psi _m^{(\pm )k}`$ and $`\varphi _n^{(\pm )k}`$ that we have derived, eqs (28), (31), (33) and (34), on the right hand side of eq (36) and verify the completeness relation by explicit evaluation of the integral. There are nine cases to consider corresponding to $`n=0,n>0,n<0`$ and $`m=0,m<0,m>0`$. For illustration we analyse the case of $`n=0,m=0`$. We must evaluate $$\frac{2}{\pi }_0^\pi 𝑑k𝒜^{}(k)A(k)$$ (A8) where $`𝒜(k)`$ and $`A(k)`$ are as given in eqs (30) and (35). Since the integrand is symmetric in $`k`$ we extend the range of integration from $`\pi `$ to $`\pi `$ and substitute $`ze^{ik}`$ to obtain a contour integral about the unit circle $$\frac{dz}{2\pi i}\frac{1}{z}\frac{(1ϵ)(z+1)^2}{ϵ^2(z1)^2(1ϵ)^2(z+1)^2}.$$ (A9) Evaluation via Cauchy’s theorem reveals that the integral equals one as required for completeness. The remaining eight cases also succumb to this method of analysis. ## B Inverse z-transform Consider the series $`f(t)`$, $`t=0,1,2,\mathrm{}`$ Its $`z`$-transform is defined as $$f(z)=\underset{t=0}{\overset{\mathrm{}}{}}f(t)z^t.$$ (B1) Some $`z`$-transforms can be inverted by inspection. For example the inverse transform of $`(1\alpha z)^1`$ is evidently $$(1\alpha z)^1f(t)=\alpha ^t.$$ (B2) In other cases the inverse transform can be found by performing the complex integral $$f(t)=_C\frac{dz}{2\pi i}\frac{f(z)}{z^{t+1}}.$$ (B3) The contour $`C`$ must enclose the origin but no singularities of $`f(z)`$. For illustration let us analyse $$f(z)=(1z)^{1/2}(1\alpha z)^1$$ (B4) needed to go from eq (43) to (44) in section IIB. Here $`\alpha >1`$. $`f(z)`$ has a pole at $`1/\alpha `$ and a branch cut at $`z=1`$ (see fig 10). We deform the contour $`C`$ that encloses the origin to contours $`C_1`$ and $`C_2`$ that encircle the pole and pass above and below the branch cut. Hence obtain $$f(t)=\sqrt{\frac{\alpha }{\alpha 1}}\alpha ^t\frac{1}{\alpha \pi }_1^{\mathrm{}}𝑑x(x1)^{1/2}\frac{1}{x^{t+1}}\left(x\frac{1}{\alpha }\right)^1.$$ (B5) The first term is the contribution of the pole; the second, of the branch cut. ## C Asymptotics of $`\mathrm{\Phi }_n(u)`$ The asymptotics of the functions $`\mathrm{\Phi }_n(u)`$ defined by eq (103) are needed to obtain the asymptotic behaviour of the scaling functions in sections IIC and IVB. The large $`u`$ behaviour poses no difficulty. Evidently $$\mathrm{\Phi }_n(u)\frac{\sqrt{\pi }}{2}\frac{1}{\sqrt{u}}\mathrm{as}u\mathrm{}$$ (C1) for all $`n`$. The small $`u`$ behaviour is a bit more subtle. Moreover, it turns out that due to cancellations we will need as many as five or six terms in the small $`u`$ series for $`\mathrm{\Phi }_n`$ to obtain the leading behaviour of the scaling functions. For definiteness consider the small $`u`$ behaviour of $$\mathrm{\Phi }_1(u)=_0^{\mathrm{}}𝑑s\frac{e^{us^2}}{(1+s^2)}.$$ (C2) The leading term is obtained by setting $`u=0`$, $$\mathrm{\Phi }_1(0)=\frac{\pi }{2}.$$ (C3) To obtain the next term it is tempting to expand the integrand in powers of $`u`$ but this leads to divergent integrals. The divergence signals that the asymptotic series is not a simple power series in $`u`$. It turns out the next term goes as $`\sqrt{u}`$. To show this, and to efficiently obtain many more terms in the series, consider $$g(x)=_0^{\mathrm{}}𝑑s\frac{e^{x^2s^2}}{(1+s^2)}.$$ (C4) We will show that $`g(x)`$ is regular about $`x=0`$ and that its asymptotic behaviour is a simple power series. To this end we observe that $`g(x)`$ obeys the first order differential equation $$\frac{d}{dx}g2xg(x)+\sqrt{\pi }=0.$$ (C5) $`x=0`$ is a regular point for this equation; hence we attempt a series solution $$g(x)=b_0+b_1x+b_2x^2+\mathrm{}$$ (C6) We find $`b_1=\sqrt{\pi }`$ and the simple recurrence relation $$b_n=\frac{2}{n}b_{n2}.$$ (C7) Evidently $`b_0=g(0)=\pi /2`$. Hence we obtain the asymptotic series $$g(x)=\frac{\pi }{2}\sqrt{\pi }x+\frac{\pi }{2}x^2\frac{2}{3}\sqrt{\pi }x^3+\frac{\pi }{4}x^4\frac{4}{15}\sqrt{\pi }x^5+\mathrm{}$$ (C8) Substituting $`x\sqrt{u}`$ we conclude $`\mathrm{\Phi }_1(u)`$ $`=`$ $`{\displaystyle \frac{\pi }{2}}\sqrt{\pi }u^{1/2}+{\displaystyle \frac{\pi }{2}}u{\displaystyle \frac{2}{3}}\sqrt{\pi }u^{3/2}+{\displaystyle \frac{\pi }{4}}u^2`$ (C10) $`{\displaystyle \frac{4}{15}}\sqrt{\pi }u^{5/2}+\mathrm{}`$ for small $`u`$. Similarly $`\mathrm{\Phi }_2(u)`$ $`=`$ $`{\displaystyle \frac{\pi }{4}}{\displaystyle \frac{\pi }{4}}u+{\displaystyle \frac{2}{3}}\sqrt{\pi }u^{3/2}{\displaystyle \frac{3}{8}}\pi u^2+{\displaystyle \frac{8}{15}}\sqrt{\pi }u^{5/2}+\mathrm{}`$ (C11) $`\mathrm{\Phi }_3(u)`$ $`=`$ $`{\displaystyle \frac{3\pi }{16}}{\displaystyle \frac{\pi }{16}}u+{\displaystyle \frac{3\pi }{32}}u^2{\displaystyle \frac{4}{15}}\sqrt{\pi }u^{5/2}+\mathrm{}`$ (C12) ## D Lattice Green’s Function ### 1 Two Dimensions Consider the Green’s function in two dimensions for the lattice Schrödinger equation discussed in section VA. The real space Green’s function at the origin is given by $$𝒢(E)=_\pi ^\pi \frac{dk}{2\pi }_\pi ^\pi \frac{dp}{2\pi }\left\{E\frac{1}{4}(1+\mathrm{cos}p)(1+\mathrm{cos}k)\right\}^1$$ (D1) \[cf. eq (121) and (128)\]. We consider real $`E>1`$. In this appendix we show that $$𝒢(E)=\frac{2}{\pi E}K\left(\frac{1}{\sqrt{E}}\right);$$ (D2) Here $`K`$ is a complete elliptic integral of the first kind. From the well-documented properties of these integrals or by direct analysis of eq (D8) below it follows that as $`E1^+`$ $$𝒢(E)\frac{1}{\pi }\mathrm{ln}\frac{1}{E1}.$$ (D3) In section V C we are interested in the behaviour of $`G(z)`$, eq (131), as the real variable $`z1^{}`$. Comparing eq (131) to (D1) we see that $$G(z)=\frac{1}{z}𝒢\left(E\frac{1}{z}\right).$$ (D4) Hence the singularity of $`G(z)`$ as $`z1`$ is $$G(z)=\frac{1}{\pi }\mathrm{ln}(1z).$$ (D5) Eq (D5) is the main result of this section of the Appendix. To demonstrate eq (D2) we regard $`p`$ as a complex variable $`px+iy`$. The integral over $`p`$ in eq (D1) may be regarded as an integral around the contour sketched in Fig 11 since the two vertical segments cancel by the periodicity of the integrand and the horizontal segment at infinity makes no contribution because the integrand vanishes along it. The integrand in eq (D1) has a simple pole at $`p=iy`$, where $`y`$ satisfies $$\mathrm{cosh}\left(\frac{y}{2}\right)=\frac{\sqrt{E}}{\mathrm{cos}(k/2)},$$ (D6) with residue $$\left(iE\sqrt{1\frac{1}{E}\mathrm{cos}^2\frac{k}{2}}\right)^1.$$ (D7) Hence by Cauchy’s theorem $$𝒢(E)=\frac{1}{E}_\pi ^\pi \frac{dk}{2\pi }\left(1\frac{1}{E}\mathrm{cos}^2\frac{k}{2}\right)^{1/2}.$$ (D8) Comparing to the definition of the elliptic integral of the first kind $$K(k)=_0^{\pi /2}𝑑\theta (1k^2\mathrm{sin}^2\theta )^{1/2}$$ (D9) we obtain eq (D2). ### 2 Below Two Dimensions In this section we analyse the singular behaviour as $`z1`$ of $`G(z)`$ in less than two dimensions. The approximate long wavelength expression for $`G`$, eq (133), provides a useful starting point. To analyse the divergence in $`D=1`$ we would note that the integrand in (133) is a sharply peaked Lorentzian. This justifies working to quadratic order in $`S(\stackrel{}{p})`$ and extending the range of integration (strictly confined to the Brillouin zone, $`\pi <p<\pi `$ in one dimension) to $`\pm \mathrm{}`$. Result: $`G(z)=(1z)^{1/2}`$. To continue this result to non-integral $`D`$ we use ’t Hooft and Veltman’s dimensional regularization trick . We write $$G(z)_0^{\mathrm{}}𝑑s\frac{d\stackrel{}{p}}{(2\pi )^D}\mathrm{exp}s[(1z)+\stackrel{}{p}^2];$$ (D10) extend the range of integration, outside the Brillouin zone and over all $`\stackrel{}{p}`$-space; and replace $$\frac{d\stackrel{}{p}}{(2\pi )^D}\frac{\mathrm{\Omega }_D}{(2\pi )^D}_0^{\mathrm{}}𝑑pp^{D1},$$ (D11) since the integrand in eq (D10) is isotropic in $`\stackrel{}{p}`$. Here $`\mathrm{\Omega }_D=2\sqrt{\pi }^D/\mathrm{\Gamma }(D/2)`$ is the total solid angle in $`D`$ dimensions (some familiar special cases: $`\mathrm{\Omega }_1=2,\mathrm{\Omega }_2=2\pi ,\mathrm{\Omega }_3=4\pi ,\mathrm{\Omega }_4=2\pi ^2`$.) The result is $$G(z)=\frac{\mathrm{\Gamma }(1D/2)}{\sqrt{\pi }^D}(1z)^{D/21}$$ (D12) for $`D<2`$. This analysis breaks down in two dimensions and higher because the integrand diverges as $`\stackrel{}{p}\mathrm{}`$. The divergence is an artifact of the quadratic approximation in eq (133) and of extending the integral outside the Brillouin zone. The spurious divergence is revealed in eq (D12) as a pole in the Gamma function factor as $`D2`$. ## E Analysis of Partition Function The purpose of this appendix is to show that the partition functions for bosons and fermions cancel. Thus $$\mathrm{Tr}[S_F^R(T)]\mathrm{Tr}[S_B^R(T)]=1.$$ (E1) A similar relation holds for the advanced bosons and fermions. We discuss the retarded case explicitly. For brevity the superscript $`R`$ will be omitted. We write the fermion S-matrix as $$S_F(t)=\mathrm{exp}\left(i\frac{\pi t}{T}\underset{n}{}c_n^{}c_n\right)𝒮_F(t).$$ (E2) $`𝒮_F(t)`$ is then governed by the Hamiltonian eq (146) without the extra term included in eq (153). To make further progress we introduce $`e_l^n(t)`$, the solution to the Schrödinger eq (144) $$i\frac{}{t}e_l^n(t)=m_l(t)e_{l+1}^n(t)+m_{l1}^{}(t)e_{l1}^n(t)$$ (E3) subject to $`e_l^n(t0)=\delta _{nl}`$. The scattering formula $$c_n𝒮_F(t)=\underset{l}{}e_n^l(t)𝒮_F(t)c_l$$ (E4) will prove very useful. To derive it, rewrite eq (E4) as $$𝒮_F(t)^1c_n𝒮_F(t)=\underset{l}{}e_n^l(t)c_l$$ (E5) and regard it as an ansatz with the functions $`e_n^l(t)`$ unspecified. Making use of eqs (147), (148) and (151), the $`t`$ derivative of the left hand side is $`𝒮_F(t)^1\{m_n(t)c_{n+1}+m_{n1}^{}(t)c_{n1}\}𝒮_F(t)`$ (E6) $`={\displaystyle \underset{l}{}}c_l\{m_l(t)e_{n+1}^l(t)+m_{l1}^{}(t)e_{n1}^l(t)\}.`$ (E7) To obtain the second line we have made use of the ansatz (E5). Comparing eq (E6) to the $`t`$ derivative of the right hand side of eq (E5) we conclude that $`e_n^l(t)`$ does obey the Schrödinger eq (E3). This completes the proof of the scattering formula (E4). Another relation that will prove useful is $$𝒮_F(t)\|0=\|0.$$ (E8) This follows because the Hamiltonian (eq 146) annhilates the vacuum; $`𝒮_F`$ is the (chronologically ordered) exponential of the Hamiltonian. Equipped with these results we write the fermion partition function as $`Z_F(T)`$ $`=`$ $`\mathrm{Tr}\{\mathrm{exp}[i\pi {\displaystyle \underset{n}{}}c_n^{}c_n]𝒮_F(T)\}`$ (E9) $`=`$ $`0\|𝒮_F(T)\|0`$ (E14) $`{\displaystyle \underset{n}{}}0\|c_n𝒮_F(T)c_n^{}\|0`$ $`+{\displaystyle \frac{1}{2}}{\displaystyle \underset{n_1,n_2}{}}0\|c_{n_1}c_{n_2}𝒮_F(T)c_{n_2}^{}c_{n_1}^{}\|0`$ $`{\displaystyle \frac{1}{3!}}{\displaystyle \underset{n_1,n_2,n_3}{}}0\|c_{n_1}c_{n_2}c_{n_3}𝒮_F(T)c_{n_3}^{}c_{n_2}^{}c_{n_1}\|0`$ $`+\mathrm{}`$ The trace is taken over the entire Fock space including states with different total numbers of fermions. The alternating signs are due to the factor $`\mathrm{exp}[i\pi _nc_n^{}c_n]`$ in the trace. The factorials are because the sums over the site indices $`n_i`$ are unrestricted; hence each state gets counted a multiple number of times. We now shift the S-matrix to the left using the scattering formula (E4), make use of the adjoint of (E7) and calculate the vacuum expectations of the fermion operators (Wick’s theorem). The result for the second-order term is $`Z_F(T)`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left[{\displaystyle \underset{n}{}}e_n^n(T)\right]^2{\displaystyle \frac{1}{2}}{\displaystyle \underset{n_1,n_2}{}}e_{n_1}^{n_2}(T)e_{n_2}^{n_1}(T)`$ (E16) $`+\mathrm{others}.`$ Fig 12(b) shows a diagrammatic representation of this term. Note that the diagram series for the partition function $`Z_F(T)`$ contains both connected and unconnected graphs. By familiar arguments we can write $$Z_F(T)=\mathrm{exp}[\mathrm{\Omega }(T)]$$ (E17) where the “free energy” $`\mathrm{\Omega }(T)`$ has the linked diagram expansion shown in fig 12(c). We turn now to the boson partition function. The boson scattering formula $$b_nS_B(t)=\underset{l}{}e_n^l(t)S_B(t)b_l$$ (E18) can be proved in the same way as eq (E4). Eq (E7) remains true when we replace $`𝒮_FS_B`$ The boson partition function is therefore given by $`Z_B(T)`$ $`=`$ $`\mathrm{Tr}\{S_B(T)\}`$ (E19) $`=`$ $`0\|S_B(T)\|0`$ (E23) $`+{\displaystyle \underset{n}{}}0\|b_nS_B(T)b_n^{}\|0`$ $`+{\displaystyle \frac{1}{2!}}{\displaystyle \underset{n_1,n_2}{}}0\|b_{n_1}b_{n_2}S_B(T)b_{n_2}^{}b_{n_1}^{}\|0`$ $`+\mathrm{}`$ This equation resembles eq (E8) but there is an extra subtlety in the combinatoric factors. In the two-boson case for example, for the offdiagonal terms $`(n_1n_2)`$ the factor $`(1/2!)`$ is to offset double counting as in the fermion case. For the diagonal terms, that vanish in the fermion case, there is no double counting but the factor $`(1/2!)`$ is needed for normalization. By shifting $`S_B(T)`$ to the left by use of the scattering formula we see that the series for $`Z_B(T)`$ is the same as for $`Z_F(T)`$ except for the minus signs. Hence $$Z_B(T)=\mathrm{exp}\left(+\mathrm{\Omega }(T)\right)$$ (E24) where $`\mathrm{\Omega }(T)`$ is defined by the diagram series in fig 12(c). Eq (E10) and E(13) together lead to (E1), the result we sought to prove here.
warning/0007/math0007046.html	ar5iv	text	# A simple proof of Bailey’s very-well-poised {_𝟔}𝝍_𝟔 summation ## 1. Introduction The theories of unilateral (or one-sided) hypergeometric and basic ($`q`$-)hypergeometric series have quite a rich history dating back to at least Euler. Formulae for bilateral (basic) hypergeometric series were not discovered until 1907 when Dougall , using residue calculus, derived summations for the bilateral $`{}_{2}{}^{}H_{2}^{}`$ and very-well-poised $`{}_{5}{}^{}H_{5}^{}`$ series. Ramanujan extended the $`q`$-binomial theorem by finding a summation formula for the bilateral $`{}_{1}{}^{}\psi _{1}^{}`$ series. Later, Bailey , carried out systematical investigations of summations and transformations for bilateral basic hypergeometric series. Further significant contributions were made by Slater ,, a student of Bailey. See and for an excellent survey of the above classical material. Bailey’s \[6, Eq. (4.7)\] very-well-poised $`{}_{6}{}^{}\psi _{6}^{}`$ summation (cf. \[11, Eq. (5.3.1)\]) is a very powerful identity, as it stands at the top of the classical hierarchy of summation formulae for bilateral series. Some of the applications of the $`{}_{6}{}^{}\psi _{6}^{}`$ summation to partitions and number theory are given in Andrews . Though several proofs of Bailey’s $`{}_{6}{}^{}\psi _{6}^{}`$ summation are already known (see, e. g., Bailey , Slater and Lakin , Andrews , Askey and Ismail , and Askey ), none of them is entirely elementary. Here we provide a new simple proof of the very-well-poised $`{}_{6}{}^{}\psi _{6}^{}`$ summation formula, directly from three applications of Rogers’ \[22, p. 29, second eq.\] nonterminating $`{}_{6}{}^{}\varphi _{5}^{}`$ summation (cf. \[11, Eq. (2.7.1)\]) and elementary manipulations of series. The method of proof we apply extends that already used by M. Jackson \[19, Sec. 4\] in her first elementary proof (as pointed out to us by George Andrews ) of Ramanujan’s $`{}_{1}{}^{}\psi _{1}^{}`$ summation formula (cf. \[11, Eq. (5.2.1)\]). Jackson’s proof essentially derives the $`{}_{1}{}^{}\psi _{1}^{}`$ summation from the $`q`$-Gauß summation, by manipulation of series. In view of this background, it is surprising that this method has not been further applied for half a century. A possible explanation is that the applicability of her method was viewed as too limited. In fact, only after changing the order of steps in Jackson’s proof, we were able to extend her proof to a “method”. Indeed, the method can also be applied to derive other summations. After recalling some notation for (basic) hypergeometric series in Section 2, we review Jackson’s elementary proof of the $`{}_{1}{}^{}\psi _{1}^{}`$ summation in Section 3. In Section 4, we apply our extension of Jackson’s method to give an elementary proof of Dougall’s $`{}_{2}{}^{}H_{2}^{}`$ summation. Finally, in Section 5, we give an elementary derivation of Bailey’s very-well-poised $`{}_{6}{}^{}\psi _{6}^{}`$ summation. We want to point out that by using a similar but slightly different method, the author has found elementary derivations of transformations for bilateral basic hypergeometric series. In fact, in we use Bailey’s nonterminating very-well-poised $`{}_{8}{}^{}\varphi _{7}^{}`$ summation theorem combined with bilateral series identities to derive a very-well-poised $`{}_{8}{}^{}\psi _{8}^{}`$ transformation, a very-well-poised $`{}_{10}{}^{}\psi _{10}^{}`$ transformation, and by induction, Slater’s general transformation for very-well-poised $`{}_{2r}{}^{}\psi _{2r}^{}`$ series. Similarly, some other bilateral series identities are also elementarily derived in . In the near future, we plan to apply the methods of this article and of to the settings of multiple basic hypergeometric series. See Milne , Gustafson , v. Diejen , and Schlosser , for several of these different settings. We are quite confident that we may not only get simpler proofs for already known results but should also obtain derivations of new formulae. Finally, we wish to gratefully acknowledge the helpful comments and suggestions of George Andrews, Mourad Ismail, and Stephen Milne. ## 2. Background and notation Here we recall some notation for hypergeometric series (cf. ), and basic hypergeometric series (cf. ). We define the shifted factorial for all integers $`k`$ by the following quotient of Gamma functions (cf. \[3, Sec. 1.1\]), $$(a)_k:=\frac{\mathrm{\Gamma }(a+k)}{\mathrm{\Gamma }(a)}.$$ Further, the (ordinary) hypergeometric $`{}_{r}{}^{}F_{s}^{}`$ series is defined as $${}_{r}{}^{}F_{s}^{}[\begin{array}{c}a_1,a_2,\mathrm{},a_r\\ b_1,b_2,\mathrm{},b_s\end{array};z]:=\underset{k=0}{\overset{\mathrm{}}{}}\frac{(a_1)_k\mathrm{}(a_r)_k}{(b_1)_k\mathrm{}(b_s)_k}\frac{z^k}{k!},$$ (2.1) and the bilateral hypergeometric $`{}_{r}{}^{}H_{s}^{}`$ series as $${}_{r}{}^{}H_{s}^{}[\begin{array}{c}a_1,a_2,\mathrm{},a_r\\ b_1,b_2,\mathrm{},b_s\end{array};z]:=\underset{k=\mathrm{}}{\overset{\mathrm{}}{}}\frac{(a_1)_k\mathrm{}(a_r)_k}{(b_1)_k\mathrm{}(b_s)_k}z^k.$$ (2.2) See \[26, p. 45 and p. 181\] for the criteria of when these series terminate, or, if not, when they converge. Let $`q`$ be a complex number such that $`0<\|q\|<1`$. We define the $`q`$-shifted factorial for all integers $`k`$ by $$(a;q)_{\mathrm{}}:=\underset{j=0}{\overset{\mathrm{}}{}}(1aq^j)\text{and}(a;q)_k:=\frac{(a;q)_{\mathrm{}}}{(aq^k;q)_{\mathrm{}}}.$$ For brevity, we employ the usual notation $$(a_1,\mathrm{},a_m;q)_k(a_1;q)_k\mathrm{}(a_m;q)_k$$ where $`k`$ is an integer or infinity. Further, we utilize the notations $$_r\varphi _s[\begin{array}{c}a_1,a_2,\mathrm{},a_r\\ b_1,b_2,\mathrm{},b_s\end{array};q,z]:=\underset{k=0}{\overset{\mathrm{}}{}}\frac{(a_1,a_2,\mathrm{},a_r;q)_k}{(q,b_1,\mathrm{},b_s;q)_k}\left((1)^kq^{\left(\genfrac{}{}{0pt}{}{k}{2}\right)}\right)^{1+sr}z^k,$$ (2.3) and $$_r\psi _s[\begin{array}{c}a_1,a_2,\mathrm{},a_r\\ b_1,b_2,\mathrm{},b_s\end{array};q,z]:=\underset{k=\mathrm{}}{\overset{\mathrm{}}{}}\frac{(a_1,a_2,\mathrm{},a_r;q)_k}{(b_1,b_2,\mathrm{},b_s;q)_k}\left((1)^kq^{\left(\genfrac{}{}{0pt}{}{k}{2}\right)}\right)^{sr}z^k,$$ (2.4) for basic hypergeometric $`{}_{r}{}^{}\varphi _{s}^{}`$ series, and bilateral basic hypergeometric $`{}_{r}{}^{}\psi _{s}^{}`$ series, respectively. See \[11, p. 25 and p. 125\] for the criteria of when these series terminate, or, if not, when they converge. We want to point out that many theorems for $`{}_{r}{}^{}F_{s}^{}`$ or $`{}_{r}{}^{}H_{s}^{}`$ series can be obtained by considering certain “$`q1`$ limiting cases” of corresponding theorems for $`{}_{r}{}^{}\varphi _{s}^{}`$ or $`{}_{r}{}^{}\psi _{s}^{}`$ series, respectively. For instance, we describe such a $`q1`$ limiting case after stating the $`q`$-binomial theorem in (3.1). A similar $`q1`$ limiting case leads from the $`q`$-Gauß summation (3.3) to the ordinary Gauß summation (4.7). The situation is different for the $`{}_{1}{}^{}\psi _{1}^{}`$ series, though. We have a summation for the general $`{}_{1}{}^{}\psi _{1}^{}`$, but not for the $`{}_{1}{}^{}H_{1}^{}`$. On the other hand, the general $`{}_{2}{}^{}H_{2}^{}`$ with unit argument is summable but the general $`{}_{2}{}^{}\psi _{2}^{}`$ is not. Many theorems for very-well-poised $`{}_{r+1}{}^{}\varphi _{r}^{}`$ series can be specialized to theorems for very-well-poised $`{}_{r}{}^{}F_{r1}^{}`$ series. For the notion of (very-)well-poised, see \[11, Sec. 2.1\]. For detailed treatises on hypergeometric and basic hypergeometric series, we refer to Slater , and Gasper and Rahman . In our computations in the following sections, we make heavily use of some elementary identities involving ($`q`$-)shifted factorials which are listed in Slater \[26, Appendix I\], and Gasper and Rahman \[11, Appendix I\]. ## 3. M. Jackson’s proof of Ramanujan’s $`{}_{1}{}^{}\psi _{1}^{}`$ summation The $`q`$-binomial theorem, $${}_{1}{}^{}\varphi _{0}^{}[\begin{array}{c}a\\ \end{array};q,z]=\frac{(az;q)_{\mathrm{}}}{(z;q)_{\mathrm{}}},$$ (3.1) where the series either terminates, or $`\|z\|<1`$, for convergence, was first discovered by Cauchy (cf. \[11, Sec. 1.3\]). It reduces to the ordinary binomial theorem as $`aq^a`$ and $`q1^{}`$. A bilateral extension of the $`q`$-binomial theorem (3.1), the $`{}_{1}{}^{}\psi _{1}^{}`$ summation, was found by the legendary Indian mathematician Ramanujan (cf. \[11, Eq. (5.2.1)\]). It reads as follows: $${}_{1}{}^{}\psi _{1}^{}[\begin{array}{c}a\\ b\end{array};q,z]=\frac{(q,b/a,az,q/az;q)_{\mathrm{}}}{(b,q/a,z,b/az;q)_{\mathrm{}}},$$ (3.2) where the series either terminates, or $`\|b/a\|<\|z\|<1`$, for convergence. Clearly, (3.2) reduces to (3.1) when $`b=q`$. Unfortunately, Ramanujan did not provide a proof for his summation formula. Hahn \[14, $`\kappa =0`$ in Eq. (4.7)\] independently established (3.2) by considering a first order homogeneous $`q`$-difference equation. Hahn thus published the first proof of the $`{}_{1}{}^{}\psi _{1}^{}`$ summation. Not much later, M. Jackson \[19, Sec. 4\] gave the first elementary proof of (3.2). Her proof derives the $`{}_{1}{}^{}\psi _{1}^{}`$ summation from the $`q`$-Gauß summation, by manipulation of series. It turns out that Jackson’s method is effective for proving also other bilateral summation formulae. Since Jackson’s short proof of Ramanujan’s $`{}_{1}{}^{}\psi _{1}^{}`$ summation seems to be not so well known, we review her proof in the following. Before we continue, we want to point out that there are also many other nice proofs of the $`{}_{1}{}^{}\psi _{1}^{}`$ summation in the literature. A simple and elegant proof of the $`{}_{1}{}^{}\psi _{1}^{}`$ summation formula was given by Ismail who showed that the $`{}_{1}{}^{}\psi _{1}^{}`$ summation is an immediate consequence of the $`q`$-binomial theorem and analytic continuation. M. Jackson’s elementary proof of (3.2) makes use of a suitable specialization of Heine’s $`q`$-Gauß summation (cf. \[11, Eq. (II.8)\]), $${}_{2}{}^{}\varphi _{1}^{}[\begin{array}{c}a,b\\ c\end{array};q,\frac{c}{ab}]=\frac{(c/a,c/b;q)_{\mathrm{}}}{(c,c/ab;q)_{\mathrm{}}},$$ (3.3) where the series either terminates, or $`\|c/ab\|<1`$, for convergence. In (3.3), we perform the substitutions $`aaq^n`$, $`bq/b`$, and $`cq^{1+n}`$, and obtain $${}_{2}{}^{}\varphi _{1}^{}[\begin{array}{c}aq^n,q/b\\ q^{1+n}\end{array};q,\frac{b}{a}]=\frac{(q/a,bq^n;q)_{\mathrm{}}}{(q^{1+n},b/a;q)_{\mathrm{}}},$$ (3.4) provided $`\|b/a\|<1`$. Using some elementary identities for $`q`$-shifted factorials (see, e. g., Gasper and Rahman \[11, Appendix I\]) we can rewrite equation (3.4) as $$\frac{(q,b/a;q)_{\mathrm{}}}{(q/a,b;q)_{\mathrm{}}}\underset{k=0}{\overset{\mathrm{}}{}}\frac{(q/b;q)_k(a;q)_{n+k}}{(q;q)_k(q;q)_{n+k}}\left(\frac{b}{a}\right)^k=\frac{(a;q)_n}{(b;q)_n}.$$ (3.5) In this identity, we multiply both sides by $`z^n`$ and sum over all integers $`n`$. On the right side we obtain $${}_{1}{}^{}\psi _{1}^{}[\begin{array}{c}a\\ b\end{array};q,z].$$ On the left side we obtain $$\frac{(q,b/a;q)_{\mathrm{}}}{(q/a,b;q)_{\mathrm{}}}\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}z^n\underset{k=0}{\overset{\mathrm{}}{}}\frac{(q/b;q)_k(a;q)_{n+k}}{(q;q)_k(q;q)_{n+k}}\left(\frac{b}{a}\right)^k.$$ (3.6) Next, we interchange summations in (3.6) and shift the inner index $`nnk`$. (Observe that the sum over $`n`$ is terminated by the term $`(q;q)_{n+k}^1`$ from below.) We obtain $$\frac{(q,b/a;q)_{\mathrm{}}}{(q/a,b;q)_{\mathrm{}}}\underset{k=0}{\overset{\mathrm{}}{}}\frac{(q/b;q)_k}{(q;q)_k}\left(\frac{b}{az}\right)^k\underset{n=0}{\overset{\mathrm{}}{}}\frac{(a;q)_n}{(q;q)_n}z^n.$$ Now, twice application of the $`q`$-binomial theorem (3.1) gives us the right side of (3.2), as desired. Now, we have to admit that M. Jackson did not give her proof in the above precise order. In fact, her proof in \[19, Sec. 4\] goes backwards. (This is also how the author originally rediscovered Jackson’s proof.) She started with the $`{}_{1}{}^{}\psi _{1}^{}`$ summation (3.2) and equated coefficients of $`z^n`$ on both sides. The resulting identity is true by the $`q`$-Gauß summation. A reason why M. Jackson’s method of proof has so far not been used to prove other bilateral summations could be that the applicability of her derivation was viewed as too limited. Equating coefficients of a power of a Laurent series variable in a bilateral basic hypergeometric series identity is easy if, as in (3.2), there is an argument $`z`$ which is independent of the other parameters. But this seems to be more particluar to the $`{}_{1}{}^{}\psi _{1}^{}`$ summation, as not in all bilateral series there is such an independent argument. The starting point for making M. Jackson’s proof to a “method” is to read the proof backwards, as displayed above. The essence here is that a unilateral series identity, (3.3), is specialized such that there is the factor $`(q;q)_{n+k}^1`$ in the series, see (3.5), so that summing over all $`n`$ again gives a (summable) unilateral series. In the next two sections, we use the method to give proofs of two other important bilateral hypergeometric and basic hypergeometric summation theorems. In particular, in Section 4, we give a simple proof of Dougall’s $`{}_{2}{}^{}H_{2}^{}`$ summation, whereas in Section 5, we give a simple proof of Bailey’s very-well-poised $`{}_{6}{}^{}\psi _{6}^{}`$ summation. ## 4. Dougall’s $`{}_{2}{}^{}H_{2}^{}`$ summation In Section 3, we multiplied both sides of the identity (3.5) by a suitable factor depending on $`n`$ and summed over all integers $`n`$. On one side, we interchanged sums and found that the inner sum was summable by the $`q`$-binomial theorem. Now, what if we start with a different factor following a similar procedure such that we can evaluate the inner sum by, say, the $`q`$-Gauß summation? If the analysis works out we may end up with an evaluation for a $`{}_{2}{}^{}\psi _{2}^{}`$ series. Let us see what happens: In identity (3.5), let us first replace $`b`$ by $`c`$. Then we multiply both sides by $$\frac{(b;q)_n}{(d;q)_n}\left(\frac{d}{ab}\right)^n$$ and sum over all integers $`n`$. On the right side we obtain $${}_{2}{}^{}\psi _{2}^{}[\begin{array}{c}a,b\\ c,d\end{array};q,\frac{d}{ab}].$$ (4.1) On the left side we obtain $$\frac{(q,c/a;q)_{\mathrm{}}}{(q/a,c;q)_{\mathrm{}}}\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}\frac{(b;q)_n}{(d;q)_n}\left(\frac{d}{ab}\right)^n\underset{k=0}{\overset{\mathrm{}}{}}\frac{(q/c;q)_k(a;q)_{n+k}}{(q;q)_k(q;q)_{n+k}}\left(\frac{c}{a}\right)^k.$$ (4.2) Next, we interchange summations in (4.2) and shift the inner index $`nnk`$. We obtain, again using some elementary identities for $`q`$-shifted factorials, $$\frac{(q,c/a;q)_{\mathrm{}}}{(q/a,c;q)_{\mathrm{}}}\underset{k=0}{\overset{\mathrm{}}{}}\frac{(q/c;q)_k(b;q)_k}{(q;q)_k(d;q)_k}\left(\frac{bc}{d}\right)^k\underset{n=0}{\overset{\mathrm{}}{}}\frac{(a,bq^k;q)_n}{(q,dq^k;q)_n}\left(\frac{d}{ab}\right)^n.$$ Now the inner sum, provided $`\|d/ab\|<1`$, can be evaluated by (3.3) and we obtain $$\frac{(q,c/a;q)_{\mathrm{}}}{(q/a,c;q)_{\mathrm{}}}\underset{k=0}{\overset{\mathrm{}}{}}\frac{(q/c;q)_k(b;q)_k}{(q;q)_k(d;q)_k}\left(\frac{bc}{d}\right)^k\frac{(dq^k/a,d/b;q)_{\mathrm{}}}{(dq^k,d/ab;q)_{\mathrm{}}},$$ which can be simplied to $$\frac{(q,c/a,d/a,d/b;q)_{\mathrm{}}}{(q/a,c,d,d/ab;q)_{\mathrm{}}}\underset{k=0}{\overset{\mathrm{}}{}}\frac{(q/c,aq/d;q)_k}{(q,q/b;q)_k}\left(\frac{c}{a}\right)^k.$$ (4.3) Hence, equating (4.1) and (4.3), we have derived the transformation $${}_{2}{}^{}\psi _{2}^{}[\begin{array}{c}a,b\\ c,d\end{array};q,\frac{d}{ab}]=\frac{(q,c/a,d/a,d/b;q)_{\mathrm{}}}{(q/a,c,d,d/ab;q)_{\mathrm{}}}{}_{2}{}^{}\varphi _{1}^{}[\begin{array}{c}q/c,aq/d\\ q/b\end{array};q,\frac{c}{a}],$$ (4.4) where the series terminate, or $`\mathrm{max}(\|d/ab\|,\|c\|,\|c/a\|)<1`$, for convergence. Unfortunately, the $`{}_{2}{}^{}\varphi _{1}^{}`$ on the right side of (4.4) simplifies only in special cases. If $`d=aq`$, then the $`{}_{2}{}^{}\varphi _{1}^{}`$ sum reduces just to the first term, 1, and we have the summation $${}_{2}{}^{}\psi _{2}^{}[\begin{array}{c}a,b\\ aq,c\end{array};q,\frac{q}{b}]=\frac{(q,q,aq/b,c/a;q)_{\mathrm{}}}{(aq,q/a,q/b,c;q)_{\mathrm{}}},$$ (4.5) where the series terminates, or $`\mathrm{max}(\|q/b\|,\|c\|)<1`$, for convergence. We want to add that the transformation in (4.4) is a special case of Bailey’s \[7, Eq. (2.3)\] $`{}_{2}{}^{}\psi _{2}^{}`$ transformation, $${}_{2}{}^{}\psi _{2}^{}[\begin{array}{c}a,b\\ c,d\end{array};q,z]=\frac{(az,d/a,c/b,dq/abz;q)_{\mathrm{}}}{(z,d,q/b,cd/abz;q)_{\mathrm{}}}{}_{2}{}^{}\psi _{2}^{}[\begin{array}{c}a,abz/d\\ az,c\end{array};q,\frac{d}{a}],$$ (4.6) where the series terminate, or $`\mathrm{max}(\|z\|,\|cd/abz\|,\|d/a\|,\|c/b\|)<1`$, for convergence. Namely, if we perform in (4.6) the simultaneous substitutions $`ab`$, $`ba`$, and $`zd/ab`$, and reverse the order of summation in the truncated series on the right side, we obtain (4.4). In Section 3, we found, following M. Jackson, a sum for a general $`{}_{1}{}^{}\psi _{1}^{}`$ series. So far in this section, we applied her method to obtain a transformation for a particular $`{}_{2}{}^{}\psi _{2}^{}`$ into a (multiple of a) $`{}_{2}{}^{}\varphi _{1}^{}`$ series. As a matter of fact, there is no closed form (as a product of linear factors) for the summation of a general $`{}_{2}{}^{}\psi _{2}^{}`$ series. The situation is different in the $`q1`$ case, though. In the following, we review the classical $`{}_{2}{}^{}F_{1}^{}`$ and $`{}_{2}{}^{}H_{2}^{}`$ summations and then prove the latter by our elementary method. In his doctoral dissertation , Gauß showed that $${}_{2}{}^{}F_{1}^{}[\begin{array}{c}a,b\\ c\end{array};1]=\frac{\mathrm{\Gamma }(c)\mathrm{\Gamma }(cab)}{\mathrm{\Gamma }(ca)\mathrm{\Gamma }(cb)},$$ (4.7) where the series either terminates, or $`\mathrm{}(cab)>0`$, for convergence. Dougall \[10, Sec. 13\] extended this result to $${}_{2}{}^{}H_{2}^{}[\begin{array}{c}a,b\\ c,d\end{array};1]=\frac{\mathrm{\Gamma }(1a)\mathrm{\Gamma }(1b)\mathrm{\Gamma }(c)\mathrm{\Gamma }(d)\mathrm{\Gamma }(c+dab1)}{\mathrm{\Gamma }(ca)\mathrm{\Gamma }(cb)\mathrm{\Gamma }(da)\mathrm{\Gamma }(db)},$$ (4.8) where the series either terminates, or $`\mathrm{}(c+dab1)>0`$, for convergence. Clearly, the $`d1`$ case of (4.8) is (4.7). We are ready to derive (4.8) from (4.7): In (4.7), we perform the simultaneous substitutions $`aa+n`$, $`b1c`$, and $`c1+n`$, and obtain $${}_{2}{}^{}F_{1}^{}[\begin{array}{c}a+n,1c\\ 1+n\end{array};1]=\frac{\mathrm{\Gamma }(1+n)\mathrm{\Gamma }(ca)}{\mathrm{\Gamma }(1a)\mathrm{\Gamma }(c+n)},$$ (4.9) provided $`\mathrm{}(ca)>0`$. Using some elementary identities for shifted factorials (see, e. g., Slater \[26, Appendix I\]) we can rewrite equation (4.9) as $$\frac{\mathrm{\Gamma }(1a)\mathrm{\Gamma }(c)}{\mathrm{\Gamma }(ca)}\underset{k=0}{\overset{\mathrm{}}{}}\frac{(1c)_k(a)_{n+k}}{(1)_k(1)_{n+k}}=\frac{(a)_n}{(c)_n}.$$ (4.10) Alternatively, we could have used (3.5) with the substitutions $`aq^a`$ and $`bq^c`$, and then let $`q1^{}`$, to arrive directly at (4.10). In (4.10), we multiply both sides by $`(b)_n/(d)_n`$ and sum over all integers $`n`$. On the right side we obtain $${}_{2}{}^{}H_{2}^{}[\begin{array}{c}a,b\\ c,d\end{array};1].$$ On the left side we obtain $$\frac{\mathrm{\Gamma }(1a)\mathrm{\Gamma }(c)}{\mathrm{\Gamma }(ca)}\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}\frac{(b)_n}{(d)_n}\underset{k=0}{\overset{\mathrm{}}{}}\frac{(1c)_k(a)_{n+k}}{(1)_k(1)_{n+k}}.$$ (4.11) Next, we interchange summations in (4.11) and shift the inner index $`nnk`$. (Observe that the sum over $`n`$ is terminated by the term $`(1)_{n+k}^1`$ from below.) We obtain $$\frac{\mathrm{\Gamma }(1a)\mathrm{\Gamma }(c)}{\mathrm{\Gamma }(ca)}\underset{k=0}{\overset{\mathrm{}}{}}\frac{(1c)_k(b)_k}{(1)_k(d)_k}\underset{n=0}{\overset{\mathrm{}}{}}\frac{(a)_n(bk)_n}{(1)_n(dk)_n}.$$ Now, the inner sum, provided $`\mathrm{}(dab)>0`$, can be evaluated by (4.7) and we obtain $$\frac{\mathrm{\Gamma }(1a)\mathrm{\Gamma }(c)}{\mathrm{\Gamma }(ca)}\underset{k=0}{\overset{\mathrm{}}{}}\frac{(1c)_k(b)_k}{(1)_k(d)_k}\frac{\mathrm{\Gamma }(dk)\mathrm{\Gamma }(dab)}{\mathrm{\Gamma }(dak)\mathrm{\Gamma }(db)},$$ which can be simplified to $$\frac{\mathrm{\Gamma }(1a)\mathrm{\Gamma }(c)\mathrm{\Gamma }(d)\mathrm{\Gamma }(dab)}{\mathrm{\Gamma }(ca)\mathrm{\Gamma }(da)\mathrm{\Gamma }(db)}\underset{k=0}{\overset{\mathrm{}}{}}\frac{(1c)_k(1+ad)_k}{(1)_k(1b)_k}.$$ To the last inner sum, provided $`\mathrm{}(c+dab1)>0`$, we can again apply (4.7) and eventually obtain the right side of (4.8), as desired. We note here that to apply Gauß’ $`{}_{2}{}^{}F_{1}^{}`$ summation theorem three times, we needed certain conditions of the parameters, for convergence. These were $`\mathrm{}(ca)>0`$, $`\mathrm{}(dab)>0`$, and $`\mathrm{}(c+dab1)>0`$. But in the end the first two of these conditions may be removed by analytic continuation. In particular, both sides of identity (4.8) are analytic in $`a`$ for $`\mathrm{}(a)<\mathrm{}(c+db1)`$ (and excluding some poles). In the course of our derivation, we have shown the identity for $`\mathrm{}(a)<\mathrm{min}(\mathrm{}(c),\mathrm{}(db),\mathrm{}(c+db1))`$ (excluding some poles). By analytic continuation, we extend the identity, when defined, to be valid for $`\mathrm{}(a)<\mathrm{}(c+db1)`$, the region of convergence of the series. ## 5. Bailey’s very-well-poised $`{}_{6}{}^{}\psi _{6}^{}`$ summation One of the most powerful identities for bilateral basic hypergeometric series is Bailey’s very-well-poised $`{}_{6}{}^{}\psi _{6}^{}`$ summation: $$\begin{array}{c}{}_{6}{}^{}\psi _{6}^{}[\begin{array}{c}q\sqrt{a},q\sqrt{a},b,c,d,e\\ \sqrt{a},\sqrt{a},aq/b,aq/c,aq/d,aq/e\end{array};q,\frac{a^2q}{bcde}]\hfill \\ \hfill =\frac{(aq,aq/bc,aq/bd,aq/be,aq/cd,aq/ce,aq/de,q,q/a;q)_{\mathrm{}}}{(aq/b,aq/c,aq/d,aq/e,q/b,q/c,q/d,q/e,a^2q/bcde;q)_{\mathrm{}}},\end{array}$$ (5.1) provided the series either terminates, or $`\|a^2q/bcde\|<1`$, for convergence. To prove Bailey’s $`{}_{6}{}^{}\psi _{6}^{}`$ summation, we start with a suitable specialization of Rogers’ $`{}_{6}{}^{}\varphi _{5}^{}`$ summation: $${}_{6}{}^{}\varphi _{5}^{}[\begin{array}{c}a,q\sqrt{a},q\sqrt{a},b,c,d\\ \sqrt{a},\sqrt{a},aq/b,aq/c,aq/d\end{array};q,\frac{aq}{bcd}]=\frac{(aq,aq/bc,aq/bd,aq/cd;q)_{\mathrm{}}}{(aq/b,aq/c,aq/d,aq/bcd;q)_{\mathrm{}}},$$ (5.2) provided the series either terminates, or $`\|aq/bcd\|<1`$, for convergence. Note that (5.2) is just the special case $`ea`$ of (5.1). In (5.2), we perform the simultaneous substitutions $`ac/a`$, $`bb/a`$, $`ccq^n`$ and $`dcq^n/a`$, and obtain $$\begin{array}{c}{}_{6}{}^{}\varphi _{5}^{}[\begin{array}{c}c/a,q\sqrt{c/a},q\sqrt{c/a},b/a,cq^n,cq^n/a\\ \sqrt{c/a},\sqrt{c/a},cq/b,q^{1n}/a,q^{1+n}\end{array};q,\frac{aq}{bc}]\hfill \\ \hfill =\frac{(cq/a,q^{1n}/b,aq^{1+n}/b,q/c;q)_{\mathrm{}}}{(cq/b,q^{1n}/a,q^{1+n},aq/bc;q)_{\mathrm{}}},\end{array}$$ (5.3) where $`\|aq/bc\|<1`$. Using some elementary identities for $`q`$-shifted factorials (see, e. g., Gasper and Rahman \[11, Appendix I\]) we can rewrite equation (5.3) as $$\begin{array}{c}\frac{(cq/b,q/a,q,aq/bc;q)_{\mathrm{}}}{(cq/a,q/b,aq/b,q/c;q)_{\mathrm{}}}\underset{k=0}{\overset{\mathrm{}}{}}\frac{(1cq^{2k}/a)}{(1c/a)}\frac{(c/a,b/a;q)_k(c;q)_{n+k}(a;q)_{nk}}{(q,cq/b;q)_k(q;q)_{n+k}(aq/c;q)_{nk}}\left(\frac{a}{b}\right)^k\hfill \\ \hfill =\frac{(b,c;q)_n}{(aq/b,aq/c;q)_n}\left(\frac{a}{b}\right)^n.\end{array}$$ (5.4) In this identity, we multiply both sides by $$\frac{(1aq^{2n})}{(1a)}\frac{(d,e;q)_n}{(aq/d,aq/e;q)_n}\left(\frac{aq}{cde}\right)^n$$ and sum over all integers $`n`$. On the right side we obtain $${}_{6}{}^{}\psi _{6}^{}[\begin{array}{c}q\sqrt{a},q\sqrt{a},b,c,d,e\\ \sqrt{a},\sqrt{a},aq/b,aq/c,aq/d,aq/e\end{array};q,\frac{a^2q}{bcde}].$$ On the left side we obtain $$\begin{array}{c}\frac{(cq/b,q/a,q,aq/bc;q)_{\mathrm{}}}{(cq/a,q/b,aq/b,q/c;q)_{\mathrm{}}}\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}\frac{(1aq^{2n})}{(1a)}\frac{(d,e;q)_n}{(aq/d,aq/e;q)_n}\left(\frac{aq}{cde}\right)^n\hfill \\ \hfill \times \underset{k=0}{\overset{\mathrm{}}{}}\frac{(1cq^{2k}/a)}{(1c/a)}\frac{(c/a,b/a;q)_k(c;q)_{n+k}(a;q)_{nk}}{(q,cq/b;q)_k(q;q)_{n+k}(aq/c;q)_{nk}}\left(\frac{a}{b}\right)^k.\end{array}$$ (5.5) Next, we interchange summations in (5.5) and shift the inner index $`nnk`$. (Observe that the sum over $`n`$ is terminated by the term $`(q;q)_{n+k}^1`$ from below.) We obtain, again using some elementary identities for $`q`$-shifted factorials, $$\begin{array}{c}\frac{(cq/b,q/a,q,aq/bc;q)_{\mathrm{}}}{(cq/a,q/b,aq/b,q/c;q)_{\mathrm{}}}\underset{k=0}{\overset{\mathrm{}}{}}\frac{(1cq^{2k}/a)}{(1c/a)}\frac{(c/a,b/a;q)_k}{(q,cq/b;q)_k}\hfill \\ \hfill \times \frac{(1aq^{2k})}{(1a)}\frac{(a;q)_{2k}(d,e;q)_k}{(aq/c;q)_{2k}(aq/d,aq/e;q)_k}\left(\frac{cde}{bq}\right)^k\\ \hfill \times \underset{n=0}{\overset{\mathrm{}}{}}\frac{(1aq^{2k+2n})}{(1aq^{2k})}\frac{(aq^{2k},c,dq^k,eq^k;q)_n}{(q,aq^{12k}/c,aq^{1k}/d,aq^{1k}/e;q)_n}\left(\frac{aq}{cde}\right)^n.\end{array}$$ Now the inner sum, provided $`\|aq/cde\|<1`$, can be evaluated by (5.2) and we obtain $$\begin{array}{c}\frac{(cq/b,q/a,q,aq/bc;q)_{\mathrm{}}}{(cq/a,q/b,aq/b,q/c;q)_{\mathrm{}}}\underset{k=0}{\overset{\mathrm{}}{}}\frac{(1cq^{2k}/a)}{(1c/a)}\frac{(c/a,b/a;q)_k(aq;q)_{2k}}{(q,cq/b;q)_k(aq/c;q)_{2k}}\hfill \\ \hfill \times \frac{(d,e;q)_k}{(aq/d,aq/e;q)_k}\left(\frac{cde}{bq}\right)^k\frac{(aq^{12k},aq^{1k}/cd,aq^{1k}/ce,aq/de;q)_{\mathrm{}}}{(aq^{12k}/c,aq^{1k}/d,aq^{1k}/e,aq/cde;q)_{\mathrm{}}},\end{array}$$ which can be simplified to $$\begin{array}{c}\frac{(cq/b,q/a,q,aq/bc,aq,aq/cd,aq/ce,aq/de;q)_{\mathrm{}}}{(cq/a,q/b,aq/b,q/c,aq/c,aq/d,aq/e,aq/cde;q)_{\mathrm{}}}\hfill \\ \hfill \times \underset{k=0}{\overset{\mathrm{}}{}}\frac{(1cq^{2k}/a)}{(1c/a)}\frac{(c/a,b/a,cd/a,ce/a;q)_k}{(q,cq/b,q/d,q/e;q)_k}\left(\frac{a^2q}{bcde}\right)^k.\end{array}$$ To the last sum, provided $`\|a^2q/bcde\|<1`$, we can again apply (5.2) and after some simplifications we finally obtain the right side of (5.1), as desired. Our derivation of the $`{}_{6}{}^{}\psi _{6}^{}`$ summation (5.1) is simple once the nonterminating $`{}_{6}{}^{}\varphi _{5}^{}`$ summation (5.2) is given. But the latter summation follows by an elementary computation from F. H. Jackson’s terminating $`{}_{8}{}^{}\varphi _{7}^{}`$ summation (cf. \[11, Eq. (2.6.2)\]) $$\begin{array}{c}_8\varphi _7[\begin{array}{c}a,q\sqrt{a},q\sqrt{a},b,c,d,a^2q^{1+n}/bcd,q^n\\ \sqrt{a},\sqrt{a},aq/b,aq/c,aq/d,bcdq^n/a,aq^{1+n}\end{array};q,q]\hfill \\ \hfill =\frac{(aq,aq/bc,aq/bd,aq/cd;q)_n}{(aq/b,aq/c,aq/d,aq/bcd;q)_n}\end{array}$$ (5.6) as $`n\mathrm{}`$. Jackson’s terminating $`{}_{8}{}^{}\varphi _{7}^{}`$ summation itself can be proved by various ways. An algorithmic approach uses the $`q`$-Zeilberger algorithm, see Koornwinder . For an inductive proof, see Slater \[26, Sec. 3.3.1\]. For another elementary classical proof, see Gasper and Rahman \[11, Sec. 2.6\]. Concluding this section, we would like to add another thought, kindly initiated by an anonymous referee. It is worth comparing our proof with Askey and Ismail’s elegant (and now classical) proof of Bailey’s $`{}_{6}{}^{}\psi _{6}^{}`$ summation. Their proof uses a method in this context often referred to as “Ismail’s argument” since Ismail was apparently the first to apply Liouville’s standard analytic continuation argument in the context of bilateral basic hypergeometric series. Askey and Ismail use Rogers’ $`{}_{6}{}^{}\varphi _{5}^{}`$ summation once to evaluate the $`{}_{6}{}^{}\psi _{6}^{}`$ series at an infinite sequence and then apply analytic continuation. Here, we evaluate the $`{}_{6}{}^{}\psi _{6}^{}`$ series on a domain, and, for the full theorem, we also need analytic continuation. In fact, we need, in addition to $`\|a^2q/bcde\|<1`$ two other inequalities on $`a,b,c,d,e`$, namely $`\|aq/bc\|<1`$ and $`\|aq/cde\|<1`$, in order to apply the $`{}_{6}{}^{}\varphi _{5}^{}`$ summation theorem. In the end, these additional conditions can be removed. In particular, both sides of identity (5.1) are analytic in $`1/c`$ around the origin. So far, we have shown the identity for $`\|1/c\|<\mathrm{min}(\|b/aq\|,\|de/aq\|,\|bde/a^2q\|)`$. By analytic continuation, we extend the identity to be valid for $`\|1/c\|<\|bde/a^2q\|`$, the radius of convergence of the series.
warning/0007/cond-mat0007139.html	ar5iv	text	# Dynamics of a vortex in a trapped Bose-Einstein condensate ## I Introduction The experimental achievement of Bose-Einstein condensation in confined alkali-atom gases has stimulated great interest in the generation and observation of vortices in such systems . Rotating a totally anisotropic harmonic trap at an angular frequency $`\mathrm{\Omega }`$ can, in principle, generate vortices; they are energetically stable for $`\mathrm{\Omega }>\mathrm{\Omega }_c`$ . There are several other ideas to create vortices in a trapped Bose-Einstein condensate (BEC) . Vortex formation in BEC was recently observed experimentally . In general, a vortex line in a trapped Bose-Einstein condensate is nonstationary. The vortex line can move as a whole, undergo deformation of its shape or perform oscillatory motion like helical waves . An extensive literature exists on vortex dynamics in superfluids . The nonlinear Schrödinger equation (Gross-Pitaevskii model) has served to study the dynamics and reconnection of vortices, their time evolution, and scattering interactions of superfluid vortex rings. Vortex precession in a nonuniform light beam has recently been observed and discussed in terms of the nonlinear Schrödinger equation . The dynamics of a vortex line in a spatially inhomogeneous two-dimension (2D) condensate was considered in , while the problem of curvature-driven motion of a vortex line in a homogeneous superfluid in three dimensions (3D) was studied in . A normal mode with negative frequency that corresponds to a vortex precession was found numerically and analytically for a large 3D disk-shape BEC , and for a small BEC . The motion of vortex lines and rings in Bose-Einstein condensates in harmonic traps was studied in 2D and 3D by numerical solution of the Gross-Pitaevskii equation . Minimum energy configurations of vortices in a rotating trap were considered in . In a nonrotating trap, the vortex state has a higher energy than the ground-state Bose condensate, so that the vortex is thermodynamically unstable . However, the vortex (with unit circulation quantum) is dynamically stable and can decay only in the presence of dissipation. Dissipative dynamics and the decay time of the vortex state (due to the interaction of the vortex with the thermal cloud) in a trapped Bose-condensed gas are discussed in , where the friction coefficient is found to be proportional to the temperature. At temperatures relevant to current experiments, one can neglect dissipation in studying the normal modes of the vortex because the vortex decay rate is much smaller than the frequencies of the normal modes. In this paper we consider the dynamics of a vortex line in a zero-temperature condensate in the Thomas-Fermi (TF) limit, when the vortex core radius $`\xi d^2/R`$ is small compared to the mean oscillator length $`d`$ and the mean dimension $`R`$ of the condensate \[here, $`d=(d_xd_yd_z)^{1/3}`$ with $`d_i=\sqrt{\mathrm{}/M\omega _i}`$ and trap frequencies $`\omega _i`$ ($`i=x,y,z`$)\]. We derive a general nonlinear equation for the motion of the vortex that includes the effects of the trap potential, the vortex curvature and the angular velocity of the trap rotation \[see Eq. (38) below\]. Linearization of this equation around stationary configurations gives rise to the equation for the normal modes of the vortex line. We investigate normal modes of the vortex in 2D and 3D condensates. For a 2D condensate there are solutions in the form of helical waves. For nonrotating trap some of the solutions have negative eigenfrequencies (these modes are formally unstable); furthermore, in a nonaxisymmetric trap, some solutions can have imaginary eigenfrequencies, implying that a straight central vortex line is unstable with respect to finite self-induced curvature. In a 3D condensate the spectrum of normal modes becomes discrete. For a vortex near the $`z`$ axis, the number of normal modes with negative frequency depends on the aspect ratio $`R_{}/R_z`$. A vortex in a disc-shape condensate ($`R_z<R_{}`$) has only one mode with negative frequency. However, if we change the aspect ratio to a cigar-shape condensate with $`R_{}<R_z`$, more modes with negative frequency appear. Thus it is more difficult to stabilize a vortex in a cigar-shape condensate rather than in a disc-shape one. The plan of the paper is the following. In Sec. II we derive a general equation of vortex dynamics using the method of matched asymptotic expansions. In Secs. III and IV, we discuss the normal modes of a vortex line for 2D and 3D condensates. In Sec. V we investigate normal modes with imaginary frequencies that appear for a vortex in a nonaxisymmetric condensate. In the last section we study the motion of a straight vortex line in a slightly nonspherical trap. ## II General equation of the vortex dynamics Consider a condensate in a nonaxisymmetric trap that rotates with an angular velocity $`𝛀`$. At zero temperature in a frame rotating with the angular velocity $`𝛀`$, the trap potential $`V_{\mathrm{tr}}`$ is time-independent, and the evolution of the condensate wave function $`\mathrm{\Psi }`$ is described by the time-dependent Gross-Pitaevskii (GP) equation: $$\left(\frac{\mathrm{}^2}{2M}^2+V_{\mathrm{tr}}+g\|\mathrm{\Psi }\|^2\mu (\mathrm{\Omega })+i\mathrm{}𝛀\left(𝒓\times \mathbf{}\right)\right)\mathrm{\Psi }=i\mathrm{}\frac{\mathrm{\Psi }}{t},$$ (1) where $`V_{\mathrm{tr}}=\frac{1}{2}M\left(\omega _x^2x^2+\omega _y^2y^2+\omega _z^2z^2\right)`$ is the external trap potential, $`g=4\pi \mathrm{}^2a/M>0`$ is the effective interparticle interaction strength, and $`\mu (\mathrm{\Omega })`$ is the chemical potential in the rotating frame. We assume that the condensate contains a $`q`$-fold quantized vortex with the position vector $`𝒓_0(z,t)`$. In this section we use the method of matched asymptotic expansions to determine the vortex velocity as a function of the local gradient of the trap potential $`V_{\mathrm{tr}}`$, the vortex curvature $`k`$ and the angular velocity $`𝛀`$, generalizing the two-dimensional results obtained by Rubinstein and Pismen to the case of a three-dimensional rotating potential. The method applies when the external potential does not change significantly on distances comparable with the core size $`\|q\|\xi R_{}`$ (this is the TF limit) and when the curvature is not too large ($`k1/\|q\|\xi `$); it matches the outer asymptotic form of the solution of Eq. (1) in the vortex-core region ($`\|𝝆𝝆_0\|\|q\|\xi `$) with the short-distance behavior of the solution in the region far from the vortex core ($`\|𝝆𝝆_0\|\|q\|\xi `$). To find the solution in the vortex-core region, one may consider Eq. (1) in a local coordinate frame centered at the point $`𝒓_0`$ of the vortex line that moves with the vortex velocity $`𝐕`$. In the general case, the vortex line has a curvature $`k`$ that depends on the specific element in question. We introduce a local coordinate system $`(x,y,z)`$, so that the $`x`$ axis is directed along the vortex normal, the $`y`$ axis is along the binormal $`\widehat{b}`$ and $`z`$ axis is along the tangent $`\widehat{t}`$ (see Fig. 1). Fig. 1. Local coordinate system associated with the vortex line. The solution is assumed to be stationary in the comoving frame and satisfies the equation (in the local coordinates): $$\left(\frac{\mathrm{}^2}{2M}\left(^2k_x\right)+V_{\mathrm{tr}}(𝒓_0)+g\|\mathrm{\Psi }\|^2\mu (\mathrm{\Omega })+i\mathrm{}(𝛀\times 𝒓_0)\mathbf{}\right)\mathrm{\Psi }=$$ $$=i\mathrm{}𝐕\mathbf{}\mathrm{\Psi },$$ (2) where the term $`k_x`$ arises from the transformation to local coordinates. One can remove $`𝛀`$ from this equation by a shift $`𝐕𝐕(𝛀\times 𝒓_0)`$. In the vortex core region we may seek a solution in the form of an expansion in the small parameters $`\xi /R_{}`$ and $`k\xi `$: $$\mathrm{\Psi }=\mathrm{\Psi }_0(\rho )+\mathrm{\Psi }_1=\left[\|\mathrm{\Psi }_0(\rho )\|\chi (\rho ,z)\mathrm{cos}\varphi \right]e^{iq\varphi i\eta (\rho ,z)\mathrm{sin}\varphi },$$ (3) where $`\mathrm{\Psi }_0`$ is the condensate wave function with $`V_{\mathrm{tr}}`$ replaced by $`V_{\mathrm{tr}}(𝒓_0)`$; it satisfies a zero-order equation $$\left(\frac{\mathrm{}^2}{2M}^2+V_{\mathrm{tr}}(𝒓_0)+g\|\mathrm{\Psi }_0\|^2\mu (\mathrm{\Omega })\right)\mathrm{\Psi }_0=0,$$ (4) and $`\chi `$, $`\eta `$ characterize the perturbation in the absolute value and phase. Physically $`\mathrm{\Psi }_0`$ is the analogous wave function for a laterally unbounded condensate with chemical potential $`\mu (\mathrm{\Omega })V_{\mathrm{tr}}(𝝆_0)`$. The polar angle $`\varphi `$ is measured from the direction of the vortex normal ($`\widehat{n}\widehat{x}`$) and $`\rho `$ is the radial cylindrical coordinate in the local frame. The perturbation $`\mathrm{\Psi }_1`$ obeys the following equation $$L(\mathrm{\Psi }_1,\mathrm{\Psi }_1^{})=\frac{2Mi}{\mathrm{}}𝐕\mathbf{}\mathrm{\Psi }_0+\frac{2M}{\mathrm{}^2}\mathrm{\Psi }_0𝝆\mathbf{}_{}V_{\mathrm{tr}}(𝒓_0)+k_x\mathrm{\Psi }_0,$$ (5) where $$L(\mathrm{\Psi }_1,\mathrm{\Psi }_1^{})^2\mathrm{\Psi }_1+\frac{2M}{\mathrm{}^2}\left[\left(\mu (\mathrm{\Omega })V_{\mathrm{tr}}(𝒓_0)2g\|\mathrm{\Psi }_0\|^2\right)\mathrm{\Psi }_1g\mathrm{\Psi }_0^2\mathrm{\Psi }_1^{}\right]$$ (6) is a self-conjugate operator and $`\mathbf{}_{}`$ is the gradient operator in a plane perpendicular to the vortex line. This equation is linear in $`\mathrm{\Psi }_1`$ and in $`𝐕`$; it contains $`\mathbf{}_{}V_{\mathrm{tr}}`$ and $`k`$ as independent sources, so that the velocity of the vortex line is a sum of independent contributions due to $`\mathbf{}_{}V_{\mathrm{tr}}`$ and $`k`$. Also, the function $`\mathrm{\Psi }_0`$ depends only on the coordinates in the direction perpendicular to the vortex line; therefore, in the dot product $`𝐕\mathbf{}\mathrm{\Psi }_0`$ only the component of the velocity perpendicular to the vortex line is relevant. We also assume that $`𝐕`$ has no component along the line. For simplicity, one can assume that $`\mathbf{}_{}V_{\mathrm{tr}}`$ lies along $`\widehat{n}`$ and derive the vortex velocity as a sum of two independent contributions. The final result written in vector form remains valid for arbitrary directions of $`\mathbf{}_{}V_{\mathrm{tr}}`$ and $`\widehat{n}`$. Under this assumption, we have $`𝐕\mathbf{}\mathrm{\Psi }_0=V_y\mathrm{\Psi }_0=V\left(\mathrm{sin}\varphi _\rho +\rho ^1\mathrm{cos}\varphi _\varphi \right)\mathrm{\Psi }_0`$ in polar coordinates. Then, writing $`\mathrm{\Psi }_1=(\chi \mathrm{cos}\varphi +i\eta \|\mathrm{\Psi }_0\|\mathrm{sin}\varphi )e^{iq\varphi }`$ in terms of the small perturbations $`\chi `$ and $`\eta `$, Eq. (5) has the form: $$\left(_{\rho \rho }^2+\frac{1}{\rho }_\rho +_{zz}^2\right)\chi +\frac{2M}{\mathrm{}^2}\left(\mu (\mathrm{\Omega })V_{\mathrm{tr}}(𝒓_0)3g\|\mathrm{\Psi }_0\|^2\right)\chi \frac{q^2+1}{\rho ^2}\chi \frac{2q}{\rho ^2}\|\mathrm{\Psi }_0\|\eta $$ $$=\frac{2M}{\mathrm{}^2}\|\mathrm{\Psi }_0\|\left(\frac{\mathrm{}q}{\rho }V\rho \|\mathbf{}_{}V_{\mathrm{tr}}(𝒓_0)\|\right)k_\rho \|\mathrm{\Psi }_0\|$$ (7) $$\left(_{\rho \rho }^2+\frac{1}{\rho }_\rho +_{zz}^2\frac{1}{\rho ^2}\right)\eta +\frac{2}{\|\mathrm{\Psi }_0\|}\left(_\rho \|\mathrm{\Psi }_0\|_\rho \eta +_z\|\mathrm{\Psi }_0\|_z\eta \frac{q}{\rho ^2}\chi \right)$$ $$=\frac{2M}{\mathrm{}}V\frac{_\rho \|\mathrm{\Psi }_0\|}{\|\mathrm{\Psi }_0\|}+\frac{kq}{\rho }$$ (8) We can remove $`V`$ from these equations with the following gauge transformation: $$\eta =\stackrel{~}{\eta }\frac{M}{\mathrm{}}\rho V$$ (9) Further, for large distances $`\|q\|\xi \rho R_{}`$, we can use $`g\|\mathrm{\Psi }_0\|^2g\|\mathrm{\Psi }_{TF}\|^2\mu (\mathrm{\Omega })V_{\mathrm{tr}}(𝒓_0)`$ and rewrite Eqs. (7) and (8) as follows $$2g\|\mathrm{\Psi }_{TF}\|\chi =\rho \|\mathbf{}_{}V_{\mathrm{tr}}(𝝆_0)\|+\frac{k\mathrm{}^2}{2M\|\mathrm{\Psi }_{TF}\|}_\rho \|\mathrm{\Psi }_0\|,$$ (10) $$\left(_{\rho \rho }^2+\frac{1}{\rho }_\rho \frac{1}{\rho ^2}\right)\stackrel{~}{\eta }\frac{\chi }{\|\mathrm{\Psi }_{TF}\|}\frac{2q}{\rho ^2}=\frac{kq}{\rho };$$ (11) equivalently, $$\chi =\frac{\rho }{2g\|\mathrm{\Psi }_{TF}\|}\|\mathbf{}_{}V_{\mathrm{tr}}(𝒓_0)\|+\frac{k\mathrm{}^2}{4Mg\|\mathrm{\Psi }_{TF}\|^2}_\rho \|\mathrm{\Psi }_0\|,$$ (12) $$\left(\rho ^2_{\rho \rho }^2+\rho _\rho 1\right)\stackrel{~}{\eta }\frac{q\rho \|\mathbf{}_{}V_{\mathrm{tr}}(𝒓_0)\|}{g\|\mathrm{\Psi }_{TF}\|^2}=kq\rho +\frac{qk\mathrm{}^2_\rho \|\mathrm{\Psi }_0\|}{2Mg\|\mathrm{\Psi }_{TF}\|^3}.$$ (13) In Eq. (13), we can omit the last term, which is smaller with respect to the term $`kq\rho `$ by the factor $`\xi ^4/\rho ^4`$. As a result, for $`\rho \|q\|\xi `$ the perturbations have the following asymptotic form: $$\eta \frac{q}{2}\left(\frac{\|\mathbf{}_{}V_{\mathrm{tr}}(𝒓_0)\|}{g\|\mathrm{\Psi }_{TF}\|^2}+k\right)\rho \mathrm{ln}\left(A\rho \right)\frac{M}{\mathrm{}}\rho \left(𝑽+𝛀\times 𝒓_0\right)\widehat{y},$$ (14) $$\chi \frac{\|\mathbf{}_{}V_{\mathrm{tr}}(𝒓_0)\|}{2g\|\mathrm{\Psi }_{TF}\|}\rho +\frac{k\mathrm{}^2}{4Mg\|\mathrm{\Psi }_{TF}\|^2}_\rho \|\mathrm{\Psi }_0\|.$$ (15) In terms of the phase $`S`$, the solution (14) (the inner expansion in the coordinate frame centered at the vortex line) has the form: $$S=q\varphi \frac{q}{2}\left(\frac{\|\mathbf{}_{}V_{\mathrm{tr}}(𝒓_0)\|}{g\|\mathrm{\Psi }_{TF}\|^2}+k\right)\mathrm{ln}\left(A\rho \right)y+\frac{M}{\mathrm{}}\left(𝑽+𝛀\times 𝒓_0\right)𝒓$$ (16) The parameters $`A`$ and $`𝑽`$ must be determined by matching the solution (16) with that far from the vortex core. To the lowest order in the small parameter $`\xi /R_{}`$, Eq. (1) far from the vortex core reduces to an equation for the condensate phase only $$\|\mathrm{\Psi }_{TF}\|^2^2S+\mathbf{}\|\mathrm{\Psi }_{TF}\|^2\mathbf{}S\frac{M}{\mathrm{}}𝛀\left(𝒓\times \mathbf{}\right)\|\mathrm{\Psi }_{TF}\|^2=0,$$ (17) where $`\mathrm{\Psi }=\|\mathrm{\Psi }\|e^{iS}`$. In the frame rotating with the trap and for $`𝛀=\mathrm{\Omega }\widehat{z}`$, the phase has the form $$S=S_0\frac{M}{\mathrm{}}\frac{(\omega _x^2\omega _y^2)}{(\omega _x^2+\omega _y^2)}\mathrm{\Omega }xy,$$ (18) where $`S_0`$ is independent of $`\mathrm{\Omega }`$ . Under a shift of coordinates $`𝒓𝒓_0+𝒓`$, we have $$SS_0+\frac{M}{\mathrm{}}\left((𝛀\times 𝒓_0)+\frac{2}{M(\omega _x^2+\omega _y^2)}\left(\mathbf{}V_{\mathrm{tr}}(𝒓_0)\times 𝛀\right)\right)𝒓.$$ (19) Comparison of (16) and (19) allows us to find the contribution to the vortex velocity due to the trap rotation $$𝐕=𝐕_0+\frac{2}{M(\omega _x^2+\omega _y^2)}\left(\mathbf{}V_{\mathrm{tr}}(𝒓_0)\times 𝛀\right),$$ (20) where $`𝐕_0`$ is the velocity for a nonrotating trap. It is next necessary to find the asymptotic form of $`S_0`$ far from the vortex core. This function $`S_0`$ satisfies the following equation (in the shifted frame): $$\|\mathrm{\Psi }_{TF}\|^2^2S_0+\mathbf{}\|\mathrm{\Psi }_{TF}\|^2\mathbf{}S_0=0.$$ (21) Introduce a function $`\mathrm{\Phi }`$ such that $$S_{0x}=q\left(\mathrm{\Phi }_y+\mathrm{\Phi }_y\mathrm{ln}\|\mathrm{\Psi }_{TF}\|^2\right),$$ (22) $$S_{0y}=q\left(\mathrm{\Phi }_x+\mathrm{\Phi }_x\mathrm{ln}\|\mathrm{\Psi }_{TF}\|^2\right),$$ (23) where $`S_{0x}=_xS_0`$ and $`S_{0y}=_yS_0`$. This representation satisfies Eq. (21) automatically. In addition, the condition $$\widehat{z}\mathbf{}\times (\mathbf{}_{}S)=S_{0yx}S_{0xy}=\mathbf{}_{}\left(S_{0y}\widehat{x}S_{0x}\widehat{y}\right)=2\pi q\delta ^{(2)}\left(𝝆\right)$$ (24) gives an equation for $`\mathrm{\Phi }`$ containing a point source at the vortex location $$\mathbf{}_{}\left[\mathbf{}_{}\mathrm{\Phi }+\mathrm{\Phi }\mathbf{}_{}\mathrm{ln}\|\mathrm{\Psi }_{TF}\|^2\right]=2\pi \delta ^{(2)}\left(𝝆\right).$$ (25) Hence $$\mathbf{}_{}\left[e^{\mathrm{ln}\|\mathrm{\Psi }_{TF}\|}\mathbf{}_{}\left(\mathrm{\Phi }e^{\mathrm{ln}\|\mathrm{\Psi }_{TF}\|}\right)\mathbf{}_{}\left(e^{\mathrm{ln}\|\mathrm{\Psi }_{TF}\|}\right)e^{\mathrm{ln}\|\mathrm{\Psi }_{TF}\|}\mathrm{\Phi }\right]=2\pi \delta ^{(2)}\left(𝝆\right),$$ (26) or $$e^{\mathrm{ln}\|\mathrm{\Psi }_{TF}\|}_{}^2\left(\mathrm{\Phi }e^{\mathrm{ln}\|\mathrm{\Psi }_{TF}\|}\right)_{}^2\left(e^{\mathrm{ln}\|\mathrm{\Psi }_{TF}\|}\right)e^{\mathrm{ln}\|\mathrm{\Psi }_{TF}\|}\mathrm{\Phi }=2\pi \delta ^{(2)}\left(𝝆\right).$$ (27) We can put $`_{}^2\left(e^{\mathrm{ln}\|\mathrm{\Psi }_{TF}\|}\right)e^{\mathrm{ln}\|\mathrm{\Psi }_{TF}\|}_{}^2V_{\mathrm{tr}}/2g\|\mathrm{\Psi }_{TF}\|^2`$, so that Eq. (27) becomes $$_{}^2\left(\mathrm{\Phi }e^{\mathrm{ln}\|\mathrm{\Psi }_{TF}\|}\right)\frac{_{}^2V_{\mathrm{tr}}}{2g\|\mathrm{\Psi }_{TF}\|^2}\mathrm{\Phi }e^{\mathrm{ln}\|\mathrm{\Psi }_{TF}\|}=2\pi \delta ^{(2)}\left(𝝆\right)e^{\mathrm{ln}\|\mathrm{\Psi }_{TF}\|}$$ (28) To find the solution, we rewrite Eq. (28) in the local coordinate frame associated with the vortex line, taking into account the effect of curvature: $$\left(_{}^2k_x\right)\left(\mathrm{\Phi }e^{\mathrm{ln}\|\mathrm{\Psi }_{TF}\|}\right)\frac{_{}^2V_{\mathrm{tr}}}{2g\|\mathrm{\Psi }_{TF}\|^2}\mathrm{\Phi }e^{\mathrm{ln}\|\mathrm{\Psi }_{TF}\|}=2\pi \delta ^{(2)}\left(𝝆\right)e^{\mathrm{ln}\|\mathrm{\Psi }_{TF}\|},$$ (29) or $$_{}^2\left(\mathrm{\Phi }e^{\mathrm{ln}\|\mathrm{\Psi }_{TF}\|kx/2}\right)\left(\frac{_{}^2V_{\mathrm{tr}}}{2g\|\mathrm{\Psi }_{TF}\|^2}+\frac{k^2}{4}\right)\mathrm{\Phi }e^{\mathrm{ln}\|\mathrm{\Psi }_{TF}\|kx/2}=2\pi \delta ^{(2)}\left(𝝆\right)e^{\mathrm{ln}\|\mathrm{\Psi }_{TF}\|}.$$ (30) The solution of this inhomogeneous equation is (note that an additional solution of the homogeneous equation does not satisfy the boundary conditions at large $`\rho `$ and should be omitted): $$\mathrm{\Phi }=e^{kx/2}K_0\left(\sqrt{\frac{_{}^2V_{\mathrm{tr}}}{2g\|\mathrm{\Psi }_{TF}\|^2}+\frac{k^2}{4}}\rho \right),$$ (31) where $`K_0`$ is a modified Bessel function \[for small $`x`$, we have $`K_0\mathrm{ln}\left(e^Cx/2\right)`$ where $`C=0.577\mathrm{}`$ is the Euler constant\]. Further, under the logarithm we may put $`_{}^2V_{\mathrm{tr}}/4g\|\mathrm{\Psi }_{TF}\|^21/R_{}^2`$. Hence $`\mathrm{\Phi }`$ (in the local coordinate frame centered at the vortex line) has the short-distance form $$\mathrm{\Phi }\mathrm{ln}\left(\frac{e^C}{\sqrt{2}}\sqrt{\frac{1}{R_{}^2}+\frac{k^2}{8}}\rho \right).$$ (32) To make the asymptotic matching, we can express the solution (16) in the vortex-core region in terms of $`\mathrm{\Phi }`$ and then compare with the formula (32). Using the definitions (22), (23) of the function $`\mathrm{\Phi }`$, one can show that the expression (16) (for $`\mathrm{\Omega }=0`$) corresponds to the following function $`\mathrm{\Phi }`$ (in the coordinate frame centered at the vortex line): $$\mathrm{\Phi }\left[1+\frac{x}{2}\left(\frac{\|\mathbf{}_{}V_{\mathrm{tr}}(𝒓_0)\|}{g\|\mathrm{\Psi }_{TF}\|^2}+k\right)\right]\mathrm{ln}\left(A\rho \right)+\frac{MV_0}{\mathrm{}q}x.$$ (33) To verify this expression, we use $`_x\mathrm{ln}\|\mathrm{\Psi }_{TF}\|^2\|\mathbf{}_{}V_{\mathrm{tr}}(𝒓_0)\|/g\|\mathrm{\Psi }_{TF}\|^2`$ and $`_y\mathrm{ln}\|\mathrm{\Psi }_{TF}\|^20`$ in the local coordinate frame where $`\mathrm{\Phi }_x\left(_xk\right)\mathrm{\Phi }`$. Substituting (33) into (22) and (23), we obtain $$S_{0x}=\frac{qy}{\rho ^2},$$ (34) $$S_{0y}=\frac{qx}{\rho ^2}\frac{q}{2}\left(\frac{\|\mathbf{}_{}V_{\mathrm{tr}}(𝒓_0)\|}{g\|\mathrm{\Psi }_{TF}\|^2}+k\right)\mathrm{ln}\left(A\rho \right)+\frac{MV_0}{\mathrm{}}.$$ (35) Therefore, in the coordinate frame centered at the vortex line, the remaining contribution to the phase is $$S_0=q\varphi \frac{q}{2}\left(\frac{\|\mathbf{}_{}V_{\mathrm{tr}}(𝒓_0)\|}{g\|\mathrm{\Psi }_{TF}\|^2}+k\right)\mathrm{ln}\left(A\rho \right)\rho \mathrm{sin}\varphi +\frac{M}{\mathrm{}}V_0\rho \mathrm{sin}\varphi $$ (36) which coincides with (16) (for $`\mathrm{\Omega }=0`$). Matching (33) and (32) (at $`\rho \|q\|\xi `$) gives an expression for the constant $`A`$ (with logarithmic accuracy): $$\mathrm{ln}\left(Ae\right)=\mathrm{ln}\sqrt{\frac{1}{R_{}^2}+\frac{k^2}{8}},$$ (37) where $`R_{}`$ is the mean transverse dimension of the condensate, and for the velocity $`V_0`$. Finally, in general vector form (in the frame rotating with the trap), the vortex velocity is: $$𝐕(𝒓_\mathrm{𝟎})=\frac{q\mathrm{}}{2M}\left(\frac{\widehat{t}\times \mathbf{}V_{\mathrm{tr}}(𝒓_0)}{g\|\mathrm{\Psi }_{TF}\|^2}+k\widehat{b}\right)\mathrm{ln}\left(\|q\|\xi \sqrt{\frac{1}{R_{}^2}+\frac{k^2}{8}}\right)+\frac{2\left(\mathbf{}V_{\mathrm{tr}}(𝒓_0)\times 𝛀\right)}{\mathrm{\Delta }_{}V_{\mathrm{tr}}(𝒓_0)},$$ (38) where $`\widehat{b}`$ is a unit vector in the direction to the vortex binormal, $`\widehat{t}`$ is a tangent vector to the vortex line and $`\mathrm{\Delta }_{}`$ is the Laplacian operator in the plane perpendicular to $`𝛀`$. This formula is valid for arbitrary directions of the local gradient of the trap potential, the normal to the vortex line and $`𝛀`$. Near the condensate boundary the denominator of the first term in this formula goes to zero. Therefore, $`\widehat{t}\times \mathbf{}V_{\mathrm{tr}}`$ must also vanish near the boundary, implying that $`\widehat{t}`$ is parallel to$`\mathbf{}V_{\mathrm{tr}}`$; as a result, the vortex line obeys the boundary condition that its axis is perpendicular to the boundary. ## III Normal modes in two dimensions To understand the implications of Eq. (38), it is valuable to consider first the case of a 2D condensate with $`𝛀=\mathrm{\Omega }\widehat{z}`$ and $`\omega _z=0`$ (hence no confinement in the $`z`$ direction). Let the vector $`𝝆_0(z,t)=(x(z,t),y(z,t))`$ describe the time-dependent position of the vortex line during its motion, and $`𝐤`$ be the vector of the principal curvature; then $`k\widehat{b}=\widehat{t}\times 𝐤`$ and $`d^2𝝆_0/ds^2=𝐤`$, where $`s`$ is the length measured along the vortex line. For small displacements of the line from the $`z`$ axis, we have $`sz`$, $`\widehat{t}\widehat{z}`$ and $`𝐤d^2𝝆_0/dz^2`$. Then using $$\widehat{z}\times \mathbf{}V_{\mathrm{tr}}=M\omega _y^2y\widehat{x}+M\omega _x^2x\widehat{y}$$ (39) and $$k\widehat{b}=\widehat{t}\times 𝐤\widehat{x}\frac{^2y}{z^2}+\widehat{y}\frac{^2x}{z^2},$$ (40) we obtain the following coupled differential equations for $`x(z,t)`$ and $`y(z,t)`$ $$\frac{x}{t}=\frac{q\mathrm{}}{2M}\left(\frac{2y}{R_y^2}+\frac{^2y}{z^2}\right)\mathrm{ln}\left(\|q\|\xi \sqrt{\frac{1}{R_{}^2}+\frac{k^2}{8}}\right)+\frac{4\mathrm{\Omega }\mu }{M(\omega _x^2+\omega _y^2)}\frac{y}{R_y^2},$$ (41) $$\frac{y}{t}=\frac{q\mathrm{}}{2M}\left(\frac{2x}{R_x^2}+\frac{^2x}{z^2}\right)\mathrm{ln}\left(\|q\|\xi \sqrt{\frac{1}{R_{}^2}+\frac{k^2}{8}}\right)\frac{4\mathrm{\Omega }\mu }{M(\omega _x^2+\omega _y^2)}\frac{x}{R_x^2}.$$ (42) These equations have solutions in the form of helical waves $$x=\epsilon _x\mathrm{sin}(\omega t+\kappa z+\phi _0),y=\epsilon _y\mathrm{cos}(\omega t+\kappa z+\phi _0),$$ (43) with the following dispersion relation between $`\omega `$ and $`\kappa `$ (under the logarithm we take $`k\|\kappa \|`$): $$\omega =\pm \frac{q\mathrm{}}{2MR_xR_y}\sqrt{\left(2\kappa ^2R_x^2\stackrel{~}{\mathrm{\Omega }}\right)\left(2\kappa ^2R_y^2\stackrel{~}{\mathrm{\Omega }}\right)}\mathrm{ln}\left(\|q\|\xi \sqrt{\frac{1}{R_{}^2}+\frac{\|\kappa \|^2}{8}}\right).$$ (44) The associated amplitudes obey the relation $$\epsilon _y=\pm \frac{R_y}{R_x}\epsilon _x\sqrt{\frac{2\kappa ^2R_x^2\stackrel{~}{\mathrm{\Omega }}}{2\kappa ^2R_y^2\stackrel{~}{\mathrm{\Omega }}}},$$ (45) where $$\stackrel{~}{\mathrm{\Omega }}=\frac{4MR_x^2R_y^2\mathrm{\Omega }}{q\mathrm{}(R_x^2+R_y^2)\mathrm{ln}(\|q\|\xi \sqrt{\frac{1}{R_{}^2}+\frac{\|\kappa \|^2}{8}})^1},$$ (46) is a dimensionless rotation speed. For a nonaxisymmetric trap (for example, $`R_x>R_y`$) the oscillation frequency becomes imaginary if $`\sqrt{(2\stackrel{~}{\mathrm{\Omega }})}/R_x<\|\kappa \|<\sqrt{(2\stackrel{~}{\mathrm{\Omega }})}/R_y`$, in which case the initial orientation of the vortex line along the $`z`$ axis is unstable with respect to the formation of finite curvature. For $`\stackrel{~}{\mathrm{\Omega }}>\stackrel{~}{\mathrm{\Omega }}_m=2`$, however, the oscillation frequency is real (and positive) for any $`\kappa `$. Thus the trap rotation stabilizes a vortex line that initially lies along the $`z`$ axis. For a straight vortex line ($`\kappa =0`$), the frequency is always real. This unstable normal mode has the most negative frequency \[we choose the sign in (44) that corresponds to positive-norm solution\] with $$\omega =\frac{q\mathrm{}}{MR_xR_y}\left(\mathrm{ln}\left(\frac{R_{}}{\|q\|\xi }\right)\frac{4\mu \mathrm{\Omega }}{q\mathrm{}(\omega _x^2+\omega _y^2)}\right).$$ (47) For $`q>0`$ and $`\mathrm{\Omega }=0`$, the vortex moves (around the $`z`$ axis) counterclockwise in the positive sense. With increasing rotation frequency $`\mathrm{\Omega }`$ of the trap, the vortex velocity (as seen in the rotating frame) decreases towards zero and vanishes at $`\mathrm{\Omega }=\mathrm{\Omega }_m`$, where the metastable angular velocity $`\mathrm{\Omega }_m`$ of trap rotation is given by $$\mathrm{\Omega }_m=\frac{\|q\|\mathrm{}(\omega _x^2+\omega _y^2)}{4\mu }\mathrm{ln}\left(\frac{R_{}}{\|q\|\xi }\right).$$ (48) This value $`\mathrm{\Omega }_m`$ corresponds to the angular velocity of trap rotation at which a straight vortex line at the trap center first becomes a local minimum of energy (for $`\mathrm{\Omega }<\mathrm{\Omega }_m`$, the central position is a local maximum) . For $`\mathrm{\Omega }>\mathrm{\Omega }_m`$ the apparent motion of the vortex becomes clockwise. For $`\kappa =0`$, this straight vortex follows an elliptic trajectory along the line $`V_{\mathrm{tr}}=const`$, as expected from the dissipationless character of the GP equation. For a uniform condensate ($`R_x,`$ $`R_y\mathrm{}`$), Eq. (44) coincides with the well-known dispersion law of small oscillations of a straight vortex line (Kelvin modes): $$\omega =\pm \frac{q\mathrm{}}{2M}k^2\mathrm{ln}(\|q\|\xi k).$$ (49) Note that one can represent the helical wave solution (43) as a sum of two plane-wave solutions: $$x_1=\epsilon _x\mathrm{cos}(kz)\mathrm{sin}(\omega t+\phi _0),y_1=\epsilon _y\mathrm{cos}(kz)\mathrm{cos}(\omega t+\phi _0),$$ (50) $$x_2=\epsilon _x\mathrm{sin}(kz)\mathrm{sin}(\omega t+\phi _0+\pi /2),y_2=\epsilon _y\mathrm{sin}(kz)\mathrm{cos}(\omega t+\phi _0+\pi /2).$$ (51) One can easily see that Eqs. (50) and (51) are indeed solutions of Eqs. (41) and (42) with the same dispersion law $`\omega =\omega (k)`$ as the helical wave (44). In fact, the general motion of the vortex line can be represented as a combination of plane-wave solutions; helical waves are just one of the possible combinations and hence do not represent a different set of solutions. Solutions (50), (51) have different parity, but the same eigenfrequency (in 2D there is degeneracy). In 3D case, the plane-wave solutions are not degenerate and, therefore, in 3D it is impossible to construct a simple analog of helical waves. The general vortex motion in 3D is a combination of plane waves (plane-wave solutions in 3D exist, at least, for an axisymmetric trap) with different numbers of nodes along the symmetry axis and hence different frequencies. ## IV Dynamics of a vortex in three dimensions Let us consider small displacements of the vortex from the $`z`$ axis and $`𝛀=\mathrm{\Omega }\widehat{z}`$. The vortex curvature is proportional to the vortex displacement, so one can put $`k0`$ under the logarithm in (38). Further, for small displacements $$\widehat{t}\times \mathbf{}V_{\mathrm{tr}}M\left(\widehat{x}\left(\omega _z^2zy^{}\omega _y^2y\right)+\widehat{y}\left(\omega _x^2x\omega _z^2zx^{}\right)\right),$$ (52) where a prime denotes derivative with respect to $`z`$. Then in dimensionless coordinates $`xR_xx`$, $`yR_yy`$, $`zR_zz`$, Eq. (38) becomes: $$\dot{x}=\frac{q\mathrm{}}{2MR_xR_y}\left(\frac{2\left(\beta zy^{}y\right)}{\left(1z^2\right)}\beta y^{\prime \prime }\right)\mathrm{ln}\left(\frac{R_{}}{\|q\|\xi }\right)+\frac{4\mathrm{\Omega }\mu }{M(\omega _x^2+\omega _y^2)}\frac{y}{R_xR_y},$$ (53) $$\dot{y}=\frac{q\mathrm{}}{2MR_xR_y}\left(\frac{2\left(\alpha zx^{}x\right)}{\left(1z^2\right)}\alpha x^{\prime \prime }\right)\mathrm{ln}\left(\frac{R_{}}{\|q\|\xi }\right)\frac{4\mathrm{\Omega }\mu }{M(\omega _x^2+\omega _y^2)}\frac{x}{R_xR_y},$$ (54) where $$\alpha =\frac{R_x^2}{R_z^2},\beta =\frac{R_y^2}{R_z^2}$$ are parameters characterizing the trap anisotropy. One can seek solution of these equations in the form $$x=x(z)\mathrm{sin}(\omega t+\phi _0),y=y(z)\mathrm{cos}(\omega t+\phi _0)$$ and obtain the following ordinary differential equations for $`x(z)`$, $`y(z)`$ and $`\omega `$ $$\stackrel{~}{\omega }x=\frac{2\left(\beta zy^{}y\right)}{\left(1z^2\right)}\beta y^{\prime \prime }+\stackrel{~}{\mathrm{\Omega }}y,$$ (55) $$\stackrel{~}{\omega }y=\frac{2\left(\alpha zx^{}x\right)}{\left(1z^2\right)}\alpha x^{\prime \prime }+\stackrel{~}{\mathrm{\Omega }}x,$$ (56) where we introduce dimensionless angular velocities $$\stackrel{~}{\omega }=\frac{2MR_xR_y\omega }{q\mathrm{}\mathrm{ln}\left(\frac{R_{}}{\|q\|\xi }\right)},\stackrel{~}{\mathrm{\Omega }}=\frac{4MR_x^2R_y^2\mathrm{\Omega }}{q\mathrm{}(R_x^2+R_y^2)\mathrm{ln}\left(\frac{R_{}}{\|q\|\xi }\right)}$$ ### A Stationary configurations Consider a nonrotating trap with $`\stackrel{~}{\mathrm{\Omega }}=0`$. In this section, we seek stationary configurations in which the vortex line remains at rest (in this case, the contributions to the velocity from the vortex curvature and the trap potential compensate each other). To find the stationary configurations we need to solve Eqs. (55) and (56) with the condition $`\stackrel{~}{\omega }=0`$. The resulting equations for $`x`$ and $`y`$ uncouple; for example, the equation for $`x(z)`$ has the form: $$\left(1z^2\right)x^{\prime \prime }2zx^{}+\frac{2}{\alpha }x=0.$$ (57) The general solution of Eq. (57) can be expressed in terms of hypergeometric functions, but it is impossible to satisfy the boundary conditions that $`x(z)`$ should be finite at $`z=\pm 1`$ unless $`2/\alpha =n(n+1)`$, where $`n`$ is an integer ($`n0`$). In this case, the solutions reduce to Legendre polynomials $$xP_n(z).$$ (58) For example, the first three physical solutions are (we ignore $`n=0`$ which corresponds to $`\alpha =\mathrm{}`$): $$x_1=Cz,\alpha =1$$ (59) $$x_2=\epsilon (13z^2),\alpha =\frac{1}{3}$$ (60) $$x_3=C\left(z\frac{5}{3}z^3\right),\alpha =\frac{1}{6}$$ (61) If $`2/\alpha n(n+1)`$, the only possible solution of Eq. (57) is the trivial one with $`x=0`$. The equation for the $`y`$ coordinate has the same solutions $`yP_m(z)`$, if $`2/\beta =m(m+1)`$, and $`y=0`$, if $`2/\beta m(m+1)`$. That is, if $`2/\alpha n(n+1)`$ and $`2/\beta m(m+1)`$, there are no stationary configurations of the vortex line apart from the straight orientation along the $`z`$ axis. One should note that the integer $`n`$ (or $`m`$) enumerates the solutions not only for $`\stackrel{~}{\omega }=0`$; in particular, $`n`$ represents to number of times that the vortex line (precessing with angular velocity $`\stackrel{~}{\omega }_n`$) crosses $`z`$ axis. For an axisymmetric trap ($`\alpha =\beta `$), we can consider $`\stackrel{~}{\omega }_n`$ as a function of $`\alpha `$; the function $`\stackrel{~}{\omega }_n`$ changes sign at $`\alpha =\alpha _n=2/n(n+1)`$. This observation allows us to find the number of normal modes with negative frequency at a fixed value of anisotropy parameter $`\alpha `$. If $`\alpha >1`$ there is only one mode with negative frequency. If $`\frac{1}{3}<\alpha <1`$ there are 2 such normal modes. If $`\frac{1}{6}<\alpha <\frac{1}{3}`$ there are 3 modes, and so on. If $`\alpha _n<\alpha <\alpha _{n1}`$ there are $`n`$ normal modes with negative frequency. ### B Dynamics of a vortex in a disk-shape condensate $`R_zR_{}`$: Investigation of unstable mode In the limit $`\alpha `$, $`\beta 1`$ the approximate solution of Eqs. (55) and (56) that corresponds to the unstable mode is $$x=\epsilon \left(1+\frac{z^2}{2\alpha }\right),$$ (62) $$y=\epsilon \left(1+\frac{z^2}{2\beta }\right),$$ (63) with the corresponding eigenvalue $$\stackrel{~}{\omega }=\stackrel{~}{\mathrm{\Omega }}3\frac{1}{10}\left(\frac{1}{\alpha }+\frac{1}{\beta }\right).$$ (64) For $`\stackrel{~}{\mathrm{\Omega }}=0`$ the excitation energy is negative and hence formally unstable. If the trap rotates, the solution (64) becomes stable at $`\|\mathrm{\Omega }\|`$ $`\mathrm{\Omega }_m`$, where $$\mathrm{\Omega }_m=\frac{\|q\|\mathrm{}(\omega _x^2+\omega _y^2)}{8\mu }\left(3+\frac{1}{10}\left(\frac{1}{\alpha }+\frac{1}{\beta }\right)\right)\mathrm{ln}\left(\frac{R_{}}{\|q\|\xi }\right).$$ (65) This expression generalizes that for the angular velocity at which a straight vortex at the center of a thin disk-shape condensate becomes metastable , including the corrections of order $`\alpha ^1`$ and $`\beta ^1`$. ### C Dynamics of a vortex in a cigar-shape condensate $`R_zR_{}`$: Investigation of unstable modes In the opposite limit $`\alpha `$, $`\beta 1`$, the unstable-mode solution of Eqs. (55) and (56) corresponds to exponential growth of the vortex displacement as a function of $`z`$. Such a solution is possible in 3D because the condensate is bounded along the $`z`$ axis. For simplicity, consider an axisymmetric trap, so that $`\alpha =\beta `$. In this case, we have only one equation because one can seek a solution in the form $`x(z)=y(z)`$: $$\stackrel{~}{\omega }x=\frac{2\left(\alpha zx^{}x\right)}{\left(1z^2\right)}\alpha x^{\prime \prime }+\stackrel{~}{\mathrm{\Omega }}x,$$ (66) In the limit $`\alpha 1`$ Eq. (66) has the following approximate solution $$x=y=\epsilon e^{\|z\|/\alpha },$$ (67) $$\stackrel{~}{\omega }=\frac{1}{\alpha }+\stackrel{~}{\mathrm{\Omega }}.$$ (68) The approximate solution (67) has a nonanalytic behavior at $`z=0`$, where a thin boundary layer appears. The actual solution differs from (67) in a small vicinity of $`z=0`$ and represents a smooth crossover from the region $`z<0`$ into the region $`z>0`$. Equation (68) yields the metastable angular velocity $`\mathrm{\Omega }_m`$ in an elongated cigar-shape condensate $$\mathrm{\Omega }_m=\frac{\|q\|\mathrm{}(\omega _x^2+\omega _y^2)}{8\mu }\frac{R_z^2}{R_{}^2}\mathrm{ln}\left(\frac{R_{}}{\|q\|\xi }\right).$$ (69) In contrast to Eq. (65) for a flattened condensate, this expression becomes very large for highly elongated trap. Consequently, it is significantly more difficult to stabilize a vortex in a cigar-shape condensate than in one with a disk shape. ### D Numerical results for 3D We have used Eq. (66) to evaluate the eigenvalues and eigenfunctions for an axisymmetric trap ($`\alpha =\beta `$). A finite trap rotation produces only a shift of eigenvalues by $`\stackrel{~}{\mathrm{\Omega }}`$, so we can set $`\stackrel{~}{\mathrm{\Omega }}=0`$. In addition, solutions of Eq. (66) can be classified as even or odd functions of $`z`$. One can enumerate the solutions by the number $`m`$ of times that the vortex line crosses $`z`$ axis, $`m=0,1,2,\mathrm{}`$. The lowest (most negative) eigenvalue corresponds to $`m=0`$. In Fig. 2 we plot the angular velocity of the vortex precession $`\stackrel{~}{\omega }`$ as a function of the trap anisotropy $`\alpha =R_{}^2/R_z^2`$ for $`m=0,1,2`$. In appropriate limits, the numerical solution $`\stackrel{~}{\omega }_0`$ coincides with those found analytically: $`\stackrel{~}{\omega }_03\frac{1}{5}\alpha `$ for $`\alpha 1`$, and $`\stackrel{~}{\omega }_01/\alpha `$ for $`\alpha 1`$. The next two solutions $`\stackrel{~}{\omega }_1`$ and $`\stackrel{~}{\omega }_2`$ are proportional to $`\alpha `$ for large $`\alpha `$ and diverge like $`1/\alpha `$ for small $`\alpha `$. If $`\alpha >1`$ only one mode has a negative frequency (namely $`\stackrel{~}{\omega }_0`$). If $`\frac{1}{3}<\alpha <1`$ there are 2 such normal modes; if $`\frac{1}{6}<\alpha <\frac{1}{3}`$ there are 3 modes, etc. (these numerical results coincide with those found analytically). The more elongated the trap, the larger the number of modes with negative frequencies (see also ). This conclusion represents one of our main findings. For a disk-shape condensate, the angular velocity $`\mathrm{\Omega }_m`$ for the onset of metastability is smaller than the thermodynamic critical angular velocity $`\mathrm{\Omega }_c`$, with $`\mathrm{\Omega }_m=\frac{3}{5}\mathrm{\Omega }_c`$ . The situation is completely different for a cigar-shape condensate with $`\alpha =R_{}^2/R_z^2<0.26`$, because $`\mathrm{\Omega }_m`$ then becomes larger than $`\mathrm{\Omega }_c`$. For comparison, Fig. 2 includes the line $`\stackrel{~}{\omega }=\stackrel{~}{\mathrm{\Omega }}_c=5`$. In Fig. 3 we plot the shape of the vortex line for the lowest (unstable) normal mode ($`m=0`$) for different values of trap anisotropy $`\alpha `$. The function $`x(z)`$ is an even function of $`z`$ without nodes and $`\epsilon =x(z=0)`$. In Fig. 4 and 5 we plot the shape of the vortex line for normal modes with one and two nodes for different values of trap anisotropy $`\alpha `$. In Fig. 4 $`x_{\mathrm{max}}=\|x(z=R_z)\|`$. Fig. 2. Dimentionless frequencies $`\stackrel{~}{\omega }`$ of the first three normal modes of the vortex as a function of the trap anisotropy $`\alpha =R_{}^2/R_z^2`$. The lower horizontal line represents the dimensionless thermodynamic critical angular velocity. Fig. 3. Shape of the vortex line for the normal mode with the lowest frequency (the most unstable mode) for different values of the trap anisotropy $`\alpha `$. Fig. 4. Shape of the vortex line for the first odd normal mode for different $`\alpha `$. Fig. 5. Shape of the vortex line for the second even normal mode for different $`\alpha `$. ## V Normal modes of a vortex with imaginary frequencies for a nonaxisymmetric 3D condensate For an axisymmetric trap all small-amplitude normal modes have real frequencies. For a nonaxisymmetric trap ($`\alpha \beta `$), however, Eqs. (55), and (56) can have solutions with imaginary frequencies (in certain regimes of trap anisotropy). As an example, we recall that a spherical trap has a special normal mode with $`m=1`$ and $`\stackrel{~}{\omega }=0`$. Let us now consider a nearly spherical trap (this geometry is relevant to the JILA experiments ) $$\|\alpha 1\|1,\|\beta 1\|1.$$ One can rewrite Eqs. (55) and (56) as follows: $$\stackrel{~}{\omega }(1z^2)\left(\begin{array}{c}x\\ y\end{array}\right)=\widehat{H}_0\left(\begin{array}{c}x\\ y\end{array}\right)+\widehat{V}\left(\begin{array}{c}x\\ y\end{array}\right),$$ (70) where $$\widehat{H}_0=\{2+_z[(1z^2)_z]\}\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right),$$ $$\widehat{V}=_z[(1z^2)_z]\left(\begin{array}{cc}0& \beta 1\\ \alpha 1& 0\end{array}\right)+(1z^2)\stackrel{~}{\mathrm{\Omega }}\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right),$$ and $`\widehat{V}`$ is a small perturbation. The unperturbed equation corresponds to the equation for the normal modes of a vortex in a spherical nonrotating trap and, therefore, all eigenfrequencies of the unperturbed equation are real. The eigenvalue $`\stackrel{~}{\omega }=0`$ of the unperturbed equation is degenerate, and there are two solutions that correspond to $`\stackrel{~}{\omega }=0`$: $$\left(\begin{array}{c}x_1\\ y_1\end{array}\right)=\left(\begin{array}{c}z\\ 0\end{array}\right),\left(\begin{array}{c}x_2\\ y_2\end{array}\right)=\left(\begin{array}{c}0\\ z\end{array}\right).$$ (71) The eigenfunctions of $`\widehat{H}_0`$ are Legendre polynomials, and these eigenfunctions form a complete basis. Therefore, one can apply the usual perturbation theory to solve Eq. (70). The matrix elements are given by: $$V_{11}=V_{22}=0,$$ $$V_{12}=\frac{4}{3}\left(\beta 1+\frac{1}{5}\stackrel{~}{\mathrm{\Omega }}\right),$$ $$V_{21}=\frac{4}{3}\left(\alpha 1+\frac{1}{5}\stackrel{~}{\mathrm{\Omega }}\right).$$ To first order in the perturbation, the eigenfrequencies have the form $$\stackrel{~}{\omega }=\pm \frac{\sqrt{V_{12}V_{21}}}{_1^1(1z^2)z^2𝑑z}=\pm 5\sqrt{\left(\alpha 1+\frac{1}{5}\stackrel{~}{\mathrm{\Omega }}\right)\left(\beta 1+\frac{1}{5}\stackrel{~}{\mathrm{\Omega }}\right)}.$$ (72) If, for example, $`\alpha <1<\beta `$ (namely $`R_x<R_z<R_y`$), then the solutions have imaginary frequencies for $`\stackrel{~}{\mathrm{\Omega }}<5(1\alpha )`$: $$\stackrel{~}{\omega }=\pm i\gamma ,\gamma =5\sqrt{\left(1\alpha \frac{1}{5}\stackrel{~}{\mathrm{\Omega }}\right)\left(\beta 1+\frac{1}{5}\stackrel{~}{\mathrm{\Omega }}\right)}>0.$$ (73) Here, the imaginary frequency means a vortex line oriented along the $`z`$ axis (which is the intermediate principal axis in our example) corresponds to an unstable equilibrium. In contrast, a vortex line oriented along the other principal axes is stable. If the angular velocity of the trap rotation increases, however, then the solution (72) becomes real for $`\stackrel{~}{\mathrm{\Omega }}>5(1\alpha )`$ and the vortex line along the $`z`$ axis becomes stable because of the rotation. It is straightforward to consider a general anisotropic trap (not necessarily close to a spherical shape). The result is the following: if the parameters $`\alpha `$ and $`\beta `$ satisfy the inequality: $$\alpha <\frac{2}{n(n+1)}<\beta ,$$ (74) where $`n`$ is a non-negative integer, then there is a normal-mode solution with imaginary frequency that corresponds to the $`n`$th Legendre polynomial. Moreover, if $$\alpha <\frac{2}{n(n+1)}<\frac{2}{m(m+1)}<\beta ,$$ (75) then there are $`(nm+1)`$ solutions with imaginary frequencies. Increasing the external trap rotation sequentially eliminates such solutions. ## VI Motion of a straight vortex line in a slightly nonspherical 3D condensate In the previous sections, we studied the motion of the vortex line for small displacements of the vortex from equilibrium position (normal modes). In this section we solve the general nonlinear equation of the vortex dynamics (38) for a slightly nonspherical trap. This problem is directly related to a recent JILA experiment involving the evolution of an initially straight vortex line in a nearly spherical condensate . In practice, the trap slightly deviates from the spherical shape ($`R_xR_yR_z`$). For a strictly spherical trap, Eq. (38) has a solution representing a motionless straight vortex line ($`\widehat{t}V_{\mathrm{tr}}`$ ) that passes through the center of the trap, with the shape $$x=\gamma _xs,y=\gamma _ys,z=\gamma _zs.$$ (76) Here $`s`$ is the length measured along the vortex line starting from the trap center and $`\gamma _x`$, $`\gamma _y`$, $`\gamma _z`$ are the direction cosines of the angles between the vortex line and the principal axes $`x`$, $`y`$, $`z`$ \[so that $`\gamma _x^2+\gamma _y^2+\gamma _z^2=1`$, $`\widehat{t}=(\gamma _x,\gamma _y,\gamma _z)`$\]. For a slightly anisotropic trap, however, the solution has approximately the same form as (76), but the coefficients $`\gamma _x`$, $`\gamma _y`$, $`\gamma _z`$ become time-dependent. To find a solution to first order in the trap anisotropy, one can omit the curvature of the vortex line and put $`k=0`$ in (38). Also we should use the vortex-free condensate density $`\|\mathrm{\Psi }_{TF}\|^2`$ and take $`R_{}`$ to be equal to the value for a spherical trap. Using the standard perturbation theory, we obtain the following equations for the coefficients $`\gamma _x`$, $`\gamma _y`$, $`\gamma _z`$ (here, we assume that there is no trap rotation): $$\left(\begin{array}{c}\dot{\gamma }_x\\ \\ \dot{\gamma }_y\\ \\ \dot{\gamma }_z\end{array}\right)=\frac{q\mathrm{}}{2\mu }\frac{_R^Rs^2𝑑s}{_R^Rs^2\left(1\frac{s^2}{R^2}\right)𝑑s}\mathrm{ln}\left(\frac{R}{\|q\|\xi }\right)\left(\begin{array}{c}\gamma _y\gamma _z(\omega _z^2\omega _y^2)\\ \\ \gamma _x\gamma _z(\omega _x^2\omega _z^2)\\ \\ \gamma _x\gamma _y(\omega _y^2\omega _x^2)\end{array}\right),$$ (77) or, evaluating the integrals, $$\dot{\gamma }_x=\frac{5q\mathrm{}}{4\mu }\mathrm{ln}\left(\frac{R}{\|q\|\xi }\right)(\omega _z^2\omega _y^2)\gamma _y\gamma _z,$$ (78) $$\dot{\gamma }_y=\frac{5q\mathrm{}}{4\mu }\mathrm{ln}\left(\frac{R}{\|q\|\xi }\right)(\omega _x^2\omega _z^2)\gamma _x\gamma _z,$$ (79) $$\dot{\gamma }_z=\frac{5q\mathrm{}}{4\mu }\mathrm{ln}\left(\frac{R}{\|q\|\xi }\right)(\omega _y^2\omega _x^2)\gamma _x\gamma _y.$$ (80) Equations of this type are known in classical mechanics as Euler’s equations of rigid-body motion; they describe the evolution of the angular velocity of free motion of a body with different principal moments of inertia as seen in the body-fixed frame . Equations (78)-(80) have the following integrals of motion: $$\gamma _x^2+\gamma _y^2+\gamma _z^2=1,$$ (81) $$\omega _x^2\gamma _x^2+\omega _y^2\gamma _y^2+\omega _z^2\gamma _z^2=const.$$ (82) That is, the ends of the straight vortex move along trajectories that correspond to the intersection of a sphere and an ellipsoid with principal axes proportional to $`R_x`$, $`R_y`$, $`R_z`$. In fact, Eq. (82) is the equation of energy conservation during the vortex motion. In the particular case of an axisymmetric trap (for example, $`\omega _x=\omega _y=\omega _{}`$), Eqs. (78)-(80) have the following solution $$\gamma _z=\gamma _z(0)=const,$$ (83) $$\gamma _x=\gamma _x(0)\mathrm{cos}(\omega t)+\gamma _y(0)\mathrm{sin}(\omega t),$$ (84) $$\gamma _y=\gamma _y(0)\mathrm{cos}(\omega t)\gamma _x(0)\mathrm{sin}(\omega t),$$ (85) where $$\omega =\frac{5q\mathrm{}(\omega _z^2\omega _{}^2)}{4\mu }\gamma _z(0)\mathrm{ln}\left(\frac{R}{\|q\|\xi }\right)=\frac{5q\mathrm{}}{2M}\left(\frac{1}{R_z^2}\frac{1}{R_{}^2}\right)\gamma _z(0)\mathrm{ln}\left(\frac{R}{\|q\|\xi }\right),$$ (86) and \[$`\gamma _x(0),\gamma _y(0),\gamma _z(0)`$\] fixes the initial orientation of the vortex line. The line precesses around the $`z`$ axis (the axis of symmetry) at a fixed angle of inclination with respect to the $`z`$ axis. The frequency of this precession depends on the inclination and vanishes if the vortex line is perpendicular to the $`z`$ axis with $`\gamma _z(0)=0`$. We now consider a general nonaxisymmetric trap. To be specific, we assume that $`\omega _x>\omega _y>\omega _z`$ and introduce new scaled functions $$\delta _x=\frac{5q\mathrm{}}{4\mu }\mathrm{ln}\left(\frac{R}{\|q\|\xi }\right)\sqrt{(\omega _x^2\omega _z^2)(\omega _x^2\omega _y^2)}\gamma _x,$$ (87) $$\delta _y=\frac{5q\mathrm{}}{4\mu }\mathrm{ln}\left(\frac{R}{\|q\|\xi }\right)\sqrt{(\omega _y^2\omega _z^2)(\omega _x^2\omega _y^2)}\gamma _y,$$ (88) $$\delta _z=\frac{5q\mathrm{}}{4\mu }\mathrm{ln}\left(\frac{R}{\|q\|\xi }\right)\sqrt{(\omega _y^2\omega _z^2)(\omega _x^2\omega _z^2)}\gamma _z.$$ (89) Then one can rewrite Eqs. (78)-(80) as follows: $$\dot{\delta }_x=\delta _y\delta _z,$$ (90) $$\dot{\delta }_y=\delta _x\delta _z,$$ (91) $$\dot{\delta }_z=\delta _x\delta _y.$$ (92) These equations have the following property: if $`\delta _x`$, $`\delta _y`$, $`\delta _z`$ is a solution of these equations, then if we change sign of any two functions (for example, $`\delta _x\delta _x`$, $`\delta _y\delta _y`$, $`\delta _z\delta _z`$), we obtain another solution. This property can serve to construct solutions that satisfy specific initial conditions. Equations (90)-(92) have three stationary solutions. Two of them (the vortex line parallel to the $`x`$ or $`z`$ axis) correspond to a stable equilibrium, while the third solution (the vortex parallel to the $`y`$ axis) is an unstable equilibrium. If $`\|\delta _x(0)\|<\|\delta _z(0)\|`$, the vortex line oscillates around $`z`$ axis (this is one of the equilibrium orientations), and one can express the solution of Eqs. (90)-(92) in terms of Jacobian elliptic functions as follows: $$\delta _x=\sqrt{\delta _x^2(0)+\delta _y^2(0)}\mathrm{cn}(\sqrt{\delta _z^2(0)+\delta _y^2(0)}t+C,k),$$ (93) $$\delta _y=\pm \sqrt{\delta _x^2(0)+\delta _y^2(0)}\mathrm{sn}(\sqrt{\delta _z^2(0)+\delta _y^2(0)}t+C,k),$$ (94) $$\delta _z=\pm \sqrt{\delta _z^2(0)+\delta _y^2(0)}\mathrm{dn}(\sqrt{\delta _z^2(0)+\delta _y^2(0)}t+C,k),$$ (95) where the modulus $`k=\sqrt{\delta _x^2(0)+\delta _y^2(0)}/\sqrt{\delta _z^2(0)+\delta _y^2(0)}`$ is less than one, and $`C`$ is a constant that must be chosen to satisfy the initial conditions. The solution is a periodic function of time with the period $$T=\frac{4}{\sqrt{\delta _z^2(0)+\delta _y^2(0)}}_0^{\pi /2}\frac{d\phi }{\sqrt{1k^2\mathrm{sin}^2\phi }}=\frac{4}{\sqrt{\delta _z^2(0)+\delta _y^2(0)}}K(k),$$ (96) where $`K(k)`$ is the complete elliptical integral of the first kind. The solution (93)-(95) represents a superposition of a nonuniform circular motion in a plane perpendicular to the $`z`$ axis and oscillations along the $`z`$ axis: $`\delta _z`$ oscillates within the following segment: $$\sqrt{\delta _z^2(0)\delta _x^2(0)}\|\delta _z\|\sqrt{\delta _z^2(0)+\delta _y^2(0)}$$ (97) If $`\|\delta _x(0)\|>\|\delta _z(0)\|`$ (the modulus $`k`$ is greater than one), the vortex line oscillates around $`x`$ axis. In this case one can use reciprocal modulus transformation ($`k\mathrm{sn}(u,k)=\mathrm{sn}(ku,1/k)`$, $`\mathrm{cn}(u,k)=\mathrm{dn}(ku,1/k)`$, $`\mathrm{dn}(u,k)=\mathrm{cn}(ku,1/k)`$) and rewrite the solution (93)-(95) as follows $$\delta _x=\pm \sqrt{\delta _x^2(0)+\delta _y^2(0)}\mathrm{dn}(\sqrt{\delta _x^2(0)+\delta _y^2(0)}t+\stackrel{~}{C},1/k),$$ (98) $$\delta _y=\pm \sqrt{\delta _z^2(0)+\delta _y^2(0)}\mathrm{sn}(\sqrt{\delta _x^2(0)+\delta _y^2(0)}t+\stackrel{~}{C},1/k),$$ (99) $$\delta _z=\sqrt{\delta _z^2(0)+\delta _y^2(0)}\mathrm{cn}(\sqrt{\delta _x^2(0)+\delta _y^2(0)}t+\stackrel{~}{C},1/k).$$ (100) The solution is a periodic function of time with the period $$T=\frac{4}{\sqrt{\delta _z^2(0)+\delta _y^2(0)}}_0^{\pi /2}\frac{d\phi }{\sqrt{k^2\mathrm{sin}^2\phi }}=\frac{4}{\sqrt{\delta _x^2(0)+\delta _y^2(0)}}K\left(1/k\right).$$ (101) If $`\|\delta _x(0)\|=\|\delta _z(0)\|`$ ($`k=1`$), the solution reduces to $$\delta _x=\pm \delta _z=\frac{\sqrt{\delta _x^2(0)+\delta _y^2(0)}}{\mathrm{cosh}\left(\sqrt{\delta _x^2(0)+\delta _y^2(0)}t+C\right)},$$ (102) $$\delta _y=\pm \sqrt{\delta _x^2(0)+\delta _y^2(0)}\mathrm{tanh}\left(\sqrt{\delta _x^2(0)+\delta _y^2(0)}t+C\right);$$ (103) during the motion, the vortex remains on a plane through the $`y`$ axis, oriented along $`\delta _x=\pm \delta _z`$ (see Fig. 6); it eventually lines up along the $`y`$ axis (which is a direction of unstable equilibrium for the geometry $`\omega _x>\omega _y>\omega _z`$ that we are considering). Finally, one can rewrite these solutions directly in terms of the parameters $`\gamma _x`$, $`\gamma _y`$, $`\gamma _z`$ that describe the orientation of the vortex line. For example, instead of (93)-(96), we have $$\gamma _x=\sqrt{\gamma _x^2(0)+\frac{(\omega _y^2\omega _z^2)}{(\omega _x^2\omega _z^2)}\gamma _y^2(0)}\mathrm{cn}(\omega t+C,k),$$ (104) $$\gamma _y=\pm \sqrt{\gamma _y^2(0)+\frac{(\omega _x^2\omega _z^2)}{(\omega _y^2\omega _z^2)}\gamma _y^2(0)}\mathrm{sn}(\omega t+C,k),$$ (105) $$\gamma _z=\pm \sqrt{\gamma _z^2(0)+\frac{(\omega _x^2\omega _y^2)}{(\omega _x^2\omega _z^2)}\gamma _y^2(0)}\mathrm{dn}(\omega t+C,k),$$ (106) where $$k=\sqrt{\frac{(\omega _x^2\omega _y^2)}{(\omega _y^2\omega _z^2)}\frac{\left[(\omega _x^2\omega _z^2)\gamma _x^2\left(0\right)+(\omega _y^2\omega _z^2)\gamma _y^2\left(0\right)\right]}{\left[(\omega _x^2\omega _z^2)\gamma _z^2\left(0\right)+(\omega _x^2\omega _y^2)\gamma _y^2\left(0\right)\right]}},$$ (107) $$\omega =\frac{5q\mathrm{}\sqrt{\omega _y^2\omega _z^2}}{4\mu }\sqrt{(\omega _x^2\omega _z^2)\gamma _z^2\left(0\right)+(\omega _x^2\omega _y^2)\gamma _y^2\left(0\right)}\mathrm{ln}\left(\frac{R}{\|q\|\xi }\right)$$ $$=\frac{5q\mathrm{}}{2M}\sqrt{\frac{1}{R_y^2}\frac{1}{R_z^2}}\sqrt{\left(\frac{1}{R_x^2}\frac{1}{R_z^2}\right)\gamma _z^2\left(0\right)+\left(\frac{1}{R_x^2}\frac{1}{R_y^2}\right)\gamma _y^2\left(0\right)}\mathrm{ln}\left(\frac{R}{\|q\|\xi }\right).$$ (108) The period of the motion is given by $$T=\frac{4}{\omega }_0^{\pi /2}\frac{d\phi }{\sqrt{1k^2\mathrm{sin}^2\phi }}=\frac{4}{\omega }K(k).$$ (109) For a nonaxisymmetric trap we plot trajectories of the end of the vortex line in Fig. 6. The plot corresponds to the geometry $`R_x<R_y<R_z`$. There are two stable orientations (along the $`x`$ and $`z`$ axes) and one unstable (along the $`y`$ axis). The shape of the trajectories strongly depends on the initial orientation of the vortex. Fig. 6. Typical trajectories of the end of a straight vortex line (that passes through the condensate center) during its motion in a slightly nonspherical trap with $`R_x<R_y<R_z`$. For an axisymmetric trap, we have $`\omega _x=\omega _y`$ so that $`k=0`$. Then, using the properties $`\mathrm{sn}(z,0)=\mathrm{sin}z`$, $`\mathrm{cn}(z,0)=\mathrm{cos}z`$, $`\mathrm{dn}(z,0)=1`$, we can reproduce formulas (83)-(86). Because the motion is periodic, one can anticipate one or more revivals of the vortex image in the recent JILA experiments. The revival time depends on the orientation of the vortex with respect to the symmetry axes and the deviations from sphericity. Recently, such revivals of the vortex image (as well as vortex precession ) was seen by the JILA group, and their results agree quantitatively with our theory . One should note that our analytical results for the normal-mode frequencies are valid with the logarithmic accuracy, namely when $`\mathrm{ln}(R/\|q\|\xi )1`$. It is, however, straightforward to go beyond logarithmic accuracy and obtain the numerical correction to the logarithm (see Ref. ). This correction modifies our formulas for the normal-mode frequencies as follows: instead of $`\mathrm{ln}(R/\|q\|\xi )`$, one should use $`\mathrm{ln}(R/\|q\|\xi )+0.675=\mathrm{ln}(1.96R/\|q\|\xi )`$. In the JILA experiments, $`\mathrm{ln}(R/\xi )3.5`$, and the inclusion of the correction improves the quantitative agreement with the experimental observations. ## VII Conclusions In this paper, we consider the dynamics of a vortex line in 2D and 3D condensates in the TF limit. We took into account the nonuniform nature of the system (namely the trap potential), the vortex curvature and a possible trap rotation. We derived a general equation of vortex dynamics and investigated various normal modes of the vortex line. For an axisymmetric trap, all eigenvalues are real and the motion of the vortex line can be represented as a superposition of planar normal modes with different frequencies. In a 2D condensate, the normal modes are degenerate, and, as a result, a superposition of planar waves can produce helical waves. In 3D, there is no such degeneracy and there is no simple analog of the helical waves. An externally applied trap rotation $`\mathrm{\Omega }`$ shifts the normal-mode frequencies and makes the vortex locally stable for sufficiently large $`\mathrm{\Omega }`$. For a cigar-shape condensate, the vortex curvature has a significant effect on the frequency of the most unstable normal mode (that with the most negative frequency), and additional modes with (less) negative frequencies appear. As a result, it is more difficult to stabilize central vortex in a cigar-shape condensate than in a disc-shape one. Normal modes with imaginary frequencies can exist for a nonaxisymmetric condensate (both in 2D and 3D). This means that the corresponding equilibrium orientation of the vortex is unstable. As an example of the solution of the general nonlinear problem of vortex dynamics, we considered the motion of a straight vortex line in a slightly nonspherical condensate. The vortex line changes its orientation in space at a rate proportional to the trap anisotropy. ###### Acknowledgements. This work was supported in part by the National Science Foundation, Grant No. DMR 99-71518, and by Stanford University (A.A.S.). We are grateful to B. Anderson, J. Anglin, E. Cornell and D. Feder for valuable correspondence and discussions. This work benefitted from our participation in recent workshops at the Lorentz Center, Leiden, The Netherlands and at ECT European Centre for Theoretical Studies in Nuclear Physics and Related Areas, Trento, Italy; we thank H. Stoof and S. Stringari for organizing these workshops and for their hospitality.
warning/0007/hep-th0007068.html	ar5iv	text	# 1 Introduction ## 1 Introduction The result that this paper presents is about WZ couplings for generalized sigma-orbifold fixed-points. The usual orbifold fixed-points do not involve any kind of gauge fields, the corresponding GS 6-form that enconding the complete gauge-gravitational anomaly and its opposite inflow only involves the standard gauge field of the D-branes. From the other side, the sigma-orbifold fixed-points have one sigma-gauge field and the corresponding GS 6-form that enconding the complete sigma-gauge-gravitational anomaly and its opposite inflow involves both the standard gauge field of the D-branes and such sigma-gauge field. The generalized sigma-orbifold fixed-points that this paper considers have two gauge fields: the sigma-gauge field and one non-standard SO(2n)-gauge field, and then, the corresponding generalized GS 6-form that encoding the complete sigma-standard gauge-gravitational-non standard gauge anomaly and its opposite inflow, involves three gauge fields: the standard gauge field of the D-branes, the sigma-gauge field and the non-standard SO(2n)-gauge field. The aim of the present paper is to display the Wess-Zumino couplings to the RR forms for such generalized sigma-orbifold fixed-points. For the usual orbifold fixed-points the Wess-Zumino couplings have the following form, which can be derived both from index theorems and from direct string computation by factorization of the one-loop partition functions: (Scrucca and Serone,hep-th/9912108) $`𝑺_{𝑭𝒌}^{\mathbf{(}\mathrm{𝟐}𝒌\mathbf{)}}\mathbf{=}\sqrt{\mathrm{𝟐}𝝅}\mathbf{}_{𝒊_𝒌\mathbf{=}\mathrm{𝟏}}^{𝑵_𝒌}𝑪_{\mathbf{(}\mathrm{𝟐}𝒌\mathbf{)}}^{𝒊_𝒌}\mathbf{}𝒁_{\mathbf{(}\mathrm{𝟐}𝒌\mathbf{)}}`$ Where the Mukai vector of RR charges for the usual orbifold fixed-points is given by: $`𝒁_{\mathbf{(}\mathrm{𝟐}𝒌\mathbf{)}}\mathbf{(}\frac{𝐑}{\mathrm{𝟒}}\mathbf{)}\mathbf{=}\mathbf{}\frac{\mathrm{𝟒}}{\sqrt{𝐍}}\mathit{ϵ}_𝒌\sqrt{\frac{\mathbf{\|}𝐜_𝐤^\mathrm{𝟏}𝐜_𝐤^\mathrm{𝟐}𝐜_𝐤^\mathrm{𝟑}\mathbf{\|}}{\mathbf{\|}𝐬_𝐤^\mathrm{𝟏}𝐬_𝐤^\mathrm{𝟐}𝐬_𝐤^\mathrm{𝟑}\mathbf{\|}}}\sqrt{𝑳\mathbf{(}\frac{𝐑}{\mathrm{𝟒}}\mathbf{)}}`$ In these formulaes C is the vector of the RR potential forms. L is the Hirzebruch genus that generates the Hirzebruch polynomials which are given in terms of Pontryaguin classes for real bundles. The Pontryaguin classes are given in terms of the 2-form curvature of the corresponding real bundle. The formula for Z involves only the gravitational curvature. For each k-twisted sector, s and c are the vectors of sines and cosines respectively of the twist vector v, $`ϵ`$ is the sign of the product of the components of the vector s. $`N_k`$ is the number of fixed points. Here the relevant group is $`Z_N`$. From the other side,for the usual sigma-orbifold fixed-points the Wess-Zumino couplings have the following form, which can be derived both from equivariant index theorems and from direct string computation by factorization of the one-loop partition functions: (Scrucca and Serone,hep-th/0006201) $`𝑺_𝑭\mathbf{=}\sqrt{\mathrm{𝟐}𝝅}\mathbf{}_{𝒌\mathbf{=}\mathrm{𝟏}}^{\frac{\mathrm{𝟏}}{\mathrm{𝟐}}\mathbf{(}𝑵\mathbf{}\mathrm{𝟏}\mathbf{)}}\mathbf{}_{𝒊_𝒌\mathbf{=}\mathrm{𝟏}}^{𝑵_𝒌}\mathbf{}𝑪_{\mathbf{(}\mathrm{𝟐}𝒌\mathbf{)}}^{𝒊_𝒌}\mathbf{}𝒁_{\mathbf{(}\mathrm{𝟐}𝒌\mathbf{)}}`$ where, $`𝒁_{\mathbf{(}\mathrm{𝟐}𝒌\mathbf{)}}\mathbf{(}\frac{𝐑}{\mathrm{𝟒}}\mathbf{,}\frac{𝐆}{\mathrm{𝟒}}\mathbf{)}\mathbf{=}\mathbf{}\frac{\mathrm{𝟒}}{\sqrt{𝐍}}\mathit{ϵ}_𝒌\sqrt{\frac{\mathbf{\|}𝐂_{\mathrm{𝟐}𝐤}\mathbf{\|}}{\mathbf{\|}𝐂_𝐤^\mathrm{𝟐}\mathbf{\|}}}\sqrt{𝑳_𝒌\mathbf{(}\frac{𝐆}{\mathrm{𝟒}}\mathit{ϵ}_{\mathrm{𝟐}𝒌}\mathbf{)}}\sqrt{𝑳\mathbf{(}\frac{𝐑}{\mathrm{𝟒}}\mathbf{)}}`$ For this case, $`L_k`$ is the Hirzebruch equivariant factor, G is the sigma gauge field and $`C_k`$ is the product of the components of the vector s. In this paper is presented the formula for the WZ couplings of the generalized sigma-orbifold fixed-points which also have one aditional SO(2n) Yang-Mills gauge field . Such WZ coupling is given by the following formula: $`𝑺_𝑭\mathbf{=}\sqrt{\mathrm{𝟐}𝝅}\mathbf{}_{𝒌\mathbf{=}\mathrm{𝟏}}^{\frac{\mathrm{𝟏}}{\mathrm{𝟐}}\mathbf{(}𝑵\mathbf{}\mathrm{𝟏}\mathbf{)}}\mathbf{}_{𝒊_𝒌\mathbf{=}\mathrm{𝟏}}^{𝑵_𝒌}\mathbf{}𝑪_{\mathbf{(}\mathrm{𝟐}𝒌\mathbf{)}}^{𝒊_𝒌}\mathbf{}𝒁_{\mathbf{(}\mathrm{𝟐}𝒌\mathbf{)}}`$ where, $`𝒁_{\mathbf{(}\mathrm{𝟐}𝒌\mathbf{)}}\mathbf{(}𝑹\mathbf{,}𝑮\mathbf{,}𝒀\mathbf{)}\mathbf{=}\mathbf{}\frac{\mathrm{𝟒}}{\sqrt{𝐍}}\mathit{ϵ}_𝒌\sqrt{\frac{\mathbf{\|}𝐂_{\mathrm{𝟐}𝐤}\mathbf{\|}}{\mathbf{\|}𝐂_𝐤^\mathrm{𝟐}\mathbf{\|}}}\sqrt{\mathbf{}_{𝒊\mathbf{=}\mathrm{𝟏}}^\mathrm{𝟑}\frac{\mathrm{𝐬𝐢𝐧}\mathbf{(}\mathrm{𝐩𝐢𝐤𝐯}_𝐢\mathbf{)}}{\mathrm{𝐬𝐢𝐧}\mathbf{(}\mathrm{𝐩𝐢𝐤𝐯}_𝐢\mathbf{+}\frac{𝐆_𝐢}{\mathrm{𝟒}𝐩𝐢}\mathbf{)}}}\sqrt{𝑨\mathbf{(}\frac{𝐑}{\mathrm{𝟐}}\mathbf{)}\mathbf{}_{𝒂\mathbf{=}\mathrm{𝟏}}^𝒏\frac{\mathrm{𝐜𝐨𝐬}\mathbf{(}\mathrm{𝐩𝐢𝐤𝐮}_𝐚\mathbf{+}\frac{𝐘_𝐚}{\mathrm{𝟒}𝐩𝐢}\mathbf{)}}{\mathrm{𝐜𝐨𝐬}\mathbf{(}\mathrm{𝐩𝐢𝐤𝐮}_𝐚\mathbf{)}}}`$ In this last formula, the productory that involves the sines is the equivariant Dirac-roof factor, the productory that involves the cosines is the equivariant Mayer class, v is the twist vector corresponding to the action of the group $`Z_N`$ over the normal bundle with respect to space-time of the fixed submanifold, u is the twist vector corresponding to the action of the group $`Z_N`$ over the SO(2n) bundle on the space-time, Y is the vector of eigenvalues of the 2-form curvature of the SO(2n) bundle, A is the Dirac-roof genus for the tangent bundle of the fixed submanifold. When the SO(2n) bundle over the space-time is the tangent bundle of the space-time, one obtains the usual formula for Z corresponding to the usual sigma-orbifold fixed-points. For this, the equivariant Mayer class for the tangent bundle of the space-time is factorized as the product of the usual Mayer class for the tangent bundle of the fixed submanifold with the equivariant Mayer factor for the normal bundle of the fixed submanifold. In such case, the vector u is reduced to the vector v and the vector Y is reduced to the sum of the vector of eigenvalues for the 2-form curvature of the tangent bundle of fixed manifold, with the vector G. Finally one can to use the following identity: $`𝑨\mathbf{(}\frac{𝐑}{\mathrm{𝟐}}\mathbf{)}𝑴𝒂𝒚𝒆𝒓\mathbf{(}\frac{𝐑}{\mathrm{𝟐}}\mathbf{)}\mathbf{=}𝑳\mathbf{(}\frac{𝐑}{\mathrm{𝟒}}\mathbf{)}`$ In the following section the generalized GS 6-form that enconding the complete sigma-standard gauge-gravitational-non standard gauge anomaly and its opposite inflow, will be given . In the final third section some conclutions are presented. ## 2 The generalized GS anomaly-inflow 6-form Let E be a SO(2n)-bundle over the space-time and consider a formal factorisation for the total Pontryaguin classs of the real bundle E, which has the following form: $`𝒑\mathbf{(}𝑬\mathbf{)}\mathbf{=}\mathbf{}_{𝒂\mathbf{=}\mathrm{𝟏}}^𝒏\mathbf{(}\mathrm{𝟏}\mathbf{+}𝒚_𝒂^\mathrm{𝟐}\mathbf{)}`$ The total Pontryaguin classs of the real bundle E,has the following formal sumarisation in terms of the corresponding Pontryaguin classes: $`𝒑\mathbf{(}𝑬\mathbf{)}\mathbf{=}\mathbf{}_{𝒋\mathbf{=}\mathrm{𝟎}}^{\mathbf{}}𝒑_𝒋\mathbf{(}𝑬\mathbf{)}`$ The total Mayer class for the real bundle E has the following formal factorisation: $`𝑴𝒂𝒚𝒆𝒓\mathbf{(}𝑬\mathbf{)}\mathbf{=}\mathbf{}_{𝒂\mathbf{=}\mathrm{𝟏}}^𝒏𝒄𝒐𝒔𝒉\mathbf{(}\frac{𝐲_𝐚}{\mathrm{𝟐}}\mathbf{)}`$ The total Mayer class for the real bundle E has the following formal sumarisation in terms of the Mayer polynomials which are formed from the corresponding Pontryaguin classes : $`𝑴𝒂𝒚𝒆𝒓\mathbf{(}𝑬\mathbf{)}\mathbf{=}\mathbf{}_{𝒋\mathbf{=}\mathrm{𝟎}}^{\mathbf{}}𝑴𝒂𝒚𝒆𝒓_𝒋\mathbf{(}𝒑_\mathrm{𝟏}\mathbf{(}𝑬\mathbf{)}\mathbf{,}\mathbf{}\mathbf{,}𝒑_𝒋\mathbf{(}𝑬\mathbf{)}\mathbf{)}`$ The Mayer polynomials are given by: $`𝑴𝒂𝒚𝒆𝒓_\mathrm{𝟎}\mathbf{(}𝒑_\mathrm{𝟎}\mathbf{(}𝑬\mathbf{)}\mathbf{)}\mathbf{=}𝑴𝒂𝒚𝒆𝒓_\mathrm{𝟎}\mathbf{(}\mathrm{𝟏}\mathbf{)}\mathbf{=}\mathrm{𝟏}`$ $`𝑴𝒂𝒚𝒆𝒓_\mathrm{𝟏}\mathbf{(}𝒑_\mathrm{𝟏}\mathbf{(}𝑬\mathbf{)}\mathbf{)}\mathbf{=}\frac{𝐩_\mathrm{𝟏}\mathbf{(}𝐄\mathbf{)}}{\mathrm{𝟖}}`$ $`𝑴𝒂𝒚𝒆𝒓_\mathrm{𝟐}\mathbf{(}𝒑_\mathrm{𝟏}\mathbf{(}𝑬\mathbf{)}\mathbf{,}𝒑_\mathrm{𝟐}\mathbf{(}𝑬\mathbf{)}\mathbf{)}\mathbf{=}\frac{𝐩_\mathrm{𝟏}\mathbf{(}𝐄\mathbf{)}^\mathrm{𝟐}\mathbf{+}\mathrm{𝟒}𝐩_\mathrm{𝟐}\mathbf{(}𝐄\mathbf{)}}{\mathrm{𝟑𝟖𝟒}}`$ $`𝑴𝒂𝒚𝒆𝒓_\mathrm{𝟑}\mathbf{(}𝒑_\mathrm{𝟏}\mathbf{(}𝑬\mathbf{)}\mathbf{,}𝒑_\mathrm{𝟐}\mathbf{(}𝑬\mathbf{)}\mathbf{,}𝒑_\mathrm{𝟑}\mathbf{(}𝑬\mathbf{)}\mathbf{)}\mathbf{=}\frac{𝐩_\mathrm{𝟏}\mathbf{(}𝐄\mathbf{)}^\mathrm{𝟑}\mathbf{+}\mathrm{𝟏𝟐}𝐩_\mathrm{𝟏}\mathbf{(}𝐄\mathbf{)}𝐩_\mathrm{𝟐}\mathbf{(}𝐄\mathbf{)}\mathbf{+}\mathrm{𝟒𝟖}𝐩_\mathrm{𝟑}\mathbf{(}𝐄\mathbf{)}}{\mathrm{𝟒𝟔𝟎𝟖𝟎}}`$ The pontryaguin classes of the real bundle E have the following realizations in terms of the powers of the 2-form curvature for such bundle. For this curvature the y’s are the eigenvalues: $`𝒑_\mathrm{𝟏}\mathbf{(}𝑬\mathbf{)}\mathbf{=}𝒑_\mathrm{𝟏}\mathbf{(}𝑹_𝑬\mathbf{)}\mathbf{=}\mathbf{}\frac{\mathrm{𝟏}}{\mathrm{𝟖}𝐩𝐢^\mathrm{𝟐}}𝒕𝒓𝑹_𝑬^\mathrm{𝟐}`$ $`𝒑_\mathrm{𝟐}\mathbf{(}𝑬\mathbf{)}\mathbf{=}𝒑_\mathrm{𝟐}\mathbf{(}𝑹_𝑬\mathbf{)}\mathbf{=}\frac{\mathrm{𝟏}}{\mathrm{𝟏𝟔}𝐩𝐢^\mathrm{𝟒}}\mathbf{[}\frac{\mathrm{𝟏}}{\mathrm{𝟖}}\mathbf{(}𝒕𝒓𝑹_𝑬^\mathrm{𝟐}\mathbf{)}^\mathrm{𝟐}\mathbf{}\frac{\mathrm{𝟏}}{\mathrm{𝟒}}𝒕𝒓𝑹_𝑬^\mathrm{𝟒}\mathbf{]}`$ $`𝒑_\mathrm{𝟑}\mathbf{(}𝑬\mathbf{)}\mathbf{=}𝒑_\mathrm{𝟑}\mathbf{(}𝑹_𝑬\mathbf{)}\mathbf{=}\frac{\mathrm{𝟏}}{\mathrm{𝟔𝟒}𝐩𝐢^\mathrm{𝟔}}\mathbf{[}\mathbf{}\frac{\mathrm{𝟏}}{\mathrm{𝟒𝟖}}\mathbf{(}𝒕𝒓𝑹_𝑬^\mathrm{𝟐}\mathbf{)}^\mathrm{𝟑}\mathbf{}\frac{\mathrm{𝟏}}{\mathrm{𝟔}}𝒕𝒓𝑹_𝑬^\mathrm{𝟔}\mathbf{+}\frac{\mathrm{𝟏}}{\mathrm{𝟖}}𝒕𝒓𝑹_𝑬^\mathrm{𝟐}𝒕𝒓𝑹_𝑬^\mathrm{𝟒}\mathbf{]}`$ Using all these expretions one can to obtain the following expantion: $`𝑴𝒂𝒚𝒆𝒓\mathbf{(}\frac{𝐑_𝐄}{\mathrm{𝟐}}\mathbf{)}\mathbf{=}\mathrm{𝟏}\mathbf{+}\frac{𝐩_\mathrm{𝟏}\mathbf{(}𝐑_𝐄\mathbf{)}}{\mathrm{𝟑𝟐}}\mathbf{+}\frac{𝐩_\mathrm{𝟏}\mathbf{(}𝐑_𝐄\mathbf{)}^\mathrm{𝟐}\mathbf{+}\mathrm{𝟒}𝐩_\mathrm{𝟐}\mathbf{(}𝐑_𝐄\mathbf{)}}{\mathrm{𝟔𝟏𝟒𝟒}}\mathbf{+}\frac{𝐩_\mathrm{𝟏}\mathbf{(}𝐑_𝐄\mathbf{)}^\mathrm{𝟑}\mathbf{+}\mathrm{𝟏𝟐}𝐩_\mathrm{𝟏}\mathbf{(}𝐑_𝐄\mathbf{)}𝐩_\mathrm{𝟐}\mathbf{(}𝐑_𝐄\mathbf{)}\mathbf{+}\mathrm{𝟒𝟖}𝐩_\mathrm{𝟑}\mathbf{(}𝐑_𝐄\mathbf{)}}{\mathrm{𝟐𝟗𝟒𝟗𝟏𝟐𝟎}}\mathbf{+}\mathbf{}`$ Now one has the following expantions: $`𝑨\mathbf{(}\frac{𝐑}{\mathrm{𝟐}}\mathbf{)}\mathbf{=}\mathrm{𝟏}\mathbf{}\frac{𝐩_\mathrm{𝟏}\mathbf{(}𝐑\mathbf{)}}{\mathrm{𝟗𝟔}}\mathbf{+}\frac{\mathrm{𝟕}𝐩_\mathrm{𝟏}\mathbf{(}𝐑\mathbf{)}^\mathrm{𝟐}\mathbf{}\mathrm{𝟒}𝐩_\mathrm{𝟐}\mathbf{(}𝐑\mathbf{)}}{\mathrm{𝟗𝟐𝟏𝟔𝟎}}\mathbf{+}\mathbf{}`$ $`𝑳\mathbf{(}\frac{𝐑}{\mathrm{𝟒}}\mathbf{)}\mathbf{=}\mathrm{𝟏}\mathbf{+}\frac{𝐩_\mathrm{𝟏}\mathbf{(}𝐑\mathbf{)}}{\mathrm{𝟒𝟖}}\mathbf{+}\frac{\mathbf{}𝐩_\mathrm{𝟏}\mathbf{(}𝐑\mathbf{)}^\mathrm{𝟐}\mathbf{+}\mathrm{𝟕}𝐩_\mathrm{𝟐}\mathbf{(}𝐑\mathbf{)}}{\mathrm{𝟏𝟏𝟓𝟐𝟎}}\mathbf{+}\mathbf{}`$ Using these three expantions it is easy to obtain the following identities: $`𝑨\mathbf{(}\frac{𝐑}{\mathrm{𝟐}}\mathbf{)}𝑴𝒂𝒚𝒆𝒓\mathbf{(}\frac{𝐑}{\mathrm{𝟐}}\mathbf{)}\mathbf{=}𝑳\mathbf{(}\frac{𝐑}{\mathrm{𝟒}}\mathbf{)}`$ $`𝑨\mathbf{(}𝑹\mathbf{)}𝑴𝒂𝒚𝒆𝒓\mathbf{(}𝑹\mathbf{)}\mathbf{=}𝑳\mathbf{(}\frac{𝐑}{\mathrm{𝟐}}\mathbf{)}`$ $`𝑨\mathbf{(}\mathrm{𝟐}𝑹\mathbf{)}𝑴𝒂𝒚𝒆𝒓\mathbf{(}\mathrm{𝟐}𝑹\mathbf{)}\mathbf{=}𝑳\mathbf{(}𝑹\mathbf{)}`$ $`𝑨\mathbf{(}\mathrm{𝟐}^𝒒𝑹\mathbf{)}𝑴𝒂𝒚𝒆𝒓\mathbf{(}\mathrm{𝟐}^𝒒𝑹\mathbf{)}\mathbf{=}𝑳\mathbf{(}\mathrm{𝟐}^{𝒒\mathbf{}\mathrm{𝟏}}𝑹\mathbf{)}`$ $`\mathbf{[}𝑨\mathbf{(}𝑹\mathbf{)}\mathrm{𝟐}^𝒌𝑴𝒂𝒚𝒆𝒓\mathbf{(}𝑹\mathbf{)}\mathbf{]}_{𝒕𝒐𝒑𝒇𝒐𝒓𝒎}\mathbf{=}𝑳\mathbf{(}𝑹\mathbf{)}_{𝒕𝒐𝒑𝒇𝒐𝒓𝒎}`$ Now, the group $`Z_N`$ can be thought to act like the automorphism group of the tangent bundle over space-time, of the standard gauge bundle and of the SO(2n)-bundle E. For the normal bundle respect to the space-time of the fixed submanifold, the group $`Z_N`$ acts according to the twist vector v. For the SO(2n)-bundle E over the space-time, the group $`Z_N`$ acts according to the twist vector u. Then, the equivariant Mayer class for E has the following factorization: $`𝑴𝒂𝒚𝒆𝒓_𝒌\mathbf{(}\frac{𝐲}{\mathrm{𝟐}}\mathbf{)}\mathbf{=}\mathbf{}_{𝒂\mathbf{=}\mathrm{𝟏}}^𝒏\frac{\mathrm{𝐜𝐨𝐬𝐡}\mathbf{(}\mathrm{𝐢𝐩𝐢𝐤𝐮}_𝐚\mathbf{+}\frac{𝐲_𝐚}{\mathrm{𝟒}}\mathbf{)}}{\mathrm{𝐜𝐨𝐬𝐡}\mathbf{(}\mathrm{𝐢𝐩𝐢𝐤𝐮}_𝐚\mathbf{)}}`$ When one defines: $`𝒚_𝒂\mathbf{=}\frac{\mathrm{𝐢𝐘}_𝐚}{\mathrm{𝐩𝐢}}`$ one obtains: $`𝑴𝒂𝒚𝒆𝒓_𝒌\mathbf{(}\frac{𝐘}{\mathrm{𝟐}}\mathbf{)}\mathbf{=}\mathbf{}_{𝒂\mathbf{=}\mathrm{𝟏}}^𝒏\frac{\mathrm{𝐜𝐨𝐬}\mathbf{(}\mathrm{𝐩𝐢𝐤𝐮}_𝐚\mathbf{+}\frac{𝐘_𝐚}{\mathrm{𝟒}𝐩𝐢}\mathbf{)}}{\mathrm{𝐜𝐨𝐬}\mathbf{(}\mathrm{𝐩𝐢𝐤𝐮}_𝐚\mathbf{)}}`$ When the ten dimentional space-time is of the form $`X=R^{1,3}xT^6`$ and the fixed submanifold by the action of the group $`Z_N`$ is $`N_k`$ copies of $`R^{1,3}`$ , then the equivariant Dirac-roof factor has the following factorization, where the twist vector v corresponds to the action of the group $`Z_N`$ over the normal bundle with respect to the space-time of the fixed submanifold and G is the sigma gauge field: $`𝑨_𝒌\mathbf{(}\frac{𝐆}{\mathrm{𝟐}}\mathbf{)}\mathbf{=}\mathbf{}_{𝒊\mathbf{=}\mathrm{𝟏}}^\mathrm{𝟑}\frac{\mathrm{𝐬𝐢𝐧}\mathbf{(}\mathrm{𝐩𝐢𝐤𝐯}_𝐢\mathbf{+}\frac{𝐆_𝐢}{\mathrm{𝟒}𝐩𝐢}\mathbf{)}}{\mathrm{𝐬𝐢𝐧}\mathbf{(}\mathrm{𝐩𝐢𝐤𝐯}_𝐢\mathbf{)}}`$ Using all these notations, the Mukay vector of RR charges for the generalized sigma orbifol fixed-points can be writen as follows: $`𝒁_{\mathbf{(}\mathrm{𝟐}𝒌\mathbf{)}}\mathbf{(}𝑹\mathbf{,}𝑮\mathbf{,}𝒀\mathbf{)}\mathbf{=}\mathbf{}\frac{\mathrm{𝟒}}{\sqrt{𝐍}}\mathit{ϵ}_𝒌\sqrt{\frac{\mathbf{\|}𝐂_{\mathrm{𝟐}𝐤}\mathbf{\|}}{\mathbf{\|}𝐂_𝐤^\mathrm{𝟐}\mathbf{\|}}}\sqrt{𝑨_𝒌\mathbf{(}\frac{𝐆}{\mathrm{𝟐}}\mathbf{)}}\sqrt{𝑨\mathbf{(}\frac{𝐑}{\mathrm{𝟐}}\mathbf{)}𝑴𝒂𝒚𝒆𝒓_𝒌\mathbf{(}\frac{𝐘}{\mathrm{𝟐}}\mathbf{)}}`$ For the particular case when the SO(2n) gauge bundle E is the SO(10) tangent bundle of the space-time X, the equivariant Mayer class for E has the following factorization, where x is the vector of eigenvalues of the 2-form curvatur R of the tangent bunlde of the fixed submanifold: $`𝑴𝒂𝒚𝒆𝒓_𝒌\mathbf{(}\frac{𝐑_𝐗}{\mathrm{𝟐}}\mathbf{)}\mathbf{=}\mathbf{}_{𝒂\mathbf{=}\mathrm{𝟏}}^\mathrm{𝟓}\frac{\mathrm{𝐜𝐨𝐬}\mathbf{(}\mathrm{𝐩𝐢𝐤𝐮}_𝐚\mathbf{+}\frac{𝐘_𝐚}{\mathrm{𝟒}𝐩𝐢}\mathbf{)}}{\mathrm{𝐜𝐨𝐬}\mathbf{(}\mathrm{𝐩𝐢𝐤𝐮}_𝐚\mathbf{)}}\mathbf{=}\mathbf{(}\mathbf{}_{𝒊\mathbf{=}\mathrm{𝟏}}^\mathrm{𝟑}\frac{\mathrm{𝐜𝐨𝐬}\mathbf{(}\mathrm{𝐩𝐢𝐤𝐯}_𝐢\mathbf{+}\frac{𝐆_𝐢}{\mathrm{𝟒}𝐩𝐢}\mathbf{)}}{\mathrm{𝐜𝐨𝐬}\mathbf{(}\mathrm{𝐩𝐢𝐤𝐯}_𝐢\mathbf{)}}\mathbf{)}\mathbf{(}\mathbf{}_{𝒋\mathbf{=}\mathrm{𝟏}}^\mathrm{𝟐}\mathrm{𝐜𝐨𝐬}\mathbf{(}\frac{\mathrm{𝐢𝐩𝐢𝐱}_𝐣}{\mathrm{𝟒}𝐩𝐢}\mathbf{)}\mathbf{)}\mathbf{=}𝑴𝒂𝒚𝒆𝒓_𝒌\mathbf{(}\frac{𝐆}{\mathrm{𝟐}}\mathbf{)}𝑴𝒂𝒚𝒆𝒓\mathbf{(}\frac{𝐑}{\mathrm{𝟐}}\mathbf{)}`$ Using this last factorization it is easy to obtain the usual formula for the Mukay vector of the usual sigma orbifold fixed-points: $`𝒁_{\mathbf{(}\mathrm{𝟐}𝒌\mathbf{)}}\mathbf{(}𝑹\mathbf{,}𝑮\mathbf{,}𝒀\mathbf{)}\mathbf{=}\mathbf{}\frac{\mathrm{𝟒}}{\sqrt{𝐍}}\mathit{ϵ}_𝒌\sqrt{\frac{\mathbf{\|}𝐂_{\mathrm{𝟐}𝐤}\mathbf{\|}}{\mathbf{\|}𝐂_𝐤^\mathrm{𝟐}\mathbf{\|}}}\sqrt{𝑨_𝒌\mathbf{(}\frac{𝐆}{\mathrm{𝟐}}\mathbf{)}}\sqrt{𝑨\mathbf{(}\frac{𝐑}{\mathrm{𝟐}}\mathbf{)}𝑴𝒂𝒚𝒆𝒓_𝒌\mathbf{(}\frac{𝐆}{\mathrm{𝟐}}\mathbf{)}𝑴𝒂𝒚𝒆𝒓\mathbf{(}\frac{𝐑}{\mathrm{𝟐}}\mathbf{)}}\mathbf{=}\mathbf{}\frac{\mathrm{𝟒}}{\sqrt{𝐍}}\mathit{ϵ}_𝒌\sqrt{\frac{\mathbf{\|}𝐂_{\mathrm{𝟐}𝐤}\mathbf{\|}}{\mathbf{\|}𝐂_𝐤^\mathrm{𝟐}\mathbf{\|}}}\sqrt{𝑨_𝒌\mathbf{(}\frac{𝐆}{\mathrm{𝟐}}\mathbf{)}𝑴𝒂𝒚𝒆𝒓_𝒌\mathbf{(}\frac{𝐆}{\mathrm{𝟐}}\mathbf{)}}\sqrt{𝑨\mathbf{(}\frac{𝐑}{\mathrm{𝟐}}\mathbf{)}𝑴𝒂𝒚𝒆𝒓\mathbf{(}\frac{𝐑}{\mathrm{𝟐}}\mathbf{)}}\mathbf{=}\mathbf{}\frac{\mathrm{𝟒}}{\sqrt{𝐍}}\mathit{ϵ}_𝒌\sqrt{\frac{\mathbf{\|}𝐂_{\mathrm{𝟐}𝐤}\mathbf{\|}}{\mathbf{\|}𝐂_𝐤^\mathrm{𝟐}\mathbf{\|}}}\sqrt{𝑳_𝒌\mathbf{(}\frac{𝐆}{\mathrm{𝟒}}\mathbf{)}}\sqrt{𝑳\mathbf{(}\frac{𝐑}{\mathrm{𝟒}}\mathbf{)}}\mathbf{=}𝒁_{\mathbf{(}\mathrm{𝟐}𝒌\mathbf{)}}\mathbf{(}\frac{𝐑}{\mathrm{𝟒}}\mathbf{,}\frac{𝐆}{\mathrm{𝟒}}\mathbf{)}`$ Here are used the following identities: $`𝑨\mathbf{(}\frac{𝐑}{\mathrm{𝟐}}\mathbf{)}𝑴𝒂𝒚𝒆𝒓\mathbf{(}\frac{𝐑}{\mathrm{𝟐}}\mathbf{)}\mathbf{=}𝑳\mathbf{(}\frac{𝐑}{\mathrm{𝟒}}\mathbf{)}`$ $`𝑨_𝒌\mathbf{(}\frac{𝐆}{\mathrm{𝟐}}\mathbf{)}𝑴𝒂𝒚𝒆𝒓_𝒌\mathbf{(}\frac{𝐆}{\mathrm{𝟐}}\mathbf{)}\mathbf{=}𝑳_𝒌\mathbf{(}\frac{𝐆}{\mathrm{𝟒}}\mathbf{)}`$ Now, the total generalized GS couplings can be obtained by summing the D-brane and generalized sigma orbifold fixed-points contributions. For the D-brane the relevant coupling has the following form (Scrucca and Serone, hep-th/0006201): $`𝒀_{\mathbf{(}\mathrm{𝟐}𝒌\mathbf{)}}\mathbf{(}𝑹\mathbf{,}𝑮\mathbf{,}𝑭\mathbf{)}\mathbf{=}\frac{\mathrm{𝟏}}{\sqrt{𝐍}}\mathit{ϵ}_{\mathrm{𝟐}𝒌}\sqrt{\frac{\mathrm{𝟏}}{\mathbf{\|}𝐂_{\mathrm{𝟐}𝐤}\mathbf{\|}}}𝒄𝒉_{\mathrm{𝟐}𝒌}\mathbf{(}\mathit{ϵ}_{\mathrm{𝟐}𝒌}𝑭\mathbf{)}\sqrt{𝑨_𝒌\mathbf{(}𝑮\mathit{ϵ}_𝒌\mathbf{)}}\sqrt{𝑨\mathbf{(}𝑹\mathbf{)}}`$ Defining the quantities $`X_{(2k)}=Y_{(2k)}+Z_{(2k)}`$, one has: $`𝑺_{𝑮𝑺}\mathbf{=}\sqrt{\mathrm{𝟐}𝝅}\mathbf{}_{𝒌\mathbf{=}\mathrm{𝟏}}^{\frac{\mathrm{𝟏}}{\mathrm{𝟐}}\mathbf{(}𝑵\mathbf{}\mathrm{𝟏}\mathbf{)}}\mathbf{}_{𝒊_𝒌\mathbf{=}\mathrm{𝟏}}^{𝑵_𝒌}\mathbf{}𝑪_{\mathbf{(}\mathrm{𝟐}𝒌\mathbf{)}}^{𝒊_𝒌}\mathbf{}𝑿_{\mathbf{(}\mathrm{𝟐}𝒌\mathbf{)}}`$ Using the explicit forms of the Mukay vectors of RR charges for the D-branes and generalized sigma orbifol fixed-points and the tadpole cancellation condition, one can to check that the total RR charges $`X_{(2k)}^{(0)}`$ with respect to the RR 4-forms are zero, and the following results for the total RR charges $`X_{(2k)}^{((2)}`$ and $`X_{(2k)}^{((4)}`$ with respect to the RR 2-forms and the RR 0-forms are found: $`𝑿_{\mathbf{(}\mathrm{𝟐}𝒌\mathbf{)}}^{\mathbf{(}\mathrm{𝟐}\mathbf{)}}\mathbf{=}\frac{\mathrm{𝟏}}{\sqrt{𝐍}\mathrm{𝟐}𝐩𝐢𝐍_𝐤^{\frac{\mathrm{𝟏}}{\mathrm{𝟒}}}}\mathbf{[}𝒊𝒕𝒓\mathbf{(}𝜸_{\mathrm{𝟐}𝒌}𝑭\mathbf{)}\mathbf{+}\frac{\mathrm{𝟏}}{\mathrm{𝟒}}𝒕𝒓\mathbf{(}𝜸_{\mathrm{𝟐}𝒌}\mathbf{)}\mathbf{(}\mathbf{}_{𝒊\mathbf{=}\mathrm{𝟏}}^\mathrm{𝟑}𝑮_𝒊𝒕𝒂𝒏\mathbf{(}𝒑𝒊𝒌𝒗_𝒊\mathbf{)}\mathbf{+}\mathbf{}_{𝒂\mathbf{=}\mathrm{𝟏}}^𝒏𝒀_𝒂𝒕𝒂𝒏\mathbf{(}𝒑𝒊𝒌𝒖_𝒂\mathbf{)}\mathbf{)}\mathbf{]}`$ $`𝑿_{\mathbf{(}\mathrm{𝟐}𝒌\mathbf{)}}^{\mathbf{(}\mathrm{𝟒}\mathbf{)}}\mathbf{=}\mathit{ϵ}_{\mathrm{𝟐}𝒌}\frac{\mathbf{}\mathrm{𝟏}}{\mathrm{𝟐}\sqrt{𝐍}\mathbf{(}\mathrm{𝟐}𝐩𝐢\mathbf{)}^\mathrm{𝟐}𝐍_𝐤^{\frac{\mathrm{𝟏}}{\mathrm{𝟒}}}}\mathbf{\{}𝒕𝒓\mathbf{(}𝜸_{\mathrm{𝟐}𝒌}𝑭^\mathrm{𝟐}\mathbf{)}\mathbf{}\frac{\mathrm{𝟏}}{\mathrm{𝟔𝟒}}𝒕𝒓\mathbf{(}𝜸_{\mathrm{𝟐}𝒌}\mathbf{)}𝒕𝒓\mathbf{(}𝑹^\mathrm{𝟐}\mathbf{)}\mathbf{+}𝒊𝒕𝒓\mathbf{(}𝜸_{\mathrm{𝟐}𝒌}𝑭\mathbf{)}\mathbf{}_{𝒊\mathbf{=}\mathrm{𝟏}}^\mathrm{𝟑}𝑮_𝒊𝒄𝒐𝒕\mathbf{(}\mathrm{𝟐}𝒑𝒊𝒌𝒗_𝒊\mathbf{)}\mathbf{}𝒕𝒓\mathbf{(}𝜸_{\mathrm{𝟐}𝒌}\mathbf{)}\mathbf{[}\mathbf{}_{𝒊\mathbf{=}\mathrm{𝟏}}^\mathrm{𝟑}\frac{\mathrm{𝟑}}{\mathrm{𝟏𝟔}}𝑮_𝒊^\mathrm{𝟐}𝒕𝒂𝒏^\mathrm{𝟐}\mathbf{(}𝒑𝒊𝒌𝒗_𝒊\mathbf{)}\mathbf{+}\mathbf{}_{𝒊\mathbf{}𝒋\mathbf{=}\mathrm{𝟏}}^\mathrm{𝟑}\frac{\mathrm{𝐜𝐨𝐬}\mathbf{(}\mathrm{𝟐}𝐩𝐢𝐤𝐯_𝐢\mathbf{)}\mathrm{𝐜𝐨𝐬}\mathbf{(}\mathrm{𝟐}𝐩𝐢𝐤𝐯_𝐣\mathbf{)}\mathbf{}\mathrm{𝐜𝐨𝐬}^\mathrm{𝟐}\mathbf{(}\mathrm{𝐩𝐢𝐤𝐯}_𝐢\mathbf{)}\mathrm{𝐜𝐨𝐬}^\mathrm{𝟐}\mathbf{(}\mathrm{𝐩𝐢𝐤𝐯}_𝐣\mathbf{)}}{\mathrm{𝟐}𝐬𝐢𝐧\mathbf{(}\mathrm{𝟐}𝐩𝐢𝐤𝐯_𝐢\mathbf{)}\mathrm{𝐬𝐢𝐧}\mathbf{(}\mathrm{𝟐}𝐩𝐢𝐤𝐯_𝐣\mathbf{)}}𝑮_𝒊𝑮_𝒋\mathbf{+}\mathbf{}_{𝒂\mathbf{=}\mathrm{𝟏}}^𝒏\frac{𝐘_𝐚^\mathrm{𝟐}\mathbf{(}\mathrm{𝟐}𝐜𝐨𝐬^\mathrm{𝟐}\mathbf{(}\mathrm{𝐩𝐢𝐤𝐮}_𝐚\mathbf{)}\mathbf{+}\mathrm{𝐬𝐢𝐧}^\mathrm{𝟐}\mathbf{(}\mathrm{𝐩𝐢𝐤𝐮}_𝐚\mathbf{)}\mathbf{)}}{\mathrm{𝟏𝟔}𝐜𝐨𝐬^\mathrm{𝟐}\mathbf{(}\mathrm{𝐩𝐢𝐤𝐮}_𝐚\mathbf{)}}\mathbf{}\mathbf{}_{𝒂\mathbf{=}\mathrm{𝟏}}^𝒏\mathbf{}_{𝒊\mathbf{=}\mathrm{𝟏}}^\mathrm{𝟑}\frac{\mathrm{𝐬𝐢𝐧}\mathbf{(}\mathrm{𝐩𝐢𝐤𝐮}_𝐚\mathbf{)}\mathrm{𝐜𝐨𝐬}\mathbf{(}\mathrm{𝐩𝐢𝐤𝐯}_𝐢\mathbf{)}}{\mathrm{𝟖}𝐜𝐨𝐬\mathbf{(}\mathrm{𝐩𝐢𝐤𝐮}_𝐚\mathbf{)}\mathrm{𝐬𝐢𝐧}\mathbf{(}\mathrm{𝐩𝐢𝐤𝐯}_𝐢\mathbf{)}}𝒀_𝒂𝑮_𝒊\mathbf{}\mathbf{}_{𝒂\mathbf{}𝒃\mathbf{=}\mathrm{𝟏}}^𝒏\frac{\mathrm{𝐬𝐢𝐧}\mathbf{(}\mathrm{𝐩𝐢𝐤𝐮}_𝐚\mathbf{)}\mathrm{𝐬𝐢𝐧}\mathbf{(}\mathrm{𝐩𝐢𝐤𝐮}_𝐛\mathbf{)}}{\mathrm{𝟖}𝐜𝐨𝐬\mathbf{(}\mathrm{𝐩𝐢𝐤𝐮}_𝐚\mathbf{)}\mathrm{𝐜𝐨𝐬}\mathbf{(}\mathrm{𝐩𝐢𝐤𝐮}_𝐛\mathbf{)}}𝒀_𝒂𝒀_𝒃\mathbf{]}\mathbf{\}}`$ These results are generalizations of the equations (6.10) and (6.11) in hep-th/0006201. When the SO(2n) bundle is the SO(10) tangent bundle, the equations in this paper are reduced to the equations (6.10) and (6.11) in hep-th/0006201. For such case the vector u is reduced to the vector v and the vector Y is reduced to $`G(i\pi x)`$ Finally, one arrives at a very simple factorized expression for the 6-form that are encoding the complete sigma- standard gauge-gravitational-non standard gauge anomaly and its opposite inflow(Scrucca and Serone, hep-th/0006201): $`𝑨^{\mathbf{(}\mathrm{𝟔}\mathbf{)}}\mathbf{=}𝑰^{\mathbf{(}\mathrm{𝟔}\mathbf{)}}\mathbf{=}𝒊\mathbf{}_{𝒌\mathbf{=}\mathrm{𝟏}}^{\frac{\mathrm{𝟏}}{\mathrm{𝟐}}\mathbf{(}𝑵\mathbf{}\mathrm{𝟏}\mathbf{)}}𝑵_𝒌𝑿_{\mathbf{(}\mathrm{𝟐}𝒌\mathbf{)}}^{\mathbf{(}\mathrm{𝟐}\mathbf{)}}\mathbf{}𝑿_{\mathbf{(}\mathrm{𝟐}𝒌\mathbf{)}}^{\mathbf{(}\mathrm{𝟒}\mathbf{)}}`$ ## 3 Conclutions Using both Mayer class and equivariant Mayer class it is possible to write the WZ couplings for certain generalized sigma orbifold fixed-points. This involves a new non standard SO(2n) gauge bundle. When such new bundle is the SO(10) tangent bundle of the ten dimensional space-time of the superstrings theories, then one can to obtain the WZ couplings for the usual sigma orbifold fixed-points. Finally when the new WZ coupling for the such generalized sigma orbifold fixed-points is combined with the usual WZ coupling for the usual Dp-brane, on can to obtain the generalized 6 form that are encoding the complete anomaly and its opposite inflow. ## 4 References ### 4.1 About WZ couplings for Dp-branes and Op-planes J. Morales, C. Scrucca and M. Serone, Anomalous couplings for D-branes and O-planes, hep-th/9812071 B.Stefanski,Jr., Gravitational Couplings of D-branes and O-planes, hep-th/9812088 K. Dasgupta, D. Jatkar and S. Mukhi, Gravitational couplings and Z2 orientifolds, Nucl. Phys. B523 (1998) 465, hep-th/9707224. K. Dasgupta and S. Mukhi, Anomaly inflow on orientifold planes, J. High Energy Phys. 3 (1998) 4, hep-th/9709219. Ben Craps and Frederik Roose, (Non-)Anomalous D-brane and O-plane couplings:the normal bundle, hep-th/9812149. Sunil Mukhi and Nemani V. Suryanarayana, Gravitational Couplings, Orientifolds and M-Planes, hep-th/9907215 ### 4.2 About WZ couplings for generalized Op-planes J.F. Ospina, Gravitation couplings for generalized Op-planes, hep-th/0006076,hep-th/0006095, hep-th/0006149 ### 4.3 About WZ couplings for sigma orbifold fixed points C. Scrucca and M. Serone, Gauge and gravitational anomalies in D=4 N=1 orientifolds, hep-th/9912108 C. Scrucca and M. Serone, hep-th/0006201 ### 4.4 About Mayer classes F. Hirzebruch, Topological Methods in Algebraic Geometry, 1978 Christian Bar, Elliptic Symbols, december 1995, Math. Nachr. 201, 7-35 (1999) Keke Li, Character-valued Index Theorems In Supersymmetric String Theories, CTP 1460, May 1987 J.F. Ospina, A Heterotic Susy version of Mayer integrality theorem, hep-th/9606186
warning/0007/astro-ph0007148.html	ar5iv	text	# Infrared Spectroscopy of Molecular Supernova Remnants ## 1 Introduction Infrared spectroscopy is a particularly useful tool for studying supernova remnants (and other shocks) in dense gas for the following reasons. (1) Infrared atomic fine-structure lines are produced by gas with the physical conditions expected behind radiative shocks. Dense, shocked gas rapidly cools to temperatures of $`10^2`$$`10^3`$ K, where infrared transitions are excited. This cool layer dominates the shocked column density, while the layer that emits optical to X-ray lines is extremely narrow. The kinetic temperature of the cooling post-shock gas is perfectly suited to the energy levels of infrared fine-structure lines from atoms and ions and infrared vibrational, ro-vibrational, and quadrupole vibrational (H<sub>2</sub>) lines from molecules. (2) Infrared fine-structure lines trace the ground-state populations of the all of the abundant elements that have electronic ground-states with non-zero angular momentum. By measuring the brightness of transitions directly to the ground-state or within a narrow multiplet of energy levels including the ground-state, a reasonably accurate measure of the column density of each species can be determined. Using pairs of lines from different elements due to energy levels with similar spacings, it is possible to measure relative abundances with confidence. (3) Far-infrared fine-structure lines are not affected by extinction. The light produced by post-shock gas is significantly affected by extinction from dust in the pre-shock cloud and the intervening medium along the line of sight. Extinction effects are severe in the ultraviolet and also in the optical and near-infrared for inner-galaxy supernova remnants, where the line-of-sight column density often exceeds $`10^{22}`$ cm<sup>-2</sup>, but it is negligible in the far-infrared. Many supernova remnants interacting with clouds are known only from their radio or hard X-ray emission, despite being copious producers of soft X-rays and optical line emission, because they are too distant or they are within or behind a molecular cloud. Recent advances in infrared astronomy have opened the far-infrared window to spectroscopy of faint emission. The observations used in this paper were made with the Long-Wavelength Spectrometer (LWS; Clegg et al. (1996)) aboard the Infrared Space Observatory (ISO; Kessler et al. (1996)), whose mission lasted from 1996–1997. A complete spectrum with the LWS, covering the 40–190 $`\mu `$m wavelength region with a resolution of $`0.29`$$`0.6`$ $`\mu `$m, took less than half an hour to reach a sensitivity adequate to detect atomic fine-structure lines and molecular rotational lines for a variety of astronomical sources. Instruments planned for the Stratospheric Observatory for Infrared Astronomy (SOFIA; Becklin (1997); expected to begin observing in 2002) will open opportunities for far-infrared spectroscopic observations over most of the far-infrared wavelength range—specifically, the part not affected by absorption and emission by very abundant molecules in the residual atmosphere such as H<sub>2</sub>O. In the somewhat-more-distant future, instruments aboard the Far-Infrared Space Telescope (FIRST; Poglitsch (1998); anticipated launch 2007) will allow continuous spectral coverage to even deeper levels for extended sources. The goal of this paper is to present our ISO LWS and SWS spectral observations of the molecular supernova remnants W 44, W 28, and 3C 391, and to illustrate from basic physical principles why the spectral lines we detected are bright. Shocked gas in supernova remnants emits virtually all fine-structure lines from the abundant elements. We show the inferred elemental abundances after impact from the supernova shock, and we discuss implication of grain destruction for elements such as Fe, Si, and P which were locked in grains. This paper continues our investigation of infrared emission from molecular supernova remnants. The brightness of the \[O I\] 63.2 $`\mu `$m line for W 44 and 3C 391 was presented in our first paper (Reach & Rho (1996)). The far-infrared molecular emission of OH, CO, and H<sub>2</sub>O was presented in the second paper (Reach & Rho (1998)). The millimeter-wave molecular emission of CO, CS, HCO<sup>+</sup> was presented in the third paper for 3C 391 (Reach & Rho (1999)), and in this paper we also present some molecular results for W 44 from new IRAM 30-m observations. Two recent and closely-related studies of molecular emission were presented by Seta et al. (1998) for W 44 and Arikawa et al. (1999) for W 28. The millimeter-wave observations show beyond doubt that all three of the supernova remnants discussed in this paper are interacting with molecular clouds. ## 2 Observations ### 2.1 Far-infrared spectral lines We obtained complete ISO LWS spectra of 1 position in each of three supernova remnants as part of our program of observing supernova-molecular cloud interactions (Reach & Rho 1996, 1998, 1999). The positions observed were the brightest OH masers (Frail et al. (1996)) in W 28 ($`18^h01^m52.3^s`$ $`23^{}19^{}25^{\prime \prime }`$) and W 44 ($`18^h56^m28.4^s`$ $`+1^{}29^{}59^{\prime \prime }`$) and the molecular peak in 3C 391:BML (18<sup>h</sup>49<sup>m</sup>21.9<sup>s</sup> -05722<sup>′′</sup>); all coordinates are J2000. Figure 1 shows the ISO LWS spectrum from 42–188 $`\mu `$m for all three remnants. For each of the 10 LWS detectors, we removed the fringes using the ISO Spectral Analysis Package<sup>1</sup><sup>1</sup>1The ISO Spectral Analysis Package (ISAP) is a joint development by the LWS and SWS Instrument Teams and Data Centers. Contributing institutes are CESR, IAS, IPAC, MPE, RAL and SRON. and subtracted a constant (due to residual dark current). Table 1 shows a list of the identified atomic fine structure lines and their measured brightness. In order of atomic number, we detected atomic fine-structure lines from C<sup>+</sup>, N<sup>+</sup>, N<sup>++</sup>, O<sup>0</sup>, O<sup>++</sup>, O<sup>+++</sup>, and P<sup>+</sup> in the LWS spectra. Part of the bright continuum and \[C II\] 157.7 $`\mu `$m line in Fig. 1 are due to dust and gas along the line of sight, unrelated to the supernova remnant. At least part of the continuum is likely to be due to dust grains surviving the shocks. For example, for W 44 we estimated 30% (Reach & Rho (1996)) based on comparison to a reference position; a spectral analysis of the continuum is given in §6 of this paper. To estimate the line-of-sight contamination of the spectral line brightnesses, we calculated the galactic line-to-continuum ratio using the COBE Far-Infrared Absolute Spectrophotometer spectrum of the galactic plane at longitude $`30^{}`$ (Reach et al. (1995)). At the ISO LWS resolution, the expected line-to-continuum ratio due to unrelated gas for \[C II\] 157.7 $`\mu `$m is 2.2, while that for \[N II\] at 121.9 $`\mu `$m is 0.05. Comparing to our spectra, the 157.7 $`\mu `$m line to continuum ratios are 2.4, 1.2, and 1.9, and the 121.9 $`\mu `$m line to continuum ratios are 0.19, 0.09, and 0.16, for W 44, W 28, and 3C 391, respectively. Thus much of the \[C II\] $`\mu `$m line emission could be due to line-of-sight gas, while only a small fraction the \[N II\] (and essentially none of the other lines) are due to line-of-sight gas. The subject of \[C II\] emission being related to the shock is addressed again in §4, focusing on the abundance of C. An improved separation of the source and background brightness might be possible by analyzing the spatial variation of the brightness, which is beyond the scope of this paper. ### 2.2 Mid-infrared Si, Fe, and H<sub>2</sub> lines We used the ISO SWS (de Graauw et al. (1996)) to observe narrow ranges centered 25.98 and 34.82 $`\mu `$m for the ground-state transitions of \[Fe II\] and \[Si II\]. The spectra are shown in Figures 2.1 and 2.1, respectively. At the same time, we obtained spectra of the H<sub>2</sub> S(9) and S(3) lines; the spectra of these lines are shown in Figures 2.1 and 2.1, respectively. The brightnesses of the lines are summarized in Table 2. The SWS observations were toward the three lines of sight listed above, plus one extra position—toward the bright, northwestern radio ridge and a bright peak in the \[O I\] 63 $`\mu `$m strip map (position b in Fig. 3 of Reach & Rho (1996)): 3C 391:radio (18<sup>h</sup>49<sup>m</sup>19.4<sup>s</sup> -05505<sup>′′</sup>). The two positions in 3C 391 have similar spectra of ionic fine structure lines: the ratio of \[Fe II\]/\[Si II\] is nearly identical, although the \[O IV\] is relatively brighter toward 3C 391:radio by at least a factor of 3. In contrast, the brightness of the H<sub>2</sub> lines, relative to the ionic lines, is more than a factor of 10 lower for 3C 391:radio than for 3C 391:BML. This is consistent with the millimeter-wave molecular lines revealing no evidence for shocked molecular gas toward 3C 391:radio, while bright and broad shocked molecular lines were detected toward 3C 391:BML (Figs. 2 and 4 of Reach & Rho (1999)). The spectra of these two positions serve a valuable comparison and diagnostic tool, because one of them is dominated by ionic lines, while the other has significant molecular emission. ### 2.3 High-resolution \[O I\] and CO spectra We used the Fabry-Perot in the LWS to obtain high-resolution spectral profiles of the \[O I\] 63 $`\mu `$m lines for 3C 391:BML and W 44. The positions are the same as those where the complete LWS spectra and the SWS line spectra were taken, and they correspond to peaks in the \[O I\] strip maps presented in paper 1<sup>2</sup><sup>2</sup>2 Please note that the apparent 400 km s<sup>-1</sup> shift in line center in 3C 391 from one side of the remnant to the other, mentioned in paper 1, is apparently an instrumental artifact due to mapping in the dispersion direction of the grating. The observations presented in this paper, with the Fabry-Perot, are considered reliable measurements of the line profile.. The instrumental resolution of the LWS Fabry-Perot is 36 km s<sup>-1</sup>. The two spectra are shown in Figures 2.2 and 2.2. The lines are clearly resolved, with a full width at half maximum (FWHM) of 100 km s<sup>-1</sup>. The lower panels of Figs. 2.2 and 2.2 show the CO($`21`$) and <sup>13</sup>CO($`10`$) line profiles for the same lines of sight as the O I spectra, obtained with the IRAM 30-m telescope (Wild (1995)). The observing techniques are described in our previous paper, where the 3C 391 observations were presented (Reach & Rho (1999)); the W 44 observations are new. The line profiles are all clearly, very different. The <sup>13</sup>CO lines have a narrow (1.5 km s<sup>-1</sup> FWHM) velocity dispersion, and they trace the column density, so they arises from large clouds of cold, unshocked or unrelated gas. Although the CO($`21`$) lines also have narrow components, identical to the <sup>13</sup>CO lines, from the unshocked or unrelated gas, the CO($`21`$) lines are dominated by a broad component (20–50 km s<sup>-1</sup> FWHM). The broad component was seen in the CO, CS, and HCO<sup>+</sup> observations of 3C 391, with relatively more emission from the higher-excitation and higher dipole moment lines (Reach & Rho (1999)). Similarly, the CO observations of W 28 showed that the broad component was much brighter than the narrow component in the $`32`$ line as compared to the $`10`$ line (Arikawa et al. (1999)). The broad components of the molecular lines are still not as broad as the \[O I\] line, suggesting that each of the spectral lines in Figs. 2.2 and 2.2 trace different regions. Inspecting the \[O I\] line profiles, there is only slight evidence for structure. The lines from both remnants can be reasonably represented by gaussians with FWHM $`100`$ km s<sup>-1</sup>. For 3C 391:BML, the line profile is somewhat better represented by multiple components: an unresolved component centered near 50 km s<sup>-1</sup> and a broad, non-Gaussian component centered near 0 km s<sup>-1</sup> with 120–150 km s<sup>-1</sup>. The integrated line brightness is dominated by the wider component. The unresolved component contributes only about 10% of the line integral, and it may be associated with preshock or unrelated gas, which would agree with the brightness of the ‘off-remnant’ positions from the \[O I\] strip map (Reach & Rho (1996)). Neither 3C 391 nor W 44 show evidence for self-absorption in the line profiles. This was a particular concern in interpreting the line brightnesses for molecular supernova remnants, both because of the possibility of dense pre-shock clouds and the large column densities of cold gas along the lines of sight to the remnants. A single cold cloud might be expected to yield an essentially unresolved (Gaussian) absorption dip, which does not occur. The foreground gas would be expected to absorb in approximate proportion to the line of sight column density. After inspecting the H I 21-cm profile for these lines of sight, we found no correlation that could be due to absorption dips. Instead, it is likely that the line-of-sight gas contributes slightly to the emission. The level of \[O I\] emission from the line-of-sight gas is most easily estimated from the reference positions just outside the remnant, where it was found to be quite faint (Reach & Rho (1996)). Therefore, we interpret the \[O I\] line profiles in Figs. 2.2 and 2.2 as emission from the shocked gas. This inference yields a direct measure of the shock velocity of 100–150 km s<sup>-1</sup>, albeit in projection along the line of sight. There are likely to be multiple shocks within the beam, and the FWHM of the line profile is approximately equal to the ‘typical’ shock velocity. Comparing the \[O I\] and CO line profiles, we can partially describe the three-dimensional geometry of the shocks. For 3C 391, note that the centroids of the \[O I\] and CO lines are very different. The CO line occurs at nearly the systemic velocity at the remnant’s distance, with the broad component due to the shocked gas very slightly blueshifted from the narrow line from the preshock gas. The \[O I\] line, on the other hand, is significantly red-shifted from the pre-shock gas. Thus the shock that produces the \[O I\] line has a significant velocity component toward us; that is, the shock is propagating partially along the line of sight. This suggests that the 3C 391 progenitor exploded on the far side of its parent molecular cloud, driving the dense blast wave partially in our direction. For W 44, the centroids of the \[O I\] and CO lines are similar, suggesting that the shocks are either perpendicular to the line of sight, or that the lines come from similar geometries. One straightforward conclusion from comparing the line profiles, both in centroid and width, is that the \[O I\] line does not arise exclusively from slow shocks into very dense gas such as produces the CS, OH, and H<sub>2</sub>O emission (Reach & Rho (1999)). Instead, most of the \[O I\] emission must arise from faster shocks into moderate-density gas. ## 3 Comparison of detected lines to the periodic table ### 3.1 Atomic fine-structure lines: Basic principles The three basic principles that determine which spectral lines will be bright in the spectrum of a parcel of gas are abundance, ionization, and excitation. The abundances of the elements are given by Anders & Grevesse (1989) for the Solar System. In the interstellar medium, since part of the abundance will be locked in solids, the gas phase abundance will be smaller than Solar System abundance for some elements (Savage & Sembach (1996)). It is clear that we need only consider the top three rows of the Periodic Table and the middle of the fourth row; all other elements have negligible abundance. The ionization state of each element depends on where the gas is located. We can consider three typical locations: the diffuse interstellar medium (where all electrons bound by less than 13.6 eV are removed), dark clouds (where elements are neutral), and hot plasmas (where the higher ionization states are present). The excitation of the atomic levels will depend on the density and temperature of the region, but a ‘zero-th’ order effect determines whether there are any atomic fine structure lines in the ground electronic state. Elements and ions with zero angular momentum ($`L=0`$) in their ground state electron configuration will not have the fine-structure lines considered here. For example, in spectroscopic notation, S configurations have no lines, while doublet <sup>2</sup>P configurations have 1 line and triplet <sup>3</sup>P configurations have 2 lines. The fine-structure lines are magnetic dipole transitions, so they are technically ‘forbidden lines’, but with spontaneous radiative decay rates of order $`10^5`$ s<sup>-1</sup> and collisional de-excitation rates of order $`10^7`$ cm<sup>-3</sup> s<sup>-1</sup>, they are easily excited under astronomical conditions and are thermalized at densities $`10^2`$ cm<sup>-3</sup> (for far-infrared lines) to $`10^4`$ cm<sup>-3</sup> (some mid-infrared lines). In order to understand and predict which fine-structure lines are present in astronomical spectra, we constructed a ‘periodic table for astronomical fine structure lines’. Table 3 contains the electron configurations and wavelengths for the p-shell ions of the second and third rows of the Periodic Table (C, N, O, Fe, Ne, Al, Si, P, S, Cl, and Ar), while Table 4 has the same information for the fourth row of the Periodic Table (Cr, Mn, Fe, CO, and Ni). In both tables, the dominant ionization state in the diffuse interstellar medium is shown in bold. Spectra of regions behind shock fronts or with very high radiation fields will contain the dominant ionization state as well as and the ions to the left of the dominant ionization state. Spectra of regions that are shielded from the interstellar radiation field will contain neutral forms and ions to the right of the dominant ionization state. Some elements with negligible astronomical abundance are not listed; also, to reduce clutter, unlikely ionization states of rare elements are not listed. Alkaline metals (in low ionization states) have no $`p`$-shell electrons to make infrared fine-structure lines, so they are not listed. These tables can be used both to predict which lines will be present for a given gas and to decide which lines should be present when another line of the same ion is detected. In all cases, the transitions are simple cascades down a ladder, with the ground state at the bottom; therefore, the upper energy level for a line can be calculated by summing the photon energies from a given line to the bottom. ### 3.2 Comparison of detected lines to the Periodic Table In the ‘periodic table for atomic fine-structure lines’ (Tables 3 and 4), the lines we detected from molecular supernova remnants observations are in italics. It appears that we detect all of the available ground-state fine-structure lines of the astronomically abundant molecules; the undetected ground-state lines are generally outside our wavelength range. All abundant elements, C, N, O and Si, are detected as shown in Tables 1 and 2. Non-detections of F and Co lines are due to low abundances. Lines from high-ionization states were also not detected. The only element with high abundance that we did not detect is Ne, because our observations did not cover the wavelengths of its ground-state fine-structure lines. Our own observations contain the wavelengths of 12 fine-structure transitions among energy levels within 600 K of the ground state, from elements with abundance $`>10^7`$ relative to H in their expected ionization states. We detect essentially all of these lines, as well as lines from higher ionization states of the abundant elements: O III, O IV, and N III. In addition to these,the ground-state line of \[P II\] is also detected for W28 and 3C391. Some of the other lines were detected in the younger and brighter supernova remnant RCW 103 (Oliva et al. 1999b ); upon reanalysis of these data from the ISO archive, we find that the detected lines include the ground-state line of \[P II\]. We can use these general principles to predict other bright lines, outside of our observed spectral range. For example, Ne is a very abundant element with lines in the mid-infrared; given that we detected ionized states of O, we expect bright emission from \[Ne II\] at 12.8 $`\mu `$m and perhaps \[Ne III\] at 15.5 $`\mu `$m. The dashed contour running approximately from the upper left to lower right of each panel shows the constraint placed on the density and temperature by requiring that the emitting region is smaller than the entire supernova remnant, while still producing lines as bright as observed. The vertical dashed line in the upper panel at $`T=10^4`$ K shows where O should be collisionally ionized. The dashed contour near the top of the lower panel shows the constraint that the O III region be at least as hot and no more dense than the O I region. Hatched regions in both panels show the allowed solutions. The asterisk in the upper panel shows the solution that we chose for detailed calculations. The upper-left asterisk in the lower panel shows this same solution if the gas is fully ionized, and the asterisks connect to it show the set of solutions in thermal pressure equilibrium with it. The lower set of asterisks shows the same solutions assuming the ionization fraction in the O III region is 1%. ## 4 Abundances of the shocked gas ### 4.1 Excitation of Oxygen Based on the arrangement of the energy levels of the ground-state atomic fine structure, the ratios the upper and lower \[O I\] and \[O III\] lines constrain the density and temperature of the emitting region, independent of the geometry. The upper panel of Figure 3.2 shows the line ratio predicted for \[O I\] as a function of the total H density and temperature. At each density and temperature, the level populations were calculated and the emergent intensity predicted using the escape probability (Hollenbach & McKee (1979)). The column density, $`N_\mathrm{H}`$, was adjusted iteratively for each model such that the model predicts the correct brightness of the \[O I\] 63 $`\mu `$m line. Solid lines show the locus of models that predict the observed line ratios from Table 1. The locus of models for which the path length through the emitting region, $`z=N_\mathrm{H}/n_\mathrm{H}`$, is equal to 3 pc (i.e. comparable to the radius of the supernova remnant) is shown as a dashed contour, and the region below this contour is labeled ‘region to large.’ Low-temperature and low-density models are not capable of producing the observed line brightness within the confines of the observed region. Another constraint is drawn as a vertical line at $`T=10^4`$ K, where the neutral O atom would be collisionally ionized; higher temperature models should therefore be excluded. The lower panel of Figure 3.2 shows the line ratio predicted for \[O III\] as a function of density and temperature, with each model normalized to predict the correct brightness of the \[O III\] 88 $`\mu `$m line. The \[O III\] and \[O I\] emitting regions are not necessarily coincident, so we consider the density, temperature, and ionization in the two regions separately. One relative constraint between the two emitting regions is practically guaranteed: the temperature will be at least as high, the density will be at least as low (and the ionization will be at least as high) in the \[O III\] emitting region as compared to the \[O I\] emitting region. This rules out the very dense models for the \[O III\] emitting region, as indicated by the upper dashed curve in the lower panel of Figure 3.2. To determine the physical conditions where the lines are produced, and to determine the abundances of the elements in the shocked gas, we define two emitting regions. Region M (molecular) is the \[O I\] emitting region, with density, temperature, and column density set to give the observed brightnesses of the 63 and 145 $`\mu `$m lines. Region A (atomic) is the \[O III\] emitting region, with density, temperature, and column density set to give the observed brightnesses of the 52 and 88 $`\mu `$m lines. Such a separation of emitting regions is anticipated because we expect that the neutral and doubly-ionized states of O would not be co-spatial—although there could in principle be overlap, if the gas is far from ionization equilibrium. Combining all of the constraints, the plausible range of density and temperature for region M has density $`2000<n_\mathrm{H}^{(M)}`$ (cm<sup>-3</sup>)$`<5000`$ and temperature $`100<T^{(M)}`$ (K)$`<1000`$ K. For convenience, we have chosen a nominal solution of $`n_\mathrm{H}^{(M)}=3000`$ cm<sup>-3</sup> and $`T^{(M)}=600`$ K, which is marked with an asterisk ($``$) in Figure 3.2 and listed in Table 5, for detailed calculations. In region A, there is a wider range of solutions within the constraints already mentioned. In particular, models with constant $`n_e=n_\mathrm{H}x`$, where $`x`$ is the fraction of H that is ionized, produce the same excitation for the ions in region A. If region A and region M were layers behind the same shock front, then we expect region A to have a comparable or somewhat higher pressure than region M (lest the shock move backwards). For the nominal solution for region M, the region A solutions with the same pressure but density lower by factors of 1, 2, 4, and 8 are shown as two lines connecting asterisks in the lower panel of Figure 3.2. The upper line applies if region A is completely ionized, and the upper left asterisk is just a direct copy of the region M solution. It is evident that region A cannot be completely ionized and be a shock layer upstream from region M: the region A pressure would be so low that the shock would move backwards. The lower line applies if region A has an ionization fraction $`x=10^2`$, and the upper-left asterisk of the lower line is the region M solution with $`n_e=n_H/100`$. This would be a very surprising result, because it means that O<sup>++</sup> coexists with a significant amount of H<sup>0</sup>. Very low density gas can be far from ionization equilibrium, with higher ionization states coexisting with neutral states (Schmutzler & Tscharnuter (1993)), but at the high densities required to get the observed brightness of the 63 and 88 $`\mu `$m lines the ionization should be reasonably close to equilibrium. For this reason, we consider the \[O III\] emitting region (region A) separate from the \[O I\] emitting region. For detailed calculations, the region A model has density $`n_e^{(A)}=8`$ cm<sup>-3</sup> and temperature $`T^{(A)}=3000`$ K; the properties are listed in Table 5. In region A, the abundance of O<sup>0</sup> is set to zero, and the abundance of O<sup>++</sup> and N<sup>++</sup> are set to the entire cosmic abundance of O and N, respectively. ### 4.2 Predicted line brightnesses for different regions The impact of a supernova blast wave with a molecular cloud leads to a complex interaction, with regions of different density reacting in very different ways. Based on the excitation of \[O I\] and \[O III\] (discussed in the previous section), and the excitation of H<sub>2</sub>O and OH (Reach & Rho (1998)), we identify three types of post-shock gas and their infrared emitting lines. Figure 4.2 is a cartoon showing these shocks, and Table 5 lists their properties. The lowest-density pre-shock gas is region A (for ‘atomic’) with pre-shock density $`n_0^{(A)}<1`$ cm<sup>-3</sup>. (If the pre-shock density of region A were higher than 1 cm<sup>-3</sup>, the post-shock density would be much higher than the 8 cm<sup>-3</sup> and would be inconsistent with the 52/88 $`\mu `$m line ratio.) The post-shock gas in region A gives rise to bright H II recombination lines and infrared and optical O III and N III (and other high ion) forbidden lines. The moderate-density pre-shock gas is region M (for ‘molecular’) with pre-shock density $`n_0^{(M)}10^2`$ cm<sup>-3</sup>. The post-shock gas in region M gives rise to bright O I and most of the other bright infrared forbidden lines we observed, as well as some optical emission (which is, however, highly obscured). The highest-density pre-shock gas is in the dense clumps, with pre-shock density $`n_0^{(C)}10^4`$ cm<sup>-3</sup>. The post-shock gas in region C gives rise to bright H<sub>2</sub>, OH, CO, CS, and H<sub>2</sub>O emission, with a possible contribution to the O I emission. We do not mean to imply that there are three sharply-defined regions with fixed densities; we do suggest that the full range of pre-shock densities is needed, such that no single region is consistent with all of the observed spectral lines. The shock velocities in the different regions are not well known, but reasonable estimates can be made. If region A is the lowest-density pre-shock gas, then it should represent the leading edge of the supernova remnant. The Sedov solutions for these remnants are consistent with the density of region A for explosions with energy of order $`10^{51}`$ erg (Rho (1995)). Thus the shock velocity should be of order $`2R/5t`$, where $`R`$ is the radius and $`t`$ is the age of the remnant. This works out to $`V_s^{(A)}500`$–600 km s<sup>-1</sup> for the remnants considered in this paper (Rho (1995)). For region M, we have some direct evidence of the shock velocity from the width of the O I line (§2.3) that $`V_s^{(M)}100`$ km s<sup>-1</sup>. This is also consistent with the expanding H I shell observed toward W 44 (Koo & Heiles (1995)). For region C, there is some direct evidence of the shock velocity from the width of the millimeter-wave CO and CS lines (Reach & Rho (1999)) that $`V_s^{(C)}30`$ km s<sup>-1</sup>. The shock ram pressure $`p_{ram}n_0V_s^2`$ should be comparable in these regions as they are all driven by the same blast wave. Indeed we find $`p_{ram}`$ is progressively higher for the denser regions, but the range of $`p_{ram}`$ is only a factor of 36. A pressure enhancement in the denser regions has been theoretically justified by Chevalier (1999) as due to the shock into the denser gas (regions M and C) being driven from the radiative shell (region A) as opposed to being isolated shocks driven by the blast wave. In Chevalier’s (1999) paper, two regimes of density were considered for the pre-shock gas; these densities correspond roughly to our regions A and M (cf. Table 5). Chevalier explained the \[O I\] 63 $`\mu `$m luminosity of the remnant using only relatively low-density shocks into region A. Based on our new results for the line ratio and excitation of O I and the high-resolution spectrum of the 63 $`\mu `$m line showing a width of only 100 km s<sup>-1</sup>, we suspect that the 63 $`\mu `$m arises in relatively denser regions, like region M. Indeed, it is impossible for a low-density emitting region to produce 63 $`\mu `$m lines as bright as we observed, within the confines of the supernova remnant. In the model of Hollenbach & McKee (1989), the 63 $`\mu `$m line brightness is proportional to $`n_0V_s`$. For the cartoon model parameters in Table 5, it is clear that the denser regions will dominate the line brightness if they fill an appreciable portion of the beam. In region C, however, the O is probably largely tied up in molecular form; furthermore, region C is somewhat beam diluted in the ISO LWS beam. Therefore, we believe that the 63 $`\mu `$m lines presented here arise from shocks into the moderate-density molecular gas of region M. For regions A and M, we calculated the model brightness of all of the lines that we observed with ISO. Each element was assumed to have its entire solar abundance in the gas phase and in the ionization state that gives rise to the line, with the following exceptions. In region M, oxygen was assumed to be neutral (as we observed), and nitrogen was assumed to be singly ionized (since the only detected line was from N II, not N I) in order to obtain an upper limit to the predicted N II brightness.. In region A, oxygen was assumed to be ionized. Our purpose in these calculations is to measure the richness or depletion relative to solar abundances. Our calculations do not yield absolute abundances, because in reality the total abundance of each element is actually divided among the ionization states that are present. In the next section, we go through the list of elements and discuss their expected ionization state and the inferred total abundances in the post-shock gas. Table 6 shows the brightness predictions for regions A and M. We can conclude the following, without considering the ionization balance. First, it is possible to reproduce all of the observed fine-structure lines, except N III and O III using region M alone. (In fact, it is possible to get the N III and O III lines, if we could accept the multiply ionized species existing in region M, where H is mostly neutral.) Therefore, we expect to be able to derive abundances in the ‘region M’ post-shock gas. For region C, we discuss details in §5 below. ### 4.3 Post-shock gas abundances We now step through each element and derive the total abundance of that element in the post-shock gas. Table 7 summarizes the abundances inferred for the post-shock gas. For each element $`X`$, $`[X/H]_{}`$ is the solar abundance (Anders & Grevesse (1989)) and $`[X/H]_{obs}`$ is the total abundance of that element from our observations and model. Carbon—The only carbon line we observed was from C II, which is likely to be the dominant ionization state in region M. The observed line brightness compared to the region M prediction (Table 6) requires essentially all of the carbon to be in the gas phase. However, a significant fraction of the 157 $`\mu `$m line could come from line-of-sight gas, so this result is inconclusive. The carbon abundance in the post-shock gas could be better addressed if we knew the brightness of the line just outside of the remnant. Using our current knowledge, we can only say that the ISO observations are consistent with complete destruction of the C-bearing dust grains if all of the emission comes from the remnant, or partial or no destruction if less than half of the emission comes from the remnant. In Table 7 we list the inferred depletion of carbon in the post-shock gas assuming all of the 157 $`\mu `$m line comes from the remnant; this yields an upper limit to the gas-phase abundance, and therefore lower limit to the depletion. Nitrogen—We observed nitrogen in its singly- and doubly-ionized states. The dominant state in region M should be neutral, because the ionization potential of N I is greater than 13.6 eV. Comparing the observed to predicted brightness (Table 6), we require only about 1/3 of the nitrogen to be N II, which is reasonable. When all nitrogen is assumed to be singly ionized, the observed to predicted brightness ratio of \[N II\] is only 0.3 (Table 6), which implies either nitrogen has lower abundance or the predicted N II is too high. Since the former is unlikely, the ratio implies that only about 1/3 of the nitrogen is N II and the rest is neutral. In region A, which was normalized to give the correct 88 $`\mu `$m line brightness, we require about 1/3 of the N III, which is also plausible. In Table 7 we list the depletion of nitrogen assuming all of the nitrogen is ionized; this is a lower limit to the gas-phase abundance, and an upper limit to the depletion, because it neglects the neutral state. Oxygen—The column density of our model was adjusted to match the observed brightness of the oxygen lines, so most of our measurements are, effectively, relative to the O abundance. We cannot easily determine the total abundance of O because we lack an accurate measurement of the total column density of H. Silicon and Iron—We find that a significant fraction of the total abundance of Si and Fe, which are normally highly depleted into dust grains in the interstellar medium, is required to produce the observed line brightness. Both Si and Fe were observed in their dominant ionization state in region M. Region A does not have enough column density to contribute significantly to the observed line brightness. (Recall that the region A column density is set by the \[O III\] line brightness, so increasing its column density would make too much \[O III\] emission.) For region M, we find that about 1/2 of the Si and 1/4 of the Fe are required to be in the gas phase. These abundances of Fe and Si are interestingly comparable to those of gas-phase abundances in intermediate-velocity interstellar clouds which are recently measured using Goddard High Resolution Spectrograph, suggesting that resilient cores of grains are not easily destroyed in shocks (Fitzpatrick (1996)). We will discuss the implications for grain destruction more in §6.2 below. Phosphorus—The abundance of P is known for only a few lines of sight in the interstellar medium. Our abundance estimate in the post-shock gas uses the observed brightness of the 60.6 $`\mu `$m line, for which this is the first reported detection (at the time of writing). Figure 4.3 shows the spectrum of this newly-detected line. We used the radiative decay rates are as given by Mendoza & Zeippen (1982), and we use the collision strength $`\mathrm{\Omega }(^3\mathrm{P}_10)=1.46`$ (Krueger & Zsyzak (1970)). The P abundance relative to O is very well constrained because the energy levels of the two observed ions are similar, and both elements were observed in their dominant ionization state. We find that the observed line brightness is consistent with solar abundance of P. Along the line of sight toward $`\zeta `$ Oph, it appears that about half of interstellar P is normally locked in grains (Savage & Sembach (1996)). Our observation suggests that some or all of the P-bearing solid material was vaporized by the strong shock wave, as we find for Fe and Si (see §6.2). ## 5 Shock-excited molecules ### 5.1 H<sub>2</sub> The S(3) and S(9) lines of H<sub>2</sub> were both detected from each of the supernova remnants we observed, and the line brightnesses are listed in Table 2. These H<sub>2</sub> lines arise from shocks in the molecular gas, like the CO, OH, and H<sub>2</sub>O lines (Reach & Rho (1998)). Using our cartoon model, we can now localize which region(s) give rise to which lines and further constrain the properties of the dense clumps (region C). First, for region M, if we assume that the H<sub>2</sub> rotational lines are collisionally excited with rotational excitation temperature equal to the kinetic temperature, we predict that the brightest H<sub>2</sub> line is the S(3) line. The brightness of the S(3) line from region M, assuming all of the H is in H<sub>2</sub>, is $`2\times 10^4`$ erg cm<sup>-2</sup> s<sup>-1</sup> sr<sup>-1</sup>. This is in very good agreement with the observed line brightness, suggesting that the H is largely molecular behind the region M shocks. The column density of H<sub>2</sub> is just enough that H<sub>2</sub> molecules are self-shielded from dissociating photons in the interstellar radiation field (Hollenbach, Werner, & Salpeter (1971)). We cannot explain the brightness of the higher H<sub>2</sub> S(9) rotational line, or the higher-lying ro-vibrational lines in the near-infrared, using the region M shocks. These lines most likely arise in molecular regions that are much warmer than the recombined molecules behind a fast, dissociative shock. Instead, the S(3) lines probably arise from non-dissociative shocks into denser gas (discussed below). In a previous paper (Reach & Rho (1998)), we derived the physical properties of region C using the excitation of CO, OH, and H<sub>2</sub>O. In that paper, we found the puzzling result that abundance of H<sub>2</sub>O was lower than expected in a region with ‘warm’ chemistry, where O is efficiently converted into molecules and OH is converted into H<sub>2</sub>O. Two papers have shown that OH can be formed from dissociation of H<sub>2</sub>O (Lockett et al. (1998), Wardle (1999)), so the abundance OH is no longer such a puzzle. We now reexamine the H<sub>2</sub>O abundance relative to H<sub>2</sub>. The column density of H<sub>2</sub>O was found to be $`3\times 10^{17}`$ cm<sup>-2</sup> using the far-infrared line brightnesses. The column density of H<sub>2</sub> was taken from the volume density times path length, yielding $`8\times 10^{23}`$ cm<sup>-2</sup>. We suggested this column density was consistent with the brightness of the H<sub>2</sub> S(3) line, but we now suspect that much of the S(3) line arises from region M and not the dense shock of region C. If this is true, then the H<sub>2</sub> column density of region C could be much lower than we estimated before. Let us start with some assumptions about the molecular abundances and infer the H<sub>2</sub> abundance. First, CO:OH:H<sub>2</sub>O is 100:1:15 from our observations (Reach & Rho (1998)). Assume that the C and O are largely in molecules. If a fraction $`f_{CO}`$ of the C is locked in CO, then the abundance ratio of H<sub>2</sub>O/H<sub>2</sub> is $`X_\mathrm{O}f_{CO}`$, where $`X_\mathrm{O}`$ is the abundance of O. The column density of H<sub>2</sub> is then $`N(\mathrm{H}_2)^{(C)}=3\times 10^{21}f_{CO}^1`$. Such a column density for region C seems reasonable, but it is lower than our previous estimate based on the observed size of the CO emitting region in the millimeter observations (Reach & Rho (1999)). This could be explained by a low filling factor for the dense gas, of order $`0.007f_{CO}^1`$ in the ISO SWS beam. Taking $`f_{CO}=0.5`$ as an example, the size of the emitting region would be of order $`1^{\prime \prime }`$, and it would be highly structured in future, higher-angular-resolution observations. If true, this new interpretation would change our estimate of the mass of the 3C 391:BML post-shock gas, making the self-gravity of the shocked gas negligible, and weakening our argument that the dense clump might imminently collapse to form stars (Reach & Rho (1999)). Table 5 shows the cartoon model for region C assuming $`f_{CO}=0.5`$. The H<sub>2</sub> 0-0 S(9) line cannot be assigned to any of the emitting regions that we have defined so far. It’s energy level is too high above the ground state to be excited in the cooling region where the far-infrared lines are produced. It is unlikely that the S(9) line can be produced in the region M shock, even upstream from the cooling region, because the H<sub>2</sub> is probably still dissociated. Thus we suspect that the S(9) line is produced upstream from the region C shock. To get the observed line S(9) brightness, in an emitting region with rotational excitation temperature of order 1500 K, a column density of order $`10^{19}`$ cm<sup>-2</sup> is needed, which could easily fit in the upstream region considering the cooling region has a column density 600 times larger. The 1500 K excitation temperature seems high, but high excitation is predicted by shock models (Draine et al. (1983)) and is required to explain observed line ratios in IC 443 (Rho et al. (2000), Richter, Graham, & Wright (1995)). Future observations should be able to distinguish which shocks produce the S(3) and S(9) lines, by measuring the width of the lines. According to our cartoon model, much of the S(3) line comes from region M and would have a width of 100 km s<sup>-1</sup>. The S(9) line comes upstream from region C and would have a width of 30 km s<sup>-1</sup>. The line profiles of progressively higher excitation lines would evolve from wide to narrow lines. ### 5.2 Unidentified line at 74.3 $`\mu `$m One relatively bright line remains unidentified. Figure 5.1 shows the spectrum of the 74.3 $`\mu `$m line for 3C 391:BML. The line was present in our complete LWS spectra of W 28, W 44, and 3C 391, and it has also been seen in the remnant RCW 103 (Oliva et al. 1999b ) and the planetary nebula NGC 7027 (Liu et al. (1996)). Table 8 lists the brightnesses and central line wavelengths for our remnants. The wavelengths are all consistent with 74.26 $`\mu `$m. All of the spectra that contain this line contain also contain both atomic fine structure and molecular lines. Based on our list of plausible lines compiled above in Tables 3 and 4, and a careful inspection of all atomic fine structure lines, we can find no candidate atomic or ionic source for the line. (The only fine structure line that matches the observed wavelength is a transition among high energy levels of \[Ti III\], which is extremely unlikely based both on abundance and energetics.) Therefore, we suspect this line is molecular in origin. Despite searching the JPL spectral line database (Pickett et al. (1996)), we cannot identify the origin of this line, which is apparently caused by an unusual molecule or meta-stable state not yet found in laboratories or planetary atmospheres. ## 6 Continuum and Grain Destruction ### 6.1 Continuum spectrum The moderate-resolution far-infrared spectra are dominated by a bright continuum underlying the spectral lines. Because the remnants we discuss here are in the galactic plane, much of this continuum emission is likely to be unrelated to the remnant, produced instead by clouds along the line of sight both in front of and behind the remnant, as well as the molecular cloud close to the remnant. We have attempted to extract the remnant contribution to the continuum emission by searching for distinct spatial and spectral variations. In this paper we present a spectral separation of the remnant emission from the unrelated emission; in paper 1 and in future work, we use the reference positions and spatial correlations with the O I line brightness<sup>3</sup><sup>3</sup>3 We would like to clarify an issue from paper 1. Despite a statement by Oliva et al. (1999a) that our continuum for W 44 is 10 times too bright compared to IRAS data, the continuum brightness of the remnant as observed with ISO appears real. The IRAS data suffer from lower angular resolution, which makes it extremely difficult to separate the remnant from unrelated galactic emission. The continuum emission will be discussed in more detail in a forthcoming paper.. The continuum spectrum of 3C 391:BML is shown in Figure 6.1. Even by inspection, it is evident that the spectrum cannot be fitted with a modified blackbody: there is a significant positive excess at wavelengths shorter than 120 $`\mu `$m. It is not surprising for the continuum spectrum to require multiple components. In the COBE FIRAS observations of the galactic plane in a $`7^{}`$ beam, the spectrum of the inner galaxy at wavelengths of 100–200 $`\mu `$m has contributions from two components: 20 K dust associated with atomic gas, and 13 K dust associated with molecular gas (Reach et al. (1995), Lagache et al. (1998)). At wavelengths shorter than 100 $`\mu `$m, continuing all the way to the PAH features in the mid-infrared, emission from small, transiently-heated dust grains becomes important (Désert, Boulanger, & Puget (1990)). For our supernova remnant observations in an $`80^{\prime \prime }`$ beam, we expect all components to be brighter than in the COBE/FIRAS data, because (1) the atomic gas column density peaks sharply in the galactic plane, (2) the lines of sight pass directly through giant molecular clouds, and (3) the emission from the remnants themselves are completely beam-diluted in the COBE FIRAS beam. To estimate the remnant contribution to the spectrum, we used a nominal spectrum for dust in the local interstellar medium, which big grains, at a temperature of 17.5 K, and small grains with a brightness that matches the COBE Diffuse Infrared Background Experiment (Reach & Boulanger (1998) and references therein). We use two components, for the atomic and molecular gas, with radiation fields equal to $`\chi `$ times the local interstellar radiation field. For each component, the temperature of big grains scales as $`\chi ^{\frac{1}{6}}`$, and the brightness of the small grain emission scales as $`\chi `$. The atomic regions were assumed to have $`\chi =2`$ (to account for the enhanced radiation field in the inner Galaxy), and the molecular regions were assumed to have $`\chi =0.5`$ (to account for extinction of the ultraviolet radiation field inside molecular clouds). The scale of the model and the ratio of the two components was set to match the far-infrared spectrum at wavelengths longward of 120 $`\mu `$m. The continuum spectral models are shown for each remnant in Fig. 6.1. The models with a single diffuse ISM component cannot explain the brightness of the continuum at wavelengths shortward of 110 $`\mu `$m; there is a clear excess at 40–110 $`\mu `$m. The amount of excess emission depends rather sensitively on the assumed spectrum of the unrelated interstellar material, so the present results must be used with caution. There is a good correlation of the 74 $`\mu `$m continuum with the O I line brightness (paper 1), supporting our claim that a significant portion of the 40–110 $`\mu `$m emission is from the remnants. There is also a tentative detection of 3C 391 at 60 $`\mu `$m using IRAS data (Arendt (1996)). This excess emission has a color temperature of $`40`$ K, corresponding to an enhancement of the radiation field by more than 2 orders of magnitude. Such a high radiation does not occur in the diffuse interstellar medium over many-pc scales such as observed here. Instead, we must be seeing dust that is either physically different from that in the diffuse interstellar medium, or dust that is excited locally, for example behind the shock fronts. The optical depth of the remnant emission at 100 $`\mu `$m is $`\tau 2\times 10^4`$, and the total surface brightness of the far-infrared continuum, assuming a $`\nu ^2B_\nu (42\mathrm{K})`$ spectrum, is $`2\times 10^2`$ erg s<sup>-1</sup> cm<sup>-2</sup> sr<sup>-1</sup>. Thus the infrared continuum is significantly brighter than any of the individual spectral lines. The luminosity of the excess emission, within the single LWS beam that we observed for each remnant, ranges from $`600L_{}`$ for 3C 391:BML to $`10L_{}`$ for W 44 and W 28; however, these numbers are rather uncertain. Using the optical depth required to produce the excess emission, the opacity per unit mass for ‘normal’ interstellar dust (Draine & Lee (1984)), the mass of shock-excited dust within the single LWS beam that we observed for 3C 391:BML, is $`1M_{}`$. ### 6.2 Grain destruction In order for the Fe and Si emission lines to be as bright as observed, at least part of the refractory dust mass must have been vaporized. In the cool diffuse cloud toward $`\zeta `$ Oph, 95% of Si and $`>99`$% of Fe are locked in dust grains (Savage & Sembach (1996)). Using our single-slab model, we found that 50% of the Si and 30% of the Fe are required to be in the post-shock gas, which would mean that about half of the dust mass is destroyed by the shock. The single-slab model will, however, give less accurate measures of the abundances for Fe and Si relative to O (for which the model parameters were determined), because the Fe and Si lines we observed are from higher energy levels than those of O. There is likely to be a stratified post-shock emitting region, with somewhat warmer regions producing some of the Si and Fe lines. Thus, the exact fraction of refractory dust mass that was vaporized is highly uncertain. We have experimented with a multiple-slab model, using the approximate density and temperature profiles behind a 100 km s<sup>-1</sup> J-shock into $`10^3`$ cm<sup>-3</sup> gas (Hollenbach & McKee (1989)), to see how much of the observed 35.2 and 26 $`\mu `$m line brightnesses could be produced by warmer layers closer to the shock. The multiple slab model predicts relatively brighter 26 $`\mu `$m lines for a given column density, by up to a factor of two in extreme cases. Thus it is possible that the gas-phase abundance of Fe was overestimated by up to a factor of 2 using the single-slab model, meaning that as little as $`15`$% of the Fe is required to produce the observed lines. An identical result is obtained for Si. Even with the lower gas-phase abundance of Si and Fe, we can conclude that a non-trivial fraction of the Fe-bearing dust mass was vaporized. It is essentially impossible to produce all of the observed brightness of the 26 $`\mu `$m line from normally-depleted pre-shock interstellar matter. Using the observed traces of the vaporized and surviving dust grains, we can work backwards to estimate the fraction of dust that was destroyed. The total mass of Si and Fe vapor in 3C 391:BML is about 0.11 $`M_{}`$, within the LWS beam. Based on the dust-phase elemental abundances in a cold interstellar cloud, the total mass of elements (mainly O, Fe, C, Si, and Mg) in dust is 3.6–5.0 times the mass of Si and Fe in dust, with the lower number coming from B-star reference abundances and the higher number coming from Solar reference abundances (Savage & Sembach (1996)). Therefore, we infer that about 0.5 $`M_{}`$ of dust was vaporized by shocks into 3C 391:BML. For comparison, our far-infrared conntinuum observations indicated at least 1 $`M_{}`$ of dust remaining in solid form after the shock. Therefore, we estimate that approximately 1/3 of the dust mass was vaporized by the shocks. How does our inferred fraction of dust vaporization compare to theoretical predictions? Jones, Tielens, & Hollenbach (1996) have calculated the destruction probabilities for dust grains in shock fronts with a range of shock velocities and pre-shock densities. Their models included sputtering and shattering: sputtering removes material from the grain via collisions with the hot gas, while shattering destroys large grains and creates small grains from them. The models for $`v_s=100`$ km s<sup>-1</sup> are most relevant for the shocks we discuss in this paper. Clearly there are faster shocks that create the X-ray-emitting plasma, and slower shocks that create the molecular line emission. But the ionic lines from which we determined the abundances arise from a layer with a post-shock density of order $`10^3`$ cm<sup>-3</sup>4) into which shocks with a velocity of about 100 km s<sup>-1</sup> have been driven (§2.3). The model of Jones et al. (1996) with parameters most similar to our present case has pre-shock density $`n_0=25`$ cm<sup>-3</sup>. While their models do not extend to higher $`n_0`$, the dependence on $`n_0`$ is surprisingly mild, with a slight trend toward more dust destruction for denser shocks. The models of Jones et al. (1996) predict that about 37% of silicate dust and 12% of graphite dust are destroyed, mostly via non-thermal sputtering. These predictions are in excellent agreement with our results. ## 7 Nature of the shock fronts powering the emission lines There are at least two significant theoretical models that predict the structure (density and temperature) and emission behind radiative shocks comparable to those that seem to be occurring in the molecular supernova remnants we observed. One model, by Hollenbach & McKee (1989; hereafter HM89), includes an extensive network of chemical reactions and detailed treatment of the dynamics of the post-shock gas (from Hollenbach & McKee (1979)). The other model, by Hartigan, Raymond, & Hartmann (1987; hereafter HRH87), includes more ionization states and more ultraviolet and optical lines. We will use the HM89 model as a benchmark, and we will use the \[O I\] 63 $`\mu `$m line as the normalizing factor. In the HM89 model, the 63 $`\mu `$m line is a significant coolant and should reliably indicate the total energy passing through the shock. Based on the observed width of the line (§2.3), we consider the models with shock velocity $`V_s=100`$ km s<sup>-1</sup>. The HM89 model at that velocity that matches the observed \[O I\] 63 $`\mu `$m brightness has pre-shock density $`n_0=10^3`$ cm<sup>-3</sup>. Such a pre-shock density is reasonable considering that we are observing shocks into molecular clouds. This shock corresponds to what we called ‘region M’ (§4.2), although we suggested a somewhat lower pre-shock density of $`10^2`$ cm<sup>-3</sup> for this shock. In fact we do not directly measure $`n_0`$ and instead estimated it from the density in the emitting region, so a higher $`n_0`$ could be consistent with the observations. The compression by the shock depends on many unknown factors, including the strength of the preshock magnetic field. Before stepping through the comparison of more observed results to the models, we should clarify that the models are themselves significantly different. The physics of shock fronts, especially into dense gas and dust clouds, with so many modes of energy release and interaction, is exceedingly complex and cannot be expected to be accurately predicted by a simple theoretical model. Let us compare the HM89 and HRH87 models for $`n_0=10^3`$ cm<sup>-3</sup> and $`V_s=100`$ km s<sup>-1</sup>. The HRH87 model is called D100 in their paper. The magnetic field is negligible in the D100 model, but comparing the E100 and B100 models it appears that a stronger magnetic field makes relatively little difference in predicted line brightness within the context of the HRH87 model. Comparing the line brightnesses between the HM89 and HRH87 models, it is clear that there are very significant differences. An extreme case is the \[O I\] 6300 Å line, which is very bright in the HM89 model but weak in the HRH87 model; the difference is a factor of about 8. A significant part of the difference could be neglect of the higher ionization states in the HM89 model; we suspect the \[O II\] and \[O III\] optical lines are much brighter than \[O I\]. Other significant differences are the importance of Ly$`\alpha `$, many high-excitation ultraviolet lines, and the 2-photon continuum; all of these are very bright in the HRH87 model. Part of the difference could be neglect of the radiative transfer through dust in the HRH87 model; we expect most of the ultraviolet photons are absorbed by the cooler layers of the post-shock gas and dust or by the remaining unshocked gas and dust. Comparing our observations to the theoretical shock models, we reach the following conclusions: Dust emission—We observe a far-infrared continuum toward the shock fronts, in excess of what we would expect from dust in the quiescent pre-shock gas. The dust has a color temperature distinctly warmer than dust in molecular or atomic clouds. This could be due to significant collisional or radiative heating in the post-shock gas raising the equilibrium temperature of the grains, or it could be due to a significant difference in the size distribution of the post-shock grains. Indeed the post-shock size distribution is predicted to be significantly different because the larger grains are shattered (Jones, Tielens, & Hollenbach (1996)). The HM89 model predicts too little \[Fe II\] 26 $`\mu `$m and \[Si II\] 35 $`\mu `$m emission by about a factor of 10. This is probably due to treatment of dust destruction: the elemental abundances in the HM89 model are set at depleted cosmic abundances typical of what may occur in the pre-shock gas, while dust destruction is now predicted to be significant at these shock velocities (Jones, Tielens, & Hollenbach (1996)). The HM89 model predicts some infrared continuum, but its flux is weaker than the \[O I\] 63 $`\mu `$m line, which does not agree with our observations: we find that the continuum flux is much greater than that of any of the emission lines. The HRH87 is on the opposite extreme: they use the entire cosmic abundance in the gas phase, implicitly assuming complete destruction of the dust and predicting no infrared continuum. Infrared versus ultraviolet—The infrared emission escapes from the shocks and arrives at our telescopes, while the ultraviolet (and optical) emission either never escapes the shock or is absorbed before reaching us. The HRH87 model predicts copious ultraviolet emission, with the strongest emission (in order of brightness) arising from Ly$`\alpha `$, 2-photon continuum, C III 977 Å, and He II 304 Å. The sum of the 5 infrared lines they report is still 20 times weaker than Ly$`\alpha `$ and 3 times weaker than just the He II 304 Å line. Thus ultraviolet lines (and 2-photon continuum) are by far the dominant cooling behind radiative shocks according to the HRH87 model. On the other hand, the HM89 model predicts a significantly different spectrum, with relatively brighter optical and infrared lines. The brightest line in the HM89 model is the optical \[O I\] 6300 Å line, followed by H$`\alpha `$ and comparable brightnesses of infrared \[O I\] 63 $`\mu `$m, optical \[S II\] 6731 Å and \[N II\] 6560 Å, and ultraviolet C II\] 2326 Å lines. Remarkably, the brightness of the C II\] 2326 Å and \[O II\] 3726+3729 Å lines, and H$`\alpha `$ and H$`\beta `$, agree for the two models, but there are many brighter ultraviolet lines in the HRH87 model. The difference between the two models is at least partially due to the different assumptions about dust grains. The HRH87 model has no dust grains, so that ultraviolet transitions that are optically thick can scatter from ion to ion and eventually escape the shock or have a quantum decay into two photons. If dust grains were included, the ultraviolet photons would not last long before being absorbed. The ultraviolet photons are produced in a thin layer behind the shock, and their energy would heat the surrounding gas, as a radiative precursor into the pre-shock gas and as an enhanced radiation field in the cooler layers of the post-shock gas. In the HM89 model, the radiative transfer does take into account absorption of the ultraviolet lines by dust grains (Hollenbach & McKee (1979)), which explains why this model predicts much weaker ultraviolet emission than the HRH87 model. Our detection of a far-infrared continuum suggests that the models must include the effects of dust on the ultraviolet photons. Molecules—We detected both the S(3) and S(9) lines of H<sub>2</sub>2.2), as well as OH and weak H<sub>2</sub>O emission (Reach & Rho (1999)). The model of HRH87 does not include molecules, so we will not discuss it here. The model of HM89 includes a chemical reaction network, and they predict the brightness of the S(3) and S(9) lines to be $`9\times 10^6`$ and $`8\times 10^5`$ erg s<sup>-1</sup> cm<sup>-2</sup> sr<sup>-1</sup>, respectively. The predicted brightness of the S(9) line is about half of what we observed, suggesting that a significant fraction of the S(9) line could be produced by a J-type shock. However, the predicted brightness of the S(3) line is far lower than we observed. Indeed, for all 4 of the molecular shocks we observed, the S(3) line is brighter than the S(9) line; this is contrary to the model, which predicts S(3) about 10 times weaker than S(9). Instead, it appears that a denser, non-dissociating shock is required to explain the S(9) line brightness. This picture agrees with presence of a preshock cloud with a range of densities, as shown in our cartoon model. Therefore, it appears that the shock models cannot simultaneously predict both the ionic lines and (all of) the molecular lines with $`n_0=10^3`$ cm<sup>-3</sup> and $`V_s=100`$ km s<sup>-1</sup>. Instead, it appears that a denser, non-dissociating shock is required to explain the S(3) line brightness. To explain the brightness and excitation of the OH and H<sub>2</sub>O emission, and the $`30`$ km s<sup>-1</sup> width of the CO and CS lines, we suggested C-type shocks into gas with density $`>10^4`$ cm<sup>-3</sup> (Reach & Rho 98, 99). Models for such shocks were developed by Draine, Roberge, & Dalgarno (1983; hereafter DRD). For $`n_0=10^4`$ cm<sup>-3</sup> and $`V_s=30`$ km s<sup>-1</sup>, the DRD models predict that H<sub>2</sub> is the dominant coolant, with the brightest lines being S(3) and S(5). In the DRD model, the S(9) line is about 10 times weaker than the S(3) line. Given that we observe S(3) and S(9) lines of comparable brightness, it appears that we are required to have the bulk of the the S(3) line arise from a C-type shock and part of the S(9) line arise from a J-type shock. The observed surface brightness of the S(3) line is lower than the predicted brightness from the DRD model ($`I_{pred}=3\times 10^3`$), but that could be due to beam dilution of small, dense clumps that produce the C-type shock. Comparison to a model of W 44— Two recent papers by Cox et al. (1999) and Shelton et al. (1999) presented a model for the remnant W 44. In this model, all observed properties of the remnant, from radio emission to $`\gamma `$-rays, and specifically including the infrared lines, are produced by a blast wave into a smooth medium with a density of 6 cm<sup>-3</sup>. This model is in marked contrast to the one we presented here, with pre-shock densities ranging from $`<1`$ to $`>10^2`$ cm<sup>-3</sup>, and some discussion may help to resolve confusion about how they could both be explaining the same data. First, we note that the evidence for molecular cloud interaction is very strong. Shocked CO line with widths of $`30`$ km s<sup>-1</sup> were revealed by Seta et al. (1998) and ourselves (this paper, Fig. 2.2). Second, we note that the uniform pre-shock conditions in the Cox et al. model are not realistic in a large region of the multiple-phase interstellar medium. In our model, we attribute the molecular line observations of W 44, including H<sub>2</sub> near-infrared emission (Seta et al. (1998)), infrared CO lines (Reach & Rho (1998)), wide millimeter CO and other lines (Seta et al. (1998)), and radio OH 1720 MHz masers (Claussen et al. (1997)), to dense, shocked gas. These lines cannot be produced by an ionized low-density post-shock region. The analogy to IC 443 is useful—W 44 is very similar in many regards, with similar brightnesses of the 63 $`\mu `$m line and brightnesses and widths of the millimeter-wave CO (and other) lines—and few would deny that IC 443 is interacting with a molecular cloud. Finally, the observed brightness of the \[O I\] 63 $`\mu `$m line cannot easily be explained by a low-density medium, in contrast to the discussion claimed by Cox et al. (1999). The excitation of \[O I\] evidenced by the 145/63 $`\mu `$m ratio requires a higher density region (Fig. 3.2). Obviously, \[O I\] also requires a neutral region, for if the H is ionized the O will also be ionized. The observed lines were so ‘shockingly’ bright that we concluded immediately that they were due to a dense shock (Reach & Rho (1996)). In this paper we considered a wide range of densities and ruled out the low densities because (1) they are inconsistent with the line ratios, and (2) the path length required through a low-density gas would be longer than the remnant size in order to build up enough column density to produce the observed line brightness. Cox et al. reduced the brightness of the 63 $`\mu `$m by a factor of 100 to match their models. One of the reduction factors is a a factor of 10 to match the estimated average over the entire remnant. Another factor applied by Cox et al. was to enhance their prediction because of limb brightening and local rippling of the shell. While this is perfectly reasonable, this correction was applied to the low-resolution 63 $`\mu `$m line observations, which trace relatively cool gas, and not to the higher-resolution H$`\alpha `$ observations (Giacani et al. (1997)), which trace warmer gas that might be even more edge brightened. All of these factors probably cancel out to first order when the average over our large beam is made, so that the \[O I\]/H$`\alpha `$ ratio is not matched by their model. The truth is probably that there were both high and low density regions in the interstellar medium around the W 44 progenitor star (as considered in this paper and by Chevalier (1999)). The particular lines of sight studied in this paper are centered on a special regions, where evidence for a dense shocks are overwhelming. Away from these positions, it is likely that the lower-density shocks contribute some of the low-level 63 $`\mu `$m emission. In Figure 4.2, we are pointed right at a region C clump, while the average over a large region would include different regions in different proportions. Because the 63 $`\mu `$m line brightness is expected to scale approximately as $`n_0V_s`$ (HM89), the denser regions contribute relatively more to the observed brightness, when they are present. Further observational work may be able to separate the dense shocked gas from the more rarefied and widespread interclump gas, using imaging to resolve the clumpy spatial distribution and spectroscopy to separate the narrower line widths. ## 8 Conclusions Based on ISO spectroscopic observations of supernova remnants interacting with molecular clouds, we found evidence for a wide range of pre- and post-shock conditions; these properties are summarized in Table 5 and illustrated in Figure 4.2. Of the three density regimes that we identified, most of the infrared emission arose from the shocks into moderate-density molecular gas (‘region M’). The energy of these shocks is radiated via continuum emission from surviving dust grains and \[O I\] 63 $`\mu `$m and \[Si II\] and other infrared atomic fine structure lines. The dust continuum was difficult to separate from the far-infrared emission from cold, unrelated dust, but the shocked dust was evident at 80–100 $`\mu `$m as a warmer component of the spectrum. Some of the dust grains must have been destroyed, because we observe bright emission lines from dust ‘vapors,’ including the \[Fe II\] 26 $`\mu `$m line. To produce the observed line dust vapor lines, we require that 1/3 of the Si and Fe-bearing dust mass was vaporized in the shocks. Using the Fabry-Perot observations, the width of the 63 $`\mu `$m line was found to be $`100`$ km s<sup>-1</sup>. Theoretical models can reproduce the brightness of the 63 $`\mu `$m line with such shock velocities and the properties of ‘region M’ (Hollenbach & McKee (1989)); theoretical models can also reproduce the observed amount of grain destruction (Jones, Tielens, & Hollenbach (1996)). Higher-density gas was required to explain the bright H<sub>2</sub> line emission, leading to our ‘region C,’ for clumps. Lower-density gas was required to explain the higher-ionization lines, leading to our ‘region A,’ for atomic gas. A general conclusion from these observations is that molecular shock fronts are copious producers of infrared emission. We anticipate many applications of these types of observations for further understanding the nature of molecular supernova remnants, such as the fate of dense clumps, the fate of dust grains, and the effect of cooling on the remnant evolution. In turn, these studies can elucidate, to some extent, the nature of the molecular clouds before the shocks. If molecular clouds were uniform, or nearly so, then we would expect an orderly progression of pre-shock to post-shock gas, with radiative coolants characteristic of a single shock velocity. On the contrary, we observe emission from coolants as diverse as multi-atom molecules, multiply-ionized atoms, and dust grains and their vapors. We would like to thank Pierre-Olivier Lagage for helpful discussions in interpreting atomic fine structure lines, and we thank Emmanuel Caux, Steve Lord, and Sergio Molinari for clarifying details of the LWS instrument and calibration.
warning/0007/cond-mat0007354.html	ar5iv	text	# Morphology of Ion-Sputtered Surfaces ## I Introduction Sputtering is the removal of material from the surface of solids through the impact of energetic particles . It is a widespread experimental technique, used in a large number of applications with a remarkable level of sophistication. It is a basic tool in surface analysis, depth profiling, sputter cleaning, micromachining, and sputter deposition. Perhaps the largest community of users is in the thin film and semiconductor fabrication areas, sputter erosion being routinely used for etching patterns important to the production of integrated circuits and device packaging. To have a better control over this important tool, we need to understand the effect of the sputtering process on the surface morphology. In many cases sputtering is routinely used to smooth out surface features. On the other hand, some investigations indicate that sputtering can also roughen the surface. Consequently, sputter erosion may have different effects on the surface, depending on many factors, such as incident ion energy, mass, angle of incidence, sputtered substrate temperature and material composition. The experimental results on the effect of sputter erosion on the surface morphology can be classified in two main classes. There exists ample experimental evidence that ion sputtering can lead to the development of periodic ripples on the surface . Depending on the sputtered substrate and the sputtering conditions these ripples can be surprisingly straight and ordered. However, a number of recent investigations have provided rather detailed and convincing experimental evidence that under certain experimental conditions ion eroded surfaces become rough, and the roughness follows the predictions of various scaling theories . Moreover, these investigations did not find any evidence of ripple formation on the surface. Up to recently these two morphological features were treated separately and no unified theoretical framework describing these morphologies was available. The first widely accepted theoretical approach describing the process of ripple formation on amorphous substrates was developed by Bradley and Harper (BH) . This theory is rather successful in predicting the ripple wavelength and orientation in agreement with numerous experimental observations. However, a number of experimental results have systematically eluded this theory. For example, the BH theory predicts an unlimited exponential increase in ripple amplitude in contrast with the observed saturation of the surface width. Similarly, it cannot account for surface roughening, or for ripple orientations different from those defined by the incoming ion direction or perpendicular to it. Finally, recent experiments have observed ripples whose wavelength is independent of the substrate temperature, and linear in the ion energy, in contrast with the BH prediction of a ripple wavelength which depends exponentially with temperature and decreases with ion energy. In the light of the accumulated experimental results, it is clear that a theory going beyond the BH approach is required, motivating the results described in this paper. Thus here we investigate the morphology of ion-sputtered amorphous surfaces aiming to describe in an unified framework the dynamic and scaling behavior of the experimentally observed surface morphologies. For this we derive a nonlinear theory that describes the time evolution of the surface morphology. At short time scales the nonlinear theory predicts the development of a periodic ripple structure, while at large time scales the surface morphology may be either rough or dominated by new ripples, that are different from those existing at short time scales. We find that transitions may take place between various surface morphologies as the experimental parameters (e.g. angle of incidence, penetration depth) are varied. Usually stochastic equations describing growth and erosion models are constructed using symmetry arguments and conservation laws. In contrast, here we show that for sputter eroded surfaces the growth equation can be derived directly from a microscopic model of the elementary processes taking place in the system. A particular case of our theory was presented in . In addition, we show that the presented theory can be extended to describe low temperature ripple formation as well. We demonstrate that under certain conditions ion-sputtering can lead to preferential erosion that appears as a surface diffusion term in the equation of motion even though no mass transport along the surface takes place in the system. To distinguish it from ordinary surface diffusion, in the following we refer to this phenomenon as effective smoothing (ES). We calculate analytically an effective surface diffusion constant accounting for the ES effect, and study its dependence on the ion energy, flux, angle of incidence, and penetration depth. The effect of ES on the morphology of ion-sputtered surfaces is summarized in a morphological phase diagram, allowing for direct experimental verification of our predictions. A restricted study along these lines appeared in . The paper is organized as follows. In Section II we review the recent advances in the scaling theory of rough (self-affine) interfaces. Section III is dedicated to a brief overview of the experimental results on surface morphology development under ion sputtering. A short summary of the theoretical approaches developed to describe the morphology of ion sputtered surfaces is presented in Section IV. This section also contains a short description of Sigmund’s theory of sputtering, that is the basis for our calculations. In Section V we derive the nonlinear stochastic equation describing sputter erosion. Analysis of this equation is presented in Section VI, discussing separately both the high and low temperature ripple formation. We compare the predictions of our theory with experimental results on surface roughening and ripple formation in Section VII, followed by Section VIII, that summarizes our findings. ## II Scaling theory In the last decade we witnessed the development of an array of theoretical tools and techniques intended to describe and characterize the roughening of nonequilibrium surfaces and interfaces . Initiated by advances in understanding the statistical mechanics of various nonequilibrium systems, it has been observed that the roughness of many natural surfaces follows rather simple scaling laws, which can be quantified using scaling exponents. Since kinetic roughening is a common feature of ion-bombarded surfaces, before we discuss the experimental results on sputtering, we need to introduce the main quantities characterizing the surface morphology. Let us consider a two-dimensional surface that is characterized by the height function $`h(x,y,t)`$. The morphology and dynamics of a rough surface can be quantified by the interface width, defined by the rms fluctuations in the height $`h(x,y,t)`$, $$w(L,t)\sqrt{\frac{1}{L^2}\underset{x,y=1,L}{}[h(x,y)\overline{h}]^2},$$ (1) where $`L`$ is the linear size of the sample and $`\overline{h}`$ is the mean surface height of the surface $$\overline{h}(t)\frac{1}{L^2}\underset{x,y=1,L}{}h(x,y,t).$$ (2) Instead of measuring the roughness of a surface over the whole sample size $`L\times L`$, we can choose a window of size $`\mathrm{}\times \mathrm{}`$, and measure the local width, $`w(\mathrm{})`$. A general property of many rough surfaces is that the roughness depends on the length scale of observation. This can be quantified by plotting $`w(\mathrm{})`$ as a function of $`\mathrm{}`$. There are two characteristic regimes one can distinguish: (i) For length scales smaller than $`\mathrm{}_\times `$, the local width increases as $$w(\mathrm{})\mathrm{}^\alpha ,$$ (3) where $`\alpha `$ is the roughness exponent. If we are interested in surface phenomena that take place at length scales shorter than $`\mathrm{}_\times `$ then we cannot neglect the roughness of the surface. In this regime the roughness is not simply a number, but it depends on the length scale accessible to the method probing the surface. (ii) For $`\mathrm{}\mathrm{}_\times `$, $`w(\mathrm{})`$ is independent of $`\mathrm{}`$, thus, at length scales larger than $`\mathrm{}_\times `$, the surface is smooth. In this regime we can characterize the surface roughness with the saturation width $`w_{sat}=w(\mathrm{}_\times )`$. The dynamics of the roughening process can be best characterized by the time dependent total width (1). At early times the total width increases as $`w(L,t)t^\beta `$, where $`\beta `$ is the growth exponent. However, for finite systems, after a crossover time $`t_\times `$, the width saturates, following the Family-Vicsek scaling function $`w(L,t)t^\beta g\left({\displaystyle \frac{t}{L^z}}\right),`$ (4) where $`z=\alpha /\beta `$ is the dynamic exponent and $`g(u1)1`$, while $`g(u1)u^\beta `$. Scaling relations such as Eq. (4) allow us to define universality classes. The universality class concept is a product of modern statistical mechanics, and encodes the fact that there are but a few essential factors that determine the exponents characterizing the scaling behavior. Thus different systems, which at first sight may appear to have no connection between them, behave in a remarkably similar fashion. The values of the exponents $`\alpha `$ and $`\beta `$ are independent of many “details” of the system, and they are uniquely defined for a given universality class. In contrast, other quantities, such as $`A`$, $`\mathrm{}_\times `$, or $`w_{sat}`$, are non-universal, i.e. they depend on almost every detail of the system. ## III Experimental results The morphology of surfaces bombarded by energetic ions has long fascinated the experimental community. Lately, with the development of high resolution observation techniques such as atomic force (AFM) and scanning tunneling (STM) microscopies, this problem is living a new life. The various experimental investigations can be classified into two main classes. First, early investigations, corroborated by numerous recent studies, have found that sputter eroded surfaces develop a ripple morphology with a rather characteristic wavelength of the order of a few micrometers . However, a number of research groups have found no evidence of ripples, but observed the development of apparently random, rough surfaces , that were characterized using scaling theories. In the following we summarize the key experimental observations for both ripple development and kinetic roughening. ### A Ripple formation The ripple morphology of ion bombarded surfaces has been initially discovered in the mid 1970’s . Since then, a number of research groups have provided detailed quantitative results regarding the ripple characteristics and dynamics of ripple formation. It is beyond the scope of this paper to offer a comprehensive review of the vast body of the experimental literature on the subject. Thus, we selected a few experiments that offer a representative picture of the experimental features that appear to be common to different materials. Angle of incidence: An experimental parameter which is rather easy to change in sputtering is the angle of incidence $`\theta `$ of the incoming ions relative to the normal to the average surface configuration. Consequently, numerous research groups have investigated the effect of $`\theta `$ on the ripples. These results indicate that ripples appear only for a limited range of incidence angles, which, depending on materials and ions involved, typically vary between 30 and 60. Support for a well defined window in $`\theta `$ for ripple formation was offered by Stevie et al. , who observed abrupt secondary ion yield changes (correlated with the onset of ripple morphology development) in experiments on $`6`$ and $`8`$ keV O$`{}_{2}{}^{}{}_{}{}^{+}`$ sputtering of Si and $`8`$, $`5.5`$, and $`2.5`$ keV O$`{}_{2}{}^{}{}_{}{}^{+}`$ sputtering of GaAs at incidence angles between 39 and 52. These results were corroborated by Karen et al. , who investigated ripple formation on GaAs surfaces under bombardment with $`10.5`$ keV O$`{}_{}{}^{+}{}_{2}{}^{}`$ ions. They found that ripple formation takes place for angles of incidence between 30 and 60 (see Table I of Ref. ). Similarly, Wittmaack found that ripple formation occurs at incidence angles between 32 and 58 during $`10`$ keV O$`{}_{2}{}^{}{}_{}{}^{+}`$-ion bombardment of a Si target. Temperature dependence: Another parameter that has been found to influence the ripple structure, and in particular the ripple wavelength, is the temperature of the substrate, $`T`$. Two different behaviors have been documented: exponential dependence of the ripple wavelength on $`T`$ has been observed at high temperatures, while the wavelength was found to be constant at low temperatures. A series of experiments on the temperature dependence of ripple formation were reported by MacLaren et al. . They studied InP and GaAs surfaces bombarded with 17.5 KeV Cs<sup>+</sup> and 5.5 keV O$`{}_{}{}^{+}{}_{2}{}^{}`$ ion beams in the temperature range from $`50^{}`$ C to 200 C. For GaAs bombarded by Cs<sup>+</sup> ions the ripple wavelength increased from 0.89 $`\mu `$m to 2.1 $`\mu `$m as the temperature increased from 0 C to 100 C. Probably the most interesting finding of their study was that lowering the temperature, the ripple wavelength did not go continuously to zero as one would expect, since the surface diffusion constant decreases exponentially with temperature (see Sect. IV C), but rather at approximately 60 C it stabilized at a constant value. MacLaren et al. interpreted this as the emergence of radiation enhanced diffusion, that gives a constant (temperature independent) contribution to the diffusion constant. Recently, Umbach et al. have studied the sputter-induced ripple formation on SiO<sub>2</sub> surfaces using $`0.52.0`$ keV Ar ion beams. The temperature dependence of the ripple wavelength $`\mathrm{}`$ was investigated for temperatures ranging from room temperature to 800 C. It was found that for high temperatures, $`T400^{}`$ C, the ripple wavelength follows the Arrhenius law $`(1/T^{1/2})\mathrm{exp}(\mathrm{\Delta }E/2k_BT)`$, indicating the thermally activated character of the relaxation processes. However, at low temperatures the ripple wavelength was independent of temperature, indicating the presence of a temperature independent relaxation mechanism. Results indicating temperature independent non-diffusive relaxation have been reported for crystalline materials as well by Carter et al. . In these experiments Si bombarded with highly energetic $`1040`$ keV Xe<sup>+</sup>-ions led to ripple formation with wavelength $`\mathrm{}`$ 0.4 $`\mu `$m for angles of incidence close to 45. Changing the surface temperature from 100 K to 300 K the ripple wavelength and orientation did not change. This observation led the authors to conclude that the smoothing mechanism is not of thermal origin. They also found that the ripple wavevector is relatively insensitive to the primary ion energy. Flux and fluence dependence: The effect of the flux on the surface dynamics was studied by Chason et al. . In these experiments a $`1`$ keV Xe ion beam was directed towards a SiO<sub>2</sub> sample with an angle of incidence of 55. The typical incoming flux was 10<sup>13</sup> cm<sup>-2</sup>s<sup>-1</sup> and fluence (the number of incoming atoms per surface area, playing the role of time) was up to $`10\times 10^{15}`$ cm<sup>-2</sup>. The surface was analyzed using in situ energy dispersive X-ray reflectivity and ex situ AFM. It was found that the interface roughness, which is proportional to the ripple amplitude, increases linearly with the fluence. Similar experiments were performed on Ge(001) surfaces using $`0.3`$, $`0.5`$, and $`1`$ keV Xe ion beams for $`T=350^{}`$ C. For flux values up to $`3`$ $`\mu `$A/cm<sup>2</sup> and fluences up to $`6\times 10^{16}`$ cm<sup>-2</sup>, the roughness is seen to increase as the square of the flux. Ion energy: The ripple wavelength dependence on the incident ion energy and the angle of incidence was reported in Refs. . The experiments indicate that the ripple wavelength is linear in the energy, following $`\mathrm{}ϵ\mathrm{cos}\theta `$. Similar results were obtained in Ref. , providing an extensive study of ripple formation by secondary ion spectrometry and scanning electron microscopy. The ripple topography was observed during O$`{}_{}{}^{+}{}_{2}{}^{}`$ primary ion bombardment of a Si(100) substrate with ion energies ranging between $`1.5`$ keV and $`9`$ keV. No ripples were observed for energies less than $`1.5`$ keV or for high energies, such as $`1.5`$ keV and $`7`$ keV, when Ar<sup>+</sup> primary ions were used. The experiments indicate that the ripple wavelength increases linearly from 100 to 400 nm when the energy of the primary ion changes from $`1`$ to $`9`$ keV. Furthermore, the experimental data indicated that the primary ion penetration depth $`a`$ and the ripple wavelength $`\mathrm{}`$ are related by the empirical relation $`\mathrm{}=40a`$. The wavelength of the ripples is found to be independent of the primary ion flux and dependent merely on fluence, i.e. sputtered depth. The recent results by Umbach et al. provided further strong evidence for the linear relationship between the ion energy and the ripple wavelength for SiO<sub>2</sub> substrates (see below). Ripple amplitude: Indirect results on the ripple amplitude were presented by Vajo et al. in their study of the surface topography induced secondary ion yield changes on SiO<sub>2</sub> surfaces bombarded by O$`{}_{2}{}^{}{}_{}{}^{+}`$ ions. The authors have found that the yield changes exponentially in the first stages of ripple development and saturates for large sputtered depth. Direct evidence on ripple amplitude saturation was obtained by Erlebacher et al. , who measured the time evolution of the ripple amplitude in experiments bombarding Si(100) surfaces with $`0.75`$ keV Ar<sup>+</sup> ions. They found that, while at short times the ripple amplitude increases exponentially, it saturates Material Ion Angle Ion energy Ripple Ref. type of (keV) wavelength incidence ($`\mu `$m) GaAs(100) O$`{}_{2}{}^{}{}_{}{}^{+}`$ $`39^{}`$ 8 0.2 GaAs(100) O$`{}_{2}{}^{}{}_{}{}^{+}`$ $`42^{}`$ 5.5 0.1 GaAs(100) O$`{}_{2}{}^{}{}_{}{}^{+}`$ $`37^{}`$ 10.5 0.23 GaAs(100) O$`{}_{2}{}^{}{}_{}{}^{+}`$ $`42^{}`$ 5.5 0.13 GaAs(100) O$`{}_{2}{}^{}{}_{}{}^{+}`$ $`39^{}`$ 8.0 0.21 GaAs(100) O$`{}_{2}{}^{}{}_{}{}^{+}`$ $`37^{}`$ 10.5 0.27 GaAs(100) O$`{}_{2}{}^{}{}_{}{}^{+}`$ $`57^{}`$ 13 0.33 GaAs O$`{}_{2}{}^{}{}_{}{}^{+}`$ $`40^{}`$ 3.0 0.075 GaAs O$`{}_{2}{}^{}{}_{}{}^{+}`$ $`40^{}`$ 7.0 0.130 Ge(001) Xe<sup>+</sup> $`55^{}`$ 1 0.2 Si(001) O$`{}_{}{}^{+}{}_{2}{}^{}`$ $`41^{}`$ 6 0.4 Si(001) O$`{}_{}{}^{+}{}_{2}{}^{}`$ $`42^{}`$ 5.5 0.5 Si(100) O$`{}_{}{}^{+}{}_{2}{}^{}`$ $`39^{}`$ 8 0.5 Si(100) O$`{}_{2}{}^{}{}_{}{}^{+}`$ $`40^{}`$ 3 0.198 Si(100) O$`{}_{2}{}^{}{}_{}{}^{+}`$ $`40^{}`$ 5 0.302 Si(100) O$`{}_{2}{}^{}{}_{}{}^{+}`$ $`40^{}`$ 9 0.408 Si(100) Ar<sup>+</sup> $`67.5^{}`$ 0.75 0.57 Si Xe<sup>+</sup> $`45^{}`$ 40 0.4 Si O$`{}_{}{}^{+}{}_{2}{}^{}`$ $`37^{}`$ 12.5 0.35 SiO<sub>2</sub> Ar<sup>+</sup> $`45^{}`$ 0.5-2 0.2-0.55 SiO<sub>2</sub> Xe<sup>+</sup> $`55^{}`$ 1 0.03 TABLE I.: Summary of the ripple characteristics reported in sputter erosion experiments of non-metallic substrates. In all cases shown, the ripple wave vector is parallel to the ion beam direction. Note that a number of experiments have obtained indirect information on ripple formation from secondary ion yield changes. These have not been included in the table. after a crossover time has been reached. Furthermore, the experiments indicate that the crossover time scales with the temperature induced surface diffusion constant. Surface chemistry and other morphological features: While a number of attempts have been made to explain ripple formation based on chemical effects, such as O$`{}_{2}{}^{}{}_{}{}^{+}`$ variations , most of these studies were contradicted by subsequent investigations where such chemical component were not present. Furthermore, in Refs. it was unambiguously shown that the process of ripple formation is not caused by defects or inherited irregularities on the surface, but is determined merely by the primary ion characteristics. These results indicate that ripple formation is independent of microscopic details and the surface chemistry. Ripple formation on crystalline and metallic surfaces: As the discussed experimental results have indicated, ripple formation takes place under a variety of conditions and on surfaces of different materials, including both crystalline and amorphous materials. Despite the fact that Sigmund’s theory, the basis of all theories of ripple formation, has been developed for amorphous targets, it is worth noting that these approaches describe many features of ripple formation on crystalline surfaces as well. However, when discussing ripple formation on crystalline materials, we always have to be aware that additional effects, induced by the crystalline anisotropy, could be present. An example of ripple development on crystalline materials has been obtained for Ag(110) surfaces under low energy ($`ϵ800`$ eV) Ar<sup>+</sup> primary beam bombardment by Rusponi et al. . Ripples with wavelength of approximately 600 Å, oriented along the $`110`$ crystallographic direction, appeared in a temperature range 270 K $`T`$ $``$ 320 K at normal ion incidence. The ripple structure was found to be unstable at room temperature, i.e. substantial smoothing of the surface with time takes place. The structure depends on the ion flux and ion energy. Similar results are available for ion-sputtered Cu$`(110)`$ monocrystals using a $`1`$ keV Ar<sup>+</sup> ion beam . For normal incidence a well defined ripple structure was observed with wave vectors whose direction changes from $`001`$ to $`110`$ as the temperature of the substrate is raised. Off-normal sputtering generated ripples whose orientation depends both on the ion direction and the surface orientation. The authors suggested that this phenomenon can be explained in terms of anisotropic surface diffusion. Summary: As the presented results indicate, ripple formation on ion-sputtered surfaces has been observed by many groups in various systems (for a partial summary see Table I). The main experimental results, common to most studied materials, can be summarized as follows: $``$ Off-normal ion bombardment often produces periodically modulated structures (ripples) on the surfaces of amorphous and crystalline materials. The wavelength of the ripples $`\mathrm{}`$ is usually of the order of tenths of micrometers. $``$ For non metallic substrates, the orientation of the ripples depends on the angle of incidence $`\theta `$, and in most cases is either parallel or perpendicular to the direction of the incoming ions. $``$ At low temperatures the ripple wavelength is independent of $`T`$, while it follows the Arrhenius law $`\mathrm{}(1/T^{1/2})\mathrm{exp}(\mathrm{\Delta }E/k_BT)`$ at higher temperatures. $``$ Numerous experiments find that the ripple wavelength is proportional to the ion range, and thus to the ion energy for intermediate energies. $``$ The ripple wavelength in many cases is independent of the ion flux, but systematic flux dependence has also been reported. $``$ The amplitude of the periodic modulations grows exponentially for early times, but saturates after a typical crossover time has been reached. In many instances, the ripple wavelength $`\mathrm{}`$ is found to be independent of time. $``$ Evidence for ripple formation was obtained for different materials and different primary ions, suggesting that the mechanism responsible for ripple formation is largely independent of surface chemistry, chemical reactions on the surface, or defects. ### B Kinetic roughening Motivated by the advances in characterizing the morphology of rough surfaces, recently a number of experimental studies have focused on the scaling properties of surfaces eroded by ion bombardment . These experiments have demonstrated that under certain ion bombardment conditions ripples do not form, and the surface undergoes kinetic roughening. The goal of the present section is to review the pertinent experimental results, aiming to summarize the key features that a comprehensive theory should address. Surface roughness and dynamical exponents: In the experiments of Eklund et al. pyrolytic graphite was bombarded by 5 keV Ar ions, at an angle of incidence of 60. The experiments were carried out for two flux values, $`6.9\times 10^{13}`$ and $`3.5\times 10^{14}`$ ions s<sup>-1</sup> cm<sup>-2</sup>, the total fluences being $`10^{16}`$, $`10^{17}`$ and $`10^{18}`$ ions $``$ cm<sup>-2</sup>. STM micrographs indicated that large scale features develop with continuous bombardment, the interface becoming highly correlated and rough. The scaling properties have also been probed using the Fourier transform of the height-height correlation function, obtaining a dynamic exponent $`z`$ in the range $`1.61.8`$, and a roughness exponent in the range $`0.20.4`$. These exponents are consistent with the predictions of the continuum theory, describing kinetic roughening, proposed by Kardar, Parisi and Zhang (KPZ) , that predicts $`z1.6`$ and $`\alpha 0.38`$ (see section IV A 1). A somewhat larger roughness exponent has been measured for samples of iron bombarded with 5 keV Ar ions at an angle of incidence of 25 . The interface morphology was observed using STM, and the height–height correlation function indicated a roughness exponent $`\alpha =0.53\pm 0.02`$ . The mechanism leading to such a roughness exponent is not yet understood in terms of continuum theories, since for two dimensional surfaces the existing continuum theories predict $`\alpha `$ values of 0.38, 2/3 and 1 , far from the observed roughness exponent. Anomalous dynamic-scaling behavior of sputtered surfaces was reported by Yang et al.. The experiments performed on Si(111) surfaces with $`0.5`$ keV Ar<sup>+</sup> ions with flux $`0.2`$ $`\mu `$A/mm<sup>2</sup> in a wide range of substrate temperatures have provided evidence of scaling behavior in the limit of small distances $`r`$. The height–height correlation function has been found to follow $`C(r)=(h(r_o)h(r+r_o))^2r^{2\alpha }\mathrm{log}t`$, with $`\alpha 1.15\pm 0.08`$ for temperatures lower than 530 C. No roughening was observed for higher temperatures, demonstrating the temperature dependence of kinetic roughening. Temperature dependence: The effect of surface relaxation due to surface diffusion on roughening of GaAs(110) surfaces eroded by $`2`$ keV Ar<sup>+</sup> and Xe<sup>+</sup> was reported by Wang et al. . They found that both the height-height correlation function and the small scale roughness increase significantly faster during erosion at higher temperatures than at lower ones. The surface width in these experiments increased with $`\beta =0.3`$ at $`T=725`$ K and there was no evidence of scaling for lower temperatures, such as $`T=625`$ K. The roughness exponent has been determined as $`\alpha =0.38\pm 0.03`$. In general, Ref. concludes that on large scales the surfaces are rougher at higher temperatures, contrary to the expectation of smaller roughness due to increased relaxation by surface diffusion. Similar conclusions on the temperature dependence of the scaling properties were drawn in Ref. . A sharp transition between scaling regimes in ion-bombardment of Ge(001) surfaces with $`1`$ keV Xe ions was observed at $`T_c=488`$ K. The regimes above and below $`T_c`$ are characterized by dynamic scaling exponents $`\beta `$ with values 0.4 and 0.1, respectively. The surface roughness of Si(111) during low-energy (500 eV) ion bombardment at $`T=610`$ K was studied in Ref. using STM. It was found that the rough morphology is consistent with the early time behavior of the noisy Kuramoto-Sivashinsky (KS) equation (see Sect. IV A 3). The measured roughness exponent was $`\alpha =0.7`$ and the dynamic exponent was $`\beta =0.25`$. Low energy ion bombardment: Recently a number of experiments and simulations have focused on low energy ion bombardment (i.e., at energies $`50`$-$`500`$ eV), for which the secondary ion yields are smaller than one . In this systems, the effect of the ions is limited to the surface of the material, the collective effect created by the collision cascade being less relevant. Often, such low energy sputtering leads to layer-by-layer erosion, almost mirroring layer-by-layer growth phenomena. The effect of vacancy diffusion and Schwoebel barriers can be rather well studied in these systems, that include Ge(001) surface etching, by $`240`$ eV Xe ions , and Si(111) surfaces etched by $`100`$ eV Ar ions . In the absence of the collision cascade, ripple formation and kinetic roughening seen at higher energies, the subject of this paper, do not appear. Various experimental results on ion-bombardment induced surface roughening are summarized in Table II. These experiments demonstrate that kinetic roughening is one of the major experimental morphologies generated by ion bombardment. However, as Table II indicates, there is a considerable scattering in the scaling exponents. This scattering is not too disturbing at this point: accurate determination of the scaling exponents from experimental data is rather difficult, since often the scaling regime is masked by strong crossover effects. As we demonstrate later, due to the separation of the linear and nonlinear regimes, such crossovers are, indeed, expected in sputter erosion. Thus the main conclusion we would like to extract from this section is that numerous experiments do observe kinetic roughening, and find that scaling concepts can successfully characterize the surface morphology. It will be a major aim of the theory proposed here to account for the origin of kinetic roughening and predict the scaling exponents. ## IV Theoretical approaches The recent theoretical studies focusing on the characterization of various surface morphologies and their time evolution have revolutionized our understanding of growth and erosion (for reviews, see ). The physical understanding of the processes associated with interface roughening require the use of the modern concepts of fractal geometry, universality and scaling. In Sect. IV A we review the major theoretical contributions to this area, necessary to describe the morphology of ion-eroded surfaces. In Sects. IV B to IV E we then review the available theoretical approaches (whether through continuum equations or by the use of discrete atomistic models) that specifically describe surfaces eroded by ion-bombardment, emphasizing the procedures which allow to describe within a continuum approach some of the relevant physical processes taking place at the surface, such as surface diffusion and beam fluctuations. ### A Continuum theories of kinetic roughening The full strength of the continuum theories comes from the prediction of the asymptotic behavior of the growth process valid in the long time and large length scale limits. These limits are often beyond the experimentally or practically interesting time and length scales. A notable exception is sputter erosion, where both the short time ripple development and the asymptotic kinetic roughening have been observed experimentally. Consequently, next we discuss separately the continuum theories needed to understand sputter erosion. #### 1 Kardar-Parisi-Zhang (KPZ) equation The time evolution of a nonequilibrium interface can be described by the Kardar-Parisi-Zhang (KPZ) equation $$\frac{h}{t}=\nu ^2h+\frac{\lambda }{2}(h)^2+\eta .$$ (5) The first term on the rhs describes the relaxation of the interface due to the surface tension ($`\nu `$ is here a positive constant) and the second is a generic nonlinear term incorporating lateral growth or erosion. The noise, $`\eta (x,y,t)`$, reflects the random fluctuations in the growth process and is a set of uncorrelated random numbers with zero configurational average. For one dimensional interfaces the scaling exponents of the KPZ equation are known exactly, as $`\alpha =1/2`$, $`\beta =1/3`$, and $`z=3/2`$. However, for higher dimensions they are known only from numerical simulations. For the physically most relevant two dimensional surfaces we have $`\alpha 0.38`$ and $`\beta 0.25`$ . If $`\lambda =0`$ in Eq. (5), the remaining equation describes the equilibrium fluctuations of an interface which tries to minimize its area under the influence of the external noise. This equation, first introduced by Edwards and Wilkinson (EW) , can be solved exactly due to its linear character, giving the scaling exponents $`\alpha =(2d)/2`$ and $`\beta =(2d)/4`$. For two dimensional interfaces ($`d=2`$) we have $`\alpha =\beta =0`$, meaning logarithmic roughening of the interface, i.e., $`w(L)\mathrm{log}L`$ for saturated interfaces, and $`w(t)\mathrm{log}t`$ for early times. #### 2 Anisotropic KPZ equation The presence of anisotropy along the substrate may drastically change the scaling properties of the KPZ equation. As a physical example consider an ion bombarded surface, where the ions arrive under oblique incidence in the $`xh`$ plane. As a result, the $`x`$ and $`y`$ directions along the substrate will not be equivalent. This anisotropy is expected to appear in the erosion equation, leading to an anisotropic equation of the form ($`d=2`$) $`{\displaystyle \frac{h}{t}}`$ $`=`$ $`\nu _x_x^2h+\nu _y_y^2h+{\displaystyle \frac{\lambda _x}{2}}(_xh)^2`$ (7) $`+{\displaystyle \frac{\lambda _y}{2}}(_yh)^2+\eta (x,y,t),[\mathrm{AKPZ}]`$ where $`_xhh/x`$ and $`_yhh/y`$. The anisotropy leads to surface tension and nonlinear terms that are different in the two directions, which have been incorporated in the growth equation by considering different values for the coefficients $`\nu `$ and $`\lambda `$ (in Eq. (7), $`\nu _x`$ and $`\nu _y`$ are positive constants). Equation (7) is called the anisotropic KPZ (AKPZ) equation. It was introduced by Villain, and its nontrivial properties were studied by Wolf . We note that if $`\nu _x=\nu _y`$ and $`\lambda _x=\lambda _y`$, Eq. (7) reduces to the KPZ equation (5). The AKPZ equation has different scaling properties depending on the signs of the coefficients $`\lambda _x`$ and $`\lambda _y`$. When $`\lambda _x\lambda _y<0`$, a surface described by the AKPZ equation has the same scaling properties as the EW equation. However, when $`\lambda _x\lambda _y>0`$ the scaling properties are described by the isotropic KPZ equation (5). Thus, changing the sign of $`\lambda _x`$ or $`\lambda _y`$ can induce morphological phase transitions from power law scaling $`(wt^\beta ;w(L)L^\alpha )`$ to logarithmic scaling $`(w\mathrm{log}t;w(L)\mathrm{log}L)`$. #### 3 Kuramoto-Sivashinsky (KS) equation The Kuramoto-Sivashinsky (KS) equation, originally proposed to describe chemical waves and flame fronts , is a deterministic equation of the form: $$\frac{h}{t}=\|\nu \|^2hK^4h+\frac{\lambda }{2}(h)^2[KS].$$ (8) While it is deterministic, its unstable and highly nonlinear character gives rise to chaotic solutions. The analysis of the KS equation for one dimensional surfaces shows that in the limit of long time and length scales, the surface described by the KS equation is similar to that described by the KPZ equation, i.e. obeys self-affine scaling with exponents $`z=3/2`$ and $`\beta =1/3`$. The short time scale solution of KS equation reveals an unstable pattern-forming behavior, with a morphology reminiscent of ripples . For two dimensional surfaces, however, the results are not clear. Computer simulations are somewhat contradictory, providing evidence for both EW and KPZ scaling . The anisotropic KS equation was studied in Ref. , indicating that for some parameter values the nonlinearities cancel each other, and lead to unstable modes dominating the asymptotic morphology. At early times the surface displays a chaotic pattern, with stable domains that nucleate and grow linearly in time until ripple domains of two different orientations are formed. The pattern of domains of perpendicularly oriented ripples coarsen with time until one orientation takes over the system. There are various physical systems, including ion sputtering, in which the relevant equation for the surface height is a noisy version of the KS equation (8) . Dynamical renormalization group analysis for the surface dimensions $`d=1`$ and $`2`$ indicate that the large distance and long time behavior of such noisy generalization of Eq. (8) is the same as that of the KPZ equation, the $`d=2`$ result being only quantitative. ### B Bradley and Harper theory of ripple formation A rather successful theoretical model, capturing many features of ripple formation, was developed by Bradley and Harper (BH). They used Sigmund’s theory of sputtering (see Sect. IV E) to relate the sputter yield to the energy deposited onto the surface by the incoming ions. This work has demonstrated for the first time that the yield variation with the local surface curvature induces an instability, which leads to the formation of periodically modulated structures. This instability is caused by the different erosion rates for troughs and crests, the former being eroded faster than the latter (see Fig. 1). Consequently, any surface perturbation increases exponentially in time. Viewing the surface profile as a smooth analytical function of coordinates, BH assumed that the surface can be locally approximated by a quadratic form. Due to the erosion mechanism, described in Fig. 1, the erosion rate depends on the local curvature. Combining the curvature dependent erosion velocity with the surface smoothing mechanism due to surface diffusion (see and next section), BH derived a linear equation for surface morphology evolution $$\frac{h}{t}=v(\theta )+\nu _x(\theta )_x^2h+\nu _y(\theta )_y^2hK^4h.$$ (9) Here $`\nu _x(\theta ),\nu _y(\theta )`$ are the effective surface tensions generated by the erosion process, dependent on the angle of incidence of the ions, $`\theta `$, $`K`$ is the relaxation rate due to surface diffusion ($`K=D_s\gamma \mathrm{\Omega }^2n/k_BT\mathrm{exp}\left\{\frac{\mathrm{\Delta }E}{k_BT}\right\}`$, where $`\mathrm{\Delta }E`$ is the activation energy for surface diffusion, $`\gamma `$ is the surface free energy per unit area, $`T`$ is temperature, $`D_s`$ is the surface diffusion constant, $`\mathrm{\Omega }`$ is the atomic volume and $`n`$ is the number of molecules per unit area on the surface). The physical instability illustrated in Fig. 1 leads to the negative signs of the $`\nu _x,\nu _y`$ coefficients in Eq. (9). Eq. (9) is linearly unstable, with a Fourier mode $`k_c`$ whose amplitude exponentially dominates all the others. This mode is observed as the periodic ripple structure. Using linear stability analysis, BH derived from Eq. (9) the ripple wavelength as $$\mathrm{}_c=2\pi /k_c=2\pi \sqrt{\frac{2K}{\|\nu \|}}(JT)^{1/2}\mathrm{exp}\left\{\frac{\mathrm{\Delta }E}{k_BT}\right\},$$ (10) where $`\nu `$ is the largest in absolute value of the two negative surface tension coefficients, $`\nu _x`$ and $`\nu _y`$, and $`J`$ is the ion flux. The calculation also predicts that the ripple direction is a function of the angle of incidence: for small $`\theta `$ the ripples are parallel to the ion direction, while for large $`\theta `$ they are perpendicular to it. As subsequent experiments have demonstrated , the BH model predicts well the ripple wavelength and orientation. On the other hand, the BH equation (9) is linear, predicting unbounded exponential growth of the ripple amplitude, thus it cannot account for the stabilization of the ripples and for kinetic roughening, both phenomena being strongly supported by experiments (see Sect. III A-III B). Furthermore, the BH model cannot account for low temperature ripple formation since the only smoothing mechanism it considers is of thermal origin. At low temperatures the ion energy and flux dependence of the ripple wavelength also disagree with the BH predictions. Despite these shortcomings, the BH theory represents a major step in understanding the mechanism of surface evolution in ion sputtering since for the first time it uncovered the origin of the ion induced surface instability. Recently a generalization of BH linear theory has been successfully introduced to account for the thermally activated anisotropic surface diffusion present in metallic substrates such as Cu(110). ### C Surface diffusion and deposition noise At high temperatures surface diffusion and fluctuations in the ion beam flux are relevant physical mechanisms taking place on the surface . In this section, we discuss the standard approach to include these phenomena in continuum models. Let us consider the simplest scenario: atoms are deposited on a surface, whereupon they diffuse. If we assume that surface diffusion is the only relaxation mechanism present, the height $`h`$ obeys a continuity equation of the form $$\frac{h}{t}+\stackrel{~}{}𝒋=0,$$ (11) where $`𝒋`$ is a surface current density tangent to the surface, and $`\stackrel{~}{}`$ is calculated in a frame with axes parallel to the surface . In general, $`𝒋`$ is given by the gradient of a chemical potential $`\mu `$, $$𝒋\stackrel{~}{}\mu (𝒓,t)\stackrel{~}{}^2\frac{\delta [h]}{\delta h},$$ (12) where $`\mu `$ minimizes the free energy functional of the surface $`[h]`$ and $`\stackrel{~}{}^2`$ is the surface Laplacian or the Laplace–Beltrami operator. Taking the latter to be proportional to the total surface area $$[h]=𝑑𝒓\sqrt{g},$$ (13) with $`g`$ as defined in Appendix A, and neglecting third or higher powers of derivatives of $`h`$ one arrives at $$\frac{h}{t}=K^2(^2h)K^4h,$$ (14) where $`K`$ is a positive constant. Eq. (14) is the so-called linear MBE equation . For an amorphous solid in equilibrium with its vapor Eq. (14) was obtained in , together with the expression for the coefficient $`K`$ as in Eq. (9). In addition to the deterministic processes, there is considerable randomness in sputter erosion due to fluctuations in the intensity of the ion beam. The ion flux is defined as the number of particles arriving on the unit surface (or per lattice site) in unit time. At large length scales the beam flux is homogeneous with an average intensity $`J`$, but there are local random fluctuations, $`\eta (𝐱,t)\delta J(𝐱,t)`$, uncorrelated in space and time. We can include fluctuations in Eq. (14) by considering the ion flux to be the sum of the average flux $`J`$ and the noise $`\eta `$, which has zero average, $$\eta (𝐱,t)=0$$ (15) and is uncorrelated, $$\eta (𝐱,t)\eta (𝐱^{},t^{})=J\delta (𝐱𝐱^{})\delta (tt^{}),$$ (16) where we have assumed a Poisson distribution for the shot noise. Consequently, the stochastic growth equation describing surface diffusion and fluctuations in an erosion process has the form $$\frac{h}{t}=K^4hJ+\eta (𝐱,t).$$ (17) This variant of Eq. (14) was introduced independently by Wolf and Villain , and by Das Sarma and Tamborenea , and played a leading role in developing our understanding of MBE. We will use the methods leading to (17) to incorporate the smoothing by surface diffusion in our model of ion erosion. Note, however, that as numerous experimental studies indicate, ion bombardment leads to an enhancement of the surface adatom mobility and thus may drastically change the relaxation mechanism, as compared to regular surface diffusion. ### D Microscopic models of ripple formation and roughening Computer simulations provide invaluable insight into microscopic processes taking place in physical systems. Consequently, a number of recent studies have focused on modeling ripple formation at the microscopic level. These studies have proven useful in resolving issues related to low temperature ripple formation and provided important ideas regarding the physical mechanism governing ripple formation . Here we shortly discuss the conclusions reached in some of the most representative numerical work. Monte Carlo simulations of sputter-induced roughening were reported by Koponen et al. . Roughening of amorphous carbon surfaces bombarded by $`5`$ keV Ar<sup>+</sup> ions was studied in for incidence angles between $`0^{}`$ and $`60^{}`$. It was found that ion bombardment induces self-affine topography on the submicrometer scale, the roughness exponent being $`\alpha 0.250.47`$ depending on the angle of incidence . The growth exponent $`\beta `$ was found to be strongly dependent on the relaxation mechanism used and changed from $`\beta 0.3`$ in the model without relaxation to $`\beta 0.20.14`$ when different relaxation rules were used in the simulations. At the same time the roughness exponent $`\alpha `$ was found to be relatively insensitive to the relaxation process on the nanometer scales. Analogous results were obtained for C ions . In this Reference, the ripple wavelength was found to be relatively independent of the ion energy or the magnitude of surface diffusion. Ripple formation was observed even at zero temperature, when surface diffusion was switched off, indicating the presence of ion induced smoothing. Furthermore, these simulations led to the observation of traveling ripples, as predicted by continuum theories (see section VI A 1). Similarly, for $`5`$ keV Ar<sup>+</sup> bombardment of amorphous carbon substrates, the ripple wave vector is seen to change from parallel to normal to the beam direction as the incidence angle is increased, in agreement with BH linear theory, (see Sect. IV). The ripple structure was again observed even when no explicit relaxation mechanism was incorporated in the simulations, and ripple travelling also occurs. For length scales comparable to the cascade dimensions, self-affine topography is observed. A discrete stochastic model was introduced in Ref. to study the morphological evolution of amorphous one dimensional surfaces under ion-bombardment. This is a solid-on-solid model incorporating the erosion rate dependence of surface curvature, the local slope dependence of the sputtering yield, and thermally activated surface diffusion. Up to four different dynamical regimes have been identified. Initially the surface relaxes by surface diffusion with a growth exponent $`\beta 0.38`$, until the onset of the linear BH instability. The instability induces rapid growth ($`\beta >0.5`$). In this regime the local slopes increase rapidly, which triggers non-linear effects eventually stabilizing the surface, $`\beta `$ taking up the EW value $`\beta 1/4`$, which indicates that an effective positive value of the surface tension has been generated. Finally, in the asymptotic time limit $`\beta `$ reaches the KPZ value $`\beta 1/3`$. This behavior is consistent with that displayed by the noisy KS equation . Furthermore, the analytical study using the master equation approach to interface models has shown that the noisy KS equation indeed provides the continuum limit of the discrete stochastic model of Ref. . Conversely, the results of the simulations in support the theoretical conclusions of Ref. that the asymptotic behavior of the noisy KS equation is the same as that of the KPZ equation for one and two dimensional surfaces. In summary, Monte Carlo simulations of the sputtering process of amorphous materials have shown that intermediate and high energy ion bombardment may lead to ripple formation in a wide parameter range. Furthermore, ripple formation was observed even at zero temperature. Computer simulations have also confirmed the linear dependence of the ripple wavelength on the incident ion penetration depth and the fact that ripple formation is a process fully determined by the incident ion characteristics and not caused by any defects, irregularities or surface chemistry. The same simulations have confirmed that under some bombarding conditions the surface is rough, and obeys scaling. ### E Sigmund’s theory of sputtering The erosion rate of ion bombarded surfaces is characterized by the sputtering yield, $`Y`$, defined as the average number of atoms leaving the surface of a solid per incident particle. In order to calculate the yield and to predict the surface morphology generated by ion bombardment, we first need to understand the mechanism of sputtering, resulting from the interaction of the incident ions and the substrate . In the process of sputtering the incoming ions penetrate the surface and transfer their kinetic energy to the atoms of the substrate by inducing cascades of collisions among the substrate atoms, or through other processes such as electronic excitations. Whereas most of the sputtered atoms are located at the surface, the scattering events that might lead to sputtering take place within a certain layer of average depth $`a`$, which is the average penetration depth of the incident ion. A qualitative picture of the sputtering process is as follows: an incoming ion penetrates into the bulk of the material and undergoes a series of collisions with the atoms of the substrate. Some of the atoms undergo secondary collisions, thereby generating another generation of recoiling atoms. A vast majority of atoms will not gain enough energy to leave their lattice positions permanently. However, some of them will be permanently removed from their sites. The atoms located in the close vicinity of the surface, which can gain enough energy to break their bonds, will be sputtered. Usually the number of sputtered atoms is orders of magnitude smaller than the total number of atoms participating in the collision cascade. A quantitative description of the process of ion sputtering was developed by Sigmund . Assuming an amorphous target, Sigmund derived a set of transport equations describing the energy transfer during the sputtering process. A practically important result of Sigmund’s theory is the prediction of the deposited energy distribution: the ion deposited at a point $`P`$ inside the bulk of the material spreads its kinetic energy according to the Gaussian distribution $$E(𝒓^{\mathbf{}})=\frac{ϵ}{(2\pi )^{3/2}\sigma \mu ^2}\mathrm{exp}\left\{\frac{Z^2}{2\sigma ^2}\frac{X^2+Y^2}{2\mu ^2}\right\}.$$ (18) In (18) $`Z^{}`$ is the distance from point $`P`$ to point $`O`$ measured along the ion trajectory, and $`X^{}`$, $`Y^{}`$ are measured in the plane perpendicular to it (see Fig. 2 and the inset of Fig. 5); $`ϵ`$ denotes the total energy carried by the ion and $`\sigma `$ and $`\mu `$ are the widths of the distribution in directions parallel and perpendicular to the incoming beam respectively. Deviations of the energy distribution from Gaussian (18) occur mainly when $`M_1>M_2`$, where $`M_1`$ is the mass of the projectile and $`M_2`$ is the mass of the target atom. As shown by Sigmund and Winterbom , electronic stopping doesn’t affect much the shape of deposited-energy distribution. Subsequently, Monte Carlo simulations of the sputtering process have demonstrated that the deposited-energy distribution and damage distribution can be well approximated by Gaussian for intermediate and high energies. In general, comparison of Sigmund’s theory with experimental results has shown that it describes well the qualitative behavior of the backsputtering yield, and in many cases there is good quantitative agreement as well . A quantity of central importance is the mean path length of an ion traveling inside the bulk of the material (see Fig. 2), often called penetration depth, given by $$a(ϵ)=\frac{1m}{2m}\gamma ^{m1}\frac{ϵ^{2m}}{NC_m},$$ (19) where $`N`$ is the target atom density, $`C_m`$ is a constant dependent on the parameters of the interatomic potential and the exponent $`m=m(ϵ)`$ varies slowly from $`m=1`$ at high energies to $`m0`$ at very low energies. In the region of intermediate energies, i.e. for $`ϵ`$ between $`10`$ and $`100`$ keV, $`m1/2`$ and we can approximate the penetration depth as $`a(ϵ)ϵ`$. Eq. (18) describes the effect of a single incident ion. Actually, the sample is subject to an uniform flux $`J`$ of bombarding ions, penetrating the solid at different points simultaneously, such that the erosion velocity at an arbitrary point $`O`$ depends on the total power $`_O`$ contributed by all the ions deposited within the range of the distribution (18). If we ignore shadowing effects and redeposition of the eroded material, the normal erosion velocity at $`O`$ is given by $$V_O=p_{}𝑑𝒓\mathrm{\Phi }(𝒓)E(𝒓),$$ (20) where the integral is taken over the region $``$ of all points at which the deposited energy contributes to $`_O`$, $`\mathrm{\Phi }(𝒓)`$ is a local correction to the uniform flux $`J`$ due to variation of the local slopes, and the material constant $`p`$ depends on the surface binding energy and scattering cross-section as $$p=\frac{3}{4\pi ^2}\frac{1}{NU_oC_o},$$ (21) where $`U_o`$ is the surface binding energy and $`C_o`$ is a constant proportional to the square of the effective radius of the interatomic interaction potential. While the predictions of Sigmund’s theory have been checked on many occasions, it also has well known limitations. Next we list two, that will limit our theory on the surface morphology as well: (a) the assumption of random slowing down and arbitrary collisions works satisfactorily only at intermediate and high energies, i.e. when $`ϵ1100`$ keV, but may break down at low energies; (b) the assumption of a planar surface may influence the magnitude of the yield, since surface roughness has a tendency to increase the yield . ## V Continuum equation for the surface height Sigmund’s theory, while offering a detailed description of ion bombardment, is not able to provide direct information about the morphology of ion-sputtered surfaces. While Eq. (20) provides the erosion velocity, in the present form it cannot be used to make analytical predictions regarding the dynamical properties of surface evolution. To achieve such a predictive power, we have to eliminate the nonlocality contained in the integral (20) and derive a continuum equation describing the surface evolution depending only on the local surface morphology. The main goal of this section is to provide a detailed derivation of such an equation starting from Eq. (20). The properties of the obtained equation will be discussed in the following sections. We start by summarizing the main steps that we follow in the derivation of the equation for the surface morphology evolution: (i) Using Eq. (20), we calculate the normal component of the erosion velocity $`V_O`$ at a generic point $`O`$ of the surface. This calculation can be performed in a local frame of reference $`(\widehat{X},\widehat{Y},\widehat{Z})`$, defined as follows: the $`\widehat{Z}`$ axis is chosen to be parallel to the local normal to the surface at point $`O`$. The $`\widehat{Z}`$ axis forms a plane with the trajectory of an ion penetrating the surface at $`O`$. We choose the $`\widehat{X}`$ axis to lie in that plane and be perpendicular to $`\widehat{Z}`$. Finally, $`\widehat{Y}`$ is perpendicular to the $`(\widehat{X},\widehat{Z})`$ plane and completes the local reference frame, as shown in Fig. 3. (ii) Next we relate the quantities measured in the local frame $`(\widehat{X},\widehat{Y},\widehat{Z})`$ to those measured in the laboratory frame $`(\widehat{x},\widehat{y},\widehat{h})`$. The latter is defined by the experimental configuration as follows: $`\widehat{h}`$ is the direction normal to the uneroded flat surface. The ion direction together with the $`\widehat{h}`$ axis define the $`(\widehat{x},\widehat{h})`$ plane. Finally, the $`\widehat{y}`$ axis is perpendicular to the $`(\widehat{x},\widehat{h})`$ plane (see Fig. 3 and Appendix A). Furthermore, we have to take into account the fact that the local angle of incidence $`\phi `$, which is the angle between the ion trajectory and the local normal to the surface, changes from point to point along the surface. Consequently, $`\phi `$ is a function of the local value of the slope at $`O`$ (as measured in the laboratory frame), and the angle $`\theta `$ between the ion trajectory and the normal $`\widehat{n}`$ to the uneroded surface. (iii) Finally, to obtain the equation of motion for the surface profile $`h(x,y,t)`$, we have to project the normal component of the velocity of erosion onto the global $`\widehat{h}`$ axis. The time derivative of $`h(x,y,t)`$ at any point $`O`$ on the surface is proportional to the surface erosion velocity $`V_O`$ at that point and the local normal is defined by the gradient of the surface profile $`h(x,y,t)`$ at $`O`$. Having defined our objectives and outlined the strategy, we move on to the description of the calculations. We consider point $`O`$ to be the origin of the local system of coordinates $`(\widehat{X},\widehat{Y},\widehat{Z})`$. To describe the surface profile in a neighborhood of $`O`$ we assume that the surface can be described in terms of a smooth analytical infinitely differentiable function, i.e. there are no singularities and overhangs, and thus we can approximate the surface profile at an arbitrary point $`(X,Y,Z)`$ by $`Z(X,Y)`$ $``$ $`{\displaystyle \frac{\mathrm{\Delta }_{20}X^2}{a}}+{\displaystyle \frac{\mathrm{\Delta }_{02}Y^2}{a}}+`$ (22) $`+`$ $`{\displaystyle \underset{m,n=0,n+m=3,4}{\overset{4}{}}}{\displaystyle \frac{\mathrm{\Delta }_{nm}}{a^{m+n1}}}X^nY^m,`$ (23) where, for later convenience, we introduced the following notations: $`\mathrm{\Delta }_{nm}={\displaystyle \frac{a^{n+m1}}{n!m!}}{\displaystyle \frac{^{n+m}Z(X,Y)}{^nX^mY}}.`$ (24) Here $`\mathrm{\Delta }_{20}`$ and $`\mathrm{\Delta }_{02}`$ are proportional to the principal curvatures of the surface, i.e., to the inverses of the principal radii of curvature, $`R_X`$ and $`R_Y`$. It must be noted that, in our approximation, $`\widehat{X}`$ and $`\widehat{Y}`$ (see Fig. 3) are the two principal directions of the surface at $`O`$, along which the curvatures are extremal. This implies the absence in Eq. (23) of cross-terms of the type $`XY`$, i.e., we neglected the term $`^2Z(X,Y)/XY`$ at $`O`$. Due to its exponential nature, the deposited energy distribution (18) decays very fast and, consequently, only particles striking the surface at a point $`(X,Y,Z)`$ such that $`X/a`$, $`Y/a`$ are of order unity, contribute non-negligibly to the energy reaching the surface at $`O`$. We further assume that the surface varies slowly enough so that $`R_X`$, $`R_Y`$ and the inverses of the higher order derivatives are much larger than the penetration depth $`a`$, i.e. the surface is smooth on length scales close to $`a`$ (this fact is supported by nearly all experimental results). Now we can calculate the various factors appearing in the integral (20). To proceed with Eq. (20) we note that, with respect to local surface orientation, only the normal component of the incident flux contributes to ion erosion. Figure 4 illustrates the calculation of the normal component of the flux. In the figure we consider a point at the surface $`(X,Y,Z)`$, and two other points also on the surface, at infinitesimal distances $`L`$ and $`N`$ away from the former. We can estimate the correction to the average flux $`J`$ due to the surface slopes by projecting a square perpendicular to the ion beam with area $`n\times l`$ onto the surface area element intersected by the ion trajectories. The result is $$\mathrm{\Phi }(𝒓)J\frac{ln}{LN},$$ (25) where $`J`$ is the average flux. From Fig. 4, $$\mathrm{\Delta }\phi =\mathrm{tan}^1\left(\frac{Z}{X}\right)\frac{Z}{X},$$ (26) and $$\frac{\mathrm{}}{L}=\mathrm{cos}(\phi +\mathrm{\Delta }\phi )\mathrm{cos}\phi \frac{Z}{X}\mathrm{sin}\phi .$$ (27) On the other hand, we also have (see Fig. 4) $$\frac{n}{N}=\mathrm{cos}\alpha 1\frac{1}{2}\left(\frac{Z}{Y}\right)^21,$$ (28) so that, combining Eqs. (25)-(28), and neglecting powers of derivatives of the height, we obtain the correction to the flux $$\mathrm{\Phi }(𝒓)J(\mathrm{cos}\phi +(_XZ)\mathrm{sin}\phi ).$$ (29) Within the same approximation, the surface element $`d\text{r}`$ in Eq. (20) can be obtained in the local coordinate system $`(\widehat{X},\widehat{Y},\widehat{Z})`$ as $`d\text{r}dXdY.`$ (30) Next we determine the distances $`X^{}`$, $`Y^{}`$, $`Z^{}`$ appearing in the exponential distribution (18), evaluating them in the local coordinate system. Using Fig. 5, we have $`X^{}`$ $`=`$ $`X\mathrm{cos}\phi +Z\mathrm{sin}\phi ,`$ (31) $`Y^{}`$ $`=`$ $`Y,`$ (32) $`Z^{}`$ $`=`$ $`a+X\mathrm{sin}\phi Z\mathrm{cos}\phi .`$ (33) Using these expressions, the correction to the ion flux (29), the deposited energy distribution (18) and the expression for the surface area element $`d\text{r}`$ (30), we can calculate the integral (20) providing the erosion velocity $`V_O`$. Introducing the dimensionless variables $`\zeta _X=X/a`$, $`\zeta _Y=Y/a`$, and $`\zeta _Z=Z/a`$, and extending the integration limits to infinity, we obtain the following expression for the erosion velocity in the laboratory coordinate frame $`V_O`$ $`=`$ $`{\displaystyle \frac{ϵpJa^2}{\sigma \mu ^2(2\pi )^{3/2}}}\mathrm{exp}(a_\sigma ^2/2)`$ (35) $`\times `$ $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\zeta _X𝑑\zeta _Y(\mathrm{cos}\phi +{\displaystyle \frac{\zeta _Z}{\zeta _X}}\mathrm{sin}\phi )`$ (36) $`\times `$ $`\mathrm{exp}(\zeta _Y^2L^2)\mathrm{exp}(\zeta _XA)\mathrm{exp}({\displaystyle \frac{1}{2}}B_1\zeta _X^2)`$ (37) $`\times `$ $`\mathrm{exp}(4D\zeta _Z^2)\mathrm{exp}(2C\zeta _X\zeta _Z)\mathrm{exp}(B_2\zeta _Z),`$ (38) where we used the following notations $`A`$ $`=`$ $`a_\sigma ^2\mathrm{sin}\phi ,`$ (40) $`B_1`$ $`=`$ $`a_\sigma ^2\mathrm{sin}^2\phi +a_\mu ^2\mathrm{cos}^2\phi ,`$ (41) $`B_2`$ $`=`$ $`a_\sigma ^2\mathrm{cos}\phi ,`$ (42) $`C`$ $`=`$ $`{\displaystyle \frac{1}{2}}(a_\mu ^2a_\sigma ^2)\mathrm{sin}\phi \mathrm{cos}\phi ,`$ (43) $`D`$ $`=`$ $`{\displaystyle \frac{1}{8}}(a_\mu ^2\mathrm{sin}^2\phi +a_\sigma ^2\mathrm{cos}^2\phi ),`$ (44) $`L`$ $`=`$ $`{\displaystyle \frac{a_\mu }{\sqrt{2}}}.`$ (45) It must be noted that Eq. (38) coincides with the two-dimensional version of the local erosion velocity derived in Ref. . Now we use the approximation for the surface profile given by Eq. (23). Taking a small $`\mathrm{\Delta }_{nm}`$ (see Eq. (24)) expansion of the $`C`$ and $`B_2`$ coefficients in Eq. (38) and evaluating the Gaussian integrals over $`\zeta _X`$ and $`\zeta _Y`$, we get $`V_O`$ $`=`$ $`{\displaystyle \frac{ϵpJa^2}{\sigma \mu \sqrt{2\pi }}}\mathrm{exp}(a_\sigma ^2/2)\mathrm{exp}\left\{{\displaystyle \frac{A^2}{2B_1}}\right\}{\displaystyle \frac{1}{\sqrt{B_1}}}`$ (46) $`\times `$ $`[\mathrm{cos}\phi +\mathrm{\Gamma }_{20}\mathrm{\Delta }_{20}+\mathrm{\Gamma }_{02}\mathrm{\Delta }_{02}+\mathrm{\Gamma }_{30}\mathrm{\Delta }_{30}`$ (47) $`+`$ $`\mathrm{\Gamma }_{21}\mathrm{\Delta }_{21}\mathrm{\Gamma }_{40}\mathrm{\Delta }_{40}+\mathrm{\Gamma }_{22}\mathrm{\Delta }_{22}+\mathrm{\Gamma }_{04}\mathrm{\Delta }_{04}].`$ (48) The expressions for the coefficients $`\mathrm{\Gamma }_{nm}`$ can be found in Appendix B. Next we have to rewrite $`V_O`$ in terms of the laboratory coordinates $`(x,y,h)`$, which we perform in two steps. First, we write the angle $`\phi `$ as a function of $`\theta `$ and the slopes of the surface at $`O`$ as measured in the laboratory frame. Second, we perform the transformation between the local and the laboratory coordinates. For both steps we will have to make expansions in powers of derivatives of $`h(x,y,t)`$. In line with our earlier assumption on the smoothness of the surface, we will assume that $`h`$ varies smoothly enough so that we can neglect products of derivatives of $`h`$ for third and higher orders. In the laboratory frame, the neglect of overhangs allows us to describe a generic point at the surface, such as $`O`$, with coordinates $`(x,y,h(x,y))`$. Considering now the unit vectors $`\widehat{n}`$, $`\widehat{m}`$ shown in Fig. 3, the angle $`\phi `$ is given by $`\mathrm{cos}\phi `$ $`=`$ $`\widehat{m}\widehat{n}={\displaystyle \frac{\mathrm{cos}\theta (_xh)\mathrm{sin}\theta }{\sqrt{1+(_xh)^2+(_yh)^2}}},`$ (50) $`\mathrm{sin}\phi `$ $`=`$ $`(\mathrm{sin}^2\theta +2(_xh)\mathrm{sin}\theta \mathrm{cos}\theta (_xh)^2\mathrm{cos}^2\theta `$ (51) $`+`$ $`(_yh)^2)^{1/2}\times (1+(_xh)^2+(_yh)^2)^{1/2}.`$ (52) Thus far, expressions (50)-(52) are exact, and the values of $`_xh`$ and $`_yh`$ are already evaluated in the laboratory frame of reference. To implement our approximations, in principle we have to separate the cases for normal ($`\theta =0`$) and off–normal ($`\theta 0`$) incidence. Nevertheless, it can be shown that the former case can be obtained as a smooth limit of the latter. Therefore in the following we give the expressions pertaining to the off–normal incidence and refer the reader to Appendix C for details on the $`\theta =0`$ limit. Expanding (50) and (52) in powers of the surface height derivatives, we obtain $`\mathrm{cos}\phi `$ $``$ $`\mathrm{cos}\theta (_xh)\mathrm{sin}\theta `$ (55) $`{\displaystyle \frac{1}{2}}((_xh)^2+(_yh)^2)\mathrm{cos}\theta ,`$ $`\mathrm{sin}\phi `$ $``$ $`\mathrm{sin}\theta +(_xh)\mathrm{cos}\theta {\displaystyle \frac{1}{2}}(_xh)^2\mathrm{sin}\theta `$ (57) $`+{\displaystyle \frac{1}{2}}(_yh)^2{\displaystyle \frac{\mathrm{cos}^2\theta }{\mathrm{sin}\theta }}.`$ Note that these expressions are invariant under the coordinate transformation $`yy`$, but not under $`xx`$, a consequence of $`\theta `$ being non-zero and of our choice of coordinates. Naturally, the $`\theta 0`$ limit restores the symmetry in the $`x`$ direction. Having obtained the expressions (55) and (57), we can return to Eq. (48) to calculate the dependence of $`V_O`$ on the slopes at $`O`$. Finally, all derivatives in (24) have to be expressed in terms of the laboratory coordinates. This can be accomplished given the relation between the base vectors of the local frame $`(\widehat{X},\widehat{Y},\widehat{Z})`$ and those of the laboratory frame $`(\widehat{x},\widehat{y},\widehat{h})`$, derived in Appendix A. If the coordinates of a generic vector $`𝒓`$ are given by $`𝒓`$ $`=`$ $`X\widehat{X}+Y\widehat{Y}+Z\widehat{Z}\mathrm{local}\mathrm{frame},`$ (58) $`𝒓`$ $`=`$ $`x\widehat{x}+y\widehat{y}+z\widehat{z}\mathrm{laboratory}\mathrm{frame},`$ (59) then these quantities are related to each other through $$\left(\begin{array}{c}x\\ y\\ z\end{array}\right)=\left(\begin{array}{c}X\\ Y\\ Z\end{array}\right),$$ (60) where $``$ is a matrix which has as columns the components of the $`(\widehat{X},\widehat{Y},\widehat{Z})`$ set of vectors in terms of the $`(\widehat{x},\widehat{y},\widehat{z})`$ (see Appendix A). To obtain the expression for the erosion velocity, analogous to (48), in the laboratory frame, we use Eqs. (55), (57), and $``$ along with the chain rule for differentiation, and perform expansions in powers of derivatives of $`h(x,y,t)`$. After some algebra we obtain in the laboratory frame $`{\displaystyle \frac{^{n+m}h}{^nX^mY}}{\displaystyle \frac{^{n+m}h}{^nx^my}},`$ (61) up to fourth order in products of derivatives of $`h(x,y)`$. To summarize our results thus far, within the approximations leading to (48) and neglecting nonlinearities of cubic and higher orders in derivatives of $`h`$ in the laboratory frame, we obtained Eq. (61), providing the relation between the derivatives in the two reference frames, and relations (55) and (57) for the angle of incidence measured in the local frame as a function of the angle of incidence $`\theta `$ and of the surface slopes. Finally, to relate the velocity of erosion $`V_O`$, which is normal to the surface at $`O`$, to the velocity of erosion of the surface along the $`h`$ axis, $`h/t`$, we have to project the former onto the latter, obtaining $$\frac{h}{t}=V_O\sqrt{g},$$ (62) where the negative sign accounts for the fact that $`V_O`$ is the rate at which the surface is eroded, i.e. the average height decreases. Furthermore, taking into account surface diffusion effects, together with the fluctuations (shot noise) in the flux of the bombarding particles, as discussed in Section IV C, we complete (62) by adding these physical effects $$\frac{h}{t}=V_O\sqrt{g}K^4h+\eta (𝒓,t),$$ (63) Finally, we have to write down the contribution of the $`V_O\sqrt{g}`$ term to the evolution equation (63). Performing a small slope expansion and using Eqs. (48), (55), and (57), we obtain $`{\displaystyle \frac{h}{t}}`$ $`=`$ $`v_0+\gamma {\displaystyle \frac{h}{x}}+\xi _x\left({\displaystyle \frac{h}{x}}\right)\left({\displaystyle \frac{^2h}{x^2}}\right)+\xi _y\left({\displaystyle \frac{h}{x}}\right)\left({\displaystyle \frac{^2h}{y^2}}\right)+\nu _x{\displaystyle \frac{^2h}{x^2}}+\nu _y{\displaystyle \frac{^2h}{y^2}}+\mathrm{\Omega }_1{\displaystyle \frac{^3h}{x^3}}+\mathrm{\Omega }_2{\displaystyle \frac{^3h}{xy^2}}`$ (65) $`D_{xy}{\displaystyle \frac{^4h}{x^2y^2}}D_{xx}{\displaystyle \frac{^4h}{x^4}}D_{yy}{\displaystyle \frac{^4h}{y^4}}K^4h+{\displaystyle \frac{\lambda _x}{2}}\left({\displaystyle \frac{h}{x}}\right)^2+{\displaystyle \frac{\lambda _y}{2}}\left({\displaystyle \frac{h}{y}}\right)^2+\eta (x,y,t),`$ where the coefficients are given by the expressions $`v_0`$ $`=`$ $`Fc,`$ (66) $`\gamma `$ $`=`$ $`F{\displaystyle \frac{s}{f^2}}\left\{a_\sigma ^2a_\mu ^2c^2(a_\sigma ^21)a_\sigma ^4s^2\right\},`$ (67) $`\nu _x`$ $`=`$ $`Fa{\displaystyle \frac{a_\sigma ^2}{2f^3}}\left\{2a_\sigma ^4s^4a_\sigma ^4a_\mu ^2s^2c^2+a_\sigma ^2a_\mu ^2s^2c^2a_\mu ^4c^4\right\},`$ (68) $`\lambda _x`$ $`=`$ $`F{\displaystyle \frac{c}{2f^4}}\{a_\sigma ^8a_\mu ^2s^4(3+2c^2)+4a_\sigma ^6a_\mu ^4s^2c^4a_\sigma ^4a_\mu ^6c^4(1+2s^2)`$ (69) $``$ $`f^2(2a_\sigma ^4s^2a_\sigma ^2a_\mu ^2(1+2s^2))a_\sigma ^8a_\mu ^4s^2c^2f^4\},`$ (70) $`\nu _y`$ $`=`$ $`Fa{\displaystyle \frac{c^2a_\sigma ^2}{2f}},`$ (71) $`\lambda _y`$ $`=`$ $`F{\displaystyle \frac{c}{2f^2}}\left\{a_\sigma ^4s^2+a_\sigma ^2a_\mu ^2c^2a_\sigma ^4a_\mu ^2c^2f^2\right\},`$ (72) $`\xi _x`$ $`=`$ $`Fa{\displaystyle \frac{a_\sigma ^2sc}{2f^5}}\{6s^6a_\sigma ^8+a_\sigma ^8a_\mu ^2s^4(4+3c^2)a_\sigma ^8a_\mu ^4c^2s^2+a_\sigma ^6a_\mu ^4c^2s^2(46s^2)+a_\sigma ^6a_\mu ^2s^4(3+15s^2)`$ (73) $`+`$ $`a_\sigma ^4a_\mu ^43c^2s^2(4+3s^2)a_\sigma ^4a_\mu ^63c^4(1+s^2)+a_\sigma ^2a_\mu ^6c^4(93s^2)3a_\mu ^8c^6\},`$ (74) $`\xi _y`$ $`=`$ $`Fa{\displaystyle \frac{a_\sigma ^2sc}{2f^3}}\left\{a_\sigma ^4a_\mu ^2c^2+a_\sigma ^4s^2(2+c^2)a_\mu ^4c^4+a_\sigma ^2a_\mu ^2c^2(32s^2)\right\},`$ (75) $`\mathrm{\Omega }_1`$ $`=`$ $`Fa^2{\displaystyle \frac{3}{6}}{\displaystyle \frac{1}{f^2}}{\displaystyle \frac{s}{a_\mu ^2}}\left\{f^2fa_\sigma ^4c^2(a_\mu ^2a_\sigma ^2)c^2(f+a_\sigma ^4s^2)\right\},`$ (76) $`\mathrm{\Omega }_2`$ $`=`$ $`Fa^2{\displaystyle \frac{1}{6}}{\displaystyle \frac{1}{f^4}}\left\{3sf^2(f+a_\sigma ^4s^2)+a_\sigma ^2c^2(3a_\sigma ^2sf+a_\sigma ^6s^3)f+2(a_\mu ^2a_\sigma ^2)c^2(3f^2s+6a_\sigma ^4s^3+a_\sigma ^8s^5)\right\},`$ (77) $`D_{xx}`$ $`=`$ $`F{\displaystyle \frac{a^3}{24}}{\displaystyle \frac{1}{f^5}}\{4(3a_\sigma ^2s^2f+a_\sigma ^6s^4)f^2+a_\sigma ^2c^2(3f^2+6a_\sigma ^4s^2f+a_\sigma ^8s^4)f`$ (78) $`+`$ $`2(a_\mu ^2a_\sigma ^2)c^2(15a_\sigma ^2s^2f^2+10a_\sigma ^6s^4f+a_\sigma ^{10}s^6)\},`$ (79) $`D_{yy}`$ $`=`$ $`F{\displaystyle \frac{a^3}{24}}{\displaystyle \frac{1}{f^5}}{\displaystyle \frac{3a_\sigma ^2}{a_\mu ^2}}\left\{f^4c^2\right\},`$ (80) $`D_{xy}`$ $`=`$ $`F{\displaystyle \frac{6a^3}{24}}{\displaystyle \frac{1}{f^5}}{\displaystyle \frac{f^2}{a_\mu ^2}}\left\{2(a_\sigma ^2s^2)f^2+a_\sigma ^2c^2(f^2+a_\sigma ^4s^2f)+2(a_\mu ^2a_\sigma ^2)c^2(3a_\sigma ^2s^2f+a_\sigma ^6s^4)\right\}.`$ (81) In the above expressions, we have defined $$F\frac{Jϵpa}{\sigma \mu \sqrt{2\pi f}}e^{a_\sigma ^2a_\mu ^2c^2/2f}.$$ (82) and, as introduced in Appendix B, $`a_\sigma {\displaystyle \frac{a}{\sigma }},a_\mu {\displaystyle \frac{a}{\mu }},`$ (83) $`s\mathrm{sin}\theta ,c\mathrm{cos}\theta ,`$ (84) $`fa_\sigma ^2s^2+a_\mu ^2c^2.`$ (85) Equation (65) with the coefficients (66)-(85), fully describe the nonlinear time evolution of sputter eroded surfaces, provided that the leading relaxation mechanisms are thermally activated surface diffusion and ion-induced effective smoothing. Due to its highly nonlinear character, Eq. (65) can predict rather complex morphologies and dynamical behaviors. In the remainder of the paper we will focus on the physical interpretation of the coefficients (66)-(85), uncovering their dependence on the experimental parameters, and we discuss the morphologies described by Eq. (65). Consistent with the symmetries imposed by the geometry of the problem, the coefficients in Eq. (65) are symmetric under the transformation $`yy`$ but not under $`xx`$, while for $`\theta 0`$ the system is isotropic in the $`x`$ and $`y`$ directions, specifically $`\gamma =\xi _x=\xi _y=\mathrm{\Omega }_1=\mathrm{\Omega }_2=0`$, $`\lambda _x=\lambda _y`$, $`\nu _x=\nu _y`$, and $`D_{xx}=D_{yy}=\frac{1}{2}D_{xy}`$. ## VI Analysis of the growth equations This section is devoted to the study of the morphological properties predicted by Eq. (65). This is not a simple task, due to large number of linear and nonlinear terms, each of which influence the surface morphology. The complexity of the problem is illustrated by some special cases of Eq. (65), for which the behavior is better understood. For example, when nonlinear terms and the noise are neglected $`(\xi _x=\xi _y=\lambda _x=\lambda _y=0,\eta =0)`$, Eq. (65) reduces to a linear generalization of BH theory, which predicts ripple formation. It is also known that the isotropic KS equation, obtained by taking $`\nu _x=\nu _y,D_{xx}=D_{yy}=D_{xy}/2`$, and $`\lambda _x=\lambda _y`$, asymptotically predicts kinetic roughening, with morphology and exponents similar to those seen experimentally in ion sputtering . For positive $`\nu _x`$ and $`\nu _y`$, Eq. (65) reduces to the anisotropic KPZ equation, whose scaling behavior is controlled by the sign of the product $`\lambda _x\lambda _y`$ . Finally, recent integration by Rost and Krug of the noiseless anisotropic KS equation (i.e., when $`\eta =0`$) showed that when $`\lambda _x\lambda _y<0`$, ripples unaccounted for by the linear theory appear, their direction being rotated with respect to the ion direction. To predict the morphology of ion-sputtered surfaces, we need to gain a full understanding of the behavior predicted by (65) in the physically relevant two dimensional case, going beyond the special cases. Help is provided by the recent numerical integration of (65) by Park et al. that indicates a clear separation in time of the linear and nonlinear behaviors. The results show that before a characteristic time $`t_c`$ has been reached, the morphology is fully described by the linear theory, as if nonlinear terms were not present. However, after $`t_c`$ the nonlinear terms completely determine the surface morphology. These results offer a natural layout for our discussion. In section VI A we will limit our discussion to the linear theory. However, even in this case we have to distinguish four different cases, depending on whether the surface diffusion in the system is thermally generated or of the effective type associated with the ion erosion process. Consequently, in Sections VI A 1 \- VI A 2, we discuss the high temperature case, when relaxation is by thermal surface diffusion, treating separately the symmetric ($`\sigma =\mu `$), and asymmetric ($`\sigma \mu `$) cases. Next we turn our attention to low temperature ripple formation, when surface relaxation is dominated by erosion, and we again distinguish the symmetric and asymmetric cases (Sects. VI A 3 and VI A 4). Finally, Sections VI B 1 \- VI B 4 are devoted to the effect of the nonlinear terms, addressing such important features as ripple stabilization, rotated ripples and kinetic roughening. ### A Linear theory #### 1 Ripple formation at high temperatures: Symmetric case In this section we discuss the process of ripple formation in the symmetric case $`\sigma =\mu `$, when the relaxation is by thermally activated surface diffusion. Thus we assume that the magnitude of the thermally activated surface diffusion coefficient, $`K`$, is much larger than $`D_{xx}`$, $`D_{xy}`$, $`D_{yy}`$, generated by the ion bombardment process. This is always the case for high temperatures since $`K`$ increases as $`(1/T)\mathrm{exp}(\mathrm{\Delta }E/k_BT)`$ with $`T`$, while ion induced effective smoothing terms are independent of $`T`$. Dropping the nonlinear terms in Eq. (65), we obtain $`{\displaystyle \frac{h}{t}}`$ $`=`$ $`v_0+\gamma {\displaystyle \frac{h}{x}}+\nu _x{\displaystyle \frac{^2h}{x^2}}+\nu _y{\displaystyle \frac{^2h}{y^2}}`$ (87) $`+\mathrm{\Omega }_1{\displaystyle \frac{^3h}{x^3}}+\mathrm{\Omega }_2{\displaystyle \frac{^3h}{x^2y}}K^4h+\eta (x,y,t),`$ where the coefficients can be obtained from Eqs. (66) -(82) by taking $`\sigma =\mu `$: $`v_0`$ $`=`$ $`Fc,`$ (88) $`\gamma `$ $`=`$ $`Fs(a_\sigma ^2c^21),`$ (89) $`\nu _x`$ $`=`$ $`{\displaystyle \frac{Fa}{2}}\left\{2s^2c^2a_\sigma ^2s^2c^2\right\},`$ (90) $`\nu _y`$ $`=`$ $`{\displaystyle \frac{Fa}{2}}c^2,`$ (91) $`\mathrm{\Omega }_1`$ $`=`$ $`{\displaystyle \frac{Fa^2s}{2a_\sigma ^2}}(1a_\sigma ^2c^2),`$ (92) $`\mathrm{\Omega }_2`$ $`=`$ $`{\displaystyle \frac{Fa^2}{6a_\sigma ^2}}s\left\{a_\sigma ^2(3c^23s^2)+a_\sigma ^4c^2s^23\right\}.`$ (93) Since the surface morphology depends on the signs and absolute values of the coefficients in Eq. (87), in the following we discuss in detail their behavior as functions of the angle of incidence $`\theta `$ and the reduced penetration depth $`a_\sigma `$. Erosion velocity, $`v_0`$: The $`v_0`$ term describes the erosion velocity of a flat surface. This term does not affect the ripple characteristics, such as ripple wavelength and ripple amplitude, and can be eliminated from the surface evolution equation by the coordinate transformation $`\stackrel{~}{h}=h+v_0t`$. This corresponds to a transformation to the coordinate frame moving with the eroded surface. However, since $`v_0`$ is the largest contribution to the erosion rate and is the only one that contributes in the linear theory, it is worthwhile to investigate its dependence on $`\theta `$ and $`a_\sigma `$. Fig. 6 shows the $`v_0`$ dependence on the angle of incidence $`\theta `$ for three different values of the reduced penetration depth $`a_\sigma `$. From Eq. (93), $`v_0`$ is positive for all $`\theta `$ and $`a_\sigma `$. In experiments $`v_0(\theta )`$ corresponds to the secondary ion yield variation with the incidence angle $`\theta `$, i.e. $`v_0(\theta )=JY_{flat}(\theta )/n`$, where $`n`$ is the density of target atoms. Note that $`v_0`$ has the characteristic increasing part for small $`\theta `$, followed by saturation and decrease for large $`\theta `$, similar to the measured yield . Traveling ripples, $`\gamma `$, $`\mathrm{\Omega }_1`$, $`\mathrm{\Omega }_2`$: If we consider a periodic perturbation with wave vector $`(q_x,q_y)`$ in the form $`h`$ $`=`$ $`v_0t+A\mathrm{exp}\left[i(q_xx+q_yy\omega t)+rt\right],`$ (94) from Eq. (87) we obtain the mode velocity $`\omega `$ $`=`$ $`\gamma q_x+\mathrm{\Omega }_1q_x^3+\mathrm{\Omega }_2q_xq_y^2,`$ (95) and the growth rate $`r`$ $`=`$ $`(\nu _xq_x^2+\nu _yq_y^2+K(q_x^2+q_y^2)^2).`$ (96) Thus the coefficients $`\gamma ,\mathrm{\Omega }_1,\mathrm{\Omega }_2`$ contribute to the Fourier mode velocity $`\omega `$ in an anisotropic way that reflects the asymmetry of the $`x`$ and $`y`$ directions for oblique ($`\theta 0`$) ion incidence. The coefficients $`\nu _x,\nu _y,K`$, on the other hand, contribute to the growth rate of the mode amplitude. Carter pointed out that dispersive terms, such as $`\mathrm{\Omega }_1`$ and $`\mathrm{\Omega }_2`$, destroy the translational invariance of the periodic morphology because the different modes travel with different velocities. Note, however, that the existence of a ripple structure means that there is essentially only one Fourier mode describing the surface morphology, which will thus move across the surface with velocity $`\omega `$. The coefficient $`\gamma `$ contributes only to the velocity of the ripples along the $`x`$ direction, leaving unaffected the $`y`$ component of the ripple velocity. Thus, as expected, $`\gamma =0`$ for normal incidence ($`\theta =0`$). Similarly to the $`v_0`$ term, $`\gamma `$ does not affect the ripple characteristics and can actually be eliminated using the transformation $`\stackrel{~}{h}=h(x\gamma t,t)`$. As can be seen in Fig. 7, $`\gamma `$ can change sign with $`\theta `$, indicating that ripples travel in both positive and negative directions along the $`x`$ coordinate, depending on the angle of incidence and the penetration depth: ripples travel in the positive $`x`$ direction for small $`\theta `$ and in the negative $`x`$ direction for larger $`\theta `$. Travelling ripples were observed in numerical simulations of Koponen et al. . As discussed above, the terms $`\mathrm{\Omega }_1`$, $`\mathrm{\Omega }_2`$ also contribute to the travelling of the ripples, and thus have no further effect on the surface morphology. Fig. 8 shows the coefficients $`\mathrm{\Omega }_1`$ and $`\mathrm{\Omega }_2`$ as functions of the angle of incidence $`\theta `$. We find that the absolute value of these coefficients at small angles is small compared to $`\gamma `$ (see Fig. 7), thus the main contribution to the ripple velocity comes from the ($`\gamma h/x`$) term. On the other hand, for angles $`\theta 60^{}`$, these terms are comparable to or larger than $`\gamma `$. The coefficients $`\nu _x`$ and $`\nu _y`$: As we discussed above (see Sect. IV B) the negative surface tension coefficients are the origin of the instability responsible for ripple formation. Consequently, they play a particularly important role in determining the surface morphology. The coefficients $`\nu _x`$ and $`\nu _y`$ are not equal due to the fact that the direction of the ion beam breaks the symmetry along the surface. As seen in Eq. (93), $`\nu _y`$ is always negative, while $`\nu _x`$ can change sign as $`\theta `$ and $`a_\sigma `$ vary, as shown in Fig. 9. The sign and the magnitude of $`\nu _x`$ and $`\nu _y`$ determine both the wavelength and the orientation of the ripples. Ripple wavelength and orientation: The experimental studies on ripple formation have mainly focused on the measurement of the ripple characteristics, such as the ripple wavelength and amplitude. Thus, a successful theory must address and predict these quantities. In the following we outline the method for calculating the ripple wavelengths $`\mathrm{}_x`$ and $`\mathrm{}_y`$. Taking into account the noisy character of Eq. (87), the experimentally observed ripple wavelength corresponds to the unstable Fourier mode which yields the maximum value of the structure factor. The structure factor, $`S(𝐪,t)`$, is calculated from the Fourier transform $`h(𝐪,t)`$ of the instantaneous surface profile and is defined as $`S(𝐪,t)=h(𝐪,t)h(𝐪,t),`$ (97) where $`h(𝐪,t)={\displaystyle \frac{d𝐫}{(2\pi )^2}\mathrm{exp}(i\mathrm{𝐪𝐫})h(𝐫,t)}.`$ (98) Fourier transforming Eq. (87) and inserting the expression for the Fourier transforms of $`h(𝐫,t)`$ into (97), we obtain $`S(𝐪,t)=h(𝐪,t)h(𝐪,t)={\displaystyle \frac{J}{2}}{\displaystyle \frac{1\mathrm{exp}(r(𝐪)t)}{r(𝐪)}},`$ (99) where $`r`$ is the growth rate of the mode $`𝐪`$ given by Eq. (96) and is positive for all unstable Fourier modes in the system. We find that, depending on the sign of $`\nu _x`$ and the relative magnitude of $`\nu _x`$ and $`\nu _y`$, we can distinguish two cases: (i) For $`\nu _x<\nu _y<0`$, which, according to Eq. (93), holds when $`a_\sigma >\sqrt{{\displaystyle \frac{2}{c^2}}},`$ (100) the ripple structure is oriented in the $`x`$ direction, with ripple wavelength $`\mathrm{}_x=2\pi \sqrt{{\displaystyle \frac{2K}{\|\nu _x\|}}}.`$ (101) This means that the maximum of $`S(𝐪,t)`$ is at $`(\sqrt{\frac{\|\nu _x\|}{2K}},0)`$. To illustrate this, in Fig. 10 we show the dependence of the structure factor on the wavevectors $`q_x`$ and $`q_y`$. The contour plot indicates the existence of a global maximum at $`(\sqrt{\frac{\|\nu _x\|}{2K}},0)`$, indicating that the ripples are oriented along the $`x`$ direction. (ii) For $`\nu _x>\nu _y`$, which holds when $`a_\sigma <\sqrt{{\displaystyle \frac{2}{c^2}}},`$ (102) the ripple structure is oriented along the $`y`$-direction, with ripple wavelength $`\mathrm{}_y=2\pi \sqrt{{\displaystyle \frac{2K}{\|\nu _y\|}}}.`$ (103) Figure 11 shows an example of this regime, indicating the existence of a global maximum at point $`(0,\sqrt{\frac{\|\nu _y\|}{2K}})`$, corresponding to the ripple structure oriented along the $`y`$ direction. Phase diagram for ripple orientation— The results obtained on ripple formation can be summarized in a $`(\theta ,a_\sigma )`$ morphological phase diagram, shown in Fig. 12, which has the following regions: Region I— For small $`\theta `$ both $`\nu _x`$ and $`\nu _y`$ are negative such that $`\nu _x<\nu _y`$, thus the ripples are oriented along the $`x`$ direction. Their wavelength is $`\mathrm{}_x=2\pi \sqrt{2K/\|\nu _x\|}`$. The amplitude of the ripples is expected to be weakly modulated by the larger wavelength $`\mathrm{}_y=2\pi \sqrt{2K/\|\nu _y}\|`$. The ripple amplitude grows as $`h_0\mathrm{exp}(r_xt)`$, where $`r_x=r(\sqrt{\frac{\|\nu _x\|}{2K}},0)`$ (see Eq. (94)). The boundary of this region is defined by $`\nu _x(a_\sigma ,\theta )=\nu _y(a_\sigma ,\theta )`$, i.e. $`a_\sigma =\sqrt{{\displaystyle \frac{2}{c^2}}}.`$ (104) Region II— This region is characterized by $`\nu _y<\nu _x`$. The ripples are directed along the $`y`$ direction and have wavelength $`\mathrm{}_y=2\pi \sqrt{2K/\|\nu _y}\|`$. Note that Region II extends down to small values of the incidence angle $`\theta `$ for small enough reduced penetration depth $`a_\sigma `$. This somewhat unphysical result is a consequence of the assumption of a symmetric ($`\sigma =\mu `$) distribution of deposited energy. We will see in the next section that the more physical asymmetric case with $`\sigma >\mu `$ leads in most cases to ripples only oriented along the $`x`$ direction for small enough angles of incidence, as generally observed. Figure 13 shows the $`\theta `$ dependence of the ripple wavelengths along the $`x`$ and $`y`$ directions. In the framework of this model, where thermal surface diffusion is the only smoothing mechanism, the observed ripple orientation corresponds to the direction featuring the smallest value of $`\mathrm{}`$, and changes when $`\mathrm{}_x=\mathrm{}_y`$. The prediction for the ripple wavelength close to $`90^{}`$ is questionable since reflection and shadowing , not incorporated in the model, start to play an important role during ion-bombardment at these high angles. Summary: The dependence of $`\mathrm{}`$ on the main physical parameters characterizing the sputtering process is given by $`\mathrm{}=2\pi \sqrt{{\displaystyle \frac{2K}{\|\nu \|}}}\sqrt{{\displaystyle \frac{2K}{Fa}}}`$ (105) This prediction has a number of consequences, some of which have been verified experimentally (see Section VII): (a) Since the penetration depth, $`a`$, is proportional to $`ϵ^{2m}`$ (see IV E), and $`\sigma \mu a`$, we have $`a_\sigma `$ const, and $`F(ϵa)/(\sigma \mu )ϵ^{12m}`$, which is independent of $`ϵ`$, when $`m=1/2`$. Consequently $$\mathrm{}ϵ^{1/2},$$ (106) i.e. the ripple wavelength is expected to decrease with the ion energy. (b) Taking into account that $`K`$ is independent of the flux and $`\nu J`$, we obtain that the ripple wavelength is also a decreasing function of the incident ion flux, given by $$\mathrm{}\frac{1}{J^{1/2}}.$$ (107) (c) As was mentioned above, the negative surface tension is the origin of the instability leading to ripple formation. When both $`\nu _x`$ and $`\nu _y`$ are negative the experimentally observed ripple structure has the direction for which the growth rate $`r`$ is largest. However, in general, we expect a superposition of both wavelengths, where the long wavelength will appear as a modulation of the ripple amplitude. Indeed, such modulations have been observed both experimentally and numerically . (d) An important prediction of this model, illustrated in Fig. 13, is the existence of the critical angle $`\theta _c`$ where the ripple orientation changes. In the case when surface diffusion is thermally activated, this transition coincides with the condition $`\nu _x=\nu _y`$. #### 2 Ripple formation at high temperatures: Asymmetric case The results of the previous section were derived for the isotropic case, $`\sigma =\mu `$. While this approximation considerably simplifies our discussion, most systems present some anisotropy in the deposited energy distribution. In this section we demonstrate that the existence of anisotropy does not modify the overall qualitative result on the existence of the two parameter regions corresponding to ripples oriented along the $`x`$ or $`y`$ directions. However, anisotropy does change the numerical value of the ripple wavelength and the exact boundary between the two morphological regions: we demonstrate that, for large enough anisotropies, if the incidence angle is small only ripples oriented along the $`x`$ direction are possible. Fig. 14 shows the coefficients $`\nu _x`$ and $`\nu _y`$, given by Eq. (93), as functions of the angle of incidence $`\theta `$, for three different degrees of asymmetry $`\tau =\sigma /\mu `$ in the physical $`\tau >1`$ range. As one can observe the qualitative behavior of $`\nu _x`$ and $`\nu _y`$ is similar to that observed in the symmetric case. One interesting feature, however, must be emphasized: the increasing asymmetry leads to larger ripple wavelength, since the absolute values of $`\nu _x`$ and $`\nu _y`$ decrease. With respect to the the third order linear terms $`\mathrm{\Omega }_1`$ and $`\mathrm{\Omega }_2`$, their behavior as functions of the angle of incidence can be also seen to be qualitatively analogous to the symmetric case. Thus the asymmetry does not change our conclusions regarding the travelling ripples. Ripple wavelength: The calculation of ripple characteristics in the asymmetric case is identical to the one used in the symmetric case. Therefore, we limit ourselves to the presentation of the results. Again, there are two possible ripple directions, and the dominant one can be found from the maximum of the structure factor (99) or, as can be seen to be equivalent, from the maximum of the growth rate (96): (i) When $`\nu _x<\nu _y<0`$, i.e., $`a_\sigma >\sqrt{{\displaystyle \frac{s^2(2+c^2)+\tau ^2c^2(1+2c^2)\tau ^4c^4}{\tau ^2c^2}}},`$ (108) the ripple structure is oriented along the $`x`$-direction with ripple wavelength $`\mathrm{}_x=2\pi \sqrt{\frac{2K}{\|\nu _x\|}}`$. (ii) When $`\nu _x>\nu _y`$, i.e., $`a_\sigma <\sqrt{{\displaystyle \frac{s^2(2+c^2)+\tau ^2c^2(1+2c^2)\tau ^4c^4}{\tau ^2c^2}}},`$ (109) the ripples are oriented along the $`y`$-direction, with ripple wavelength $`\mathrm{}_y=2\pi \sqrt{\frac{2K}{\|\nu _x\|}}`$. Phase diagram for ripple orientation: To consider the effect of asymmetry on the different regimes in ion sputtering, we have studied the ripple orientation phase diagram for different values of $`\tau `$. As $`\tau `$ changes, we find a smooth evolution which does not uncover any new morphological regime. However, the topology of the phase diagram does change as $`\tau `$ increases. For $`\tau <\sqrt{3}`$ the topology of the phase diagram is similar to the symmetric case (see Fig. 12). As Fig. 15 illustrates, for $`\tau \sqrt{3}`$ the ripples oriented along the $`y`$ direction, predicted by the linear theory for small enough $`\theta `$ and $`a_\sigma `$, are absent, which is consistent with most experimental observations. The characteristics of Region $`I`$ and Region $`II`$ of the phase diagram are the same as in the symmetric case. #### 3 Ripple formation at low temperatures: Symmetric case In the previous two sections we discussed the process of ripple formation when the origin of surface smoothing is surface diffusion, described by the $`K^4h`$ term. However, in the series expansion of the erosion velocity we found linear fourth order terms of the form $`D_{xx}_x^4h`$, $`D_{xy}_x^2_y^2h`$, and $`D_{yy}_y^4h`$, which are formally equivalent to the thermally induced surface diffusion terms. These terms arise as a higher order correction to the local surface curvature, being fully determined by the process of surface erosion. Consequently, these terms do not imply actual mass transport along the surface, as thermal surface diffusion does. In this section we show that, in some parameter regions, these terms have a smoothing effect that counteracts the erosion instability, in such a way that they can also lead to ripple formation. We believe this explains the ripples observed at low temperatures both experimentally and in computer simulations . Neglecting the thermally induced relaxation terms (i.e, taking $`K=0`$), nonlinear terms and the terms $`v_0`$, $`\gamma `$, $`\mathrm{\Omega }_1`$ and $`\mathrm{\Omega }_2`$, that do not affect the ripple characteristics, from Eq. (65) we obtain the linear equation $`{\displaystyle \frac{h}{t}}`$ $`=`$ $`\nu _x{\displaystyle \frac{^2h}{x^2}}+\nu _y{\displaystyle \frac{^2h}{y^2}}D_{xx}{\displaystyle \frac{^4h}{x^4}}`$ (111) $`D_{xy}{\displaystyle \frac{^4h}{x^2y^2}}D_{yy}{\displaystyle \frac{^4h}{y^4}}+\eta (x,y,t).`$ The expressions for the coefficients of the ion-induced effective smoothing terms can be obtained from Eqs. (79)-(81) using $`\sigma =\mu `$: $`D_{xx}`$ $`=`$ $`{\displaystyle \frac{Fa^3}{24a_\sigma ^2}}\left\{a_\sigma ^4s^4c^2+a_\sigma ^2(6c^2s^24s^4)+3c^212s^2\right\},`$ (112) $`D_{xy}`$ $`=`$ $`{\displaystyle \frac{Fa^3}{24a_\sigma ^2}}6\left\{a_\sigma ^2s^2c^2+c^22s^2\right\},`$ (113) $`D_{yy}`$ $`=`$ $`{\displaystyle \frac{Fa^3}{24a_\sigma ^2}}3c^2.`$ (114) From Eq. (114), $`D_{yy}`$ is always positive, while $`D_{xy}`$ and $`D_{xx}`$ change sign with $`\theta `$. As we discuss below, the positive $`D_{xx}`$ and $`D_{yy}`$ coefficients play a role similar to thermally activated surface diffusion. Furthermore, the absolute value of $`D_{xx}`$ is comparable with the thermally activated surface diffusion coefficient even at high temperatures (see Sect. VII). Ripple wavelength and orientation: The ripple wavelength and orientation can be calculated following the arguments presented in Section VI A 1, being determined by the maxima of the structure factor $`S(𝐪,t)`$. The growth rate $`r`$ is now given by $`r(q_x,q_y)`$ $`=`$ $`(\nu _xq_x^2+\nu _yq_y^2+D_{xx}q_x^4`$ (116) $`+D_{xy}q_x^2q_y^2+D_{yy}q_y^4).`$ In principle the asymmetry of the $`D_{ij}`$ coefficients may lead to maxima of $`S(𝐪,t)`$ at nonzero $`q_x`$ and $`q_y`$ values, which correspond to ripples forming a nonzero angle with both the $`x`$ and $`y`$ directions. However, straightforward calculations indicate that the following condition holds $`D_{xy}\nu _y=2\nu _xD_{yy},`$ (117) for all values of $`a_\sigma `$ and $`\theta `$ in (93) and (114). This identity implies that no extrema $`(q_x^{},q_y^{})`$ of $`S(𝐪,t)`$ exist other than of the form $`(q_x^{},0)`$ or $`(0,q_y^{})`$. Of these two solutions the one with the largest positive value of $`r(q_x,q_y)`$ corresponds to the observed ripple structure. For small angles of incidence so that $`D_{xx}0`$ (Region I in Fig. 16), it can be easily seen that $`\nu _x<0`$, and the absolute maximum of $`S(𝐪,t)`$ is at $`(q_x^{},0)`$ with $`q_x^{}=\sqrt{\|\nu _x\|/2D_{xx}}`$, thus the ripple structure is aligned along the $`x`$ direction (see for example Fig. 17). Crossing the $`D_{xx}=0`$ line into Region II in Fig. 16 the $`(q_x^{},0)`$ solution disappears, the structure factor having an extremum at $`(0,q_y^{})`$, with $`q_y^{}=\sqrt{\|\nu _y\|/2D_{yy}}`$. However, this extremum is not a global maximum, see for example Fig. 18. The lower boundary of Region II is provided by the condition $`\nu _x=0`$. When we cross it (entering Region III in Fig. 16), we have $`D_{xx}<0`$ and $`\nu _x>0`$. Under this condition, again, there exists an extremum of $`S(𝐪,t)`$ at $`(q_x^{},0)`$, with $`q_x^{}=\sqrt{\nu _x/2\|D_{xx}\|}`$. However, the structure factor takes its absolute maximum at $`(0,q_y^{})`$. Phase diagram for ripple orientation: In summary, three different regions can be determined in the morphological phase diagram shown in Fig. 16 for the case of ion-induced effective smoothing in the symmetric $`\sigma =\mu `$ case. Region I— In this region, the surface tension coefficients $`\nu _x`$ and $`\nu _y`$ are negative, while $`D_{xx}`$ and $`D_{yy}`$ are both positive. The observed ripple structure corresponds to the maximum of $`S(𝐪,t)`$, which indicates that the ripple structure is oriented along the $`x`$ direction. The lower boundary of this region, separating it from Region II, is given by the $`D_{xx}(a_\sigma ,\theta )=0`$ line or, equivalently, by $`a_\sigma =\sqrt{{\displaystyle \frac{(2s^23c^2)+\sqrt{6c^4+4s^4}}{s^2c^2}}}.`$ (118) Region II— In this region both $`D_{xx}`$ and $`\nu _x`$ are negative. This region is bounded below by the $`\nu _x(a_\sigma ,\theta )=0`$ line, given by $`a_\sigma =\sqrt{{\displaystyle \frac{(2s^2c^2)}{s^2c^2}}}.`$ (119) In a continuum description, the maximum of $`r(𝐪)`$ is at infinite $`𝐪`$, thus our theory possibly breaks down in the sense that not even non-linear effects can be expected to stabilize the surface under such conditions. In such a case, a higher order Taylor expansion should be carried out in Sec. V in order to be able to describe our system. Additional effects, such as shadowing, could also start to play a role under such conditions. Region III — In this region $`D_{xx}`$ is negative and $`\nu _x`$ is positive. Thus the instability given by the negative $`D_{xx}`$ is smoothed out by the positive $`\nu _x`$. Since the structure factor takes on its absolute maximum at the finite wave vector $`(0,q_y^{})`$, in principle there is still a ripple structure oriented along the $`y`$ direction. However, remarks similar to those made in Region II might apply here, since we still have $`D_{xx}<0`$. Summary: In the presence of ion-induced effective smoothing the dependence of the ripple wavelength on the physical parameters is different from the case of thermal surface diffusion (see Sect. VI A 1). Here we summarize some of the differences. (a) The dependence of $`\mathrm{}`$ on the ion energy is given by $`\mathrm{}=\sqrt{{\displaystyle \frac{2D}{\|\nu \|}}}\sqrt{{\displaystyle \frac{Fa^3}{Fa}}}aϵ^{2m},`$ (120) indicating that the ripple wavelength depends linearly on the penetration depth $`a`$. This is very different from the behavior predicted by Eq. (87), derived for thermal surface diffusion, and indicates that monitoring the ripple wavelength dependence on $`ϵ`$ can be used to identify the dominant smoothing mechanism. Such a linear behavior of $`\mathrm{}`$ on $`ϵ`$ has indeed been seen experimentally (see Section III A). (b) From Eq. (120) it also follows that the ripple wavelength is independent of the incident ion flux. This prediction is again quite different from the case dominated by thermal surface diffusion, given by Eq. (107). Such a flux independent behavior has been observed experimentally (see Section III A). Finally, analogues of characteristics (c) and (d), discussed in Sect. VI A 1, apply here as well. #### 4 Ripple formation at low temperatures: Asymmetric case In this section we discuss the effect of asymmetry ($`\sigma \mu `$) of the energy distribution on the morphological regimes predicted by Eq. (111). We find that the coefficients of Eq. (111) vary slowly with the asymmetry, but this does not change the qualitative picture presented in the previous section, regarding the ripple wavelength and orientation, or the major morphological regimes found in the isotropic case, including the conditions of validity of our continuum theory. Specifically, we find that the asymmetry enlarges the region in $`\theta `$ where $`D_{xx}`$ and $`D_{yy}`$ are positive, thus shifting Region $`II`$ to larger values of $`\theta `$. Phase diagram for ripple orientation— The topology of the morphological phase diagram and the characteristics of the three main regions are not changed by the asymmetry. We find that the only effect of the asymmetry is to move the boundaries smoothly to larger values of $`\theta `$ as $`\tau `$ increases. The condition $`D(a_\sigma ,\theta )=0`$ (see Eq. (118)) now takes the form $`a_\sigma =\sqrt{(s^2+\tau ^2c^2){\displaystyle \frac{(2s^23\tau ^2c^2)+\sqrt{6\tau ^4c^4+4s^4}}{\tau ^2s^2c^2}}},`$ (121) and the condition $`\nu _x(a_\sigma ,\theta )=0`$ (see Eq. (119)) becomes $`a_\sigma =\sqrt{{\displaystyle \frac{2s^4+\tau ^2c^2s^2\tau ^4c^4}{\tau ^2s^2c^2}}}.`$ (122) ### B Nonlinear theory As we demonstrated in the previous section, linear theory can predict many features of ripple formation, such as the ripple wavelength and orientation, both at high and low temperatures. However, a number of important experimental features are incorrectly predicted by linear theory. They include the stabilization of the ripple amplitude (according to the linear theory the amplitude increases indefinitely at an exponential rate) or the presence of kinetic roughening (completely unaccounted for by the linear approach). To account for these features, we have to consider the effect of the nonlinear terms. Consequently, this section is devoted to the effect of the nonlinear terms on the surface morphology. There is an important difference between the linear and nonlinear theories: while all predictions of the linear theory can be calculated analytically (as we demonstrated in the previous section), the discussion of the nonlinear effects requires a combination of analytical and numerical tools. Even with these tools, the understanding of the nonlinear effects is far less complete than that of the linear theory. #### 1 High temperature morphology: Symmetric case In the high temperature regime, where thermal surface diffusion dominates over ion-induced effective smoothing, the nonlinear equation of the interface evolution has the form $`{\displaystyle \frac{h}{t}}`$ $`=`$ $`v_0+\gamma {\displaystyle \frac{h}{x}}+\xi _x\left({\displaystyle \frac{h}{x}}\right)\left({\displaystyle \frac{^2h}{x^2}}\right)`$ (126) $`+\xi _y\left({\displaystyle \frac{h}{x}}\right)\left({\displaystyle \frac{^2h}{y^2}}\right)+\nu _x{\displaystyle \frac{^2h}{x^2}}`$ $`+\nu _y{\displaystyle \frac{^2h}{y^2}}+{\displaystyle \frac{\lambda _x}{2}}\left({\displaystyle \frac{h}{x}}\right)^2+{\displaystyle \frac{\lambda _y}{2}}\left({\displaystyle \frac{h}{y}}\right)^2`$ $`+\mathrm{\Omega }_1{\displaystyle \frac{^3h}{x^3}}+\mathrm{\Omega }_2{\displaystyle \frac{^3h}{xy^2}}K^4h+\eta (x,y,t),`$ where the coefficients of the linear terms, $`v_0`$, $`\gamma `$, $`\nu _x`$, $`\nu _y`$, $`\mathrm{\Omega }_1`$, $`\mathrm{\Omega }_2`$, and $`K`$ have been discussed in Sections VI A 1 and VI A 2. The coefficients of the nonlinear terms in the symmetric case ($`\sigma =\mu `$) are $`\lambda _x`$ $`=`$ $`{\displaystyle \frac{F}{2}}c\left\{a_\sigma ^2(3s^2c^2)a_\sigma ^4s^2c^2\right\},`$ (127) $`\lambda _y`$ $`=`$ $`{\displaystyle \frac{F}{2}}c\{a_\sigma ^2c^2\},`$ (128) $`\xi _x`$ $`=`$ $`{\displaystyle \frac{Fa}{2}}sc\left\{a_\sigma ^4c^2s^2+a_\sigma ^2(4s^23c^2)+6\right\},`$ (129) $`\xi _y`$ $`=`$ $`{\displaystyle \frac{Fa}{2}}sc\left\{2a_\sigma ^2c^2\right\}.`$ (130) Next we discuss the physical interpretation and the behavior of these coefficients as functions of $`\theta `$ and $`a_\sigma `$. The coefficients $`\xi _x`$ and $`\xi _y`$: Fig. 19 shows the nonlinear coefficients $`\xi _x`$ and $`\xi _y`$ as functions of the angle of incidence $`\theta `$. As the numerical analysis of Eq. (126) shows, these nonlinearities are responsible for the development of overhangs on the surface . Even though the $`\xi _x`$ and $`\xi _y`$ terms are expected not to determine the asymptotic scaling behavior, they can play a relevant role at intermediate time scales after the development of the ripple structure, particularly in the region of large $`\theta `$. The precise contribution of these nonlinearities to the surface dynamics is currently under investigation . The coefficients $`\lambda _x`$ and $`\lambda _y`$: As we discussed in Sec. IV A 2, the morphology of the surface described by Eq. (126) depends on the relative signs of the the nonlinear terms $`\lambda _x`$ and $`\lambda _y`$. As it is evident from Eq. (130), $`\lambda _y`$ is negative for all angles of incidence and penetration depths. However, as shown in Fig. 20, the sign of $`\lambda _x`$ depends on $`\theta `$ and $`a_\sigma `$: $`\lambda _x`$ is negative for small $`\theta `$ and changes sign for larger angles of incidence. In principle the nonlinear terms completely determine the surface morphology. Since the nonlinear terms are always present, an important question is whether the linear regimes are relevant at all. Recent results by Park et al. indicate that, while the nonlinear effects indeed change the surface morphology, the regime described by the linear terms is still visible for a wide range of parameters. By numerically integrating Eq. (126) they have shown that there is a clear separation of the linear and nonlinear regimes in time: for times up to a crossover time $`t_c`$ the surface erodes as if the nonlinear terms would be completely absent, following the predictions of the linear theory. After $`t_c`$, however, the nonlinear terms take over and completely determine the surface morphology. The transition from the linear to the nonlinear regime can be seen either by monitoring the surface width (which is proportional to the ripple amplitude) or the erosion velocity. The simulations indicate that the width increases exponentially with time, as predicted by the linear theory, until $`t_c`$, after which the width growths at a much slower rate . This transition is typically accompanied by the disappearance of ripples predicted by the linear theory and the appearance of either kinetic roughening or of a new rotated ripple structure. The erosion velocity is constant in the linear regime (before $`t_c`$), while it increases or decreases after $`t_c`$, depending on the relative signs of the nonlinear terms. Crossover time: The crossover time $`t_c`$ from the linear to the nonlinear behavior can be estimated by comparing the strength of the linear terms with that of the nonlinear terms. Let the typical surface width at the crossover time $`t_c`$ be $`W_o=\sqrt{W^2(L,t_c)}`$. Then from the linear equation we obtain $$W_o\mathrm{exp}(\nu t_c/\mathrm{}^2),$$ (131) while from $`_th\lambda (h)^2`$ we estimate $$W_o/t_c\lambda W_o^2/\mathrm{}^2.$$ (132) Combining these relations we obtain $$t_c(K/\nu ^2)\mathrm{ln}(\nu /\lambda ).$$ (133) In this expression, $`\nu `$, $`K`$, and $`\lambda `$ correspond to the direction parallel to the ripple orientation. The predicted $`\lambda `$ dependence of $`t_c`$ has been confirmed by numerical simulations . Surface morphology: The surface morphology in the nonlinear regime depends on the relative signs of $`\nu _x`$, $`\nu _y`$, $`\lambda _x`$ and $`\lambda _y`$. The different morphological regimes can be summarized in a phase diagram, shown in Fig. 21. Next we discuss each of the phases separately. Region I— For small $`\theta `$ the nonlinear terms $`\lambda _x`$ and $`\lambda _y`$ have the same (negative) sign, the boundary of this region being given by the condition that $`\lambda _x(a_\sigma ,\theta )=0`$, or equivalently $`a_\sigma =\sqrt{{\displaystyle \frac{3s^2c^2}{c^2s^2}}}.`$ (134) In this region both $`\nu _x`$ and $`\nu _y`$ are negative, thus at short time scales $`(tt_c)`$, the linear theory (see Section VI A 1) predicts ripples oriented along the direction ($`x`$ for large $`a_\sigma `$ and $`y`$ for small $`a_\sigma `$) for which the absolute value of $`\nu `$ is largest, with ripple wavelength $`\mathrm{}=2\pi \sqrt{\frac{2K}{\|\nu \|}}`$. On the other hand, at long times $`(tt_c)`$, the ripple structure disappears and the surface undergoes kinetic roughening . Since $`\lambda _x\lambda _y>0`$, we expect the dynamics of the kinetic roughening regime to be described by the KPZ equation, i.e. the surface width follows $`WL^\alpha `$, $`Wt^\beta `$, where the scaling exponents are $`\alpha 0.38`$ and $`\beta 0.25`$ (see Sect. IV A 1). Region II— In this region both the $`\nu _x`$ and the $`\nu _y`$ coefficients are still negative, but in contrast with Region I the product $`\lambda _x\lambda _y`$ is negative. Recent studies by Park et al. have shown that in time three morphological regimes can be distinguished. For short times, $`tt_I`$, the ripple structure predicted by the linear theory (see Section IV A 1) is observed, with ripples oriented along the direction which has the largest value of $`\|\nu \|`$. For intermediate times $`t_Itt_{II}`$, the surface is rough. If this roughness would follow kinetic roughening, one would expect logarithmic scaling, described by the Edwards-Wilkinson equation (see Sections VI B 2-VI B 4), since $`\lambda _x\lambda _y<0`$. However, this transient regime is somewhat different from what we expect during kinetic roughening. The numerical simulations often show the development of individual ripples, which soon disappear, and no long-range order is present in the system. However, at a second crossover time, $`t_{II}`$, a new ripple structure suddenly forms, in which the ripples are stable and rotated an angle $`\theta _c`$ with respect to the $`x`$ direction. Rost and Krug have shown \[for the deterministic limit of Eq. (126)\] that, by defining $`\alpha _\nu =\nu _x/\nu _y`$ and $`\beta _\lambda =\lambda _x/\lambda _y`$, the fact that $`\alpha _\nu >0`$ and $`\beta _\lambda <0`$ throughout Region II implies the existence of “cancellation modes” which dominate the dynamics and lead to this rotated ripple morphology. (Note the parameter ratios $`\alpha _\nu `$ and $`\beta _\lambda `$ are not to be confused with the roughness and growth exponents $`\alpha `$ and $`\beta `$ introduced in Section II.) The angle $`\theta _c`$ calculated by Rost and Krug has the value $`\theta _c=\mathrm{tan}^1\sqrt{\lambda _x/\lambda _y}`$ (see also Appendix D), and can be obtained by moving to a rotated frame of coordinates that cancels the nonlinear terms in the transverse direction. The boundary of Region $`II`$ is given by the condition $`\nu _x(a_\sigma ,\theta )=0`$, Eq. (119). Region III — This region is characterized by a positive $`\nu _x`$ and a negative $`\nu _y`$. At short time scales, $`tt_c`$, the periodic structure associated with the instability is oriented along the $`y`$ direction and has wavelength $`\mathrm{}_y=2\pi \sqrt{\frac{2K}{\|\nu _y\|}}`$. Much less is known, however, about the nonlinear regime. Such an anisotropic and linearly unstable equation is unexplored in the context of growth equations. The analysis by Rost and Krug for the corresponding deterministic equation predicts that, given that $`\beta _\lambda <\alpha _\nu <0`$ does hold all over Region III, again cancellation modes induce a rotated ripple morphology. Summary: Even though several aspects of the scaling behavior predicted by Eq. (126) remain to be clarified, we believe that this equation contains the relevant ingredients for understanding roughening by ion bombardment. To summarize, at short time scales the morphology consists of a periodic structure oriented along the direction determined by the largest in absolute value of the negative surface tension coefficients . Modifying the values of $`a_\sigma `$ or $`\theta `$ changes the orientation of the ripples . At long time scales we expect two different morphological regimes. One is characterized by the KPZ exponents, which might be observed in Region $`I`$ in Fig. 21. Indeed, the values of the exponents reported by Eklund et al. are consistent within the experimental errors with the KPZ exponents. In Region $`II`$ kinetic roughening is not expected. Rather, nonlinear terms lead to a new ripple structure that is rotated with respect to the ion direction. Region III is less understood; analysis of the deterministic equation again predicts a rotated ripple structure. By tuning the values of $`\theta `$ and/or $`a_\sigma `$ one may induce transitions among these morphological regimes. #### 2 High temperature morphology: Asymmetric case In this section we discuss the effect of asymmetry on the scaling regimes of Eq. (126). Here again we obtain that the effect of asymmetry does not bring in new qualitative features. Specifically, we find that the qualitative behavior of $`\xi _x`$ and $`\xi _y`$ is not affected by the asymmetry. As the asymmetry grows, the absolute value of the coefficients in the region of small angles decreases and the peaks at large $`\theta `$ increase. Similarly, while the numerical values of $`\lambda _x`$ and $`\lambda _y`$ are affected by the asymmetry $`\tau `$, their qualitative features are not. Finally, we find that the morphological diagram is topologically equivalent to the phase diagram obtained in the symmetric case (see Fig. 21), the only difference being in the position of the boundaries. As $`\tau `$ changes, we find a smooth evolution of the boundaries, which does not uncover any new regimes. Since the morphological properties of the system in the three regimes are the same as those discussed in the symmetric case, we will not discuss them again. #### 3 Low temperature morphology: Symmetric case In this section we discuss the effect of the effective surface smoothing on the surface morphology in the nonlinear regime. In the absence of thermally activated surface diffusion, from Eq. (65) we obtain the following equation governing the morphology evolution $`{\displaystyle \frac{h}{t}}`$ $`=`$ $`v_0+\gamma {\displaystyle \frac{h}{x}}+\xi _x\left({\displaystyle \frac{h}{x}}\right)\left({\displaystyle \frac{^2h}{x^2}}\right)`$ (139) $`+\xi _y\left({\displaystyle \frac{h}{x}}\right)\left({\displaystyle \frac{^2h}{y^2}}\right)+\nu _x{\displaystyle \frac{^2h}{x^2}}`$ $`+\nu _y{\displaystyle \frac{^2h}{y^2}}+{\displaystyle \frac{\lambda _x}{2}}\left({\displaystyle \frac{h}{x}}\right)^2+{\displaystyle \frac{\lambda _y}{2}}\left({\displaystyle \frac{h}{y}}\right)^2`$ $`+\mathrm{\Omega }_1{\displaystyle \frac{^3h}{x^3}}+\mathrm{\Omega }_2{\displaystyle \frac{^3h}{xy^2}}D_{xx}{\displaystyle \frac{^4h}{x^4}}`$ $`D_{xy}{\displaystyle \frac{^4h}{x^2y^2}}D_{yy}{\displaystyle \frac{^4h}{y^4}}+\eta (x,y,t).`$ The terms $`\gamma `$, $`\nu _x`$, $`\nu _y`$, $`\mathrm{\Omega }_1`$, $`\mathrm{\Omega }_2`$, $`\xi _x`$, $`\xi _y`$, $`\lambda _x`$, $`\lambda _y`$, as well as the ion-induced effective smoothing coefficients $`D_{xx}`$, $`D_{xy}`$ and $`D_{yy}`$ have been discussed in the previous sections. In the following we discuss the morphological phase diagram predicted by Eq. (139) and shown in Fig. 22. The different regimes have the following characteristics: Region I: The surface tensions $`\nu _x`$ and $`\nu _y`$ are negative while $`D_{xx}`$ and $`D_{yy}`$ are positive, and $`\lambda _x`$, $`\lambda _y`$ are both negative. This regime has been previously described in Section VI B 1 (Regime I in Fig. 21), the only difference being that here the ion-induced effective surface smoothing plays the role of $`K`$. The boundary of this region is given by $`D_{xx}(a_\sigma ,\theta )=0`$, Eq. (118). Region IIa: Here $`\nu _x`$, $`\nu _y`$, $`\lambda _x`$, and $`\lambda _y`$ are still negative, $`D_{yy}`$ is positive, while $`D_{xx}<0`$. Consequently, along the $`x`$ direction the surface is unstable, all modes growing exponentially. However, the nonlinear terms $`\lambda _x`$ and $`\lambda _y`$ have the same sign. The nonlinear regime in this parameter region has not yet been explored numerically, thus its morphology is unknown. The boundary of this region is given by $`\lambda _x(a_\sigma ,\theta )=0`$, Eq. (134). Region IIb: In this region $`D_{xx}`$ is negative, $`D_{yy}`$ is positive, $`\nu _x<0`$, $`\nu _y<0`$, and $`\lambda _x>0`$, $`\lambda _y<0`$. Thus, the only difference of this region with respect to Region IIa is that the nonlinear terms have different signs. Similarly to Region IIa, nothing is known about the nonlinear behavior. The boundary of this region is given by $`\nu _x(a_\sigma ,\theta )=0`$, Eq. (119). Region III: Here we have $`\nu _x>0`$, $`\nu _y<0`$, $`D_{xx}<0`$, $`D_{yy}>0`$, $`\lambda _x>0`$ and $`\lambda _y<0`$. This region has similar features to Region III in the phase diagram of Fig. 21, except for the negative value of the $`D_{xx}`$ coefficient. #### 4 Low temperature morphology: Asymmetric case In this section we discuss the effect of asymmetry on the long distance properties of Eq. (139). The effect of the asymmetry on the coefficients appearing in the equation were discussed earlier. Therefore we concentrate here merely on the morphological phase diagram predicted by Eq. (139) for the asymmetric case, which is displayed in Fig. 23. As before, asymmetry ($`\tau 1`$) leads to a smooth shift of the boundaries of the regions provided by the lines where the coefficients $`D_{xx}(a_\sigma ,\theta )`$, $`\nu _x(a_\sigma ,\theta )`$, and $`\lambda _x(a_\sigma ,\theta )`$ change sign. In the presence of effective smoothing, however, asymmetry in the deposited energy distribution induces the appearance of a fifth morphological regime. This is caused by the smooth motion of the boundary determined by $`\lambda _x(a_\sigma ,\theta )=0`$, which intersects for some value of $`\tau `$ the boundary defined by $`D_{xx}(a_\sigma ,\theta )=0`$. Comparison of Fig. 23 and Fig. 22 illustrates how the boundaries move with $`\tau `$. Regions I, IIa, IIb and III are analogous of those shown in Fig. 22. Region IIc in the phase diagram, on the other hand, is analogous of Region II of the high temperature phase diagram, shown in Fig. 21, and all the conclusions obtained in that section regarding the morphological properties in this regime also apply here. ## VII Comparison with experiments In this section we compare the predictions of the theory presented in this paper with experimental results on ripple formation and surface roughening. For a better presentation, we choose to structure the material around well known features of the morphological evolution, present the theoretical predictions and discuss to which extent are they supported by the available experimental data. We also discuss predictions that have not been tested in sufficient detail but could offer future tests of the theory. Ripple amplitude: A key quantity in ripple formation is the time evolution of the ripple amplitude. As we have shown in Section VI A 1, at early times ($`tt_c`$) the ripple amplitude grows exponentially, following $`h\mathrm{exp}(r(q_x^{},q_y^{})t)`$, where $`r`$ is the growth rate of the most unstable mode $`(q_x^{},q_y^{})`$. According to the linear theory, this growth continues indefinitely. In contrast, the nonlinear theory predicts that the amplitude should stabilize after time $`t_c`$, where $`t_c`$ is given by Eq. (133). This is consistent with the experimental investigations . On the other hand, recently Erlebacher et al. also found that at initial stages the ripple morphology is growing exponentially. Furthermore, they observed that at some time $`t_c`$ the exponential growth stops and the ripple amplitude saturates. Measuring the temperature dependence of the saturation curves, they found that rescaling the time $`t`$ with a factor $`\nu ^2/4K`$ and the amplitude $`h`$ with $`\sqrt{\nu /2K}`$, the different curves representing the amplitude as a function of time collapse onto a single one. This result is in excellent agreement with our prediction that suggest that plotting the result in terms of the rescaled parameters, $`t/t_c`$ and $`h\sqrt{\nu /2K}`$, the different curves should collapse . They also offer direct proof that the nonlinear terms play a major role in determining the amplitude of the ripples, indicating that the incorporation of the nonlinear mechanisms in the theory of ripple formation is essential. Temperature dependence of the ripple wavelength: A key quantity that provides direct information about the nature of the relaxation mechanism is the temperature dependence of the ripple wavelength. Our results indicate that there are two mechanisms contributing to surface relaxation: thermally induced surface diffusion (Sects. VI A 1-VI A 2) and ion-induced smoothing (Sects. VI A 3-VI A 4). At high temperatures thermal surface diffusion is rather intensive, thus it is the main mechanism determining the relaxation process, the ripple wavelength being given by (101) or (103). Since the surface diffusion constant $`K`$ decreases with $`T`$ as $`(1/T)\mathrm{exp}(\mathrm{\Delta }E/k_BT)`$, the ripple wavelength is also expected to decrease exponentially with $`T`$. Indeed, such an exponential temperature dependence of $`\mathrm{}`$ has been observed by various groups . However, in Section VI A 3 we demonstrated the existence of ion induced smoothing, that is present at any temperature. Thus, up to some inessential numerical factors, the total surface diffusion constants have a form $`D^T=K+D`$, where $`D`$ is independent of temperature. Since $`K`$ decreases exponentially with $`T`$, at low enough temperatures we have $`KD`$, indicating that the main relaxation mechanism is ion-induced. Consequently, below a certain critical temperature $`T_c`$, given by $`K(T_c)=D`$, one expects the ripple wavelength to be independent of $`T`$. Support for this scenario has been provided by the experiments in and and the molecular dynamics simulations of Koponen et al. . Consequently, as the theoretical results in this paper indicate, the temperature dependence of $`\mathrm{}`$ can be used to identify the relaxation mechanism: when $`\mathrm{}`$ increases exponentially with $`T`$, we are dealing with relaxation by thermal surface diffusion, while a temperature independent wavelength is an indication of ion induced smoothing. Ripple orientation: An important feature of ripple formation is that, as the linear theory predicts, the ripple orientation depends on the angle of incidence $`\theta `$. The dependence of the ripple orientation on the experimental parameters has been summarized in the phase diagrams shown in Figs. 12, 15, 16, 21, 22, and 23. In general, for physical values of the asymmetry parameter ($`i.e.`$ for $`\tau >1`$), we find that for small angles the ripples are oriented in the $`x`$ direction (along the incoming ions), and they change orientation to the $`y`$ direction for large $`\theta `$. The boundary separating these two morphological regions depends on the parameters characterizing the sputtering process, such as the ion penetration depth and the geometry of the deposited energy distribution. Such transition in the ripple orientation has been found in the simulations of , where for $`\theta 45^{}`$ the observed ripples were oriented along the $`x`$ direction, while for $`\theta 45^{}`$, they changed their orientation to $`y`$. Furthermore, the nonlinear theory predicts that after the nonlinear terms take over, new ripples, with orientation different from both $`x`$ and $`y`$ directions, might appear (see Sect. VI B 1). To the best of our knowledge, such rotated ripples have not been observed experimentally as yet. Nevertheless, this morphology might also lead to additional effects, such as shadowing, which have been neglected in our approach. Ripple wavelength dependence on the flux: Depending on the nature of the relaxation mechanism, the linear theory has two different predictions on the flux dependence of the ripple wavelength: for high temperatures, when thermal surface diffusion dominates, one expects $`\mathrm{}J^{1/2}`$ \[see Eq. (107)\], while at low temperatures, characterized by ion induced smoothing, we expect the wavelength to be independent of flux (see Sect. VI A 3). Consequently, due to its strong dependence on the relaxation mechanism, the flux dependence of the ripple wavelength can also be used to identify the relaxation mechanism. Indeed, a number of experiments are compatible with the prediction of a flux independent wavelength. Other aspects of ripple characteristics (such as energy or temperature dependence) also lead to the same conclusion. On the other hand, we are not aware of results indicating decreasing ripple wavelength with increasing flux. However, support for the relevance of thermal surface diffusion comes from the experiments of Chason et al. , who reported that the growth rate $`r(q_x^{},q_y^{})`$ as a function of flux follows the predictions of the linear theory with thermal surface diffusion. Ripple wavelength dependence on the ion energy: The linear theory indicates that the ion energy dependence of the ripple wavelength can be used to distinguish between the two relaxation mechanisms: at high temperatures we expect $`\mathrm{}ϵ^{1/2}`$ \[see Eq. (106)\], i.e., the wavelength decreases with the energy, while at low temperature we have $`\mathrm{}ϵ^{2m}ϵ`$ \[see Eq. (120)\], i.e., the wavelength should increase with energy, strikingly different predictions. A number of experimental groups have found that the ripple wavelength increases linearly with the ion energy . However, while we are not aware of any direct observation of a decreasing ripple wavelength with increasing ion energy, the growth rate dependence on the ion energy measured by Chason et al. provided results which are in agreement with the predictions based on thermal surface diffusion. The magnitude of the effective surface diffusion constant: Since the transition between the low and high temperature regimes is determined by the relative magnitude of $`D_{xx},D_{xy},D_{yy}`$ (ion induced smoothing), and $`K`$ (thermal surface diffusion), we need to estimate the magnitude of these constants. In the following we give an order of magnitude estimate for the effective surface diffusion constant $`D_{xx}`$ and compare it to $`K`$, using data from for Si(001). Taking $`Y=2.6`$, $`J=670`$ $`\mu `$A/cm<sup>2</sup>, $`ϵ=9`$ keV, $`a=100`$ Å, $`a_\sigma =2`$, $`a_\mu =4`$, and $`\theta =40^{}`$, Eq. (79) gives $`D_{xx}12\times 10^{28}`$ cm<sup>4</sup>/s. For comparison, at T = 550 C it is estimated that $`K34\times 10^{28}`$ cm<sup>4</sup>/s. Hence, since $`K`$ decreases exponentially with temperature, ion induced smoothing can be significant at low temperatures (including room temperature), in some cases being comparable or more relevant than thermal surface diffusion. Kinetic roughening: An important feature of our theory is that it goes beyond the linear approach, handling systematically the nonlinear effects as well. As we demonstrated in Sect. VI B, the presence of the nonlinear terms can affect both the dynamics and the morphology of the surface. The first and the most dramatic consequence is the stabilization of the ripple amplitude, discussed above. Furthermore, after the stabilization of the ripple amplitude the surface morphology is rather different from the morphology predicted by the linear theory. In particular, depending on the signs of $`\lambda _x`$ and $`\lambda _y`$, different morphological features can develop. When $`\lambda _x\lambda _y`$ is positive, at large times the surface undergoes kinetic roughening, following the predictions of the KPZ equation. This behavior has, indeed, been observed experimentally and numerically , providing direct support for the predictions of the nonlinear theory. When $`\lambda _x\lambda _y`$ is negative, direct numerical integration of the nonlinear theory indicated the existence of a new, rotated ripple structure. The absence of experimental data on this phase might be due to the required large sputtering times: the simulations indicate that between the linear regime and the formation of the rotated ripple structure there is a rather long transient regime with an apparently rough surface morphology. The above predictions apply when the surface diffusion terms, ion or thermally induced, act to smooth the surface (i.e., $`D^T0`$). However, at low temperatures, when ion induced smoothing dominates, surface diffusion can generate an instability that can further modify this behavior (see Sect. VI B 3). In general, while rather detailed experimental data are available describing the linear regime of ion sputtering, explanation of the nonlinear regime is only at its beginning, hiding the possibility of new interesting phases and behaviors. ## VIII Conclusions In this paper we investigated the morphological properties of surfaces eroded by ion bombardment. Starting from the expression for the erosion velocity derived in the framework of Sigmund’s theory of sputtering of amorphous targets, we derived a stochastic partial differential equation for the surface height, which involves up to fourth order derivatives of the height, and incorporates surface diffusion and the fluctuations arising in the erosion process due to the inhomogeneities in the ion flux. In some special cases the derived nonlinear theory reduces to the much studied KS or the KPZ equations, well known descriptions of dynamically evolving surfaces. However, in contrast with these theories, which have been derived using symmetry and conservation considerations , here we derived the continuum theory directly from a microscopic model of sputtering, and thus all coefficients can be explicitly expressed in terms of the physical parameters (such as angle of incidence, ion penetration depth, etc.) characterizing the ion bombardment process. An important feature of the derived nonlinear continuum theory is that the linear and the nonlinear regimes are separated in time. As a result, they can be discussed separately, the former controlling the behavior at early times, the latter at late times. Furthermore, an important result of our calculations is that higher order effects of the sputtering process can smooth the surface. This effective mechanism was necessary to explain ripple formation at low temperatures, when thermally induced surface diffusion is not relevant. Consequently, based on these two ingredients (separation of time scales between linear and the nonlinear regimes and the existence of two different relaxation mechanisms) we have discussed four different cases. In the linear high temperature regime the equations reduce to the linear theory of Bradley and Harper, predicting ripple formation, and explaining such experimentally observed phenomenon as ripple orientation (and its change with $`a_\sigma `$ and $`\theta `$), exponential increase in ripple amplitude (valid for short times), or flux and energy dependence of the ripple wavelength. On the other hand, phenomena not explained by this approach, such as the stabilization of the ripple amplitude, can be explained by considering the nonlinear terms as well. We also show that, depending on the sign of the coefficients of the nonlinear terms, the late time morphology of the surface is either rough, or dominated by new rotated ripples. The rough phase is expected to be described by the KPZ equation, which has its own significance: while the introduction of the KPZ equation has catalyzed an explosion in the study of the morphological properties of growing surfaces, there are very few actual surfaces that are described by it (and not by one of its offsprings, such as the MBE or related equations ), and most notably none, as far as we know, that describe crystal surfaces. The sputtering problem provides one of the first systems that is convincingly described by this continuum theory. Many of the previous mysteries of low temperature ripple formation have also been solved by the present theory. By deriving the higher order ion-induced effective smoothing terms, we can explain the existence of ripples at temperatures where thermally induced surface diffusion is not active. We showed that the derived effective smoothing affects both the linear and the nonlinear regimes, governing the early time ripple formation, and the late time nonlinear behavior. The coexistence of thermal and ion induced smoothing can explain the stabilization of the ripple wavelength at low temperatures, in contrast with its exponential $`T`$ dependence at higher temperatures. On the other hand, our theory has limitations, most of which can be already identified in Sigmund’s and Bradley and Harper’s theories, with which it is related. Namely, it is devised for amorphous substrates, whereupon it neglects effects such as viscous relaxation , which might be the cause for the failure of the theory to predict the absence in many experiments of ripples at low (but non-zero) angles of incidence. This issue should constitute one of the most important extensions of our present theory. Perhaps related with this, we have seen that there exist parameter regions at low temperatures within which our theory breaks down, due to the unstable higher order derivative terms that occur. A relevant issue is thus to determine the correct continuum description of the surface under these conditions. Most of the predictions offered by the presented continuum theory have been already verified experimentally. However, many unexplained predictions remain at low temperatures both in the linear and the nonlinear regimes, as well as regarding the nonlinear regime at high temperatures. We hope that the rather precise derivations offered in this paper will guide such future experimental work. Furthermore, some of the morphologies expected in the nonlinear regime need further theoretical understanding as well, allowing for the continuation of this inquiry. With the dramatic advances in computer speed, the understanding of some of these questions, either through numerical integration of the continuum theory or through discrete models, might be not too far. ## ACKNOWLEDGMENTS We would like to acknowledge discussions with E. Chason, B. Kahng, H. Jeong, F. Ojeda, and L. Vázquez. This research was supported by NSF-DMR CAREER and ONR-YI awards (A.-L. B. and M. M.) and DGES (Spain) grant PB96-0119 (R. C.). ## A The algebraic relation between the coordinates of the laboratory frame and the local frame, depicted in Fig. 3, follows from the definitions given in point (i) of Sect. V. Accordingly, the unit vector along the $`\widehat{Z}`$ axis is the normal at point $`O`$ $$\widehat{Z}\widehat{n}=\frac{(_xh,_yh,1)}{\sqrt{g}},$$ (A1) where $`g1+(_xh)^2+(_yh)^2.`$ The vector $`\widehat{m}`$ drawn on Fig. 3 has components $$\widehat{m}=(\mathrm{sin}\theta ,0,\mathrm{cos}\theta ).$$ (A2) Therefore, the unit vector along the $`\widehat{Y}`$ axis reads: $`\widehat{Y}`$ $``$ $`{\displaystyle \frac{\widehat{n}\times \widehat{m}}{\|\widehat{n}\times \widehat{m}\|}}={\displaystyle \frac{1}{\sqrt{g}\mathrm{sin}\phi }}((_yh)\mathrm{cos}\theta ,\mathrm{sin}\theta `$ (A4) $`+(_xh)\mathrm{cos}\theta ,(_yh)\mathrm{sin}\theta ),`$ and finally (A1) and (A4) yield for the unit vector along the $`\widehat{X}`$ axis: $`\widehat{X}`$ $`=`$ $`\widehat{Y}\times \widehat{Z}`$ (A5) $`=`$ $`{\displaystyle \frac{1}{g\mathrm{sin}\phi }}(\mathrm{sin}\theta +(_xh)\mathrm{cos}\theta (_yh)^2\mathrm{sin}\theta ,`$ (A8) $`(_yh)\mathrm{cos}\theta (_xh)(_yh)\mathrm{sin}\theta ,`$ $`(_xh)\mathrm{sin}\theta +((_xh)^2+(_yh)^2)\mathrm{cos}\theta ).`$ The matrix $``$ defined in Eq. (60) and which relates the coordinates in the local and laboratory frames reads $$=\left(\begin{array}{ccc}\frac{𝓈+𝒸_𝓎𝒽+𝓈(_𝓎𝒽)^\mathcal{2}}{\sqrt{}𝓇}& \frac{𝒸_𝓎𝒽}{𝓇}& \frac{_𝓍𝒽}{\sqrt{}}\\ \frac{𝒸_𝓎𝒽𝓈_𝓍𝒽_𝓎𝒽}{\sqrt{}𝓇}& \frac{𝓈+𝒸_𝓍𝒽}{𝓇}& \frac{_𝓎𝒽}{\sqrt{}}\\ \frac{𝓈_𝓍𝒽𝒸((_𝓍𝒽)^\mathcal{2}+(_𝓎𝒽)^\mathcal{2})}{\sqrt{}𝓇}& \frac{𝓈_𝓎𝒽}{𝓇}& \frac{\mathcal{1}}{\sqrt{}}\end{array}\right),$$ (A9) where $`s=\mathrm{sin}\theta `$, $`c=\mathrm{cos}\theta `$, and $`r\sqrt{(s+c_xh)^2+(_yh)^2}`$. ## B If we perform a small $`\mathrm{\Delta }_{nm}`$ expansion in Eq. (38) we obtain $`V_O`$ $`=`$ $`{\displaystyle \frac{ϵpJa^2}{\sigma \mu (2\pi )^{3/2}}}\mathrm{exp}(a_\sigma ^2/2)`$ (B1) $`\times `$ $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\zeta _X𝑑\zeta _Y\mathrm{exp}(\zeta _YL^2)`$ (B2) $`\times `$ $`\mathrm{exp}(\zeta _XA{\displaystyle \frac{1}{2}}\zeta _X^2B_1)[\mathrm{cos}\phi (1+B_2\mathrm{\Delta }_{20}\zeta _X^2`$ (B3) $`+`$ $`B_2\mathrm{\Delta }_{02}\zeta _Y^2+B_2\mathrm{\Delta }_{30}\zeta _X^3+B_2\mathrm{\Delta }_{12}\zeta _X\zeta _Y^2`$ (B4) $`+`$ $`B_2\mathrm{\Delta }_{22}\zeta _X^2\zeta _Y^2+B_2\mathrm{\Delta }_{40}\zeta _X^4+B_2\mathrm{\Delta }_{04}\zeta _Y^4`$ (B5) $``$ $`2C\mathrm{\Delta }_{20}\zeta _X^32C\mathrm{\Delta }_{02}\zeta _Y^2\zeta _X2C\mathrm{\Delta }_{30}\zeta _X^4`$ (B6) $``$ $`2C\mathrm{\Delta }_{12}\zeta _X^2\zeta _Y^22C\mathrm{\Delta }_{22}\zeta _X^3\zeta _Y^22C\mathrm{\Delta }_{40}\zeta _X^5`$ (B7) $``$ $`2C\mathrm{\Delta }_{04}\zeta _X\zeta _Y^4)+\mathrm{sin}\phi (2\mathrm{\Delta }_{20}\zeta _X+3\mathrm{\Delta }_{30}\zeta _X^2`$ (B8) $`+`$ $`\mathrm{\Delta }_{12}\zeta _Y^2+2\mathrm{\Delta }_{22}\zeta _X\zeta _Y^2+4\mathrm{\Delta }_{40}\zeta _X^3)].`$ (B9) Evaluating the Gaussian integrals in this formula we obtain Eq. (48), where the coefficients $`\mathrm{\Gamma }_{nm}(\phi )`$ are given by $`\mathrm{\Gamma }_{20}(\phi )`$ $`=`$ $`{\displaystyle \frac{2A}{B_1}}\mathrm{sin}\phi +{\displaystyle \frac{B_2}{B_1}}\left[1+{\displaystyle \frac{A^2}{B_1}}\right]\mathrm{cos}\phi `$ (B10) $`+`$ $`{\displaystyle \frac{2AC}{B_1^2}}\left[3+{\displaystyle \frac{A^2}{B_1}}\right]\mathrm{cos}\phi ,`$ (B11) $`\mathrm{\Gamma }_{02}(\phi )`$ $`=`$ $`2{\displaystyle \frac{\mu ^2}{a^2}}\mathrm{cos}\phi \left({\displaystyle \frac{B_2}{2}}+{\displaystyle \frac{AC}{B_1}}\right),`$ (B12) $`\mathrm{\Gamma }_{30}(\phi )`$ $`=`$ $`\mathrm{sin}\phi \left({\displaystyle \frac{1}{B_1}}+{\displaystyle \frac{A^2}{B_1^2}}\right)B_2\mathrm{cos}\phi \left({\displaystyle \frac{3A}{B_1^2}}+{\displaystyle \frac{A^3}{B_1^3}}\right)`$ (B13) $``$ $`2C\mathrm{cos}\phi \left({\displaystyle \frac{3}{B_1^2}}+{\displaystyle \frac{6A^2}{B_1^3}}+{\displaystyle \frac{A^4}{B_1^4}}\right),`$ (B14) $`\mathrm{\Gamma }_{12}(\phi )`$ $`=`$ $`2{\displaystyle \frac{\mu ^2}{a^2}}\{\mathrm{sin}\phi B_2\mathrm{cos}\phi {\displaystyle \frac{A}{B_1}}`$ (B15) $``$ $`2C\mathrm{cos}\phi ({\displaystyle \frac{1}{B_1}}+{\displaystyle \frac{A^2}{B_1^2}})\},`$ (B16) $`\mathrm{\Gamma }_{40}(\phi )`$ $`=`$ $`\{4\mathrm{sin}\phi ({\displaystyle \frac{3A}{B_1^2}}{\displaystyle \frac{A^3}{B_1^3}})`$ (B17) $`+`$ $`B_2\mathrm{cos}\phi \left({\displaystyle \frac{3}{B_1^2}}+{\displaystyle \frac{6A^2}{B_1^3}}+{\displaystyle \frac{A^4}{B_1^4}}\right)`$ (B18) $`+`$ $`2C\mathrm{cos}\phi ({\displaystyle \frac{15A}{B_1^3}}+{\displaystyle \frac{10A^3}{B_1^4}}+{\displaystyle \frac{A^5}{B_1^5}})\}.`$ (B19) $`\mathrm{\Gamma }_{22}(\phi )`$ $`=`$ $`2{\displaystyle \frac{\mu ^2}{a^2}}\{2\mathrm{sin}\phi {\displaystyle \frac{A}{B_1}}+B_2\mathrm{cos}\phi ({\displaystyle \frac{1}{B_1}}+{\displaystyle \frac{A^2}{B_1^2}})`$ (B20) $`+`$ $`2C\mathrm{cos}\phi ({\displaystyle \frac{3A}{B_1^2}}+{\displaystyle \frac{A^3}{B_1^3}})\}.`$ (B21) $`\mathrm{\Gamma }_{04}(\phi )`$ $`=`$ $`3{\displaystyle \frac{\mu ^4}{a^4}}\left\{B_2\mathrm{cos}\phi +2C\mathrm{cos}\phi {\displaystyle \frac{A}{B_1}}\right\}.`$ (B22) Taking into account Eqs. (50), (52) relating the local ($`\phi `$) and the global ($`\theta `$) angles of incidence through the surface slopes $`_xh`$, $`_yh`$, a small slope approximation leads to $`e^{A^2/2B_1}`$ $``$ $`e^{a_\sigma ^4s^2/2f}\{1+{\displaystyle \frac{a_\sigma ^4}{2f}}[{\displaystyle \frac{a_\mu ^2}{f}}s(_xh)+{\displaystyle \frac{a_\mu ^2}{f}}c^2(_yh)^2`$ (B23) $`+`$ $`{\displaystyle \frac{a_\mu ^2}{f^2}}(_xh)^2(a_\mu ^2c^2(1+2s^2)a_\sigma ^2s^2(1+2c^2)`$ (B24) $`+`$ $`{\displaystyle \frac{a_\sigma ^4a_\mu ^2}{f}}s^2c^2)]\},`$ (B25) $`B_1^{1/2}`$ $``$ $`{\displaystyle \frac{1}{\sqrt{f}}}\{1{\displaystyle \frac{a_\sigma ^2a_\mu ^2}{f}}sc(_xh)`$ (B26) $`+`$ $`{\displaystyle \frac{a_\sigma ^2a_\mu ^2}{2f^2}}(_xh)^2(a_\sigma ^2s^2(1+c^2)`$ (B27) $``$ $`a_\mu ^2c^2(1+s^2)){\displaystyle \frac{a_\sigma ^2a_\mu ^2}{2f}}c^2(_yh)^2\}.`$ (B28) Also, we have $`\mathrm{\Gamma }_{20}(\theta )`$ $`=`$ $`\mathrm{\Gamma }_{20}^{(0)}(\theta )+\mathrm{\Gamma }_{20}^{(1)}(\theta )(_xh),`$ (B29) $`\mathrm{\Gamma }_{02}(\theta )`$ $`=`$ $`\mathrm{\Gamma }_{02}^{(0)}(\theta )+\mathrm{\Gamma }_{02}^{(1)}(\theta )(_xh),`$ (B30) $`\mathrm{\Gamma }_{20}^{(0)}(\theta )`$ $`=`$ $`{\displaystyle \frac{a_\sigma ^2}{2f^3}}\{2a_\sigma ^6s^4c^2+2a_\sigma ^4s^4(s^22c^2)`$ (B31) $`+`$ $`a_\sigma ^4a_\mu ^2s^2c^2(s^2c^2)+a_\sigma ^2a_\mu ^2`$ (B32) $`\times `$ $`s^2c^2(7s^25c^2)+a_\mu ^4c^4(5s^2c^2)\},`$ (B33) $`\mathrm{\Gamma }_{20}^{(1)}(\theta )`$ $`=`$ $`{\displaystyle \frac{a_\sigma ^2sc}{f^4}}\{2a_\sigma ^6s^4a_\mu ^6c^4+a_\sigma ^6a_\mu ^2`$ (B34) $`\times `$ $`(s^2+s^2c^2)a_\sigma ^4a_\mu ^4(c^2+c^2s^2)`$ (B35) $`+`$ $`a_\sigma ^4a_\mu ^2(5s^23s^2c^2)+4a_\sigma ^2a_\mu ^4c^2\},`$ (B36) $`\mathrm{\Gamma }_{02}^{(0)}(\theta )`$ $`=`$ $`{\displaystyle \frac{c^2a_\sigma ^2}{2f}},`$ (B37) $`\mathrm{\Gamma }_{02}^{(1)}(\theta )`$ $`=`$ $`{\displaystyle \frac{a_\sigma ^4cs}{f^2}}.`$ (B38) In the above expressions we used the notations $`a_\sigma {\displaystyle \frac{a}{\sigma }},a_\mu {\displaystyle \frac{a}{\mu }},s\mathrm{sin}\theta ,`$ (B39) $`c\mathrm{cos}\theta ,fa_\sigma ^2s^2+a_\mu ^2c^2.`$ (B40) ## C Equations (55)-(57) relating the incidence angle as measured in the local and laboratory reference frames apply only in the off-normal incidence case ($`\theta 0`$). In the following we derive the correct expressions for the normal incidence case ($`\theta =0`$). Indeed, if $`\theta =0`$, the vectors $`\widehat{n}`$ and $`\widehat{m}`$ shown in Fig. 3 are given by $`\widehat{n}={\displaystyle \frac{(_xh,_yh,1)}{\sqrt{g}}},\widehat{m}=(0,0,1).`$ (C1) Proceeding now as in (50)-(52), we obtain $`\mathrm{cos}\phi ={\displaystyle \frac{1}{\sqrt{g}}},\mathrm{sin}\phi =\sqrt{{\displaystyle \frac{(_xh)^2+(_yh)^2}{g}}},`$ (C2) which are the $`\theta 0`$ limit of Eqs. (50) and (52). The small gradient expansion performed on Eq. (C2) now gives $`\mathrm{cos}\phi `$ $``$ $`1{\displaystyle \frac{1}{2}}((_xh)^2+(_yh)^2),`$ (C3) $`\mathrm{sin}\phi `$ $``$ $`\sqrt{(_xh)^2+(_yh)^2}.`$ (C4) Using Eqs. (C4) in the expansions leading to Eq. (65), it can be seen that the expressions obtained for the coefficients indeed are the $`\theta 0`$ limit of Eqs. (66)-(85). ## D The solution corresponding to a rotated ripple structure follows from Eq. (126). Indeed, in the absence of the $`\xi _x`$ and $`\xi _y`$ terms, if we consider a solution of (126) of the form $`h(x,y,t)=g(xvy,t)`$ with $`v`$ an arbitrary constant, the surface morphology evolution equation takes the form $`_tg`$ $`=`$ $`v_0+\gamma _lg+(\nu _x+v^2\nu _y)_l^2g`$ (D3) $`+{\displaystyle \frac{1}{2}}(\lambda _x+v^2\lambda _y)(_lg)^2+(\mathrm{\Omega }_1+v^2\mathrm{\Omega }_2)_l^3g`$ $`K(1+v^2)^2_l^4g,`$ where $`g(l)=g(xvy)`$ is the steady wave solution . From (D3) it follows that the nonlinearity vanishes when $`\lambda _x+v^2\lambda _y=0`$, or $`v=\sqrt{\lambda _x/\lambda _y}`$. In this case we obtain an exponentially growing ripple structure with ripples forming an angle $`\theta _c=\mathrm{tan}^1(v)=\mathrm{tan}^1(\sqrt{\lambda _x/\lambda _y})`$ with respect to the $`x`$ axis.
warning/0007/cond-mat0007276.html	ar5iv	text	# Molecular mode-coupling theory applied to a liquid of diatomic molecules ## I Introduction The mode coupling theory (MCT) of the glass transition is by now an important tool to understand experiments in and simulations of supercooled liquids . For a long time most of the theoretical investigations concentrated on simple monoatomic or binary liquids. All universal and even system specific predictions of these investigations could be tested on a quantitative level in a system of hard colloids , which is an excellent realization of a hard sphere system and in computer simulations for a binary Lennard Jones system . Details of theory and tests for simple glass formers can be found in review articles and articles cited therein. Although the theory was originally formulated only for these simple systems, most of the experimental and simulation support came from research on much more complex system (e.g. tri-$`\alpha `$-naphtylbenzene , Orthoterphenyl (OTP) , $`0.4Ca(NO_3)_20.6KNO_3`$ (CKN) , Glycerol , Salol , toluene and water ). Also most of the experimental methods used, did not measure density correlation functions or their susceptibilities, for which the original theory was formulated. Even neutron scattering experiments in systems consisting of molecules whose components have different cross sections for neutrons do not measure the density correlation function exclusively, but a mixture of more complicated correlation functions involving molecular degrees of freedom (see also for a single linear molecule). Dielectric loss measurements measure directly the correlation function of a tensor of rank $`1`$. Depolarized light scattering , Kerr effect experiments , NMR and ESR (and references therein) measure correlation function of a tensor of rank $`2`$. The mentioned tensorial quantities are all related to orientational degrees of freedom (ODOF), whereas the original theory only considered translational degrees of freedom (TDOF), i.e. the center of mass motion. But of course, when comparing experimental results on complex systems with predictions of the MCT for simple liquids it was always reasonable to argue, that there are in every experiment couplings to the center of mass motion. E.g. the reorientation of dipoles measured in dielectric loss measurements can induce center of mass motion via a translation - rotation coupling. Also the reorientation of the polarizability tensor in light scattering measurements is related to a physical rotation of the molecules and will therefore be coupled to the center of mass motion of the molecules as well. A slowing down of this motion due to very slow structural relaxations can consequently also indirectly be measured in the mentioned experiments. In addition it is perfectly justified to perform tests of the universal predictions of MCT in complicated molecular and polymeric systems for $`\beta `$ scaling laws and properties of the $`\alpha `$ relaxation, like time temperature superposition principle and wave vector dependent stretching exponents, since the underlying universal features of the bifurcation scenario should also remain valid for molecular systems. But beyond the universal aspects MCT aims to be a microscopic theory of structural relaxation. This goal is to a large extent achieved for simple liquids. There it was possible to obtain quantitative agreement between experiment and theory for the full dynamic range of structural relaxation (i.e. $`\beta `$ and $`\alpha `$ \- relaxation) . Recently also a theory for anomalous high frequency oscillations (Bose peak phenomenon) was formulated within MCT . The molecular mode coupling theory (MMCT) which is under study for a few years now, intends to extend this line of research to experimentally relevant molecular systems. There are three different mode coupling theories for the description of different aspects of molecular degrees of freedom. In the motion of a single linear molecule in a liquid of spherical atoms is studied. In a site - site description is formulated, which is perfectly adapted to study neutron scattering experiments of molecular systems. In this approach the atomic structure of the molecules is considered. Finally, the MMCT , we are using in this work , is devised to investigate the dynamics of a liquid of linear molecules. For this purpose a self consistent mode coupling theory for the dynamic correlation functions of tensorial densities $`\rho _{lm}(\stackrel{}{q},t)`$ was developed. These densities are the generalized Fourier components of the microscopic density $`\rho (\stackrel{}{x},\mathrm{\Omega })`$ in an expansion in spherical harmonics with respect to the orientation $`\mathrm{\Omega }=(\theta ,\phi )`$ of the molecules and plane waves with respect to $`\stackrel{}{x}`$. An extension to arbitrary molecules is given in . A theory for arbitrary molecules was also formulated in . First results for the tensorial nonergodicity parameters (NEP) and critical amplitudes were obtained for dipolar hard spheres and , respectively. A study of the phase diagram for glass transitions of a liquid of hard ellipsoids was performed in ref. . Several aspects of the theory for general molecules were tested against simulation for water in ref. and treating water as a linear molecule in ref. . As in the MCT for simple liquids the static structure factors $`S_{ll^{}}^m(q)`$ in the q frame, i.e. the coordinate system in which the $`z`$ axis is along the wavevector, completely determines the longtime dynamics and thus the NEP and critical amplitudes. Note that the structure factors (static and dynamic) are diagonal in $`m`$ in the q -frame . The static structure factors have to be known to solve the equations of MMCT. They are either obtained by analytical theories as e.g. in ref. for ellipsoids or they have to be taken from simulations. In this work we present a detailed comparison of MMCT calculations of the NEP and critical amplitudes with the results of simulation for a liquid of diatomic molecules. A description of the system and the simulation can be found in . There also tests of the universal properties of MCT are presented. For our comparison between simulation and theory the static structure factors are taken from the simulation. In a preliminary study only the diagonal static correlators $`S_{ll}^m(q)`$ were used as input and also the dynamical correlators $`S_{ll}^m(q,t)`$ and thus the NEP were assumed to be diagonal . This severe approximation has lead to unphysical results like the existence of two different transition temperature for ODOF and TDOF. In our study we use diagonal and non diagonal elements of the static structure factor as input to calculate all components of the NEP and, in addition, of the critical amplitudes. For the calculation of the NEP, we also extend the necessary upper cut off $`l_{co}`$ for the index $`l`$ to $`l_{co}=4`$. The paper is organized as follows. In chapter (II) we review the main equations and concepts for the calculation of the critical NEP (II A) and the critical amplitudes (II B). In chapter (III) we discuss the influence of different approximation schemes on the theoretically obtained critical temperatures (III A) and its relation to simulation results. Then we present the comparison of theoretical critical NEP (III B) and critical amplitudes (III C) with simulations. Conclusions are presented in chapter IV and an appendix describes how the critical NEP can be obtained from the liquid side close the the ideal glass transition. ## II Molecular Mode Coupling Theory ### A Non ergodicity parameter The derivation of the equation of MMCT for the dynamics of linear molecules and general molecules can be found in and respectively. We only repeat the basic definitions and equations and refer the reader for details to the literature. For the present work, in which we want to calculate the critical NEP and the critical amplitudes of diatomic molecules, we only need the equations for linear molecules in the limit of time to infinity. The basic quantities are the correlation functions of tensorial densities $`\rho _{lm}(\stackrel{}{q})`$ and tensorial current densities $`j_{lm}^\alpha (\stackrel{}{q},t)`$ $`\rho _{lm}(\stackrel{}{q},t)`$ $`=`$ $`\sqrt{4\pi }i^l{\displaystyle \underset{n=1}{\overset{N}{}}}e^{i\stackrel{}{q}\stackrel{}{x}_n(t)}Y_{lm}(\stackrel{}{\mathrm{\Omega }}_n(t))\text{.}`$ (1) $`\stackrel{}{j}_{lm}^\alpha (\stackrel{}{q},t)`$ $`=`$ $`\sqrt{4\pi }i^l{\displaystyle \underset{n=1}{\overset{N}{}}}\stackrel{}{v}_n^\alpha e^{i\stackrel{}{q}\stackrel{}{x}_n(t)}Y_{lm}(\stackrel{}{\mathrm{\Omega }}_n(t)).`$ (2) The $`Y_{lm}(\stackrel{}{\mathrm{\Omega }}_n(t))`$ are the standard spherical harmonics and we follow in our notation the text book by Gray and Gubbins . $`\stackrel{}{v}_n^\alpha `$ is either the center of mass velocity $`\stackrel{}{v}_n(t)`$ of the $`n^{th}`$ molecule or its angular velocity $`\stackrel{}{\omega }_n(t)`$ depending on the index $`\alpha `$: $$\stackrel{}{v}_n^\alpha (t):=\{\begin{array}{cc}\hfill \stackrel{}{v}_n(t),& \alpha =T\hfill \\ \hfill \stackrel{}{\omega }_n(t),& \alpha =R\hfill \end{array}.$$ (3) where $`T`$ and $`R`$ stands for translational and rotational part, respectively. For the calculation of the NEP we in principal need all spatial components of the currents but here we use as in only the projection on directions defined by the wave vector $`\stackrel{}{q}`$ and the angular momentum operator $`\stackrel{}{L}`$. Taking into account also transversal currents will lead only to small correction in the NEP . We therefore define the longitudinal currents $`j_{lm}^\alpha (\stackrel{}{q},t)`$. $$j_{lm}^\alpha (\stackrel{}{q},t)=\frac{1}{q_l^\alpha }\left(\widehat{q}^\alpha \stackrel{}{j}^\alpha \right)_{lm}(\stackrel{}{q}),\alpha \{T,R\}$$ (4) with $$q_l^\alpha (q):=\{\begin{array}{cc}\hfill q,& \alpha =T\hfill \\ \hfill \sqrt{l(l+1)},& \alpha =R\hfill \end{array}$$ (5) and the definition $$\widehat{q}^\alpha :=\{\begin{array}{cc}\hfill \stackrel{}{q},& \alpha =T\hfill \\ \hfill \stackrel{}{L},& \alpha =R\hfill \end{array}$$ (6) . The quantities we are going to calculate are NEP $`F_{ll^{}}^m(q)`$ in the q -frame i.e. in a coordinate system in which the z-axis is given by the direction of the wave vector $`\stackrel{}{q}=(0,0,q)`$. In this coordinate system all correlation functions $`\rho _{lm}^{}(\stackrel{}{q},t)\rho _{l^{}m^{}}(\stackrel{}{q},0)=\delta _{mm^{}}S_{l,l^{}}^m(q,t)`$ are diagonal in $`m`$. They are real and depend on $`\|m\|`$ only. The same holds for all other tensorial quantities, we will use in the q - frame. The NEP are given by $$F_{ll^{}}^m(q)=\underset{t\mathrm{}}{lim}S_{ll^{}}^m(q,t).$$ (7) As input for the mode coupling equations we need the static structure factors $`S_{ll^{}}^m(q)`$ They are directly taken from the simulation of Kämmerer et. al. . The off diagonal elements of the static structure factors were not yet published and had to be determined from the raw data. The equation for the matrix $`𝐅(q)(F_{ll^{}}^m(q))`$ of NEP can then be written as $$𝐅(q)=\left[\mathrm{𝟏}+𝓚(q)𝐒^1(q)\right]^1𝐒(q).$$ (8) The matrix $`𝓚`$ is related to the Laplace transform of the dynamic current correlation function and can be expressed as the inverse of a memory matrix $`𝓕^m(q)(_{ll^{}}^m(q))`$ at $`t=\mathrm{}`$, $$𝒦_{ll^{}}^m(\stackrel{}{q})=\underset{\alpha \alpha ^{}}{}q_l^\alpha (q)\left(𝓕^𝒎(q)^1\right)_{ll^{}}^{\alpha \alpha ^{}}q_l^{}^\alpha ^{}(q).$$ (9) The mode coupling approximations yield: $`_{ll^{}}^{m\alpha \alpha ^{}}(q){\displaystyle \frac{1}{2N}}\left({\displaystyle \frac{\rho _0}{4\pi }}\right)^2{\displaystyle \underset{\stackrel{}{q}_1\stackrel{}{q}_2}{}^{}}{\displaystyle \underset{m_1m_2}{}}{\displaystyle \underset{l_1l_2}{}}{\displaystyle \underset{l_1^{}l_2^{}}{}}\times `$ (10) $`\times V_{ll^{}l_1l_1^{}l_2l_2^{}}^{\alpha \alpha ^{}}(q,q_1,q_2;m,m_1,m_2)F_{l_1l_1^{}}^{m_1}(q_1)F_{l_2l_2^{}}^{m_2}(q_2),`$ (11) with $`V_{ll^{}l_1l_1^{}l_2l_2^{}}^{\alpha \alpha ^{}}(q,q_1,q_2;m,m_1,m_2):=`$ (12) $`v_{ll_1l_2}^\alpha (q,q_1,q_2;m,m_1,m_2)v_{l^{}l_1^{}l_2^{}}^\alpha ^{}(q,q_1,q_2;m,m_1,m_2)^{},`$ (13) $`v_{ll_1l_2}^\alpha (q,q_1,q_2;m,m_1,m_2):=`$ (14) $`{\displaystyle \underset{l_3}{}}u_{ll_3l_2}^\alpha (q,q_1,q_2;m,m_1,m_2)c_{l_3l_1}^{m_1}(q_1)+(1)^m(12)`$ (15) and $$u_{ll_1l_2}^\alpha (q,q_1,q_2;m,m_1,m_2):=i^{l_1+l_2l}\left[\frac{(2l_1+1)(2l_2+1)}{(2l+1)}\right]^{\frac{1}{2}}\frac{1}{2}[1+(1)^{l_1+l_2+l}]\times $$ $`\times `$ $`{\displaystyle \underset{m_1^{}m_2^{}}{}}(1)^{m_2^{}}d_{m_1^{}m_1}^{l_1}(\mathrm{\Theta }_{q_1})d_{m_2^{}m_2}^{l_2}(\mathrm{\Theta }_{q_2})C(l_1l_2l;m_1^{}m_2^{}m)`$ (16) $`\times `$ $`\{\begin{array}{cc}q_1\mathrm{cos}\mathrm{\Theta }_{q_1}C(l_1l_2l;000);\hfill & \alpha =T\hfill \\ \sqrt{l_1(l_1+1)}C(l_1l_2l;101);\hfill & \alpha =R\hfill \end{array}\text{.}`$ (19) Here the functions $`C(l_1l_2l;m_1,m_2,m)`$ are the Clebsch Gordan coefficients and $`d_{m^{}m}^l(\mathrm{\Theta })`$ are related to Wigner’s rotation matrices (we follow the notation of Gray and Gubbins). For given Euler angles $`\mathrm{\Omega }=(\mathrm{\Phi },\mathrm{\Theta },\chi )`$ they are defined as ) $$D_{mm^{}}^l(\mathrm{\Omega })=e^{im\mathrm{\Phi }}d_{mm^{}}^l(\mathrm{\Theta })e^{im^{}\chi }.$$ (20) $`\mathrm{\Theta }_{q_i}`$ is the angle between $`\stackrel{}{q}`$ and $`\stackrel{}{q}_i`$. The prime at the first summation in Eq.(10) restricts $`\stackrel{}{q}_1,\stackrel{}{q}_2`$ such, that $`\stackrel{}{q}_1+\stackrel{}{q}_2=\stackrel{}{q}`$. Eqs. (8 \- 9) and (10) form a closed set of infinitely many coupled nonlinear equations for the NEP. To obtain a solvable theory we have to restrict $`l`$ to be smaller than an upper cut off, $`ll_{co}`$. The resulting equations can in principle be solved by a fixed point iteration algorithm. The physical control parameter like the temperature and the density only enter via the static structure factor. At a critical temperature or density the solution of this equations bifurcates from all functions $`F_{ll^{}}^m(q)`$ being zero to nonzero. In the simulations of Kämmerer et al. the temperature was used as control parameter. Close to the transition temperature $`T_c`$ the stability matrix of the iteration (see below) will have a largest eigenvalue $`E_0`$ approaching $`E_0=1`$ from below. Consequently the convergence of the iteration is getting very slow close to $`T_c`$. The time for one iteration depends very sensitive on the upper cut off $`l_{co}`$. For $`l_{co}=2`$ one iteration took $`10`$ minutes on a MIPS R10000, for $`l_{co}=4`$ this time increased to 6 hours. We therefore concentrated on $`l_{co}=2`$, to determine the transition point with a high accuracy and used the calculation for $`l_{co}=4`$ mostly as check for the sensitivity of our results against changing the cut off. To overcome some of the restrictions connected to the critical slowing down of the convergence close to $`T_c`$ we determined the critical NEP, i.e. the NEP at $`T_c`$, with two alternative methods. For the standard fixed point iteration we started at a temperature low enough to be in the glass state. Then the temperature is increased very slowly. At every temperature the equations for the NEP are solved by the iteration $$𝐅^{(n+1)}(q)=𝓖(𝐅^n,ϵ)$$ (21) where $`𝓖(𝐅^{(n)})`$ is the right hand side of Eq. (8) and $`ϵ=(T_cT)/T_c`$ . This iteration converges exponentially fast towards its solution as long as the temperature is not the critical temperature. The convergence rate is determined by the largest eigenvalue $`E_0`$ of the stability tensor $`𝒞_{\lambda \lambda ^{}}=𝒢_\lambda /F_\lambda ^{}`$ . The index $`\lambda `$ is an abbreviation for wavevectors $`q`$ and rotational indices $`l,l^{},m`$. The exponential convergence rate is then $`\mathrm{ln}E_0`$. Close to and below $`T_c`$ the eigenvalue $`E_0`$ can be written as $`E_0=1A\sqrt{ϵ}`$ with A being a positive constant. Therefore the convergence rate is $`A\sqrt{ϵ}`$ and the number of iterations to obtain convergence diverges inversely proportional to $`\sqrt{ϵ}`$ close to $`T_c`$. With this number of iterations the deviations of our NEP from the true critical NEP are proportional to $`\sqrt{ϵ}`$, since the NEP exhibit the well known square root singularity (cf. ). If the temperature is increased above $`T_c`$, i.e. $`ϵ<0`$, there is no non zero solution for the iteration Eq. (21). Nevertheless for $`0<ϵ1`$, the iteration is nearly stationary for a large number of iteration of the order $`\|ϵ\|^{1/2}`$ (see appendix). The approximate critical NEP is determined as the stationary point $`\widehat{𝐅}(ϵ)`$ whose change along the eigenvector with eigenvalue 1 of the critical stability matrix is minimal during iteration. But contrary to the NEP determined from the fixed point iteration for $`T<T_c`$, the stationary NEP differ only in order $`\|ϵ\|`$ instead in order $`\sqrt{ϵ}`$ from the true critical NEP. Consequently this property allows us in the following to crosscheck the very accurate results for $`l_{co}=2`$ obtained from the fixed point iteration and also to obtain the critical NEP for $`l_{co}=4`$. ### B Critical amplitudes A central prediction of the mode coupling theory of the glass transition in simple liquids is the existence of the $`\beta `$ \- relaxation regime . For the problem of a single dumbbell in an isotropic hard sphere system it was demonstrated, that the $`\beta `$ relaxation law can be detected in every quantity, which couples to the density. For a liquid of anisotropic molecules it is not very well defined which degree of freedom is driving the glass transition. The equations of MMCT couple all degrees of freedom and there are situations where the transition is not caused by the TDOF, but the ODOF . But even in these systems the factorization theorem is generically valid for all correlators. This can be proven using the standard techniques . Therefore every dynamic structure factor $`S_{ll^{}}^m(q,t;T)`$ can for $`1ϵ1`$ in the $`\beta `$ relaxation regime be written as $$S_{ll^{}}^m(q,t;T)=F_{ll^{}}^m(q)+H_{ll^{}}^m(q)G(t/t_0;\sigma )$$ (22) The function $`G(t/t_0)`$ is the same for all $`l,l^{},m,q`$. $`t_0`$ is an overall microscopic scale. $`H_{ll^{}}^m(q)`$ are the critical amplitudes determining the intensity of the asymptotic $`\beta `$ \- relaxation for a certain combination of $`l,l^{},m`$ and $`q`$. Also the correction to the asymptotics, which determine, besides the temperature, the range of validity of the law Eq. (22) depend on these amplitudes (cf. ref for simple liquids). Differences in the observability of the critical correlators between depolarized light scattering experiments and dielectric loss measurements, can be explained by differences in the amplitudes $`H_{ll^{}}^m(q)`$ involving $`l=2`$ and $`l=1`$ respectively (see ref. for a single molecule). To determine the amplitudes numerically to a high precision it is necessary to be very close to the transition point, to make sure that all correction terms of order $`ϵ`$ are small compared to the leading term of order $`\sqrt{ϵ}`$. Due to the difficulties described above we could only determine the critical amplitudes for upper cut off $`l_{co}=2`$. ## III Results ### A The critical temperatures For $`l_{co}=2`$ and in the full diagonalization approximation , in which the static structure factors $`𝐒(q)`$, the glass form factors $`𝐅(q)`$ and the memory matrix $`(q)`$ are assumed to be diagonal with respect to $`l`$, the glass transition temperature for the TDOF predicted by MMCT is below the transition temperature of the MD simulations $`T_c^{MD}=0.477`$. Note, that this temperatures are given in Lennard Jones units (cf. ). In all other known examples the MCT overestimates the tendency for vitrification. As an additional artefact of the full diagonalization the ODOF vitrify at a lower temperature than the TDOF. Since the top down symmetry of the dumbbells is broken, the full equations of MMCT (8) - (10) do generically not allow for such a scenario. As soon as we take $`𝐒(q)`$ non diagonal all degrees of freedom undergo a glass transition at the same temperature above the MD result. To study the influence of different diagonalization approximations a bit more in detail, we investigated several cases, with the main condition of $`𝐒(q)`$ being non diagonal: 1. $`𝐅(q)`$ and $`(q)`$ diagonal (dd) 2. $`𝐅(q)`$ diagonal and $`(q)`$ non diagonal (dnd) 3. $`𝐅(q)`$ and $`(q)`$ non diagonal (ndnd) Let us discuss $`l_{co}=2`$ first. For this case $`T_c`$ has been determined very accurately from the asymptotic behavior $`(1E_0(T))^2ϵ`$ (see chapter II B) for the largest eigenvalue $`E_0(T)`$. Fig. 1. demonstrates this law for case (3). The highest transition temperature is obtained for case (1). Here the transition temperature is roughly three times as large as the MD result, $`T_c^{dd}=1.4`$. In case (2) it only slightly decreases to $`T_c^{dnd}=1.38`$. If everything is taken non diagonal (case 3) the transition temperature is $`T_c=0.7521`$. Although still twice as large as the MD result, the discrepancy of this result from $`T_c^{MD}`$ is comparable to other known cases and consistent with the usual $`20\%`$ accuracy of the critical density . The equations are too complex to get a deeper theoretical understanding of this seemingly erratic jumping of the transition temperature, dependent on the approximation we are using. Particularly the fact, that the vertices (Eq. 14) are not positive anymore, makes an analytical prediction impossible. But it is at least possible to rationalize the behavior using a combination of physical and mathematical arguments. First of all it is quite clear that the full diagonalization where all matrices are assumed to be diagonal is too crude to describe the coupling of TDOF and ODOF for the system of diatomic molecules. The TDOF and ODOF are only coupled via the diagonal memory function $`_{ll^{}}^{m\alpha \alpha ^{}}(q)`$. The coupling of the equations for different $`l`$ is considerably reduced compared to the case where $`𝐒(q)`$ is taken to be non diagonal. E.g. in the $`_{00}^{mTT}(q)`$ component of the memory matrix only terms of the form (symbolically) $`_{l^{},m^{}}V_l^{}^m^{}(F_{l^{}l^{}}^m^{})^2`$ appear in the full diagonalization approximation, since the Clebsch Gordan coefficients $`C(0l^{}l^{\prime \prime },0,m^{},m^{\prime \prime })`$, which enter into the vertices (cf. Eqs. (10) \- (19)) are nonzero for $`l^{}=l^{\prime \prime }`$ only. Similarly, the memory functional $`_{11}^{m\alpha \alpha ^{}}(q)`$ only contains couplings of the form $`F_{00}^0F_{11}^m`$ and $`F_{22}^mF_{11}^m`$. For $`_{22}^{m\alpha \alpha ^{}}(q)`$ the Clebsch Gordan coefficients in the vertex allow “self” couplings $`(F_{11}^m)^2`$, $`(F_{22}^m)^2`$ and couplings to $`l=0`$ in the form $`F_{00}^0F_{22}^m`$, but no “self” couplings $`(F_{00}^0)^2`$. Due to the absence of $`(F_{00}^0)^2`$ in $`_{22}^{m\alpha \alpha ^{}}(q)`$ a freezing of the center of mass motion, i.e. $`l=l^{}=0`$ does not imply a freezing for quadrupolar dynamics $`l=l^{}=2`$. If the vertex for the coupling of the NEP with $`l=2`$ and $`l=0`$ in $`F_{11}^m`$ and $`F_{22}^m`$ are not large enough, exactly this structure of the memory matrix allows generically a separate transition of the $`l=0`$ and the $`l0`$ components of the diagonalized dynamic structure factor as observed in . But we have to stress that the approximation is not per se inadequate. In the case of water the full diagonalization approximation leads to a rather satisfactory agreement with simulations, without the artefact of separate transitions and too low transition temperatures. The pronounced angular dependence of interaction between water molecules, which is reflected in the fact that the static structure factors for $`l=2,m=0,1,2`$ and $`l=0`$ are of the same order, yield large enough vertices to produce a single transition temperature of the TDOF and ODOF. In the present case, the structure factor $`S_{00}^0`$ is clearly more dominant than $`S_{22}^m`$ (see Fig. 2). This is different from water where all $`m`$ for $`l=l^{}=2`$ are important. This leads to the a posteriori conclusion that in general the full diagonalization can only be used (if at all) for systems with “very strong” static translation rotation coupling. This statement, unfortunately, cannot be further quantified. If we now take the non diagonality of $`𝐒(q)`$ serious but leave all other matrices diagonal (case 1), additional coupling between TDOF and ODOF appear, which may lead to an effectively stronger coupling. Although the equations for the different $`l`$ components of the NEP still couple only via the diagonal memory functions $`_{ll}^{m\alpha \alpha ^{}}(q)`$, the diagonalized memory matrix contains now static couplings between all NEP. E.g., $`_{11}^{m\alpha \alpha ^{}}(q)`$ contains additional couplings between the TDOF - correlator $`F_{00}^0`$ and the correlators involving $`l=2`$ and, even more important, ”self coupling” terms $`(F_{00}^0)^2`$ due to the non vanishing structure factors $`S_{10}^0`$ and $`S_{12}^m(q)`$, respectively. This of course does not explain, but at least makes plausible the dramatic increase of the transition temperature. Any slowing down of let’s say the TDOF is immediately transferred to all other degrees of freedom and causes a further slowing down of the TDOF due to the feed back via the memory function. This enhances the tendency towards vitrification and also is responsible for the existence of a single transition temperature. The reason for the decrease of the transition temperature, when we give up the diagonalization approximations for $`(q)`$ and $`𝐅(q)`$ (case 3), is not obvious. We only note, that the transition temperature is decreasing from case 1 ($`T_c=1.4`$) over 2 ($`T_c=1.38`$) to 3 ($`T_c=0.752`$) i.e. the more off diagonal elements of the matrix $`𝓕`$ are taken into account. $`𝓕`$ is the $`t\mathrm{}`$ limit of the memory matrix i.e. the random force correlation function. Therfore it seems, that the more components of the random forces are coupled, the lower is the transition temperature. This implies, that the more the random forces can mutually influence each other, the more difficult it is to form a glass. Although we cannot proof this statement on mathematical grounds, it describes an feasible physical phenomenon. To test the sensitivity of our results to changes in the cut off, we also solved the MMCT for upper cut off $`l_{co}=4`$. The larger cut off value for $`l`$ reduces the transition temperature further towards the simulation result. For $`l_{co}=4`$, an upper and a lower bound for $`T_c`$ has been determined. The lower bound is the highest temperature for which the NEP are still nonzero after 88 iterations. The upper bound is the temperature for which the NEP are converging to zero after about 24 iterations. Since the time per iteration increases dramatically upon increasing the upper cut off the transition temperature could only be determined within $`5\%`$, $`T_c=0.61`$. It is encouraging that the real transition is approached upon increasing the cut off, but our arguments presented above, show, that this is not necessarily the case. Which of the competing mechanisms influencing the transition temperature is dominant, cannot be predicted on general grounds. ### B The non ergodicity parameters In the following we concentrate on the results for the normalized NEP $`f_{ll^{}}^m(q)=F_{ll^{}}^m(q)/\sqrt{S_{ll}^m(q)S_{l^{}l^{}}^m(q)}`$ without any diagonalization approximation. Fig. 3 shows the normalized diagonal terms of the matrix of NEP $`f_{ll}^m(q)`$ for $`(l,m)=(0,0),(1,0),(2,0),(2,1)`$. Not shown are the results for $`(l,m)=(1,1),(2,2)`$, since they do not exhibit very much structure. The corresponding simulation result is taken from . It was obtained by fitting a von Schweidler law plus corrections $`S_{ll^{}}^m(q,t)=F_{ll^{}}^m(q)H_{ll^{}}^m(q)(t/\tau _\alpha )^b+(H^{(2)})_{ll^{}}^m(q)(t/\tau _\alpha )^{2b}`$ to the simulation results for the time dependent density correlation function $`S_{ll^{}}^m(q,t)`$, where $`\tau _\alpha `$ is the $`\alpha `$ \- relaxation scale. There are three different theoretical curves. The two curves at temperatures $`T_c=0.7522,0.7521`$ are obtained with the fixed point method (on the glass side of the transition) and the quasi stability criterion (on the liquid side), respectively, as described above for upper cut off $`l_{co}=2`$. Their good agreement demonstrates the high accuracy of the solution. The third theoretical curve shows the result for upper cut off $`l_{co}=4`$ using the more accurate quasi stability criterion. Compared to the results in a clear improvement of the agreement with simulations can be observed. Especially the q - dependence of the functions are very well reproduced. Even a feature like the prepeak in $`f_{00}^0(q)`$ at $`q2.5`$ is reproduced as a shoulder in the corresponding theoretical result. This peak is not present in the static structure factors. Since $`S_{11}^0(q)`$ has a peak at about $`q2.5`$ it could appear due to a dynamic coupling of the ODOF, especially the one involving $`l=1`$ and the TDOF. Note also, that the mentioned peak exactly corresponds to the first peak in $`f_{11}^0(q)`$. There is a tendency, that the agreement is best around the wave vector, where the structure factor $`S_{00}(q)`$ has its first peak and is getting worse for large wave vectors. This might be interpreted as an indication for the glass transition being driven also for the investigated system of diatomic molecules by the TDOF. From investigations of other systems , we know that different scenarios are possible. Similar to increasing $`q`$, the agreement between simulation and theory gets worse with increasing $`l`$. This is expected due to two different reasons. First higher $`l`$ correspond to a higher angular resolution and are therefore probably much more affected by the mode coupling approximation. Second, higher $`l`$ are of course much more sensitive to the cut off $`l_{co}`$ than lower $`l`$. The curves for larger cut off increase the quality of the comparison with the MD results. But it is important to note that in our case no general rule can be given of how much the quality of the results for lower values of $`l`$ can be improved by increasing the upper cut off, as this was done in . In the case of a single dumbbell in a liquid of hard spheres the glass transition temperature is completely determined by the hard sphere liquid and does not change by increasing the cut off $`l_{co}`$. As explained above, in our case, $`T_c`$ can depend very sensitively on $`l_{co}`$. But this influences directly the amplitude of the NEP via the trivial effect of the temperature on the static structure factors. We already compensate as much as possible for this mechanism by presenting only the normalized NEP. But as in the case of hard spheres, there is still the effect, that also the normalized NEP are proportional to the static structure factor. This is a very nontrivial phenomenon, since the existence of negative vertices in the mode coupling functional $`𝓕(q)`$ could in principle lead to a violation of this correlation. But as can be inferred from Fig. 3 the NEP $`f_{00}^0(q),f_{22}^0(q)`$ for $`l_{co}=4`$ are systematically larger than the one for $`l_{co}=2`$ in a large region around the first peak of the structure factor $`S_{00}^0(q)`$ , without big differences in the functional form. This effect, especially for $`F_{00}^0(q)`$ where the mentioned trend is valid for all wavevectors, can be mainly understood as a consequence of the transition temperature being smaller for $`l_{co}=4`$ than for $`l_{co}=2`$, which causes the the first peak of $`S_{00}^0(q)`$ to increase. Additional evidence for this reasoning is presented in Fig. 4. In this figure we show the results for $`l_{co}=4`$ at temperature $`T=0.60`$, obtained with fixed point method, and $`T=0.63`$, obtained with the quasi stability criterion. The lower temperature is still in the glass (The required accuracy i.e. the maximum difference between consecutive iterates is smaller than $`10^6`$, is reached after 88 iterations). At the higher temperature the accuracy is first increasing as expected (see appendix). But after 24 iterations it begins slowly to decrease and after 42 iterations the iterates start to converge quickly towards the solution $`𝐅(q)=0`$. The results for the lower temperature agree, except for $`l=0`$, much better with the simulation, than the one for the higher temperature. But since the deviations to the true critical NEP at roughly the same number of iterations are of order $`\|TT_c\|`$ for the higher temperature compared to order $`\|\sqrt{\|TT_c\|}`$ for the lower one, we have to conclude, that the results for the higher temperature are closer to the critical NEP of the theory. The better agreement with simulations of the NEP at $`T=0.600`$ is a trivial consequence of the fact, that positive NEP increase with decreasing temperature. Due to this influence of the value of $`T_c`$ on the amplitude of the normalized critical NEP, we cannot in general conclude that increasing the cut off $`l_{co}`$ leads to a better agreement with the simulation. It might even happen, that the agreement with simulations gets worse instead of better, if increasing the cut off would lead to a larger transition temperature. This is possible due to the existence of negative vertices in the mode coupling functional $`𝓕`$. The observed trends allow the reasonable hypothesis, that the temperature effect could be the main source for the deviations between simulation and theory. In general the main structural features in the normalized non ergodicity parameter are very well represented, but they are systematically to small for nearly all wave vectors, exactly as expected, if the theoretical transition temperature is too large. Fig. 5 demonstrates, that the theory also gives good results for the off diagonal NEP. We found that $`S_{02}^0`$ and $`F_{02}^0`$ are the only important off diagonal components of the static structure factor matrix $`𝐒(q)`$ and NEP matrix $`𝐅(q)`$, respectively. In Fig. 5 we therefore show the normalized NEP $`f_{02}^0`$. The quality of the result is even better, than for the diagonal components of the NEP. ### C Critical amplitudes The critical amplitudes are determined only up to an overall scale factor. I.e. our theoretical results cannot be directly compared to the simulation results. But once we have chosen a scale factor for e.g. the amplitude $`h_{00}^0`$, all other amplitudes should be multiplied with the same scale factor to compare with the simulations. In fig 6(a) we have chosen a scale factor of 200 to obtain best agreement with the normalized critical amplitude $`h_{00}^0(q)=H_{00}^0(q)/S_{00}^0(q)`$. The features of this component are the same as in simple glass forming systems . There is a minimum at the position of the first peak of $`S_{00}^0(q)`$. Simulations and theory compare quite well for $`2q7`$ and show deviations at other wave vectors. With the chosen scale factors the other diagonal elements of the critical amplitude matrix show strong deviations from the simulation results. Especially $`h_{11}^0(q)`$ (see Fig.6(b)) does not even disagree in amplitude but also in the form, except for the minimum at $`q3`$. This is not unexpected, since the simulations have shown strong differences between the form of the dynamic correlators involving odd and even $`l`$ . Due to the weak top down anisotropy of our diatomic molecules, the dynamic of the correlators involving odd $`l`$ is only weakly coupled to the dynamics of the even components . $`180^0`$ jumps are still possible on a much faster time scale than the translational motion . Consequently there are strong corrections to the asymptotic results for the correlator $`S_{11}^0(q,t)`$ and the amplitude $`h_{11}^0(q)`$ is not very well defined. The deviations between simulation and theory in $`h_{22}^0`$ (see Fig. 6(c)) are more serious, since $`S_{22}^0(q,t)`$ exhibits a well defined $`\beta `$ \- relaxation regime. We can improve the agreement between simulation and theory by choosing a free scale factor for the simulation curves. The result is shown in Fig. 6(d) to demonstrate, in contrast to $`h_{11}^0`$, that essential structural features in $`h_{22}^0`$ are indeed reproduced by the theory. As argued above the dynamic correlators and therefore also the critical amplitudes involving $`l=2`$ are much more affected by the cut off $`l_{co}=2`$ as the one with lower $`l`$. This might be the reason for the rather large discrepancy found for $`h_{22}^0`$. To determine the critical amplitudes, it is necessary to be very close to the critical point. Restriction in computer time didn’t allow us to determine $`T_c`$ for $`l_{co}=4`$ with high enough accuracy to get reliable results for the critical amplitudes. The critical amplitudes with $`m>0`$ do not exhibit very much structure. In Fig. 7 we show for completeness the result for $`h_{22}^1(q)`$, which could in principle be measured in light scattering experiments. Again we choose an overall amplitude prefactor as a free fit parameter, but the agreement is still not very good. Much better is the agreement (similar to the NEP) for $`h_{02}^0`$ (see Fig. 8), although we still had to choose the amplitude scale of the simulation as a free fit parameter. ## IV Conclusions We have performed a quantitative test of MMCT for a liquid of diatomic molecules. The static structure factors from simulations were used as input for the MMCT to calculate the critical temperature $`T_c`$, the matrix of critical NEP $`F_{ll^{}}^m(q)`$ and the matrix of critical $`\beta `$ \- relaxation amplitudes $`H_{ll^{}}^m(q)`$. Several approximation schemes were used to test the sensitivity of the results against changing the degree of diagonalization of the $`F_{ll^{}}^m(q)`$ in $`l`$ and $`l^{}`$ and the dependence on the upper cut off $`l_{co}`$. As maximum cut off, $`l_{co}=4`$ was used. Since the computational effort increases strongly with $`l_{co}`$, we used a new, more accurate, method to determine the critical NEP from the liquid side of the transition. As expected, by giving up any diagonalization approximation for $`S_{ll^{}}^m(q)`$, $`F_{ll^{}}^m(q)`$ and the memory functional $`_{ll^{}}^{m\alpha \alpha ^{}}(q)`$ all degrees of freedom vitrify at a single temperature. In fact, to obtain a unique transition temperature, it is sufficient to keep only the static structure factor $`S_{ll^{}}^m(q)`$ nondiagonal in $`l`$ and $`l^{}`$. The strongest effect of successively applying diagonalization approximations for the NEP $`F_{ll^{}}^m(q)`$ and the memory functional $`_{ll^{}}^{m\alpha \alpha ^{}}(q)`$ is the change of the transition temperatures $`T_c`$. Also using different cut offs $`l_{co}`$ changes the transition temperature. Contrary to $`T_c`$, the overall form of NEP is much less sensitive to the different approximation schemes. The agreement with simulation is qualitatively quite good for cut off $`l_{co}`$=2 as well as for $`l_{co}`$=4. In some cases especially for the $`l=0`$, $`l^{}=0`$ component and the off diagonal $`l=0`$, $`l^{}=2`$ component of the NEP the agreement is even quantitative for wavevectors around the first peak of the structure factor $`S_{00}^0(q)`$. The comparison of the NEP with simulations for the correlator with ($`l=0`$, $`l^{}=0`$) is clearly improved using the cut off $`l_{co}=4`$ instead of $`l_{co}=2`$. Other correlators are better represented for $`l_{co}=2`$. We therefore cannot conclude that in general a further increase of the cut off will lead necessarily to still better agreement with simulations. Also the wave vector dependence of the normalized critical $`\beta `$-relaxation amplitudes agrees quite well with the one obtained from simulations. Contrary to the prediction of MMCT no common amplitude scale for the critical amplitudes with different $`l`$ and $`l^{}`$ could be found. At present it is not clear whether this failure is real or due to the fitting procedure used in to obtain the beta relaxation amplitudes. Also the restriction to cut off $`l_{co}=2`$ for the determination of the critical amplitudes could cause the found discrepancies, since also correlators with $`l>2`$ will contribute to the critical amplitudes at $`l=2`$. By neglecting them an error in the amplitude scale is possible. In summarizing our study, we may say, that the MMCT offers an overall consistent description of the glass transition in molecular liquid of diatomic molecules, at least concerning the critical NEP and the critical amplitudes. ## A Determination of the NEP using quasi stability If the temperature is chosen above the critical temperature, the only stable fixed point of the iteration (21) is $`𝐅=0`$. But there is still the possibility of determining the critical NEP with even higher accuracy than with the converging iteration for $`T<T_c`$. To implement the method, we have to initialize the iteration with a $`F_\mu ^{(0)}`$, which is close to $`F_\mu ^c`$, where $`\mu `$ is a superindex for $`q,l,l^{}`$ and $`m`$. This can be achieved by first using the fixed point iteration $$F_\mu ^{(n+1)}=𝒢_\mu (\{F_\mu ^{}^{(n)}\},ϵ).$$ (A1) below $`T_c`$ and increasing the temperatures in small steps until the liquid regime is reached. The critical temperature will eventually be missed by a small value $`0<ϵ1`$ and the iteration converges after a long transient quickly to zero. Since the iteration at a given temperature is always initialized with the final result of the previous temperature, the above assumption is fulfilled. The long transient is caused by the marginal stability of the critical point. The iteration will be dominated for a long time by the critical direction of the iteration, which is given by the eigenvector $`e^c`$, which belongs to the eigenvalue $`E_0=1`$ of the critical stability matrix $`(C_{\mu \mu ^{}}^c)=(𝒢_\mu /F_\mu ^{})^c`$ . The NEP $`\widehat{F}_\mu (ϵ)`$ at which the iteration is almost stationary is given by minimizing the distance of the projection of $`\mathrm{\Delta }C_\mu =𝒢_\mu (\{F_\mu ^{},ϵ\})F_\mu `$ on the critical direction $`\widehat{e}^c`$, where for technical reasons the left critical eigenvector $`\widehat{e}^c`$ was chosen. This condition can be written as $$\frac{_\mu \widehat{e}_\mu ^c\mathrm{\Delta }C_\mu }{F_\mu ^{}}=\underset{\mu }{}\widehat{e}_\mu ^cC_{\mu \mu ^{}}(\{\widehat{F}_{\mu ^{\prime \prime }}\},ϵ)\widehat{e}_\mu ^{}^c=0$$ (A2) Eq. (A2) determines the quasi stationary point $`\widehat{F}_\mu (ϵ)`$ such, that matrix $`(C_{\mu \mu ^{}}(\{\widehat{F}_{\mu ^{\prime \prime }},ϵ))`$ has an eigenvalue $`\widehat{E}_0(ϵ)=1`$ for $`ϵ<0`$. $`ϵ<0`$ indicates that the temperature is $`T=T_c(1ϵ)>T_c`$. This stationary point $`\widehat{F}_\mu `$ differs from $`F_\mu ^c`$ in order $`ϵ`$! This can be proven by using an argument analogous to deriving the $`\beta `$ \- relaxation equation . Defining $`\widehat{\delta }_\mu `$ by $`\widehat{F}_\mu =F_\mu ^c+\widehat{\delta }_\mu `$ and using the property that $`C_{\mu \mu ^{}}(\{\widehat{F}_{\mu ^{\prime \prime }}\},ϵ)`$ is a smooth differentiable function of $`ϵ`$, Eq. (A2) can be written, in leading order in $`\widehat{\delta }_\mu `$ and $`ϵ`$ as $$0=\underset{\mu }{}\widehat{e}_\mu ^c(\underset{\mu ^{\prime \prime }}{}C_{\mu \mu ^{}\mu ^{\prime \prime }}(𝐅^c,0)\widehat{\delta }_{\mu ^{\prime \prime }}+\frac{C_{\mu \mu ^{}}(𝐅^c,ϵ)}{ϵ}\|_{ϵ=0}ϵ)$$ (A3) here we have in addition used that $`\widehat{e}_\mu ^c`$ is the left eigenvector of the stability matrix at the critical point with eigenvalue $`E_0=1`$. The tensor $`C_{\mu \mu ^{}\mu ^{\prime \prime }}`$ is the partial derivative of the stability matrix with respect to $`F_{\mu ^{\prime \prime }}`$. Eq. (A3) immediately shows, that $`\widehat{\delta }_\mu `$ is (at least) of order $`ϵ)`$. Note, that the argument crucially depends on the well known phenomenon, that the static quantities like structure factors, which determine completely the stability matrix $`(C_{\mu \mu ^{}}(\{F_{\mu ^{\prime \prime }}\},ϵ)`$ are varying smoothly across the glass transition. We now want to estimate the number of iterations required to approach the quasi stationary point $`\widehat{𝐅}`$. For that matter the iterates $`F_\mu ^{(n)}`$ are expanded around the quasi stationary point $`F_\mu ^{(n)}=\widehat{F}_\mu +\delta _\mu ^{(n)}`$ Neglecting terms of order $`(\delta _\mu ^{(n)})^3`$ the iteration (A1) can be rewritten as $$\delta _\mu ^{(n+1)}=\mathrm{\Delta }\widehat{C}_\mu +\widehat{C}_{\mu \mu ^{}}\delta _\mu ^{}^{(n)}+\frac{1}{2}C_{\mu \mu ^{}\mu ^{\prime \prime }}\delta _\mu ^{}^{(n)}\delta _{\mu ^{\prime \prime }}^{(n)}$$ (A4) where the quantities with $`\widehat{}`$ are taken at $`(\{\widehat{F}_{\mu ^{\prime \prime }}\},ϵ)`$. The slowing down of the iteration is due to the component of $`\delta _\mu `$ along the critical direction $`e^c`$. This component can be extracted by writing $`\delta ^{(n)}=a^{(n)}e^c`$ and multiplying Eq. (A4) with the left eigenvector $`\widehat{e}^c`$ of the stability matrix $`C`$. The resulting equation for $`a^{(n)}`$ is then $$a^{(n+1)}a^{(n)}=\sigma +(\lambda 1)(a^{(n)})^2$$ (A5) Here we made use of (A2) and have normalized $`e^c,\widehat{e}^c`$ such that $`_\mu \widehat{e}_\mu ^ce_\mu ^c=1`$. The quantities $`\sigma =_\mu \widehat{e}_\mu ^c\mathrm{\Delta }C_\mu `$ and $`\lambda =1+_\mu \widehat{e}_\mu ^cC_{\mu \mu ^{}\mu ^{\prime \prime }}\widehat{e}_\mu ^{}^ce_{\mu ^{\prime \prime }}^c`$ can be identified as the separation parameter and the exponent parameter $`0<\lambda <1`$ of the $`\beta `$ \- relaxation theory, respectively . Since consecutive iterates are very close to each other in the vicinity of $`\widehat{𝐅}`$ we can rewrite (A5) as a differential equation. $$\dot{a}=(\|\sigma \|+(1\lambda )a^2)$$ (A6) Lets assume we start the iteration at $`t=0`$ with $`a(0)=a_0`$. Then the solution of (A6) is $$a(t)=\frac{\sqrt{\|\sigma \|}}{\sqrt{1\lambda }}\mathrm{tan}(\sqrt{(1\lambda )\|\sigma \|}t\mathrm{arctan}(a_0\frac{\sqrt{1\lambda }}{\sqrt{\|\sigma \|}}))$$ (A7) Since the iteration is initialized with a value for $`𝐅`$ slightly in the glass, the initial deviation $`a_0`$ from $`\widehat{𝐅}`$ is of order $`\sqrt{\|\sigma \|}=O(\sqrt{\|ϵ\|})`$. From Eq. (A7) then follows, that we need a number of iteration of the order $`1/\sqrt{\|ϵ\|}`$ to approach $`\widehat{𝐅}`$ (i.e. $`a(t)=0`$), and will stay close to $`\widehat{𝐅}`$ (i.e. $`a(t)O(\sqrt{\|\sigma \|}`$) for the same amount of iterations before the iterates decay to zero. This proves our statement that the quasi stationary solution of the iteration for $`T>T_c`$ is approached with the same amount of iterations as the stable solution slightly below $`T_c`$. But there the glass solution agrees with $`𝐅^c`$ only up to order $`\sqrt{ϵ}`$ the quasistationary solution $`\widehat{𝐅}`$ however agrees with $`𝐅^c`$ up to order $`\|ϵ\|\sqrt{\|ϵ\|}`$. Acknowledgment: It is a pleasure to thank W. Götze for a critical reading of the manuscript. We are grateful to the Sonderforschungsbereich 262 for financial support.
warning/0007/gr-qc0007003.html	ar5iv	text	# On the propagation of jump discontinuities in relativistic cosmology ## I Introduction It is well known that the mathematical nature of the evolution system within the relativistic gravitational field equations is essentially hyperbolic, with domains of dependence and influence determined by the speed of light. However, this nature is not obvious when the dynamical equations are written out in standard form, employing either a metric approach , a Hamiltonian (ADM) representation , a $`1+3`$ orthonormal frame (ONF) formulation , or a form obtained from a covariant $`(1+3)`$–decomposition . Consequently, over the years considerable effort has gone into determining ways of making this hyperbolic nature clear. Indeed, evolution systems of partial differential equations in first-order symmetric hyperbolic (FOSH) format have proven to be a main theme of research for at least the last five years. This activity was particularly motivated by the desire to carry over to the numerical investigation of relativistic effects such as generation of gravitational radiation associated with the in-spiralling of black hole and neutron star binaries the methods and expertise gained in areas of computational physics with a longer tradition like, e.g., hydrodynamics, where evolution systems of FOSH format are commonplace (see, e.g., the recent reviews and ). One of the most promising methods for obtaining a FOSH representation for the relativistic gravitational field equations is to use an extended $`1+3`$ orthonormal frame formulation that includes the once-contracted second Bianchi identities . In terms of the highest-order partial derivatives of the spacetime metric $`𝐠`$ implicitly occurring, the dynamical equations here rank at derivative level $`^3g`$. The set of geometrically defined field variables contains the components of the physically significant Weyl curvature tensor, and so we refer to this as a $`1+3`$ orthonormal frame curvature formulation of gravitational fields. It represents a completion of the physically and geometrically transparent $`1+3`$ covariant formalism for the cosmological case, in which all field variables are covariantly defined relative to a uniquely defined future-directed timelike reference congruence given by the unit 4-velocity field $`𝐮`$ of matter present . As Friedrich first showed, when the matter source for a spacetime geometry $`(,𝐠,𝐮)`$ is described phenomenologically as a perfect fluid, or one restricts to vacuum situations, upon introduction of a set of local coordinates and a specific choice to remove the gauge fixing freedom it is indeed possible to find linear combinations of the field variables and their dynamical equations that lead to an evolution system of (autonomous) partial differential equations in FOSH format, as desired . It is puzzling, then, that when this is done, as well as the expected sets of characteristic 3-surfaces associated with propagation speeds relative to $`𝐮`$ of $`\|v\|=0`$ (Coulomb-like gravitational and fluid rotational modes), $`\|v\|=c_s`$ (sound wave modes), and $`\|v\|=1`$ (transverse gravitational wave modes), the equations on the face of it indicate that the semi-longitudinal Weyl curvature characteristic eigenfields, identified long ago by Szekeres , appear related to an additional set of timelike characteristic 3-surfaces associated with propagation speed relative to $`𝐮`$ of one-half the speed of light, $`\|v\|=\frac{1}{2}`$. Friedrich assumed that these modes could be ignored, because, at least in the vacuum case, they had no invariant meaning . Two of us (HvE and GFRE), on the other hand, argued in a recent paper that in the fluid context that is relevant in relativistic cosmology, this is not the case, essentially because (i) the timelike reference congruence corresponding to $`𝐮`$ is invariantly defined by the presence of matter, and (ii) the hypothetical Weyl curvature modes cannot all be transformed to zero by using up the remaining gauge fixing freedom . Thus, these might possibly represent propagation modes that are physically meaningful. In private communications, Friedrich responded that the form of the $`^3g`$-order FOSH evolution systems discussed in Refs. and is not unique; one could obtain other reductions of the relativistic gravitational field equations that would lead to additional timelike characteristic 3-surfaces associated with propagation speeds different from $`\|v\|=\frac{1}{2}`$. And, hence, these modes cannot be physical (see also Ref. ). However, that argument is not very satisfactory. In our view, it compounds the problem, rather than solving it, by suggesting (if we take for face value the FOSH formalism and its interpretation) that still other (subluminal) propagation speeds could be associated with the transport of physical quantities, as well as $`\|v\|=\frac{1}{2}`$. The implication is that any single FOSH evolution system cannot be taken seriously without considering the set of all possible FOSH evolution systems, which is difficult to determine. What had not been fully taken into account in these discussions, even though Sec. III D 2 of Ref. made a brief reference to this aspect, is the physical role the constraint equations play in the selection of appropriate classes of initial data sets. Specifically, what kind of physically relevant jump discontinuities they do allow to be present in the derivatives of these data sets . Jump discontinuities in the derivatives of the initial data are of major dynamical importance as they can be thought of as representing the triggering or terminating events of physical generation processes.<sup></sup><sup></sup>Jürgen Ehlers (private communication). In the present paper we show that an analysis of the constraint equations leads to a resolution of the problem: while the additional timelike characteristic 3-surfaces do indeed occur in the $`^3g`$-order FOSH evolution system, the constraint equations are of such a form as to prevent the corresponding semi-longitudinal Weyl curvature modes being activated. Because jump discontinuities cannot exist in the values of the constraint equations themselves (they have to be zero everywhere), jump discontinuities in the derivatives of the initial data can only be propagated along the characteristic 3-surfaces associated with propagation speeds relative to $`𝐮`$ of $`\|v\|=0`$, $`\|v\|=c_s`$, and $`\|v\|=1`$, but not along those associated with $`\|v\|=\frac{1}{2}`$. Hence, one cannot physically send arbitrary information along the latter set of timelike characteristic 3-surfaces. The problem of additional characteristic 3-surfaces and how the constraint equations select appropriate initial data sets can be highlighted in a nice and transparent fashion by a suggestive example that derives from a set of linearised relativistic gravitational field equations given in a paper by Kind, Ehlers and Schmidt . We will briefly discuss this example in Appendix 1. A possible complementary explanation to the above viewpoint arises from the following chain of arguments: (i) For vacuum as well as perfect fluid spacetime geometries, there exist $`^2g`$-order dynamical formulations of the relativistic gravitational field equations (with harmonic coordinate gauge fixing) in which characteristic 3-surfaces only correspond to the local light cones, the local sound cones, or are generated by the fluid flow lines. (See, e.g., Refs. and .) (ii) There are dynamical formulations of order $`^2g`$ or $`^3g`$ that have additional (timelike or spacelike) characteristic 3-surfaces. Along these characteristic 3-surfaces the evolution system will propagate jump discontinuities if we put them into the related initial data (ignoring the constraint equations). (iii) Different dynamical formulations with different sets of characteristic 3-surfaces all describe for identical initial data the same spacetime geometry. (iv) Because of the existence of a dynamical formulation according to (i), jump discontinuities as described in (ii) cannot really exists (and so be physically meaningful). Hence, the only logical possibility is that in (ii) the constraint equations prevent the existence of such jump discontinuities. The present paper shows that the semi-longitudinal Weyl curvature modes that potentially exist cannot in fact carry “gravitational news”. Of course, this is expected, so this analysis simply confirms what everyone has believed all along. However, we believe it is nevertheless useful, both in terms of showing important relations between the constraint equations and the set of characteristic 3-surfaces that can occur in general FOSH evolution systems (and have not been clearly demonstrated in the literature), and because this does indeed resolve the issue at hand in the specific important case of relativistic gravitation, where the variables presented here, and hence the associated dynamical equations, have much to recommend them. Without the analysis given in this paper, the apparent ($`\|v\|=\frac{1}{2}`$)–modes remain an annoying and unresolved problem. The outline of the paper is as follows. In Sec. II we briefly discuss the mathematical concepts relevant to our analysis: we review the central ideas behind FOSH evolution systems, we address the issue of gauge fixing freedom in a $`1+3`$ orthonormal frame formulation of the relativistic gravitational field equations, we list the set of constraint equations we need to consider, and we introduce a $`(1+1+2)`$–decomposition of all geometrically defined field variables. Then, in Sec. III, we address from a purely local viewpoint the question of how the constraint equations select appropriate initial data sets for the $`^3g`$-order FOSH evolution system introduced in Ref. . To the best of our knowledge, we present in Sec. IV for the first time a fully gauge-fixed autonomous FOSH evolution system for a special class of spatially inhomogeneous perfect fluid cosmological models on the basis of a $`1+3`$ orthonormal frame formulation that is only of order $`^2g`$ in the degrees of freedom of the gravitational field.The equations of motion for the matter sources will necessarily derive from the twice-contracted second Bianchi identities and so are typically of order $`^3g`$. We refer to this as a $`1+3`$ orthonormal frame connection formulation of gravitational fields. Additionally, we derive for these cosmological models, for which there exists an Abelian $`G_2`$ isometry group, the transport equations that describe how physically relevant jump discontinuities in the derivatives of the initial data are propagated along the bicharacteristic rays of the setting. Our conclusions are contained in Sec. V. Finally, besides the above-mentioned suggestive example in Appendix 1, we give in Appendix 2 in explicit form those linear combinations of the components of the extended $`1+3`$ orthonormal frame constraint equations that prove suitable for our analysis. We will use the same conventions, units and notations as introduced in Appendix A 1 of Ref. . ## II Mathematical preliminaries ### A First-order symmetric hyperbolic evolution systems We consider evolution systems for a collection of $`k`$ real-valued field variables $`u^A=u^A(x^\mu )`$ that are composed of a set of $`k`$ quasi-linear partial differential equations of first order given by $$M^{AB\mu }(x^\nu ,u^C)_\mu u_B=N^A(x^\nu ,u^C),A,B,C=1,\mathrm{},k;$$ (1) the field variables $`u^A`$ are functions of a set of local spacetime coordinates $`\{x^\mu \}`$. Evolution systems of this form are called symmetric if the real-valued $`k\times k`$ coefficient matrices entering the principle part satisfy $`M^{AB\mu }=M^{(AB)\mu }`$; moreover, they are called hyperbolic if the contraction $`M^{AB\mu }n_\mu `$ with the coordinate components of an arbitrary past-directed timelike 1-form $`n_a`$ yields a positive-definite matrix. We remark that (i) cases with $`M^{AB\mu }=M^{AB\mu }(x^\nu )`$ are referred to as semi-linear, and (ii) cases with $`M^{AB\mu }=M^{AB\mu }(u^C)`$ and $`N^A=N^A(u^C)`$ are referred to as autonomous. In general, it proves convenient to consider a $`(1+3)`$–decomposition of Eq. (1) in the format $$M^{AB\mathrm{\hspace{0.17em}0}}(t,x^j,u^C)_tu_B+M^{ABi}(t,x^j,u^C)_iu_B=N^A(t,x^j,u^C).$$ The concept of FOSH evolution systems was first introduced by Friedrichs ; a standard reference, broughtly discussing its usefulness to applications in mathematical physics and also presenting proofs for existence and uniqueness of solutions, is the book by Courant and Hilbert . The characteristic condition $$0=Q:=det[M^{AB\mu }_\mu \varphi ]$$ (2) determines the coordinate components of the past-directed normals $`_a\varphi `$ of the set of characteristic 3-surfaces $`𝒞`$:$`\{\varphi (x^\mu )=\text{const}\}`$ associated with the FOSH evolution system (1). With $`M^{AB\mu }=M^{(AB)\mu }`$, hyperbolicity of Eq. (1) thus also corresponds to all individual roots (“eigenvalues”) $`v`$ of Eq. (2) being real-valued. Every individual $`v`$ then defines a pair of so-called left and right nullifying vectors, $`l^A`$ and $`r^A`$, by $$0=:l_A(M^{AB\mu }_\mu \varphi ),0=:(M^{AB\mu }_\mu \varphi )r_B;$$ (3) the linearly independent sets $`\{l^A\}`$ or $`\{r^A\}`$ form a basis of the $`k`$-dimensional space of field variables $`u^A`$. According to the theory discussed in Ch. VI.4.2 of Courant and Hilbert , FOSH evolution systems of the format (1) have the power to describe the physical transport along so-called bicharacteristic rays of jump discontinuities that exist in the outward first derivatives across a characteristic 3-surface $`𝒞`$:$`\{\varphi (x^\mu )=\text{const}\}`$ of the field variables $`u^A`$; the tangential first derivatives of the $`u^A`$ as well as the $`u^A`$ themselves are assumed to be continuous across $`𝒞`$:$`\{\varphi (x^\mu )=\text{const}\}`$. As is standard, we will use the notation $$[f]:=\underset{\varphi c_+}{lim}f\underset{\varphi c_{}}{lim}f=f_+f_{}$$ to symbolise a jump discontinuity (of finite magnitude) across $`𝒞`$:$`\{\varphi (x^\mu )=\text{const}\}`$ in the value of a given variable $`f`$. Under the stated assumptions, it follows from Eqs. (1) and (3) that $$0=(M^{AB\mu }_\mu \varphi )[_\varphi u_B][_\varphi u^A]=[_\varphi u]r^A,$$ (4) i.e., the jump discontinuity $`[_\varphi u^A]`$ must be proportional to a right nullifying vector $`r^A`$. The real-valued scalar of proportionality, denoted by $`[_\varphi u]`$, is assumed to have continuous first derivatives. Then, according to Chs. VI.4.2 and VI.4.9 of Ref. , for linear, semi-linear and quasi-linear FOSH evolution systems (1) the transport equation for $`[_\varphi u]`$ along bicharacteristic rays within the characteristic 3-surfaces $`𝒞`$:$`\{\varphi (x^\mu )=\text{const}\}`$ takes the effective form $$0=(l_AM^{AB\mu }r_B)_\mu [_\varphi u]+\left((l_AM^{AB\mu })_\mu r_B(l_AN^{AB}r_B)\right)[_\varphi u].$$ (5) Note, in particular, the involutive character of this relation; if $`[_\varphi u]`$ is non-zero at one point along a bicharacteristic ray, it will be non-zero everywhere along this ray, and vice versa. Note also that the present treatment of jump discontinuities breaks down when the $`𝒞`$:$`\{\varphi (x^\mu )=\text{const}\}`$ within a given family start to intersect and so prompt the formation of so-called “shocks”. Shock formation, however, cannot arise when the principal part of Eq. (1) is semi-linear. It is a special feature of the relativistic gravitational field equations that related FOSH evolution systems do have, in the branches that evolve the degrees of freedom in the gravitational field itself, principal parts which are effectively semi-linear.See, e.g., Eqs. (3.21) – (3.23) in Ref. , or Eqs. (94) – (97) in the Abelian $`G_2`$ example given in Sec. IV below. More precisely, the coefficient matrices $`M^{AB\mu }`$ in these branches depend only on those field variables $`u^A`$ which form a background for gravitational dynamics in the sense of Geroch . ### B Choice of gauge source functions and local coordinates As Friedrich emphasised in Sec. 5.2 of Ref. , there exists within a $`1+3`$ orthonormal frame representation of the relativistic gravitational field equations a set of ten so-called gauge source functions, $`G:=\{T^0,T^\alpha ,T^\alpha {}_{0}{}^{},T^\alpha {}_{\beta }{}^{}\}`$, that can be arbitrarily prescribed in any dynamical consideration (and are thus assumed to be “known”). These relate to (i) the arbitrary choice of a future-directed reference “time flow vector field” $`𝐓`$,<sup>§</sup><sup>§</sup>§The reference vector field $`𝐓`$ need not necessarily be timelike. which, in terms of the $`1+3`$ ONF basis $`\{𝐞_0,𝐞_\alpha \}`$, is expressed by $$𝐓:=T^0𝐞_0+T^\alpha 𝐞_\alpha ,T^0>0,$$ (6) and (ii) the propagation of the $`1+3`$ ONF basis $`\{𝐞_0,𝐞_\alpha \}`$ along $`𝐓`$, described by $$_𝐓𝐞_0:=T^\alpha {}_{0}{}^{}𝐞_{\alpha }^{},_𝐓𝐞_\alpha :=T^0{}_{\alpha }{}^{}𝐞_{0}^{}+T^\beta {}_{\alpha }{}^{}𝐞_{\beta }^{}.$$ (7) Parallel transport of $`\{𝐞_0,𝐞_\alpha \}`$ along $`𝐓`$, for example, thus corresponds to setting $`0=T^0{}_{\alpha }{}^{}=T^\alpha _\beta `$. Upon introduction of a dimensionless local time coordinate $`t`$ and dimensionless local spatial coordinates $`\{x^i\}`$ that comove with $`𝐓`$, the gauge conditions related to a $`1+3`$ orthonormal frame representation are made explicit by $$e_0{}_{}{}^{\mu }=\frac{1}{T^0}(M_0^1\delta ^\mu {}_{0}{}^{}T^\alpha e_\alpha {}_{}{}^{\mu }),\mathrm{\Gamma }^0{}_{\alpha 0}{}^{}=\frac{1}{T^0}(T^0{}_{\alpha }{}^{}\mathrm{\Gamma }^0{}_{\alpha \beta }{}^{}T_{}^{\beta }),\mathrm{\Gamma }^\alpha {}_{\beta 0}{}^{}=\frac{1}{T^0}(T^\alpha {}_{\beta }{}^{}\mathrm{\Gamma }^\alpha {}_{\beta \gamma }{}^{}T_{}^{\gamma });$$ (8) we keep the inverse unit of $`[\text{length}]`$, $`M_0^1`$, as a coefficient for reasons of physical dimensions. The fluid-comoving, Lagrangean perspective adopted in the discussion of Ref. by identifying the timelike reference congruence with the fluid 4-velocity field, $`𝐞_0𝐮`$, is now obtained by fixing three of the four dimensionless coordinate gauge source functions according to $`T^\alpha =0`$, resulting in an alignment $`𝐓𝐞_0`$ ($`𝐮`$). This leads to $$e_0{}_{}{}^{\mu }=M^1\delta ^\mu {}_{0}{}^{},T^0{}_{\alpha }{}^{}=T^0\mathrm{\Gamma }^0{}_{\alpha 0}{}^{}=T^0\dot{u}_\alpha ,T^\alpha {}_{\beta }{}^{}=T^0\mathrm{\Gamma }^\alpha {}_{\beta 0}{}^{}=T^0ϵ^\alpha {}_{\beta \gamma }{}^{}\mathrm{\Omega }_{}^{\gamma },$$ (9) where $`M:=T^0M_0`$. Consequently, the three frame gauge source functions $`T^0_\alpha `$ become proportional to the components of the fluid acceleration $`\dot{u}^\alpha `$, while the three frame gauge source functions $`T^\alpha _\beta `$ become proportional to the components of the rotation rate $`\mathrm{\Omega }^\alpha `$ at which the spatial frame $`\{𝐞_\alpha \}`$ fails to be Fermi-propagated along $`𝐮`$. For cosmological models $`(,𝐠,𝐮)`$ with perfect fluid matter sources, Ref. introduced proper time along $`𝐮`$ by setting $`T^0=1M=M_0`$, and derived the evolution equation for $`\dot{u}^\alpha `$ along $`𝐮`$ from the commutators on the basis of Eqs. (12) and (30) below. Additionally, Ref. set $`\mathrm{\Omega }^\alpha =0`$. The latter choice, however, is by no means compulsory, and other choices may prove equally convenient (given that $`\mathrm{\Omega }^\alpha `$ is assumed to be “known”). In order to obtain from the extended $`1+3`$ orthonormal frame relations proper partial differential equations such that the theory underlying Subsec. II A applied, Ref. expressed the coordinate components $`e_0{}_{}{}^{\mu }:=𝐞_0(x^\mu )`$ and $`e_\alpha {}_{}{}^{\mu }:=𝐞_\alpha (x^\mu )`$ of the $`1+3`$ ONF basis $`\{𝐞_0,𝐞_\alpha \}`$ in terms of the comoving local coordinate basis $`\{_t,_i\}`$ by $$𝐞_0:=M^1_t,𝐞_\alpha :=e_\alpha {}_{}{}^{i}(M_i_t+_i);$$ (10) $`M=M(t,x^i)`$ denotes the threading lapse function and $`M_idx^i=M_i(t,x^j)dx^i`$ the dimensionless threading shift 1-form. The inverse of the threading metric is $`h^{ij}:=\delta ^{\alpha \beta }e_\alpha {}_{}{}^{i}e_{\beta }^{}^j`$. See, e.g., Ref. and references therein. ### C Matter model The matter sources in the present discussion (as well as in Ref. ) are assumed to be modelled phenomenologically as a perfect fluid such that, with respect to fluid-comoving observers, $$0=q^\alpha (𝐮)=\pi _{\alpha \beta }(𝐮),$$ (11) i.e., the energy current density and the anisotropic pressure both vanish. Additionally, a baryotropic equation of state is assumed, $$p=p(\mu ),$$ (12) relating the isotropic pressure $`p(𝐮)`$ to the total energy density $`\mu (𝐮)`$. ### D Constraint equations The following relations in the set obtained from an extended $`1+3`$ orthonormal frame representation of the relativistic gravitational field equations do not contain any frame derivatives with respect to $`𝐞_0`$. Hence, it is commonplace to refer to these relations as “constraint equations”.Even though this terminology is problematic in the generic case when $`𝐞_0𝐮`$ has non-zero vorticity, $`\omega ^\alpha (𝐮)0`$, and local coordinates are introduced according to Eq. (10) above. These are $`0`$ $`=`$ $`(C_1)^\alpha :=(𝐞_\beta 3a_\beta )(\sigma ^{\alpha \beta })\frac{2}{3}\delta ^{\alpha \beta }𝐞_\beta (\mathrm{\Theta })n_\beta ^\alpha \omega ^\beta +ϵ^{\alpha \beta \gamma }[(𝐞_\beta +2\dot{u}_\beta a_\beta )(\omega _\gamma )n_{\beta \delta }\sigma _\gamma ^\delta ]`$ (13) $`0`$ $`=`$ $`(C_2):=(𝐞_\alpha \dot{u}_\alpha 2a_\alpha )(\omega ^\alpha )`$ (15) $`0`$ $`=`$ $`(C_3)^{\alpha \beta }:=(\delta ^{\gamma \alpha }𝐞_\gamma +2\dot{u}^\alpha +a^\alpha )(\omega ^\beta )\frac{1}{2}n_\gamma ^\gamma \sigma ^{\alpha \beta }+3n_\gamma ^\alpha \sigma ^{\beta \gamma }+H^{\alpha \beta }`$ (18) $`ϵ^{\gamma \delta \alpha }[(𝐞_\gamma a_\gamma )(\sigma _\delta ^\beta )+n_\gamma ^\beta \omega _\delta ]`$ $`0`$ $`=`$ $`(C_\mathrm{J})^\alpha :=(𝐞_\beta 2a_\beta )(n^{\alpha \beta })+\frac{2}{3}\mathrm{\Theta }\omega ^\alpha +2\sigma _\beta ^\alpha \omega ^\beta +ϵ^{\alpha \beta \gamma }[𝐞_\beta (a_\gamma )2\omega _\beta \mathrm{\Omega }_\gamma ]`$ (20) $`0`$ $`=`$ $`(C_\mathrm{G})^{\alpha \beta }:={}_{}{}^{}S_{}^{\alpha \beta }+\frac{1}{3}\mathrm{\Theta }\sigma ^{\alpha \beta }\sigma _\gamma ^\alpha \sigma ^{\beta \gamma }\omega ^\alpha \omega ^\beta +2\omega ^\alpha \mathrm{\Omega }^\beta E^{\alpha \beta }`$ (22) $`0`$ $`=`$ $`(C_\mathrm{G}):={}_{}{}^{}R+\frac{2}{3}\mathrm{\Theta }^2(\sigma _{\alpha \beta }\sigma ^{\alpha \beta })+2(\omega _\alpha \omega ^\alpha )4(\omega _\alpha \mathrm{\Omega }^\alpha )2\mu 2\mathrm{\Lambda }`$ (24) $`0`$ $`=`$ $`(C_4)^\alpha :=(𝐞_\beta 3a_\beta )(E^{\alpha \beta })\frac{1}{3}\delta ^{\alpha \beta }𝐞_\beta (\mu )3\omega _\beta H^{\alpha \beta }ϵ^{\alpha \beta \gamma }[\sigma _{\beta \delta }H_\gamma ^\delta +n_{\beta \delta }E_\gamma ^\delta ]`$ (26) $`0`$ $`=`$ $`(C_5)^\alpha :=(𝐞_\beta 3a_\beta )(H^{\alpha \beta })+(\mu +p)\omega ^\alpha +3\omega _\beta E^{\alpha \beta }+ϵ^{\alpha \beta \gamma }[\sigma _{\beta \delta }E_\gamma ^\delta n_{\beta \delta }H_\gamma ^\delta ]`$ (28) $`0`$ $`=`$ $`(C_{\mathrm{PF}})^\alpha :=c_s^2\delta ^{\alpha \beta }𝐞_\beta (\mu )+(\mu +p)\dot{u}^\alpha ,`$ (30) where $`{}_{}{}^{}S_{\alpha \beta }^{}`$ $`:=`$ $`𝐞_\alpha (a_\beta )+b_{\alpha \beta }ϵ^{\gamma \delta }{}_{\alpha }{}^{}(𝐞_{\|\gamma \|}2a_{\|\gamma \|})(n_{\beta \delta })`$ (31) $`{}_{}{}^{}R`$ $`:=`$ $`2(2𝐞_\alpha 3a_\alpha )(a^\alpha )\frac{1}{2}b_\alpha ^\alpha `$ (33) $`b_{\alpha \beta }`$ $`:=`$ $`2n_{\alpha \gamma }n_\beta ^\gamma n_\gamma ^\gamma n_{\alpha \beta },`$ (35) $`c_s^2(\mu ):=dp(\mu )/d\mu `$ defines the isentropic speed of sound with $`0c_s^21`$, and angle brackets denote the symmetric tracefree part.The numbering of the constraint equations we employ is based on the conventions established in $`1+3`$ covariant treatments of relativistic cosmological models $`(,𝐠,𝐮)`$ (cf. Ref. ). While Eqs. (13) – (24) derive from the Ricci identities, the Jacobi identities and the Einstein field equations and so are of order $`^2g`$, Eqs. (26) – (30) derive from the once- and twice-contracted second Bianchi identities and so are of order $`^3g`$. The divergence constraint equation (13) for the fluid rate of shear is often referred to as the “momentum constraint”. When $`0=\omega ^\alpha (𝐮)`$, such that the fluid 4-velocity field $`𝐮`$ constitutes the normals to a family of spacelike 3-surfaces $`𝒮`$:$`\{t=\text{const}\}`$, Eqs. (22) and (24) correspond to the symmetric tracefree and trace parts of the once-contracted Gauß embedding equation. In this case, one also speaks of $`(C_G)`$ as the generalised Friedmann equation, alias the “Hamiltonian constraint” or the “energy constraint”. ### E $`(1+1+2)`$–decomposition In the present discussion it proves very helpful to consider a $`(1+1+2)`$–decomposition of all geometrically defined field variables and their dynamical relations. In order to do so, we arbitrarily pick the frame basis field $`𝐞_1`$ as a second, spacelike, reference direction, in addition to $`𝐞_0𝐮`$ as a timelike one; any other spatial direction, however, would be equally acceptable. Hence, in a small isotropic neighbourhood $`𝒰`$ in the local rest 3-space of an arbitrary event $`𝒫`$, we establish the convention of regarding those spatial frame components of geometrical objects which contain the index “$`1`$” as (semi-)longitudinal with respect to $`𝐞_1`$, while regarding those which exclude the index “$`1`$” as transverse with respect to $`𝐞_1`$. Likewise, in $`𝒰`$, $`𝐞_1`$ shall constitute the outward frame derivative while $`𝐞_2`$ and $`𝐞_3`$ shall be tangential frame derivatives when we turn in Sec. III to the issue of jump discontinuities across spherical spacelike 2-surfaces $`𝒥`$:$`\{t=\text{const},\varphi (x^\mu )=\text{const}\}`$. For the frame components of spatial rank-$`2`$ symmetric tracefree tensors $`a_{\alpha \beta }=a_{\alpha \beta }`$ with squared magnitude $`a^2:=\frac{1}{2}(a_{\alpha \beta }a^{\alpha \beta })0`$, we define a new set of frame variables by $`\begin{array}{ccccc}a_+:=\frac{1}{2}(a_{22}+a_{33})=\frac{1}{2}a_{11}\hfill & & a_{}:=\frac{1}{2\sqrt{3}}(a_{22}a_{33})\hfill & & \\ & & & & \\ a_\times :=\frac{1}{\sqrt{3}}a_{23}\hfill & & a_2:=\frac{1}{\sqrt{3}}a_{31}\hfill & & a_3:=\frac{1}{\sqrt{3}}a_{12},\hfill \end{array}`$ (39) so that $$a^2=3(a_+^2+a_{}^2+a_\times ^2+a_2^2+a_3^2).$$ (40) In particular, in the present discussion we have $`a_{\alpha \beta }\{\sigma _{\alpha \beta },E_{\alpha \beta },H_{\alpha \beta }\}`$. We remark that these definitions are now adapted to the conventions of the book edited by Wainwright and Ellis , implying they differ by a factor of $`\frac{1}{3}`$ from those used in Ref. . In analogy to Eq. (39), we perform a $`(1+1+2)`$–decomposition of the spatial commutation functions $`n_{\alpha \beta }`$ by defining $`\begin{array}{ccccc}n:=n_{11}+n_{22}+n_{33}\hfill & & n_+:=n_{11}+\frac{1}{2}(n_{22}+n_{33})\hfill & & n_{}:=\frac{1}{2\sqrt{3}}(n_{22}n_{33})\hfill \\ & & & & \\ n_\times :=\frac{1}{\sqrt{3}}n_{23}\hfill & & n_2:=\frac{1}{\sqrt{3}}n_{31}\hfill & & n_3:=\frac{1}{\sqrt{3}}n_{12}.\hfill \end{array}`$ (44) The squared magnitude is then given by $$\frac{1}{2}(n_{\alpha \beta }n^{\alpha \beta })=\frac{1}{6}(n^2+2n_+^2)+3(n_{}^2+n_\times ^2+n_2^2+n_3^2).$$ (45) Note that only $`(n2n_+)`$, $`n_{}`$ and $`n_\times `$ transform as tensor components under rotations of the spatial frame $`\{𝐞_\alpha \}`$ about the reference $`𝐞_1`$-axis. By employing these conventions and definitions, we have listed in Appendix 2 certain linear combinations of the components of the constraint equations (13) – (30) that will be needed in the following sections. \********************************************************** INSERT Table I on “Conventions for $`(1+1+2)`$–decomposition” HERE. \********************************************************* ## III Physical effect of constraint equations on outward first derivatives Generic cosmological models $`(,𝐠,𝐮)`$ with a perfect fluid matter source have fluid 4-velocity fields $`𝐮`$ with non-zero vorticity, $`\omega ^\alpha (𝐮)0`$. This property makes it impossible to determine a fluid-comoving spacelike 3-surface $`𝒮`$:$`\{t=\text{const}\}`$ everywhere orthogonal to $`𝐮`$ on which initial data satisfying the constraint equations of Subsec. II D could be specified. In this case, the discussion of a well-posed Cauchy initial value problem requires that the setting of the data as well as the solution of the constraint equations be instead performed on a non-comoving spacelike 3-surface $`𝒮`$:$`\{t=\text{const}\}`$, before the data is evolved along $`𝐮`$.<sup>\</sup><sup>\</sup>\It is currently unknown whether the $`^3g`$-order FOSH evolution systems with perfect fluid matter sources in Refs. and can be generalised to a non-comoving perspective. All these complications disappear when $`0=\omega ^\alpha (𝐮)`$, and well-defined spacelike 3-surface $`𝒮`$:$`\{t=\text{const}\}`$ everywhere orthogonal to $`𝐮`$ do exist. For our purposes, however, it is sufficient to investigate the physical effect of the constraint equations on the selection of appropriate initial data sets from a purely local viewpoint, i.e., only in a small isotropic neighbourhood $`𝒰`$ in the local rest 3-space of an arbitrary event $`𝒫`$. Due to the local Minkowskian structure of all relativistic spacetime manifolds $`(,𝐠)`$, one conventionally determines (and analyses their physical properties) the set of characteristic cones $`𝒞`$:$`\{\varphi (x^\mu )=\text{const}\}`$ for a given FOSH evolution system only within the small isotropic spacetime neighbourhood $`\{\epsilon t\epsilon \}\times 𝒰`$ of $`𝒫`$. In line with this, we will consider in the following spherical spacelike 2-surfaces $`𝒥`$:$`\{t=\text{const},\varphi (x^\mu )=\text{const}\}`$ in $`𝒰`$ across which we assume to exist (i) jump discontinuities in the outward first frame derivatives of certain geometrically defined field variables $`u^A`$, and (ii) continuity of the tangential first frame derivatives of the $`u^A`$ and the $`u^A`$ themselves. As the constraint equations have to be satisfied everywhere, it is clear that across $`𝒥`$ we have $`0=[(C)^\alpha \mathrm{}]_{t,\varphi =\mathrm{const}}`$ for the value of any component in the set of Eqs. (13) – (30). Motivated by the prospect of grounding the discussion on wave-like phenomena described by the relativistic gravitational field equations on the deviation equation for a set of test particles, the dynamical considerations on the $`^3g`$-order FOSH evolution system in Ref. focused on the set of Weyl curvature characteristic eigenfields $`\{E_+,H_+,(E_3H_2),(E_2\pm H_3),(E_{}H_\times ),(E_\times \pm H_{})\}`$. There it was correctly argued that the deviation equation monitors the physical effects on the state of motion of a set of test particles of both gradual as well as sudden changes in the values of these fields (see, e.g., Refs. and ). What was overlooked in this work, however, is the fact that, on the basis of the theory underlying Subsec. II A, the $`^3g`$-order FOSH evolution system presented in Ref. can at best describe the physical transport along bicharacteristic rays of jump discontinuities in the (outward) first derivatives of these fields rather than these fields themselves, given the constraint equations do not impose any additional restrictions. It is our aim to supplement the discussion of Ref. by such a consideration in the present section. ### A Jump discontinuities at derivative level $`^3g`$ Considering in $`𝒰`$ a spherical spacelike 2-surface $`𝒥`$:$`\{t=\text{const},\varphi (x^\mu )=\text{const}\}`$, and assuming that across $`𝒥`$ all geometrically defined field variables $`u^A`$ as well as their tangential first frame derivatives are continuous, we find that the Weyl curvature divergence equations (156) – (170) amongst the constraint equations lead to the following set of jump conditions: $`\left[𝐞_1(E_+)\right]_{t,\varphi =\mathrm{const}}`$ $`=`$ $`\frac{1}{6}\left[𝐞_1(\mu )\right]_{t,\varphi =\mathrm{const}}`$ (46) $`\left[𝐞_1(H_+)\right]_{t,\varphi =\mathrm{const}}`$ $`=`$ $`0`$ (47) $`\left[𝐞_1(E_3H_2)\right]_{t,\varphi =\mathrm{const}}`$ $`=`$ $`0`$ (48) $`\left[𝐞_1(E_2\pm H_3)\right]_{t,\varphi =\mathrm{const}}`$ $`=`$ $`0`$ (49) $`\left[𝐞_1(E_{}H_\times )\right]_{t,\varphi =\mathrm{const}}`$ $`=`$ unconstrained (50) $`\left[𝐞_1(E_\times \pm H_{})\right]_{t,\varphi =\mathrm{const}}`$ $`=`$ $`\text{unconstrained}.`$ (51) The implications are three-fold: (i) Jump discontinuities in the outward first frame derivative of the Coulomb-like Weyl curvature characteristic eigenfield $`E_+`$ originate from jump discontinuities in the outward first frame derivative of the matter total energy density $`\mu `$ and are physically allowed, if the momentum conservation equation (171) is satisfied on both sides of $`𝒥`$ with $`0=[(C_{\mathrm{PF}})_1]_{t,\varphi =\mathrm{const}}`$. In more detail: assuming that all of $`p`$, $`c_s`$ and $`\dot{u}_1`$ are continuous across $`𝒥`$, Eq. (171) yields $$0=[(C_{\mathrm{PF}})_1]_{t,\varphi =\mathrm{const}}=c_s^2\left[𝐞_1(\mu )\right]_{t,\varphi =\mathrm{const}}+\left[\mu \right]_{t,\varphi =\mathrm{const}}\dot{u}_1.$$ (52) Hence, if $`p=0c_s=0`$ and $`\dot{u}_1=0`$ on both sides of $`𝒥`$, the values of $`[𝐞_1(\mu )]_{t,\varphi =\mathrm{const}}`$ and $`[\mu ]_{t,\varphi =\mathrm{const}}`$ remain unconstrained. Phenomena of the presently described kind occur, for example, across the surfaces of static, spherically symmetrical perfect fluid stars with equation of state (12) (see, e.g., Ref. ). In a $`^3g`$-order formulation, it follows that real-valued initial data for $`\mu `$ (and so for $`E_+`$) is required to be of differentiability class $`C^2(𝒰)`$ with respect to the zeroth-order derivative level of $`𝐠`$. The equations of Ref. show that $`[𝐞_1(E_+)]_{t,\varphi =\mathrm{const}}0`$ propagates with characteristic velocity $`v=0`$ relative to $`𝐮`$. Note that in the vacuum subcase $`\mu =0[𝐞_1(E_+)]_{t,\varphi =\mathrm{const}}=0`$. Jump discontinuities in $`𝐞_1(H_+)`$ are not physically allowed, and so, in a $`^3g`$-order formulation, real-valued initial data for $`H_+`$ is required to be of differentiability class $`C^3(𝒰)`$ with respect to the zeroth-order derivative level of $`𝐠`$. (ii) Jump discontinuities in the outward first frame derivatives of the semi-longitudinal Weyl curvature characteristic eigenfields $`(E_3H_2)`$ and $`(E_2\pm H_3)`$, that by restricting to the net $`^3g`$-order FOSH evolution system in Ref. (without accounting for the constraint equations) are theoretically associated with characteristic velocities $`v=\pm \frac{1}{2}`$ relative to $`𝐮`$, are not physically allowed. Hence, in a $`^3g`$-order formulation, real-valued initial data for $`(E_3H_2)`$ and $`(E_2\pm H_3)`$ needs to be of differentiability class $`C^3(𝒰)`$ (rather than $`C^2(𝒰)`$) with respect to the zeroth-order derivative level of $`𝐠`$. In short, not all initial data can be given freely. It can be easily inferred from the propagation equations along $`𝐮`$ for the Weyl curvature divergence equations (26) and (28), first published for an irrotational pressure-free fluid matter source in Ref. , and presented for a general perfect fluid matter source in Refs. and , that the characteristic velocities relative to $`𝐮`$ for the components $`(C_4)_2(C_5)_3`$ and $`(C_4)_3\pm (C_5)_2`$ are $`v=\pm \frac{1}{2}`$ too. Hence, comparing this result with Eqs. (165) and (170) in Appendix 2 and Eqs. (48) and (49) above, it becomes clear that the Weyl curvature divergence equations propagate relative to $`𝐮`$ at precisely the speed that is required to ensure that jump discontinuities in $`𝐞_1(E_3H_2)`$ and $`𝐞_1(E_2\pm H_3)`$ will always remain physically disallowed at any instant throughout the dynamical evolution of a cosmological model $`(,𝐠,𝐮)`$. It should be emphasised at this stage that this property is completely independent of the presence of matter. That is, of course the jump conditions (48) and (49) apply equally to vacuum spacetime configurations. (iii) Jump discontinuities in the outward first frame derivatives of the transverse Weyl curvature characteristic eigenfields $`(E_{}H_\times )`$ and $`(E_\times \pm H_{})`$ are physically allowed. Clearly, this situation reflects the freedom of specifying four arbitrary (non-analytic) real-valued functions $`I_{^3g}:=\{a_1(x^i),a_2(x^i),a_3(x^i),a_4(x^i)\}`$ of differentiability class $`C^2(𝒰)`$ with respect to the zeroth-order derivative level of $`𝐠`$ as the initial data for the dynamical degrees of freedom associated with the gravitational field itself. ### B Jump discontinuities at derivative level $`^2g`$ To be able to argue in terms of physical effects described by the deviation equation for a set of test particles, we have to turn our attention directly to the set of Weyl curvature characteristic eigenfields and possible discontinuous changes in their values. Such changes are driven by the dynamics of the underlying connection fields and their derivatives. Hence, to facilitate the interpretation of generic gravitational dynamics, it would be desirable to have available a $`^2g`$-order FOSH evolution system derived from a $`1+3`$ orthonormal frame connection formulation of gravitational fields (see Refs. and for reviews of the latter). Unfortunately, to date, such a formulation has not been accomplished for the generic case. The exception is the $`^2g`$-order FOSH evolution system for perfect fluid cosmological models $`(,𝐠,𝐮)`$ with an Abelian $`G_2`$ isometry group we will present in Sec. IV below. In the absence of such a generally applicable dynamical formulation, we return to our local viewpoint and investigate how, in a small isotropic neighbourhood $`𝒰`$ in the local rest 3-space of an arbitrary event $`𝒫`$, the constraint equations at derivative level $`^2g`$ restrict the occurrence of jump discontinuities in the values of the Weyl curvature characteristic eigenfields themselves. To this end, we now focus on Eqs. (13) – (24), and certain linear combinations of the components thereof provided by Eqs. (131) – (153) and (180) – (218). Again, we consider in $`𝒰`$ a spherical spacelike 2-surface $`𝒥`$:$`\{t=\text{const},\varphi (x^\mu )=\text{const}\}`$, and assume that across $`𝒥`$ all geometrically defined field variables $`u^A`$ as well as their tangential first frame derivatives are continuous. This leads to: (i) $`v=0`$ longitudinal Weyl curvature characteristic eigenfields: $`E_+`$ $`=`$ $`\frac{1}{3}𝐞_1(a_1)+\text{tangential frame derivatives/algebraic terms}`$ (53) $`H_+`$ $`=`$ $`\frac{1}{3}𝐞_1(\omega _1)+\text{tangential frame derivatives/algebraic terms},`$ (54) from Eqs. (180) and (186). Across $`𝒥`$, the generalised Gauß–Friedmann equation (153) and the fluid vorticity divergence equation (149), respectively, then impose the restrictions $`\left[E_+\right]_{t,\varphi =\mathrm{const}}`$ $`=`$ $`\frac{1}{3}\left[𝐞_1(a_1)\right]_{t,\varphi =\mathrm{const}}=\frac{1}{6}\left[\mu \right]_{t,\varphi =\mathrm{const}}`$ (55) $`\left[H_+\right]_{t,\varphi =\mathrm{const}}`$ $`=`$ $`\frac{1}{3}\left[𝐞_1(\omega _1)\right]_{t,\varphi =\mathrm{const}}=0.`$ (56) That is, jump discontinuities in the values of the Coulomb-like Weyl curvature characteristic eigenfield $`E_+`$ originate from jump discontinuities in the values of the matter total energy density $`\mu `$ and are physically allowed, if the momentum conservation equation (171) is satisfied on both sides of $`𝒥`$ with $`0=[(C_{\mathrm{PF}})_1]_{t,\varphi =\mathrm{const}}`$ \[cf. Eq. (52)\]. In a $`^2g`$-order formulation, real-valued initial data for $`\mu `$ (and so for $`E_+`$) is thus required to be of differentiability class $`C^1(𝒰)`$ with respect to the zeroth-order derivative level of $`𝐠`$. Note that in the vacuum subcase $`\mu =0[E_+]_{t,\varphi =\mathrm{const}}=0`$. Jump discontinuities in the values of $`H_+`$ are not physically allowed, and so, in a $`^2g`$-order formulation, real-valued initial data for $`H_+`$ is required to be of differentiability class $`C^2(𝒰)`$ with respect to the zeroth-order derivative level of $`𝐠`$. (ii) $`v=\pm \frac{1}{2}`$ semi-longitudinal Weyl curvature characteristic eigenfields: $`(E_3H_2)`$ $`=`$ $`\frac{1}{2}𝐞_1(\sigma _3n_2\frac{1}{\sqrt{3}}\omega _3\frac{1}{\sqrt{3}}a_2)+\text{tangential frame derivatives/algebraic terms}`$ (57) $`(E_2\pm H_3)`$ $`=`$ $`\frac{1}{2}𝐞_1(\sigma _2\pm n_3+\frac{1}{\sqrt{3}}\omega _2\frac{1}{\sqrt{3}}a_3)+\text{tangential frame derivatives/algebraic terms},`$ (58) from Eqs. (195) and (204). In this case, we find that across $`𝒥`$ the fluid shear divergence/Jacobi constraint equations (141) and (148) impose the restrictions $`\left[(E_3H_2)\right]_{t,\varphi =\mathrm{const}}`$ $`=`$ $`\frac{1}{2}[𝐞_1(\sigma _3n_2\frac{1}{\sqrt{3}}\omega _3\frac{1}{\sqrt{3}}a_2)]_{t,\varphi =\mathrm{const}}=0`$ (59) $`\left[(E_2\pm H_3)\right]_{t,\varphi =\mathrm{const}}`$ $`=`$ $`\frac{1}{2}[𝐞_1(\sigma _2\pm n_3+\frac{1}{\sqrt{3}}\omega _2\frac{1}{\sqrt{3}}a_3)]_{t,\varphi =\mathrm{const}}=0.`$ (60) That is, jump discontinuities in the values of the semi-longitudinal Weyl curvature characteristic eigenfields $`(E_3H_2)`$ and $`(E_2\pm H_3)`$ are not physically allowed. In a $`^2g`$-order formulation real-valued initial data for $`(E_3H_2)`$ and $`(E_2\pm H_3)`$ is thus required to be of differentiability class $`C^2(𝒰)`$ (rather than $`C^1(𝒰)`$) with respect to the zeroth-order derivative level of $`𝐠`$. (iii) $`v=\pm \mathrm{\hspace{0.17em}1}`$ transverse Weyl curvature characteristic eigenfields: $`(E_{}H_\times )`$ $`=`$ $`𝐞_1(\sigma _{}n_\times )+\text{tangential frame derivatives/algebraic terms}`$ (61) $`(E_\times \pm H_{})`$ $`=`$ $`𝐞_1(\sigma _\times \pm n_{})+\text{tangential frame derivatives/algebraic terms},`$ (62) from Eqs. (211) and (218). In this case, we find that across $`𝒥`$ the constraint equations impose no restrictions so that $`\left[(E_{}H_\times )\right]_{t,\varphi =\mathrm{const}}`$ $`=`$ $`\left[𝐞_1(\sigma _{}n_\times )\right]_{t,\varphi =\mathrm{const}}=\text{unconstrained}`$ (63) $`\left[(E_\times \pm H_{})\right]_{t,\varphi =\mathrm{const}}`$ $`=`$ $`\left[𝐞_1(\sigma _\times \pm n_{})\right]_{t,\varphi =\mathrm{const}}=\text{unconstrained}.`$ (64) That is, jump discontinuities in the values of the transverse Weyl curvature characteristic eigenfields $`(E_{}H_\times )`$ and $`(E_\times \pm H_{})`$ are physically allowed. Hence, they can transport arbitrary non-zero jump discontinuities $`[𝐞_1(\sigma _{}n_\times )]_{t,\varphi =\mathrm{const}}`$ and $`[𝐞_1(\sigma _\times \pm n_{})]_{t,\varphi =\mathrm{const}}`$ of finite magnitude. Again, this reflects the freedom of specifying four arbitrary (non-analytic) real-valued functions $`I_{^2g}:=\{A_1(x^i),A_2(x^i),A_3(x^i),A_4(x^i)\}`$ of differentiability class $`C^1(𝒰)`$ with respect to the zeroth-order derivative level of $`𝐠`$ as the initial data for the dynamical degrees of freedom associated with the gravitational field itself. For completeness, we now also briefly discuss the effect of the constraint equations (13) – (20), respectively, Eqs. (131) – (149), on the outward first frame derivatives of the characteristic eigenfields associated with the fluid kinematical branch of the $`^3g`$-order FOSH evolution system in Ref. \[cf. Eqs. (3.27) and (3.28)\]. The derivative level is hence one below the Weyl curvature case. For the $`v=0`$ fluid kinematical characteristic eigenfields we find the jump conditions: (iv) $`v=0`$ fluid kinematical characteristic eigenfields: $`[𝐞_1(\frac{1}{3}\mathrm{\Theta }+\sigma _+)]_{t,\varphi =\mathrm{const}}=0`$ (65) $`[𝐞_1(\omega _1)]_{t,\varphi =\mathrm{const}}=0`$ (66) $`[𝐞_1(\sigma _3+\frac{1}{\sqrt{3}}\omega _3)]_{t,\varphi =\mathrm{const}}=\text{unconstrained}`$ (67) $`[𝐞_1(\sigma _2\frac{1}{\sqrt{3}}\omega _2)]_{t,\varphi =\mathrm{const}}=\text{unconstrained}.`$ (68) The last two conditions imply the existence of two generically non-zero fluid rotational modes that were identified before by Ehlers et al in an analysis of linearised perturbations of arbitrary background dust spacetimes . Corresponding real-valued initial data for these modes is required to be of differentiability class $`C^1(𝒰)`$ with respect to the zeroth-order derivative level of $`𝐠`$. Finally, for the different parts of the $`v=\pm c_s`$ fluid kinematical characteristic eigenfields we find the jump conditions: (v) $`v=\pm c_s`$ fluid kinematical characteristic eigenfields: $`[𝐞_1(\frac{1}{3}\mathrm{\Theta }2\sigma _+)]_{t,\varphi =\mathrm{const}}=\text{unconstrained}`$ (69) $`[𝐞_1(\sigma _3\frac{1}{\sqrt{3}}\omega _3)]_{t,\varphi =\mathrm{const}}=\pm [𝐞_1(\frac{1}{\sqrt{3}}a_2+n_2)]_{t,\varphi =\mathrm{const}}=\text{unconstrained}`$ (70) $`[𝐞_1(\sigma _2+\frac{1}{\sqrt{3}}\omega _2)]_{t,\varphi =\mathrm{const}}=\pm [𝐞_1(\frac{1}{\sqrt{3}}a_3n_3)]_{t,\varphi =\mathrm{const}}=\text{unconstrained}`$ (71) $`[𝐞_1(\dot{u}_1)]_{t,\varphi =\mathrm{const}}=\text{unconstrained}`$ (72) $`[𝐞_1(\dot{u}_2)]_{t,\varphi =\mathrm{const}}=\text{unconstrained}`$ (73) $`[𝐞_1(\dot{u}_3)]_{t,\varphi =\mathrm{const}}=\text{unconstrained}.`$ (74) Again, corresponding real-valued initial data for these parts is required to be of differentiability class $`C^1(𝒰)`$ with respect to the zeroth-order derivative level of $`𝐠`$. Note especially that the jump conditions (70) and (71) are precisely of such a nature that no violations of the jump conditions (59) and (60) above for the semi-longitudinal Weyl curvature characteristic eigenfields may occur. ## IV Worked example: Cosmological models with Abelian $`G_2`$ isometry group In this section we turn to discuss in some detail a new, fully gauge-fixed, $`^2g`$-order autonomous FOSH evolution system for spatially inhomogeneous perfect fluid cosmological models $`(,𝐠,𝐮)`$ which are invariant under the transformations of an Abelian $`G_2`$ isometry group that is simply transitive on spacelike 2-surfaces. Thus, all geometrically defined field variables $`u^A`$ vary in one spatial direction only. A systematic approach to this class of cosmological models was brought forward some time ago by Wainwright in Refs. and , wherein the generic case was given the classification label “A(i)”. A number of exact, real analytic, solutions to the EFE for Abelian $`G_2`$ perfect fluid cosmological models are known, such as those listed in Refs. and or the singularity-free solution obtained by Senovilla , but most of them belong to dynamically restricted or higher-symmetry subcases. ### A Well-posed Cauchy initial value problem Choosing an orbit-aligned group-invariant $`1+3`$ ONF basis $`\{𝐞_0,𝐞_\alpha \}`$ such that commutation relations $$0=[𝝃,𝐞_0]=[𝝃,𝐞_\alpha ]=[𝜼,𝐞_0]=[𝜼,𝐞_\alpha ]$$ (75) hold between the two commuting spacelike Killing vector fields $`𝝃`$ and $`𝜼`$ and $`\{𝐞_0,𝐞_\alpha \}`$, and assuming that $`𝐞_0𝐮`$ is orthogonal to the isometry group orbits, it follows that for all solutions in the Abelian $`G_2`$ class we have $$0=𝐞_2(u^A)=𝐞_3(u^A),0=\dot{u}_2=\dot{u}_3=\omega ^\alpha (𝐮)=a_2=a_3=(n2n_+)=n_2=n_3.$$ (76) That is, $`𝐞_2`$ and $`𝐞_3`$ are tangent to the isometry group orbits. Besides $`𝐞_0𝐮`$ also the frame basis field $`𝐞_1`$ is hypersurface orthogonal . The canonical choice Wainwright proposes for Abelian $`G_2`$ perfect fluid cosmological models consists of introducing fluid-comoving local coordinates $`\{t,x,y,z\}`$ adapted to the integral curves of $`𝐮`$, $`𝝃`$ and $`𝜼`$ such that $$𝝃=_y,𝜼=_z;$$ (77) the isometry group orbits, which have vanishing Gaußian 2-curvature, are thus given by spacelike 2-surfaces $`\{t=\text{const},x=\text{const}\}`$. Additionally, the coordinate components of the $`1+3`$ ONF basis $`\{𝐞_0,𝐞_\alpha \}`$ as introduced by Eq. (10) are specialised to $`\begin{array}{ccccccc}𝐞_0& =& M^1_t,\hfill & & 𝐞_2& =& e_2{}_{}{}^{2}_{y}^{},\hfill \\ 𝐞_1& =& e_1{}_{}{}^{1}_{x}^{}+e_1{}_{}{}^{2}_{y}^{}+e_1{}_{}{}^{3}_{z}^{},\hfill & & 𝐞_3& =& e_3{}_{}{}^{2}_{y}^{}+e_3{}_{}{}^{3}_{z}^{};\hfill \end{array}`$ (80) in view of $`0=\omega ^\alpha (𝐮)`$ the choice $`0=M_i`$ is made.<sup>††</sup><sup>††</sup>††For $`0=\omega ^\alpha (𝐮)`$, the choice $`0=M_i`$ is the simplest one possible, but it is by no means compulsory. Note that Wainwright’s canonical choice establishes the property $`𝐞_2𝝃`$. All geometrically defined field variables $`u^A`$ will now only be functions of the local coordinates $`t`$ and $`x`$. From the commutators (see, e.g., Refs. and ), the canonical choice of $`1+3`$ ONF basis and local coordinates has the direct consequences $$0=(\sqrt{3}\sigma _\times +\mathrm{\Omega }_1)=(\sqrt{3}\sigma _2+\mathrm{\Omega }_2)=(\sqrt{3}\sigma _3\mathrm{\Omega }_3),$$ (81) implying the three spatial coordinate conditions $`0=e_2{}_{}{}^{1}=e_2{}_{}{}^{3}=e_3^1`$ in Eq. (80) will be automatically preserved along $`𝐮`$. Thus, the spatial frame $`\{𝐞_\alpha \}`$ will presently not be Fermi-transported along $`𝐮`$, but $`\mathrm{\Omega }^\alpha `$ (and so the three frame gauge source functions $`T^\alpha _\beta `$) will be adapted to the fluid rate of shear instead. Likewise, the canonical choice leads to $$0=(n_+\sqrt{3}n_{}).$$ (82) We now define two new frame variables for two components of the fluid rate of expansion tensor by $$\alpha :=(\frac{1}{3}\mathrm{\Theta }2\sigma _+),\beta :=(\frac{1}{3}\mathrm{\Theta }+\sigma _+),$$ (83) shadowing the variable names used earlier in a related context in Ref. (on this choice of variables, see also the remarks made on the so-called “Taub gauge” for fluid spacetime geometries in Ref. ). Substituting from Eq. (81) into the evolution equation for the semi-longitudinal fluid shear component $`\sigma _3`$ (see, e.g., Ref. ), one finds that the latter is involutive. Therefore, on performing a spatial rotation of $`\{𝐞_\alpha \}`$ about $`𝐞_1`$ at every point of a given 2-surfaces $`\{t=\text{const},x=\text{const}\}`$, one can set $$0=\sigma _3$$ (84) to hold at every event of $`(,𝐠,𝐮)`$.<sup>‡‡</sup><sup>‡‡</sup>‡‡Alternatively one could choose $`𝐞_3𝜼`$, and then set $`0=\sigma _2`$ instead. Determining the evolution along $`𝐮`$ of the frame gauge source function $`(T^0)^1T^0{}_{1}{}^{}=\dot{u}_1`$ as outlined in Subsec. II B, at this point, from a dynamical viewpoint, the coordinate and spatial frame freedom has been completely fixed (modulo coordinate reparameterisation freedom given by $`t^{}=t^{}(t)`$, $`x^{}=x^{}(x)`$, $`y^{}=y+f(x)`$ and $`z^{}=z+g(x)`$ ). It is now fairly straightforward to derive from the equations of a $`1+3`$ orthonormal frame connection formulation of gravitational fields as given in Refs. and an evolution system of autonomous partial differential equations in FOSH format for the following set of eleven geometrically defined field variables $`u^A=u^A(t,x)`$: $$u^A=(e_1{}_{}{}^{1},\beta ,\sigma _2,a_1,\mu ,\alpha ,\dot{u}_1,\sigma _{},n_\times ,\sigma _\times ,n_{})^T.$$ (85) Note that the evolution of the frame coordinate components other than $`e_1^1`$ is decoupled from this set; $`e_1^1`$ itself forms a background field for pairwise dynamical interactions between $`\{\alpha ,\dot{u}_1\}`$, $`\{\sigma _{},n_\times \}`$ and $`\{\sigma _\times ,n_{}\}`$ in the sense of Geroch . 11-dimensional autonomous first-order symmetric hyperbolic evolution system: $`M^1_te_1^1`$ $`=`$ $`\alpha e_1^1`$ (86) $`M^1_t\beta `$ $`=`$ $`\frac{3}{2}\beta ^2\frac{3}{2}(\sigma _{}^2+\sigma _\times ^2+n_{}^2+n_\times ^2)+\frac{3}{2}\sigma _2^2\frac{1}{2}(2\dot{u}_1a_1)a_1\frac{1}{2}(p\mathrm{\Lambda })`$ (87) $`M^1_t\sigma _2`$ $`=`$ $`(3\beta \sqrt{3}\sigma _{})\sigma _2`$ (88) $`M^1_ta_1`$ $`=`$ $`\beta (\dot{u}_1+a_1)+3(n_{}\sigma _\times n_\times \sigma _{})`$ (89) $`M^1_t\mu `$ $`=`$ $`(\alpha +2\beta )(\mu +p)`$ (90) $`c_s^2M^1_t\alpha c_s^2e_1{}_{}{}^{1}_{x}^{}\dot{u}_1`$ $`=`$ $`c_s^2[\alpha ^2+\beta ^23(\sigma _{}^2+\sigma _\times ^2n_{}^2n_\times ^2)9\sigma _2^2+\dot{u}_1^2a_1^2\frac{1}{2}(\mu +p)]`$ (91) $`M^1_t\dot{u}_1c_s^2e_1{}_{}{}^{1}_{x}^{}\alpha `$ $`=`$ $`\alpha \dot{u}_1(\alpha +2\beta )[c_s^2{\displaystyle \frac{d^2p}{d\mu ^2}}(\mu +p)c_s^2]\dot{u}_1`$ (93) $`c_s^2[\mathrm{\hspace{0.17em}2}a_1(\alpha \beta )+6(n_{}\sigma _\times n_\times \sigma _{})]`$ $`M^1_t\sigma _{}+e_1{}_{}{}^{1}_{x}^{}n_\times `$ $`=`$ $`(\alpha +2\beta )\sigma _{}+2\sqrt{3}\sigma _\times ^2\sqrt{3}\sigma _2^22\sqrt{3}n_{}^2(\dot{u}_12a_1)n_\times `$ (94) $`M^1_tn_\times +e_1{}_{}{}^{1}_{x}^{}\sigma _{}`$ $`=`$ $`\alpha n_\times \dot{u}_1\sigma _{}`$ (95) $`M^1_t\sigma _\times e_1{}_{}{}^{1}_{x}^{}n_{}`$ $`=`$ $`(\alpha +2\beta +2\sqrt{3}\sigma _{})\sigma _\times +(\dot{u}_12a_12\sqrt{3}n_\times )n_{}`$ (96) $`M^1_tn_{}e_1{}_{}{}^{1}_{x}^{}\sigma _\times `$ $`=`$ $`(\alpha 2\sqrt{3}\sigma _{})n_{}+(\dot{u}_1+2\sqrt{3}n_\times )\sigma _\times .`$ (97) The characteristic condition for the set (86) – (97) is invariantly given by $`0=Q=(u^a\zeta _a)^5[(u^bu^c+c_s^2h^{bc})\zeta _b\zeta _c][(u^du^e+1^2h^{de})\zeta _d\zeta _e]^2`$, where $`\zeta _a:=_a\varphi `$ are the past-directed normals to characteristic 3-surfaces $`𝒞`$:$`\{\varphi (x^\mu )=\text{const}\}`$. The connection characteristic eigenfields associated with non-zero characteristic velocities are $`\{(\alpha \pm c_s\dot{u}_1)/(1+c_s^2)^{1/2}\}`$ with $`v=\pm c_s`$ and $`\{(\sigma _{}n_\times )/\sqrt{2},(\sigma _\times \pm n_{})/\sqrt{2}\}`$ with $`v=\pm \mathrm{\hspace{0.17em}1}`$. Initial data which is invariant under the transformations of an Abelian $`G_2`$ isometry group has to satisfy the following set of constraint equations: Initial value constraint equations: $`e_1{}_{}{}^{1}M_{}^{1}_xM`$ $`=`$ $`\dot{u}_1`$ (98) $`0`$ $`=`$ $`(C_{\mathrm{com}})_1:=e_1{}_{}{}^{1}_{x}^{}e_2{}_{}{}^{2}(a_1+\sqrt{3}n_\times )e_2^2`$ (99) $`0`$ $`=`$ $`(C_{\mathrm{com}})_2:=e_1{}_{}{}^{1}_{x}^{}e_3{}_{}{}^{2}(a_1\sqrt{3}n_\times )e_3{}_{}{}^{2}+2\sqrt{3}n_{}e_2^2`$ (100) $`0`$ $`=`$ $`(C_{\mathrm{com}})_3:=e_1{}_{}{}^{1}_{x}^{}e_3{}_{}{}^{3}(a_1\sqrt{3}n_\times )e_3^3`$ (101) $`0`$ $`=`$ $`(C_\mathrm{G}):={}_{}{}^{}R+2(2\alpha +\beta )\beta 6(\sigma _{}^2+\sigma _\times ^2+\sigma _2^2)2\mu 2\mathrm{\Lambda }`$ (102) $`0`$ $`=`$ $`(C_1)_1:=e_1{}_{}{}^{1}_{x}^{}\beta +a_1(\alpha \beta )+3(n_{}\sigma _\times n_\times \sigma _{})`$ (103) $`0`$ $`=`$ $`(C_1)_3:=(e_1{}_{}{}^{1}_{x}^{}3a_1+\sqrt{3}n_\times )\sigma _2`$ (104) $`0`$ $`=`$ $`(C_{\mathrm{PF}})_1:=c_s^2e_1{}_{}{}^{1}_{x}^{}\mu +(\mu +p)\dot{u}_1;`$ (105) the 3-Ricci curvature scalar $`{}_{}{}^{}R`$ of spacelike 3-surfaces $`𝒮`$:$`\{t=\text{const}\}`$ orthogonal to $`𝐮`$ is given by $${}_{}{}^{}R:=2(2e_1{}_{}{}^{1}_{x}^{}3a_1)a_16(n_{}^2+n_\times ^2).$$ (106) The constraint equations (102) – (105) in this set are specialisations of Eqs. (153), (131), (148) and (171), respectively. The remaining ones, Eqs. (98) – (101), derive from the commutators. > Algorithm: Given an equation of state of the form $`p=p(\mu )`$, the initial data, which can be specified freely (modulo minimal differentiability requirements and the remaining coordinate reparameterisation freedom) as functions of the spatial coordinate $`x`$ on a spacelike 3-surface $`𝒮`$:$`\{t=\text{const}\}`$ orthogonal to the fluid 4-velocity field $`𝐮`$, are the values of the variables $`\{e_1{}_{}{}^{1},\alpha ,\sigma _{},n_\times ,\sigma _\times ,n_{},\mu \}`$, together with the cosmological constant $`\mathrm{\Lambda }`$. Specifying the values of the variables $`\{e_2{}_{}{}^{2},e_3{}_{}{}^{2},e_3{}_{}{}^{3},\beta ,\sigma _2,a_1\}`$ at one point on this 3-surface, their spatial distribution follows from Eqs. (99) – (102), respectively, while $`\dot{u}_1`$ is determined through Eq. (105) and $`M`$ through Eq. (98). Then all time derivatives of these variables are known. Note that one of the coordinate components $`e_1^2`$ and $`e_1^3`$ is arbitrary as functions of $`x`$, too, while the other follows from the unit-magnitude property of $`𝐞_1`$. It remains to specify boundary conditions for all variables to obtain unique solutions. Two subcases of importance are contained within the class of Abelian $`G_2`$ perfect fluid cosmological models: the orthogonally transitive subcase arises when $`0=e_1{}_{}{}^{2}=e_1{}_{}{}^{3}0=\sigma _2`$, which itself specialises to the diagonal (“polarised”) subcase when additionally $`0=e_3{}_{}{}^{2}0=\sigma _\times =n_{}`$, leading to a diagonal line element (see Refs. and ). The initial value constraint equations (99) – (105) are propagated along $`𝐮`$ via a FOSH evolution system of autonomous partial differential equations, where the characteristic speeds are $`\|v\|=0`$, according to 7-dimensional autonomous constraint evolution system: $`M^1_t(C_{\mathrm{com}})_1`$ $`=`$ $`(\alpha +\beta +\sqrt{3}\sigma _{})(C_{\mathrm{com}})_1e_2{}_{}{}^{2}(C_1)_{1}^{}`$ (107) $`M^1_t(C_{\mathrm{com}})_2`$ $`=`$ $`(\alpha +\beta \sqrt{3}\sigma _{})(C_{\mathrm{com}})_22\sqrt{3}\sigma _\times (C_{\mathrm{com}})_1e_3{}_{}{}^{2}(C_1)_{1}^{}`$ (108) $`M^1_t(C_{\mathrm{com}})_3`$ $`=`$ $`(\alpha +\beta \sqrt{3}\sigma _{})(C_{\mathrm{com}})_3e_3{}_{}{}^{3}(C_1)_{1}^{}`$ (109) $`M^1_t(C_\mathrm{G})`$ $`=`$ $`(\alpha +\beta )(C_\mathrm{G})4(\dot{u}_1+a_1)(C_1)_1`$ (110) $`M^1_t(C_1)_1`$ $`=`$ $`(\alpha +3\beta )(C_1)_1+\sigma _2(C_1)_3\frac{1}{4}(\dot{u}_1a_1)(C_\mathrm{G})\frac{1}{2}(C_{\mathrm{PF}})_1`$ (111) $`M^1_t(C_1)_3`$ $`=`$ $`(\alpha +3\beta \sqrt{3}\sigma _{})(C_1)_39\sigma _2(C_1)_1`$ (112) $`M^1_t(C_{\mathrm{PF}})_1`$ $`=`$ $`\mathrm{\hspace{0.17em}2}(\alpha +\beta )(C_{\mathrm{PF}})_1(\alpha +2\beta )[c_s^2+c_s^2{\displaystyle \frac{d^2p}{d\mu ^2}}(\mu +p)](C_{\mathrm{PF}})_12c_s^2(\mu +p)(C_1)_1.`$ (113) The full set of propagation equations for the Abelian $`G_2`$ perfect fluid cosmological models now forms a larger autonomous FOSH evolution system according to Eq. (1), with the previous 11-dimensional system as a subset. The virtue of the larger system is that it explicitly shows that the set of dynamical field equations introduced is consistent: the initial value constraint equations are preserved by the time evolution equations (if they are true initially, they remain true thereafter). This completes the discussion on a well-posed initial value problem for this class of cosmological models in the $`1+3`$ orthonormal frame connection formulation of gravitational fields. We finally list the specialisations which the expressions (180) – (218) for the Weyl curvature characteristic eigenfields undergo by the geometrical restrictions imposed by the Abelian $`G_2`$ isometry group. We use Eq. (104) to eliminate derivatives $`_x\sigma _2`$. Weyl curvature characteristic eigenfields: $`E_+`$ $`=`$ $`\frac{1}{3}e_1{}_{}{}^{1}_{x}^{}a_1\frac{1}{3}(\alpha \beta )\beta \sigma _{}^2\sigma _\times ^2+2n_{}^2+2n_\times ^2+\frac{1}{2}\sigma _2^2`$ (114) $`H_+`$ $`=`$ $`\frac{3}{2}(\sigma _{}n_\times )(\sigma _\times +n_{})+\frac{3}{2}(\sigma _{}+n_\times )(\sigma _\times n_{})`$ (115) $`(E_3H_2)`$ $`=`$ $`\sqrt{3}(\sigma _\times \pm n_{})\sigma _2`$ (116) $`(E_2\pm H_3)`$ $`=`$ $`(\beta a_1+\sqrt{3}\sigma _{}\sqrt{3}n_\times )\sigma _2`$ (117) $`(E_{}H_\times )`$ $`=`$ $`(e_1{}_{}{}^{1}_{x}^{}\alpha a_1)(\sigma _{}n_\times )\pm 2\sqrt{3}n_{}(\sigma _\times \pm n_{})\pm n_\times (\beta a_1)+\frac{\sqrt{3}}{2}\sigma _2^2`$ (118) $`(E_\times \pm H_{})`$ $`=`$ $`(e_1{}_{}{}^{1}_{x}^{}\alpha a_1)(\sigma _\times \pm n_{})2\sqrt{3}n_{}(\sigma _{}n_\times )n_{}(\beta a_1).`$ (119) Note that $`0=\sigma _20=(E_3H_2)=(E_2\pm H_3)`$ holds. ### B Transport equations for jump discontinuities in outward first derivatives To give a graphic example, which, we believe, will also be of some interest in numerical investigations of dynamical features of Abelian $`G_2`$ perfect fluid cosmological models, we conclude this section with a brief derivation of the transport equations that describe how physically relevant jump discontinuities in the outward first derivatives of the initial data are propagated along the bicharacteristic rays of the setting. We confine ourselves to modes with $`v0`$ relative to $`𝐮`$. To this end, we first turn to Eq. (4) for each of $`v\{\pm c_s,\pm \mathrm{\hspace{0.17em}1}\}`$, leading to the conditions $$[_\varphi \dot{u}_1]=\pm c_s[_\varphi \alpha ],[_\varphi n_\times ]=[_\varphi \sigma _{}],[_\varphi n_{}]=\pm [_\varphi \sigma _\times ],$$ (120) respectively. This then gives from Eq. (5), together with the evolution subsystem (91) – (97), the relations: (i) $`v=\pm c_s`$ longitudinal modes: $`0`$ $`=`$ $`(M^1_tc_se_1{}_{}{}^{1}_{x}^{}+\frac{1}{2}c_s^1(M^1_tc_sc_se_1{}_{}{}^{1}_{x}^{}c_s)`$ (122) $`+\alpha +\frac{1}{2}(\alpha +2\beta )[c_s^2{\displaystyle \frac{d^2p}{d\mu ^2}}(\mu +p)c_s^2]\frac{1}{2}c_s\dot{u}_1\pm c_sa_1)[_\varphi (\alpha \pm c_s^1\dot{u}_1)];`$ jump discontinuities in the outward first derivatives of initial data for $`\{\alpha ,\dot{u}_1\}`$, subject to Eq. (120), travel along the local sound cones. Note that, because of the general functional dependence $`c_s=c_s(\mu )`$, the evolution of sound cone initial data typically leads to the formation of “shocks”. This phenomenon becomes impossible in the special case of a linear baryotropic equation of state with $`p(\mu )=(\gamma 1)\mu `$, $`1\gamma 2`$, where $`c_s=(\gamma 1)^{1/2}=\text{const}`$. (ii) $`v=\pm \mathrm{\hspace{0.17em}1}`$ transverse modes: $`0`$ $`=`$ $`(M^1_te_1{}_{}{}^{1}_{x}^{}+\alpha +\beta \dot{u}_1\pm a_1)[_\varphi (\sigma _{}n_\times )]`$ (123) $`0`$ $`=`$ $`(M^1_te_1{}_{}{}^{1}_{x}^{}+\alpha +\beta \dot{u}_1\pm a_1)[_\varphi (\sigma _\times \pm n_{})].`$ (124) jump discontinuities in the outward first derivatives of initial data for $`\{\sigma _{},n_\times \}`$ and $`\{\sigma _\times ,n_{}\}`$, subject to Eq. (120), travel along the local light cones. ## V Conclusion The main conclusion of this work is that in examining a set of dynamical equations for a physical system such as the relativistic gravitational field equations, completed by the equations for all needed auxiliary variables, the constraint equations are crucial in determining what information can be propagated along the characteristic 3-surfaces of the evolution equations when these are expressed in FOSH format. We have explicitly shown how to examine the constraint equations to determine whether jump discontinuities in the derivatives of the initial data can be propagated along the various characteristic 3-surfaces in the case of the relativistic gravitational field equations with a baryotropic perfect fluid matter source, including the Weyl curvature variables. This makes it clear that such an investigation is needed to complement the determination of the set of characteristic 3-surfaces of any FOSH evolution system with any existing supplementary constraint equations, in order to determine which characteristic 3-surfaces in the set are physically relevant. The process outlined enables us to show why the characteristic 3-surfaces for the semi-longitudinal Weyl curvature characteristic eigenfields $`(E_3H_2)`$ and $`(E_2\pm H_3)`$ that are associated with $`\|v\|=\frac{1}{2}`$, apparent in a straightforward reduction of order $`^3g`$ of the evolution system of the relativistic gravitational field equations for a baryotropic perfect fluid matter source to FOSH format, cannot in fact be activated. It demonstrates that potentially associated semi-longitudinal gravitational radiation cannot occur, despite the occurrence of related characteristic 3-surfaces in the FOSH evolution system. It should be noted that the issue is not that the constraint equations are incompatible with the evolution equations, in the sense of not being conserved under the system’s time evolution. On the contrary, the propagation of these constraint equations is indeed compatible with the existence of these modes, as can be shown by considering an extended FOSH evolution system that includes variables representing satisfaction of the constraint equations (see also Refs. , and ). The issue is that the constraint equations do not allow the setting of jump discontinuities in (the derivatives of) the initial data for $`(E_3H_2)`$ and $`(E_2\pm H_3)`$ because the values of the components of the constraint equations themselves cannot suffer jump discontinuities: they have to be continuously zero from one spacetime event to any nearby one, i.e., everywhere. What this analysis does not do is to show in what manner the semi-longitudinal Weyl curvature characteristic eigenfields $`(E_3H_2)`$ and $`(E_2\pm H_3)`$, originally identified by Szekeres and shown there to have observable physical effects , will evolve in time, nor does it adequately characterise what freedom there is in setting initial data for these modes. It would be helpful to have some characterisation of the full freedom to assign these modes on an initial data 3-surface. The worked example for perfect fluid cosmological models with an Abelian $`G_2`$ isometry group presented in Sec. IV features neatly the conceptual and mathematical advantages one can gain from combining the idea of FOSH evolution systems with a $`1+3`$ orthonormal frame connection formulation of gravitational fields. It will be usable as a multi-facet test bed for numerical experiments of spacetime geometry evolution processes in relativistic cosmology and already provides the basis for work-in-progress on an interesting new scale-invariant, dimensionless dynamical formulation for the orthogonally transitive Abelian $`G_2`$ perfect fluid cosmological models . ###### Acknowledgements. We are grateful to Jürgen Ehlers and Claes Uggla for very helpful comments. This work was supported by the Deutsche Forschungsgemeinschaft (DFG) at Bonn, Germany (HvE), and the Foundation for Research and Development (FRD) at Pretoria, South Africa (GFRE). In parts the computer algebra packages REDUCE and CLASSI were employed. ## ### 1 Suggestive example: linearised relativistic gravitational field equations in a metric approach with Regge–Wheeler coordinate gauge fixing Rooted in tradition, a metric approach to formulating relativistic gravitational dynamics is still more frequently encountered in the literature than, e.g., the (extended) $`1+3`$ orthonormal frame formulation employed in this paper and in Ref. . This motivates the brief discussion of a metric approach based example derived from a paper by Kind, Ehlers and Schmidt that highlights the issue of the physical role of the constraint equations in the selection of appropriate initial data sets. The example arises as a special subcase of linearised relativistic gravitational field equations that describe small adiabatic, non-radial perturbations of a star in hydrostatic equilibrium. Starting from a set of local coordinates $`\{t,r,\vartheta ,\phi \}`$ and imposing Regge–Wheeler coordinate gauge fixing conditions, the Ansatz for the line element contains two real-valued metric functions $`f=f(t,r)`$ and $`g=g(t,r)`$; the angular behaviour of the line element can be given in terms of spherical harmonic functions $`Y_{lm}(\vartheta ,\phi )`$ (for more details see Sec. 2 of Ref. ). One can then obtain from the linearised Einstein field equations a second-order symmetric hyperbolic evolution system for $`f`$ and $`g`$ that takes the form $`\ddot{f}+f^{\prime \prime }`$ $`=`$ $`{\displaystyle \frac{2}{r}}f^{}{\displaystyle \frac{4}{r^2}}g^{}+{\displaystyle \frac{l(l+1)}{r^2}}f`$ (125) $`\ddot{g}+g^{\prime \prime }`$ $`=`$ $`{\displaystyle \frac{2}{r}}g^{}+{\displaystyle \frac{2}{r^2}}f+{\displaystyle \frac{1}{r^2}}(l(l+1)2)g.`$ (126) Moreover, $`f`$ and $`g`$ are bound to satisfy the constraint equation $$0=(C):=g^{\prime \prime }\frac{1}{r}f^{}+\frac{3}{r}g^{}\frac{l(l+1)}{2r^2}(f+g)\frac{1}{r}(fg).$$ (127) The characteristic 3-surfaces $`𝒞`$ determined by the principal parts of Eqs. (125) and (126) satisfy the conditions $`(tr)=\text{const}`$. If we assume that Eqs. (125) and (126) hold for $`f`$ and $`g`$, then we obtain for $`(C)`$ the evolution equation $$0=(\ddot{C})+(C)^{\prime \prime }+\frac{2}{r}(C)^{}\frac{l(l+1)}{2r^2}(C).$$ (128) Hence, if on a spacelike 3-surface $`𝒮`$:$`\{t=\text{const}\}`$ real-valued initial data $`I_{^2g}=\{f,\dot{f},g,\dot{g}\}`$ satisfy $`0=(C)=(\dot{C})`$, then any solution of Eqs. (125) and (126) satisfies $`0=(C)`$ everywhere. It is instructive to note that from Eq. (127) we can now add any arbitrary real-valued multiple $`ag^{\prime \prime }`$ to Eq. (126) to obtain a principal part of the form $$\ddot{g}+(1+a)g^{\prime \prime },$$ which instead defines characteristic 3-surfaces $`𝒞`$ that satisfy the condition $`((1+a)^{1/2}tr)=\text{const}`$. In particular, we can choose $`a=\mathrm{\hspace{0.17em}1}`$, leading to $`r=\text{const}`$, which corresponds to a propagation speed $`\|v\|=0`$. This implies that no information is transported from one event to a spatially separated nearby one. A different view considers the differentiability properties of $`f`$ and $`g`$: if one chooses real-valued initial data such that $`\{f,f^{},g,g^{}\}`$ are continuous, then the constraint equation (127) imposes no condition on the differentiability of $`f^{\prime \prime }`$ ( and so on $`\ddot{f}`$, from Eq. (125) ), while $`g^{\prime \prime }`$ ( and so $`\ddot{g}`$, from Eq. (126) ) is required to be continuous too. Hence, jump discontinuities are physically allowed only in the initial value of the former quantity and are propagated via Eq. (125). ### 2 Constraint equations component-wise Fluid shear divergence equations/Jacobi constraint equations: $`0`$ $`=`$ $`\frac{1}{2}(C_1)_1=𝐞_1(\frac{1}{3}\mathrm{\Theta }+\sigma _+)\frac{\sqrt{3}}{2}(𝐞_23a_2+\sqrt{3}n_2)(\sigma _3)\frac{\sqrt{3}}{2}(𝐞_33a_3\sqrt{3}n_3)(\sigma _2)`$ (131) $`\frac{1}{2}(𝐞_2+2\dot{u}_2a_2\sqrt{3}n_2)(\omega _3)+\frac{1}{2}(𝐞_3+2\dot{u}_3a_3+\sqrt{3}n_3)(\omega _2)`$ $`3a_1\sigma _++3(n_{}\sigma _\times n_\times \sigma _{})+\frac{1}{6}(n2n_+)\omega _1`$ $`0`$ $`=`$ $`(C_\mathrm{J})_1=\frac{1}{3}(𝐞_12a_1)(n2n_+)+(𝐞_22a_2)(a_3+\sqrt{3}n_3)+(𝐞_32a_3)(a_2\sqrt{3}n_2)`$ (134) $`+2(\frac{1}{3}\mathrm{\Theta }2\sigma _+)\omega _1+2(\sqrt{3}\sigma _2+\mathrm{\Omega }_2)\omega _3+2(\sqrt{3}\sigma _3\mathrm{\Omega }_3)\omega _2`$ $`0`$ $`=`$ $`\frac{1}{\sqrt{3}}[(C_1)_2(C_\mathrm{J})_3]=𝐞_1(\sigma _3n_2\frac{1}{\sqrt{3}}\omega _3\frac{1}{\sqrt{3}}a_2)+𝐞_2(\sigma _{}n_\times \pm \frac{1}{\sqrt{3}}a_1)+𝐞_3(\sigma _\times \pm n_{}+\frac{1}{\sqrt{3}}\omega _1)`$ (141) $`+\frac{1}{\sqrt{3}}(𝐞_23a_2+3\sqrt{3}n_2)(\sigma _+)\frac{1}{3\sqrt{3}}(𝐞_32a_3)(n+n_+)\frac{2}{3\sqrt{3}}𝐞_2(\mathrm{\Theta })`$ $`3a_1(\sigma _3n_2)3a_2(\sigma _{}n_\times )3a_3(\sigma _\times \pm n_{})(n_+\sqrt{3}n_{})\sigma _2`$ $`n_\times (\sqrt{3}\sigma _3\pm a_2)n_2(\sqrt{3}\sigma _{}\pm a_1)+\sqrt{3}n_3\sigma _\times \pm a_3n_{}`$ $`\frac{1}{\sqrt{3}}(2\sqrt{3}\sigma _22\dot{u}_32\mathrm{\Omega }_2\pm a_3\pm \sqrt{3}n_3)\omega _1\frac{1}{\sqrt{3}}(2\sqrt{3}\sigma _\times +2\mathrm{\Omega }_1\pm \frac{1}{3}n\pm \frac{1}{3}n_+\pm \sqrt{3}n_{})\omega _2`$ $`\frac{1}{\sqrt{3}}(\frac{2}{3}\mathrm{\Theta }+2\sigma _+2\sqrt{3}\sigma _{}\pm 2\dot{u}_1a_1\pm \sqrt{3}n_\times )\omega _3`$ $`0`$ $`=`$ $`\frac{1}{\sqrt{3}}[(C_1)_3\pm (C_\mathrm{J})_2]=𝐞_1(\sigma _2\pm n_3+\frac{1}{\sqrt{3}}\omega _2\frac{1}{\sqrt{3}}a_3)+𝐞_2(\sigma _\times \pm n_{}\frac{1}{\sqrt{3}}\omega _1)𝐞_3(\sigma _{}n_\times \frac{1}{\sqrt{3}}a_1)`$ (148) $`\pm \frac{1}{3\sqrt{3}}(𝐞_22a_2)(n+n_+)+\frac{1}{\sqrt{3}}(𝐞_33a_33\sqrt{3}n_3)(\sigma _+)\frac{2}{3\sqrt{3}}𝐞_3(\mathrm{\Theta })`$ $`3a_1(\sigma _2\pm n_3)3a_2(\sigma _\times \pm n_{})+3a_3(\sigma _{}n_\times )+(n_++\sqrt{3}n_{})\sigma _3`$ $`+n_\times (\sqrt{3}\sigma _2\pm a_3)n_3(\sqrt{3}\sigma _{}a_1)\sqrt{3}n_2\sigma _\times \pm a_2n_{}`$ $`\pm \frac{1}{\sqrt{3}}(2\sqrt{3}\sigma _32\dot{u}_2+2\mathrm{\Omega }_3\pm a_2\sqrt{3}n_2)\omega _1\pm \frac{1}{\sqrt{3}}(\frac{2}{3}\mathrm{\Theta }+2\sigma _++2\sqrt{3}\sigma _{}\pm 2\dot{u}_1a_1\sqrt{3}n_\times )\omega _2`$ $`\pm \frac{1}{\sqrt{3}}(2\sqrt{3}\sigma _\times 2\mathrm{\Omega }_1\frac{1}{3}n\frac{1}{3}n_+\pm \sqrt{3}n_{})\omega _3.`$ Fluid vorticity divergence equation: $$0=(C_2)=(𝐞_1\dot{u}_12a_1)(\omega _1)+(𝐞_2\dot{u}_22a_2)(\omega _2)+(𝐞_3\dot{u}_32a_3)(\omega _3).$$ (149) Generalised Gauß–Friedmann equation: $`0=(C_\mathrm{G})`$ $`=`$ $`2(2𝐞_13a_1)(a_1)+2(2𝐞_23a_2)(a_2)+2(2𝐞_33a_3)(a_3)`$ (153) $`+\frac{1}{6}n^2\frac{2}{3}n_+^26(n_{}^2+n_\times ^2+n_2^2+n_3^2)`$ $`+6(\frac{1}{3}\mathrm{\Theta }\sigma _+)(\frac{1}{3}\mathrm{\Theta }+\sigma _+)6(\sigma _{}^2+\sigma _\times ^2+\sigma _2^2+\sigma _3^2)`$ $`+2(\omega _12\mathrm{\Omega }_1)\omega _1+2(\omega _22\mathrm{\Omega }_2)\omega _2+2(\omega _32\mathrm{\Omega }_3)\omega _32\mu 2\mathrm{\Lambda }.`$ Weyl curvature divergence equations: $`0`$ $`=`$ $`\frac{1}{2}(C_4)_1=(𝐞_13a_1)(E_+)\frac{\sqrt{3}}{2}(𝐞_23a_2+\sqrt{3}n_2)(E_3)\frac{\sqrt{3}}{2}(𝐞_33a_3\sqrt{3}n_3)(E_2)`$ (156) $`+\frac{1}{6}𝐞_1(\mu )+3n_{}E_\times 3n_\times E_{}+3\sigma _{}H_\times 3\sigma _\times H_{}`$ $`3\omega _1H_+\frac{3}{2}(\sigma _2\sqrt{3}\omega _2)H_3+\frac{3}{2}(\sigma _3+\sqrt{3}\omega _3)H_2`$ $`0`$ $`=`$ $`\frac{1}{2}(C_5)_1=(𝐞_13a_1)(H_+)\frac{\sqrt{3}}{2}(𝐞_23a_2+\sqrt{3}n_2)(H_3)\frac{\sqrt{3}}{2}(𝐞_33a_3\sqrt{3}n_3)(H_2)`$ (160) $`\frac{1}{2}(\mu +p)\omega _1+3n_{}H_\times 3n_\times H_{}3\sigma _{}E_\times +3\sigma _\times E_{}`$ $`+3\omega _1E_++\frac{3}{2}(\sigma _2\sqrt{3}\omega _2)E_3\frac{3}{2}(\sigma _3+\sqrt{3}\omega _3)E_2`$ $`0`$ $`=`$ $`\frac{1}{\sqrt{3}}[(C_4)_2(C_5)_3]=(𝐞_1\pm 3\sigma _+3a_1)(E_3H_2)+(𝐞_2\sqrt{3}\sigma _3\pm 3\omega _33a_2\sqrt{3}n_2)(E_{}H_\times )`$ (165) $`+(𝐞_3\sqrt{3}\sigma _23\omega _23a_3+\sqrt{3}n_3)(E_\times \pm H_{})+\frac{1}{\sqrt{3}}(𝐞_23\sqrt{3}\sigma _33\omega _33a_2+3\sqrt{3}n_2)(E_+)`$ $`\frac{1}{\sqrt{3}}(𝐞_33\sqrt{3}\sigma _2\pm 3\omega _23a_33\sqrt{3}n_3)(H_+)\frac{1}{3\sqrt{3}}𝐞_2(\mu )\frac{1}{\sqrt{3}}(\mu +p)\omega _3`$ $`\pm \sqrt{3}(\sigma _{}n_\times )(E_3\pm H_2)\pm \sqrt{3}(\sigma _\times \pm n_{})(E_2H_3)(3\omega _1\pm n_+)(E_2\pm H_3)`$ $`0`$ $`=`$ $`\frac{1}{\sqrt{3}}[(C_4)_3\pm (C_5)_2]=(𝐞_1\pm 3\sigma _+3a_1)(E_2\pm H_3)+(𝐞_2\sqrt{3}\sigma _3\pm 3\omega _33a_2\sqrt{3}n_2)(E_\times \pm H_{})`$ (170) $`(𝐞_3\sqrt{3}\sigma _23\omega _23a_3+\sqrt{3}n_3)(E_{}H_\times )\pm \frac{1}{\sqrt{3}}(𝐞_23\sqrt{3}\sigma _33\omega _33a_2+3\sqrt{3}n_2)(H_+)`$ $`+\frac{1}{\sqrt{3}}(𝐞_33\sqrt{3}\sigma _2\pm 3\omega _23a_33\sqrt{3}n_3)(E_+)\frac{1}{3\sqrt{3}}𝐞_3(\mu )\pm \frac{1}{\sqrt{3}}(\mu +p)\omega _2`$ $`\sqrt{3}(\sigma _{}n_\times )(E_2H_3)\pm \sqrt{3}(\sigma _\times \pm n_{})(E_3\pm H_2)\pm (3\omega _1\pm n_+)(E_3H_2).`$ Momentum conservation equations: $`0=(C_{\mathrm{PF}})_1`$ $`=`$ $`c_s^2𝐞_1(\mu )+(\mu +p)\dot{u}_1`$ (171) $`0=(C_{\mathrm{PF}})_2`$ $`=`$ $`c_s^2𝐞_2(\mu )+(\mu +p)\dot{u}_2`$ (172) $`0=(C_{\mathrm{PF}})_3`$ $`=`$ $`c_s^2𝐞_3(\mu )+(\mu +p)\dot{u}_3.`$ (173) Weyl curvature characteristic eigenfields: $`E_+`$ $`=`$ $`\frac{1}{3}𝐞_1(a_1)+\frac{1}{6}(𝐞_23\sqrt{3}n_2)(a_2)+\frac{1}{6}(𝐞_3+3\sqrt{3}n_3)(a_3)`$ (180) $`+\frac{\sqrt{3}}{2}(𝐞_2a_2\frac{1}{\sqrt{3}}n_2)(n_2)\frac{\sqrt{3}}{2}(𝐞_3a_3+\frac{1}{\sqrt{3}}n_3)(n_3)`$ $`+(\frac{1}{3}\mathrm{\Theta }+\sigma _+)\sigma _++\frac{1}{3}(n2n_+)n_+`$ $`(\sigma _{}n_\times )(\sigma _{}+n_\times )(\sigma _\times +n_{})(\sigma _\times n_{})`$ $`+\frac{1}{2}(\sigma _3n_2)(\sigma _3+n_2)+\frac{1}{2}(\sigma _2+n_3)(\sigma _2n_3)`$ $`+n_{}^2+n_\times ^2+\frac{1}{3}(\omega _12\mathrm{\Omega }_1)\omega _1\frac{1}{6}(\omega _22\mathrm{\Omega }_2)\omega _2\frac{1}{6}(\omega _32\mathrm{\Omega }_3)\omega _3`$ $`+\frac{1}{2}(C_\mathrm{G})_{11}`$ $`H_+`$ $`=`$ $`\frac{\sqrt{3}}{2}(𝐞_2a_2\sqrt{3}n_2)(\sigma _2)+\frac{\sqrt{3}}{2}(𝐞_3a_3+\sqrt{3}n_3)(\sigma _3)`$ (186) $`+\frac{1}{3}(𝐞_1+2\dot{u}_1+a_1)(\omega _1)\frac{1}{6}(𝐞_2+2\dot{u}_2+a_23\sqrt{3}n_2)(\omega _2)`$ $`\frac{1}{6}(𝐞_3+2\dot{u}_3+a_3+3\sqrt{3}n_3)(\omega _3)`$ $`\frac{1}{2}(n2n_+)\sigma _+\frac{3}{2}(\sigma _{}n_\times )(\sigma _\times +n_{})+\frac{3}{2}(\sigma _{}+n_\times )(\sigma _\times n_{})`$ $`\frac{1}{2}(C_3)_{11}`$ $`(E_3H_2)`$ $`=`$ $`\frac{1}{2}𝐞_1(\sigma _3n_2\frac{1}{\sqrt{3}}\omega _3\frac{1}{\sqrt{3}}a_2)\pm \frac{1}{2}𝐞_2(\sigma _{}n_\times \pm \frac{1}{\sqrt{3}}a_1)`$ (195) $`\pm \frac{1}{2}𝐞_3(\sigma _\times \pm n_{}+\frac{1}{\sqrt{3}}\omega _1)\frac{\sqrt{3}}{2}(𝐞_2a_2+\sqrt{3}n_2)(\sigma _+)+\frac{1}{2\sqrt{3}}(𝐞_32a_3)(n_+)`$ $`\sqrt{3}(\sigma _2\pm \frac{1}{2\sqrt{3}}a_3\frac{3}{2}n_3)(\sigma _\times \pm n_{})\sqrt{3}(\sigma _3\pm \frac{1}{2\sqrt{3}}a_2\pm \frac{3}{2}n_2)(\sigma _{}n_\times )`$ $`+(\frac{1}{3}\mathrm{\Theta }+\sigma _+)\sigma _3\pm \frac{1}{2}a_1(\sigma _32n_2)\pm \frac{1}{2}(nn_+)\sigma _2+\frac{1}{3}(n2n_+)n_3`$ $`\frac{\sqrt{3}}{2}n_{}(\sigma _2n_3\pm \frac{1}{\sqrt{3}}a_3)\pm \frac{\sqrt{3}}{2}n_\times (\sigma _3\pm n_2\pm \frac{1}{\sqrt{3}}a_2)`$ $`\pm \frac{1}{2\sqrt{3}}(2\dot{u}_3\omega _2\pm 2\mathrm{\Omega }_2+a_3+\sqrt{3}n_3)\omega _1`$ $`\frac{1}{2\sqrt{3}}(\omega _12\mathrm{\Omega }_1n_+\pm \sqrt{3}n_{})\omega _2\pm \frac{1}{2\sqrt{3}}(2\dot{u}_1+a_1\sqrt{3}n_\times )\omega _3`$ $`\frac{1}{\sqrt{3}}[(C_\mathrm{G})_{12}\pm (C_3)_{31}]`$ $`(E_2\pm H_3)`$ $`=`$ $`\frac{1}{2}𝐞_1(\sigma _2\pm n_3+\frac{1}{\sqrt{3}}\omega _2\frac{1}{\sqrt{3}}a_3)\pm \frac{1}{2}𝐞_2(\sigma _\times \pm n_{}\frac{1}{\sqrt{3}}\omega _1)`$ (204) $`\frac{1}{2}𝐞_3(\sigma _{}n_\times \frac{1}{\sqrt{3}}a_1)\frac{1}{2\sqrt{3}}(𝐞_22a_2)(n_+)\frac{\sqrt{3}}{2}(𝐞_3a_3\sqrt{3}n_3)(\sigma _+)`$ $`+\sqrt{3}(\sigma _2\pm \frac{1}{2\sqrt{3}}a_3\frac{3}{2}n_3)(\sigma _{}n_\times )\sqrt{3}(\sigma _3\pm \frac{1}{2\sqrt{3}}a_2\pm \frac{3}{2}n_2)(\sigma _\times \pm n_{})`$ $`+(\frac{1}{3}\mathrm{\Theta }+\sigma _+)\sigma _2\pm \frac{1}{2}a_1(\sigma _2\pm 2n_3)\frac{1}{2}(nn_+)\sigma _3+\frac{1}{3}(n2n_+)n_2`$ $`\frac{\sqrt{3}}{2}n_{}(\sigma _3\pm n_2\pm \frac{1}{\sqrt{3}}a_2)\frac{\sqrt{3}}{2}n_\times (\sigma _2n_3\pm \frac{1}{\sqrt{3}}a_3)`$ $`\frac{1}{2\sqrt{3}}(2\dot{u}_2\pm \omega _32\mathrm{\Omega }_3+a_2\sqrt{3}n_2)\omega _1`$ $`\frac{1}{2\sqrt{3}}(2\dot{u}_1+a_1+\sqrt{3}n_\times )\omega _2\frac{1}{2\sqrt{3}}(\omega _12\mathrm{\Omega }_1n_+\sqrt{3}n_{})\omega _3`$ $`\frac{1}{\sqrt{3}}[(C_\mathrm{G})_{31}(C_3)_{12}]`$ $`(E_{}H_\times )`$ $`=`$ $`𝐞_1(\sigma _{}n_\times )\pm \frac{1}{2}𝐞_2(\sigma _3n_2+\frac{1}{\sqrt{3}}\omega _3\pm \frac{1}{\sqrt{3}}a_2)\frac{1}{2}𝐞_3(\sigma _2\pm n_3\frac{1}{\sqrt{3}}\omega _2\pm \frac{1}{\sqrt{3}}a_3)`$ (211) $`+\frac{\sqrt{3}}{2}(\sigma _2\pm \frac{1}{\sqrt{3}}a_3\pm 2n_3)(\sigma _2\pm n_3)\frac{\sqrt{3}}{2}(\sigma _3\pm \frac{1}{\sqrt{3}}a_22n_2)(\sigma _3n_2)`$ $`+(\frac{1}{3}\mathrm{\Theta }2\sigma _+)\sigma _{}\pm 3n_\times \sigma _+\pm \frac{1}{2}(n+2n_+)\sigma _\times +\frac{1}{3}(n+4n_+)n_{}`$ $`\pm a_1(\sigma _{}2n_\times )+\frac{1}{2}a_2n_2+\frac{1}{2}a_3n_3\pm n_{}\omega _1`$ $`\pm \frac{1}{2\sqrt{3}}(2\dot{u}_3\omega _2\pm 2\mathrm{\Omega }_2+a_3\sqrt{3}n_3)\omega _2\pm \frac{1}{2\sqrt{3}}(2\dot{u}_2\pm \omega _32\mathrm{\Omega }_3+a_2+\sqrt{3}n_2)\omega _3`$ $`\frac{1}{2\sqrt{3}}[(C_\mathrm{G})_{22}(C_\mathrm{G})_{33}\pm 2(C_3)_{23}]`$ $`(E_\times \pm H_{})`$ $`=`$ $`𝐞_1(\sigma _\times \pm n_{})\pm \frac{1}{2}𝐞_2(\sigma _2\pm n_3\frac{1}{\sqrt{3}}\omega _2\pm \frac{1}{\sqrt{3}}a_3)\pm \frac{1}{2}𝐞_3(\sigma _3n_2+\frac{1}{\sqrt{3}}\omega _3\pm \frac{1}{\sqrt{3}}a_2)`$ (218) $`\frac{\sqrt{3}}{2}(\sigma _2\pm \frac{1}{\sqrt{3}}a_3\pm 2n_3)(\sigma _3n_2)\frac{\sqrt{3}}{2}(\sigma _3\pm \frac{1}{\sqrt{3}}a_22n_2)(\sigma _2\pm n_3)`$ $`+(\frac{1}{3}\mathrm{\Theta }2\sigma _+)\sigma _\times 3n_{}\sigma _+\frac{1}{2}(n+2n_+)\sigma _{}+\frac{1}{3}(n+4n_+)n_\times `$ $`\pm a_1(\sigma _\times \pm 2n_{})\frac{1}{2}a_2n_3+\frac{1}{2}a_3n_2\pm n_\times \omega _1`$ $`\frac{1}{2\sqrt{3}}(2\dot{u}_2\pm \omega _32\mathrm{\Omega }_3+a_2+\sqrt{3}n_2)\omega _2\pm \frac{1}{2\sqrt{3}}(2\dot{u}_3\omega _2\pm 2\mathrm{\Omega }_2+a_3\sqrt{3}n_3)\omega _3`$ $`\frac{1}{2\sqrt{3}}[\mathrm{\hspace{0.17em}2}(C_\mathrm{G})_{23}(C_3)_{22}\pm (C_3)_{33}].`$
warning/0007/cond-mat0007436.html	ar5iv	text	# Stable and Metastable vortex states and the first order transition across the peak effect region in weakly pinned 2H-NbSe2 ## I INTRODUCTION In the presence of strong pinning, the vortex state of type II superconductors is usually characterized by the critical current density $`J_c(H,T)`$ which decreases monotonically with increasing field $`H`$ or temperature $`T`$. In the weakly pinned superconductors, on the other hand, the interplay between the intervortex interaction and the flux pinning produces an anomalous peak in $`J_c`$, as a function of both field and temperature just below the normal state boundary (usually designated as the peak effect or PE). Within the collective pinning description, this signifies that the vortex phase undergoes a transition/crossover from an ordered state to a disordered state. The detailed nature of this transition, e.g., whether it is a thermodynamic phase transition or not, remains a subject of considerable debate. One of the key issues is the detection of an anomaly in the thermodynamic quantities, such as, specific heat or equilibrium magnetization $`M_{eq}`$. $`J_c`$ and $`M_{eq}`$ can be estimated from the measured irreversible magnetization data of a superconducting sample using the relations, $$J_c(H)=[M(H)M(H)]/2g\mu _0R,$$ $`(1a)`$ $$M_{eq}(H)=[M(H)+M(H)]/2,$$ $`(1b)`$ where $`M(H)`$ and $`M(H)`$ are the magnetization in the increasing (forward) and decreasing (reverse) field cycles respectively, $`\mu _0=4\pi \times 10^7`$ W/A.m, $`R`$ is the sample dimension transverse to the applied field and $`g`$ is a factor which depends on the sample geometry. Eq. 1 implicitly assumes that $`J_c`$ is history independent and is thus uniquely determined by the local induction $`B`$. However, across the peak effect region, the above equations are not valid due to a strong history dependence in $`J_c`$ . Recently, considerable efforts have gone into ascertaining the equilibrium magnetization across the peak effect region, where an order-disorder transition occurs in the vortex matter. However, these efforts have met with ambiguous and somewhat conflicting results. For example, the construction of the equilibrium magnetization from the hysteresis loop by using two different kinds of minor magnetization curves, results in apparently different conclusions. In one case, a jump in $`M_{eq}`$ could be found at the onset of the PE, while the other case shows no increase at all. These differences apparently originate from the difficulties in establishing an unambiguous and reproducible vortex state due to a strongly history dependent configuration of the vortex matter in the PE region. The different procedures proposed to obtain $`M_{eq}`$ shall be discussed in section II. In Sec. III, we briefly discuss a recent phenomenological model, which addresses the issue of the history dependent $`J_c`$ and the metastability in the vortex state through an extension of the Bean’s critical state model. In Sec. IV, we present an experimental method based on the ideas of the model to obtain a unique ”stable” vortex state in the PE region, which is independent of the past magnetic history. We propose that this state, in effect, is the “stable” or “equilibrium”state and evaluate the critical current density and $`M_{eq}`$ of this state. We further demonstrate that, a sharp change in the equilibrium magnetization (albeit smeared) occurs across the PE region. These results imply that an underlying first order phase transition, presumably driven by a competition between elastic and pinning energies in a situation where thermal fluctuations are weak, marks the peak effect. ## II MINOR CURVES AND THE EQUILIBRIUM MAGNETIZATION ACROSS THE PEAK EFFECT In the peak effect region, the critical current density in the increasing field cycle $`J_c(H)`$ is less than that ($`J_c(H)`$) in the decreasing field cycle for $`H<H_p`$, where $`H_p`$ is the field where $`J_c`$ is maximum. However, well below the onset of the PE and at $`H>H_p`$, $`J_c`$ is independent of the magnetic history. One of its consequences is the peculiar behavior of the minor magnetization curves, which can not be reconciled within the critical state model. For instance, a typical minor magnetization curve (type I) initiated from a field $`H<H_p`$ in the PE region saturates without meeting the reverse magnetization curve, as shown in Fig. 1(a). On the other hand, the minor curves (type II) measured by increasing the field from different points on the reverse magnetization curve overshoot the forward curve as shown in Fig.1(b). The two types of anomalous behavior may be contrasted with the conventional behavior for the minor curves starting at (a) $`H`$ $`>`$ $`H_p`$ and (b) H$`<<`$H<sub>pl</sub>, i.e.,for fields well below the PE region. The latter catagories of minor curves meet the magnetization envelope, constituted by the forward and reverse curves, as expected from the Bean’s critical state model. A new procedure was proposed by Roy and Chaddah to obtain $`M_{eq}`$ from the minor magnetization curves of the type I by the relation, $$M_{eq}(H)=[M(H+\delta ,)+M_{ML}(H\delta ,)]/2,$$ $`(2)`$ where $`M(H+\delta ,)`$ is the magnetization at a field $`H+\delta `$ (denoted by point A in Fig. 1(a)) from where the minor curve is initiated on the forward curve. $`M_{ML}(H\delta ,)`$ is the magnetization on the minor curve at a field $`H\delta `$, where it saturates as indicated by the point B in Fig. 1(a). This procedure is based on the implicit assumption that the vortex state formed on the forward curve is an ”equilibrium” state. This assumption is however inconsistent with the experimental observation by Wordenweber, Kes and Tsuei, who showed that both current cycling and field cycling processes eventually establish a vortex state with a $`J_c`$ higher than that on the forward curve. Such an observation indicates that the vortex state formed on the forward curve is metastable in nature. Tenya et al have preferred a procedure given below, which is very similar to the one described above but using the minor curves of the type II described in Fig. 1(b): $$M_{eq}(H)=[M(H\delta ,)+M_{ML}(H+\delta ,)]/2,$$ $`(3)`$ where $`H\delta `$ (point C in Fig. 1(b)) is the field from where the minor curve is initiated on the reverse curve and $`H+\delta `$ (point D in Fig. 1(b)) is the field where it saturates. $`M_{ML}(H+\delta ,)`$ is the saturated magnetization value on the minor curve. This procedure too has the shortcoming similar to that in Eq. 2, viz., the vortex state on the reverse magnetization curve is actually a metastable state. Moreover, not only are these recipes deficient, they also yield different conclusions, viz., an enhancement in equilibrium magnetization is observed in one case, whereas it is absent in the other. These ambiguities point to the need to evolve a more satisfactory procedure to arrive at a unique and stable vortex state unambiguously and determine the equilibrium magnetization assuming the stable state to be the equilibrium state. ## III MODEL FOR HISTORY EFFECTS AND METASTABILITY Ravikumar et al incorporated the history dependence in the macroscopic critical current density $`J_c`$ by postulating, $$J_c(B+\mathrm{\Delta }B)=J_c(B)+(\|\mathrm{\Delta }B\|/B_r)(J_c^{st}J_c).$$ $`(4)`$ where the critical current density $`J_c(B)`$ is a macroscopic representation of a particular metastable configuration of the vortex lattice at a field $`B`$. Eq. 4 describes how the vortex state evolves from one metastable configuration to another. An important assumption of this model is the existence of a stable vortex state with a critical current density $`J_c^{st}`$, which is unique for a given field and temperature. $`B_r`$ is a macroscopic measure of metastability and describes how strongly $`J_c`$ could be history dependent. In the limit of $`B_r`$ tending to zero, however, this model reduces to the standard critical state model for which $`J_c`$ ($`=J_c^{st}`$) is independent of the magnetic history. It can be seen from Eq. 4 that a metastable vortex state with $`J_c`$ $``$ $`J_c^{st}`$, can be driven into a stable state by merely oscillating the field by a small amplitude (see Fig. 1 of the Ref. 20). In the PE region, the energy barriers between different metastable vortex configurations are much greater than the available thermal energy. The field cycling allows the vortices to move and explore the energy landscape and thereby rearrange in to a vortex configuration closer to the stable state. In the next section, we will demonstrate this experimentally and show that the stable state obtained is indeed independent of the magnetic history. In the limit $`\mathrm{\Delta }B0`$, Eq. 4 can be rewritten in the form, $$\pm dJ_c/dB=(J_c^{st}J_c)/B_r,$$ $`(5)`$ where upper and lower signs are applicable in the cases of increasing and decreasing local field $`B`$, respectively. In each case, the $`J_c(B)`$ can be obtained by solving Eq. 5, provided the functional form of $`J_c^{st}(B)`$ and $`B_r(B)`$ are known. We assume for $`J_c^{st}(B)`$ and $`B_r(B)`$ the following forms used in Ref. for calculating the minor magnetization curves: $$J_c^{st}(B)=J_{c1}(1B/\mu _0H_1)+J_{c2}e^{(B\mu _0H_p)^2/2\mu _0H_W^2}$$ $`(6)`$ and $$B_r(B)(B\mu _0H_{low})^m(\mu _0H_pB)^nforH_{low}<B/\mu _0<H_p$$ $$0otherwise$$ $`(7),`$ The first term in Eq. 6 is the field dependence of $`J_c^{st}`$ well below the peak and the second term reflects the peak in $`J_c^{st}`$ $`vs`$ $`B`$. $`B_r(B)`$ in Eq. 7 accounts for the observed history dependence in $`J_c`$ in the PE region. $`B_r=0`$ in the field ranges $`H<H_{low}`$ and $`H>H_p`$ signifies that $`J_c`$ is independent of the magnetic history and is always equal to the $`J_c^{st}`$. For the two limiting cases, $`H<H_{low}`$ and $`H>H_p`$, the intervortex interaction and the flux pinning are dominant respectively and therefore the stable state is readily accessed by the vortex lattice. The values of the different parameters used in this paper are listed in the caption of Fig. 2. $`J_c(H)`$ \[$`J_c(H)`$\] is calculated by numerically solving Eq. 5 with the upper (lower) sign with the initial condition $`J_c(H)`$ \[$`J_c(H)`$\] $`=`$ $`J_c^{st}(H)`$ at some field below $`H_{low}`$ (above $`H_p`$). In Fig. 2(a), we present an evaluation of $`J_c(H)`$ and $`J_c(H)`$ which obey the inequality $`J_c(H)<J_c^{st}(B)<J_c(H)`$. It was earlier interpreted that the vortex state formed on the decreasing field cycle is a supercooled disordered state. In other words, the vortex state formed in decreasing field (from above $`H_p`$) retains the memory of the vortex correlations from the previous fields. In analogy, we can argue that the vortex state formed on the increasing field cycle is a superheated ordered state. Both of these states are metastable in nature. As argued above, they can be driven into a stable state by oscillating the external field by a small amplitude. The magnetization hysteresis loop corresponding to $`J_c(H)`$ and $`J_c(H)`$ are shown in Fig. 2(b). Note the asymmetry in the hysteresis, usually observed in experiments. For a comparison, we also plot the magnetization hysteresis loop one would obtain within the framework of Bean’s critical state model with $`J_c=J_c^{st}`$ (applicable in the limit $`B_r0`$) which is symmetric in the forward and reverse field cycles, as shown by the dotted line in Fig. 2(b). Details of the magnetization calculation are described in Ref. 20. The minor magnetization curves of the types I and II calculated in the slab geometry, are shown in Fig. 3(a) and Fig. 3(b) respectively. They clearly mimic the behavior seen in experiments. We assumed $`M_{eq}(H)=0`$, in calculating these magnetization curves. We note that the calculated curves in Fig. 2 and Fig. 3 are not quantitative fits to experimental data, they only serve to illustrate the qualitative features of the observed data. In Fig. 3(c), we show $`M_{eq}^{}(H)`$ determined from the calculated minor curves of the type I and type II following Eq. 2 and Eq. 3, respectively. The test of the self-consistency of these procedures lies in reproducing the original form ($`M_{eq}=0`$) assumed in the calculation. $`M_{eq}^{}(H)`$ obtained from these two procedures are not only inconsistent with each other, but, also, do not conform to the original assumption that $`M_{eq}=0`$. The procedure of Eq. 2 indeed produces a peak like structure in $`M_{eq}^{}(H)`$ which has been shown earlier from an analysis of experimental data in $`2HNbSe_2`$ following the same recipe. On the other hand, the use of Eq. 3 proposed by Tenya et al yields no variation in $`M_{eq}^{}`$ vs $`H`$ across the PE region. Fig. 3(c) illustrates the unreliable and ambiguous nature of these recipes noted above and thus points to the need for a consistent approach in order to overcome their difficulties. ## IV EXPERIMENTAL RESULTS AND DISCUSSION In this section, we will show experimentally that repeated field cycling drives any metastable state into a stable state, which is unique at a given field. We study the minor hysteresis loops traced by repeated field cycling and infer from these measurements the critical current density $`J_c^{st}`$ and the equilibrium magnetization $`M_{eq}`$ of the stable state. DC magnetization measurements have been carried out using a Quantum Design (QD) Inc. SQUID magnetometer (Model MPMS5) in the peak effect region of a $`2H`$-$`NbSe_2`$ single crystal ($`T_c`$ $``$ 7.25 K) with the field applied parallel to its c-axis. The crystal is of approximate dimensions ($`a\times b\times c`$) $`4mm\times 5mm\times 0.43mm`$. As stated earlier, the peak effect in $`J_c`$ is manifested as the anomalous enhancement in the magnetization hysteresis (c.f. Fig. 1). The magnetization hysteresis has been studied at different temperatures from 6.7 to 6.95K. Magnetization hysteresis data at 6.95K was measured using a 2 cm full scan length, and the data at the other temperatures was obtained using the half-scan technique to avoid artefacts arising due to field inhomogeneity experienced by the sample along the scan length. In the temperature range investigated, $`J_c`$ at the peak field $`H_p`$ decreases with decreasing temperature (see Table 1). Fig. 4(a) depicts a part of the hysteresis loop at 6.95K, constituting $`M`$ $`vs`$ $`H`$ curves in the increasing (forward) and decreasing (reverse) field cycles measured with a 30 sec wait time at each field. We identify the onset field $`H_{pl}^+`$ of the PE on the forward curve, where $`M`$ begins to decrease sharply. The field $`H_p`$ marks the field at which magnetization hysteresis is maximum. In Fig. 4(a), we show the points A, B, C and D from where the minor hysteresis loops are initiated. A(C) and B(D) are at a field $`H<H_{pl}^+`$ ($`H>H_{pl}^+`$) on the forward and reverse curves, respectively. Minor hysteresis loops starting from both forward and reverse curves are recorded at different fields (spanning the peak region) by repeatedly cycling the field by a small amplitude $`\mathrm{\Delta }H`$. The interval $`\mathrm{\Delta }H`$ is chosen such that it is above the threshold field required to reverse the direction of the shielding currents throughout the sample. From the critical state model, we understand that magnetization values must always remain confined within the forward and reverse magnetization curves, which constitute the so called magnetization envelope. Further, the $`MH`$ loop in each field cycle must retrace itself. In Fig. 4(b), we show the minor hysteresis loops (MHLs) measured by repeatedly cycling the field, starting at point A ($`H<H_{pl}^+`$) on the forward curve. These MHLs in different field cycles retrace each other indicating that the $`J_c`$ does not change with field cycling. Therefore, we conclude that the vortex state is in a stable configuration. In contrast, the MHLs shown in Fig. 4(c) (continuous line with data points omitted) starting at B ($`H<H_{pl}^+`$) on the reverse curve, show shrinkage effects, with each successive field cycle and finally MHL collapses into the minor loop started from point A (open circles) which is replotted in Fig. 4(c). This suggests that the vortex configuration at point B is metastable with a $`J_c>J_c^{st}`$. Repeated field cycling causes the $`J_c`$ to fall towards the stable stable value as reflected in the reduction of the width of the MHL with each successive field cycle. It is remarkable that the minor loops starting from both A and B merge into precisely the same loop within the experimental accuracy. This clearly reaffirms the basic assumption of the model that there exists a unique stable state with critical current density $`J_c^{st}`$, independent of the initial vortex state from which it evolves. We now focus on the behavior of MHLs which start from a field $`H>H_{pl}^+`$. As shown in Fig. 4(d), the behavior of the minor loops starting at point C is quite different from those started at point A. The increasing field leg of the MHL moves away from the forward magnetization curve in the first field cycle itself and remains outside the magnetization envelope for subsequent field cycles. This clearly suggests that, for $`H>H_{pl}^+`$, the vortex configuration even on the forward magnetization curve is metastable. However, the behavior of the MHLs starting at point D on the reverse magnetization curve is very similar to the behavior of those that start at point B, i,e., the MHL shrinks with each successive field cycle (continuous line in Fig. 4(e)). The data in Fig. 4(d) is replotted in Fig. 4(e) (open circles connected by dotted line), which suggests that the MHLs starting from both C and D collapse into the same final loop (MHL). Firstly, the data in Fig. 4 clearly suggest the metastable nature of the vortex configuration for fields above $`H_{pl}^+`$ both on the forward and the reverse magnetization curves. Further, the eventual MHL obtained on repeated field cycling is independent of the initial vortex configuration. We note that the metastable state on the forward magnetization curve settles into the stable state much faster than that on the reverse curve. This might imply that the vortex configuration on the increasing field cycle is closer to the equilibrium configuration. The data in Fig. 4 yield the following inequalities for the critical currents in the different field ranges: (i) For $`H<H_{pl}^+`$, the vortex configuration is stable in the increasing field cycle while at the same field value, it is highly metastable in the decreasing field cycle. This can be summarized by the inequality $`J_c(H)=J_c^{st}(H)<J_c(H)`$; (ii) For $`H_{pl}^+<H<H_p`$, the vortex configurations in both increasing and decreasing field cycles are metastable, with the critical currents obeying the inequality, $`J_c(H)<J_c^{st}(H)<J_c(H)`$; (iii) For $`H>H_p`$, $`J_c(H)=J_c(H)=J_c^{st}(H)`$. These observations are in accordance with the model (cf. Fig. 2(a)). We thus assert that the Eq. 2, proposed by Roy and Chaddah is applicable only for $`H<H_{pl}^+`$. It is unsatisfactory for $`H_{pl}^+<H<H_p`$, as the vortex lattice on the forward curve is in a superheated vortex configuration which is more ordered (but metastable) than the stable configuration. Eq. 3, as proposed by Tenya et al is not appropriate in any of the field ranges because the vortex states produced on the reverse curve are supercooled vortex configurations which are more disordered than the corresponding stable states. Fig. 5 shows the M-H loop at 6.9K constituting the forward and reverse magnetization curves (dark line with data points omitted) indicating $`H_{pl}^+`$ and $`H_p`$. Note the asymmetry (also seen at 6.95K) in the forward and reverse magnetization curves which is the hall mark of the peak effect. We also measured the MHLs by repeatedly cycling the field starting at different points on the forward and reverse curves. The saturated MHLs are again found to be independent of the initial vortex state just as for 6.95K. The locus of magnetization values on the increasing and decreasing field legs of the saturated MHLs measured at different fields are also plotted in Fig. 5 (open circles connected by dotted line). This observed behavior is in excellent qualitative agreement with that expected from the model in Ref. 20 (see Fig. 2(b)). The locus of saturated magnetization values corresponds to the “table” or the “quilibrium” vortex configuration at different fields. Having established the existence of a history independent stable state we determine the critical current density $`J_c^{st}`$ and the equilibrium magnetization $`M_{eq}`$ state at each field from the saturated MHL. $`J_c^{st}`$ and $`M_{eq}`$ are given by, $$J_c^{st}(H)=[M_{st}(H)M_{st}(H)]/2g\mu _0R,$$ $`(8a)`$ $$M_{eq}(H)=[M_{st}(H)+M_{st}(H)]/2,$$ $`(8b)`$ where $`M_{st}(H)`$ and $`M_{st}(H)`$ are the magnetization values on the increasing and decreasing field legs of the saturated MHL. $`J_c^{st}`$ vs $`H`$ and $`M_{eq}`$ vs $`H`$ data at 6.95K are plotted in Fig. 6(a) and Fig. 6(b), respectively. $`M_{eq}`$ exhibits a sharp increase between $`H_{pl}^+`$ and $`H_p`$ signifying an increase in the equilibrium flux density. This is reminiscent of the characteristic of $`M_{eq}`$ across the FLL melting transition observed in cuprate superconductors. We argue that the change in $`M_{eq}`$ indicates a first order transition in the FLL from an ordered solid to a pinned amorphous state presumably analogous to a Bragg Glass to Vortex Glass/pinned liquid phase transition. The increase in $`M_{eq}`$ coincides with the increase in $`J_c^{st}`$ near the onset of the peak effect and spans the field range between $`H_{pl}^+`$ and $`H_p`$. In Fig. 7(a) and 7(b), we present the $`M_{eq}`$ vs $`H`$ and $`J_c^{st}`$ vs $`H`$ data respectively, at 6.9K. We note that the sharp change in $`M_{eq}`$ correlates with a sharp increase in $`J_c^{st}`$ between $`H_{pl}^+`$ and $`H_p`$. We also present the $`\mathrm{\Delta }M_{eq}`$ values obtained at different temperatures in Table 1. It is important to understand the nature of the vortex state in the transition region $`H_{pl}^+<H<H_p`$. One of the well known pictures is the collective pinning scenario, where the loss of long range order is expected to permeate uniformly throughout the sample. On the other hand Paltiel et al have recently proposed a picture where the disordered phase enters through surface imperfections and coexists near the surface with the ordered phase of the bulk. They argue that the boundary between the disordered region and the ordered region moves into the sample as the temperature (or field) is increased towards $`T_p`$ (or $`H_p`$). Further possibility is the coexistence of ordered and disordered phases, with an intricate geometrical connectivity of these phases. Irrespective of the particular picture used, our experiments demonstrate a specific and an unambiguous procedure, viz., subjecting the sample to a field cycling, to produce a unique stable state (in a macroscopic sense) across the peak effect region. We consider this stable state as a pinned equilibrium state, and estimate equilibrium magnetization and the free energy difference or entropy change when the vortex lattice changes from an ordered to an amorphous state. As per the Clausius-Clapeyron relation, the entropy change per vortex per inter-layer distance $`d`$ ($``$ 4 Angstroms) in the $`2HNbSe_2`$ system, $$\mathrm{\Delta }s=(\mathrm{\Delta }M_{eq}/H_p)(dH_{pl}^+/dT)(\varphi _0d/k_B),$$ where $`dH_{pl}^+/dT`$ $``$ $`dH_p/dT`$ $``$ $`0.65`$ T/K. The value of $`\mathrm{\Delta }s`$ estimated at different temperatures is tabulated in Table 1. Incidentally these values are comparable to the entropy change reported across the FLL melting transition in high $`T_c`$ cuprates. An important question that can arise is whether the entropy change can be observed in thermal measurements such as specific heat vs temperature. We recall that the metastability in the vortex state is much greater in temperature scans in a fixed magnetic field. Repeated cycling of the field by a small amplitude may be necessary to produce the ”stable” or ”equilibrium” state before a thermal measurement is carried out at each temperature. ## V CONCLUSIONS In this paper, we have presented a study of the different metastable vortex configurations occuring in the peak effect region of a weakly pinned superconductor $`2HNbSe_2`$ through magnetization measurements. Each metastable vortex configuration is characterized by a critical current density $`J_c`$ which is strongly dependent on the magnetic history. It is also shown that any metastable vortex configuration obtained under given field historys can be driven into a stable configuration by repeated field cycling. This stable configuration has a critical current density $`J_c^{st}`$, uniquely determined by field and temperature as postulated in a recent model. Field cycling appears to act as an effective temperature to drive a metastable state into the stable state, even when thermal energy itself is inadequate to sample the phase space and access the stable state. The method of recording minor hysteresis loops described here allows us to determine the pinning and equilibrium properties of the stable vortex state satisfactorily. Our equilibrium magnetization data clearly suggest that the transition of the vortex lattice from an ordered state to a disordered state is first order in nature. The smearing of the transition, i.e., the width of the transition region may be a manifestation of the spatially inhomogeneous pinning of the system. The $`J_c^{st}`$ data suggests that the loss of quasi-long range order in the vortex lattice also spans the same field window as the magnetization jump. In the collective pinning picture, this amounts to correlation volume of the vortex phase decreasing in this regime and the FLL becoming completely disordered above $`H_p`$ or $`T_p`$. The precise coincidence of the $`J_c`$ anomaly with the equilibrium magnetization anomaly further illustrates the self consistency of the procedure developed here. It would be interesting to compare the nature of this disorder-driven transition in systems with different types of pinning, e.g. high density of point pins versus low density of extended pins to further understand the nature of this presumably disorder induced phase transformation. The authors thank Dr. K. V. Bhagwat, Dr. T. V. Chandrasekhar Rao, Dr. P. K. Mishra and Mr. M. R. Singh for discussions. FIGURE CAPTIONS Fig.1. : Typical magnetization hysteresis loop observed in the peak effect region of a superconducting $`2HNbSe_2`$. In the panel (a), the minor curve obtained by decreasing the field from the point A, corresponding to a field ($`H+\delta `$) on the forward magnetization curve is shown to saturate at the point B, which corresponds to the field ($`H\delta `$). Magnetization values at A and B are $`M(H+\delta ,)`$ and $`M_{ML}(H\delta ,)`$ respectively (see text). In the panel (b) the minor curve obtained by increasing the field from the point C, corresponding to a field ($`H\delta `$) on the reverse magnetization curve, saturates at D ($`H+\delta `$) and corresponds to a magnetization value $`M_{ML}(H+\delta ,)`$. Fig. 2: (a) Calculated critical current densities $`J_c(H)`$ and $`J_c(H)`$ in the increasing and decreasing field cases, respectively. These are compared with the stable critical current density $`J_c^{st}`$ (dotted line). In this calculation, we have used $`H_{low}=0.05T`$, $`H_p=0.1T`$, $`J_{c1}`$ $`=`$ $`10^4A/m^2`$, $`J_{c2}`$ $`=`$ $`20J_{c1}`$, $`H_1`$ $`=`$ $`0.12T`$ and $`H_W`$ $`=0.008T`$ . (b) Magnetization hysteresis loop corresponding to the $`J_c`$ values shown in (a). The hysteresis loop that would be obtained within the framework of critical state model, i.e., in the limit of $`B_r0`$ is also shown in the in this panel as dotted line. The inset shows the functional form of $`B_r`$ which is non-zero in the field range $`H_{low}<H<H_p`$. Fig. 3: Calculated minor curves of type I and type II are shown in pnaels (a) and (b), respectively. In the panel (c), we show the $`M_{eq}^{}`$ vs $`H`$ obtained using Eq. 2 and Eq. 3, respectively along with the original form, $`M_{eq}=0`$, assumed in the calculation of the minor curves. Fig. 4: (a) A part of the magnetization loop (forward and reverse curves) measured at 6.95K on a $`2HNbSe_2`$ single crystal. Also indicated are the characteristic fields, $`H_{pl}^+`$ and $`H_p`$. We indicate A and B ($`H<H_{pl}^+`$) and C and D ($`H_{pl}^+<H<H_p`$) starting from which the minor hysteresis loops are measured. (b) Minor hysteresis loops started from point A (open circles). In different field cycles, they are seen to retrace the same loop. (c) The MHL started from B (continuous line) shrinks with each successive field cycle. The increasing and decreasing field legs of the first and second cycles are numbered. After five field cycles, the hysteresis loop is seen to merge with the loop shown in (b), which is replotted (open circles). (d) Minor hysteresis loops started from point C (open circles). In the first field cycle itself, increasing field leg of the MHL moves away from the forward curve and remains outside the magnetization envelope for the subsequent field cycles. (e) The minor loops starting from D (continuous line) are seen to collapse onto the loop shown in (d), which is replotted. Fig. 5: Magnetization hysteresis loop of $`2HNbSe_2`$, recorded using half scan technique at 6.9K (continuous line). The open circles are the saturated magnetization values obtained after repeated field cycling. $`H_{pl}^+`$ and $`H_p`$ are also marked. The locus of the saturated magnetization values is shown as a dotted line. Fig. 6: (a) Stable critical current density $`J_c^{st}`$ in the field range 80 mT $`<`$ $`\mu _0H`$ $`<`$ 105 mT. In the inset, we show the $`J_c^{st}`$ $`vs`$ $`H`$ in the entire field range. Filled triangles and open circles correspond to the values obtained from the MHLs intiated from the forward and reverse magnetization curves respectively. (b) Equilibrium magnetization $`M_{eq}`$ as a function of field at 6.95K. Note that the sharp change in $`M_{eq}`$ coincides with the PE onset field $`H_{pl}^+`$. Fig. 7: (a) Critical current density $`J_c^{st}`$ vs $`H`$ and (b) $`M_{eq}`$ vs $`H`$ data obtained at 6.9K. Note that the smeared jump in $`M_{eq}`$ vs $`H`$, as marked by the double sided arrow, agrees precisely with a smeared jump in critical current density $`J_c^{st}`$ in the peak regime. See text for a discussion.
warning/0007/hep-ph0007055.html	ar5iv	text	# Factorization Scheme and Parton Distributions in the Polarized Virtual Photon Target ## 1 Introduction In the two-photon process of $`e^+e^{}`$ collision experiments, we can measure the structure functions of the virtual photon (Fig.1). The advantage in studying the virtual photon target is that, in the case $$\mathrm{\Lambda }^2P^2Q^2$$ (1.1) where $`Q^2`$ ($`P^2`$) is the mass squared of the probe (target) photon, and $`\mathrm{\Lambda }`$ is the QCD scale parameter, we can calculate the whole structure function up to the next-to-leading order (NLO) in QCD by the perturbative method, in contrast to the case of the real photon target where in NLO there exist non-perturbative pieces . The spin-independent structure functions $`F_2^\gamma (x,Q^2,P^2)`$ and $`F_L^\gamma (x,Q^2,P^2)`$ as well as the parton contents were studied in the leading order (LO) and in NLO -. The target mass effect of unpolarized and polarized virtual photon structure in LO was discussed in Ref.. The information on the spin structure of the photon would be provided by the resolved photon process in polarized version of the DESY electron and proton collider HERA . More directly, polarized photon structure function can be measured by the polarized $`e^+e^{}`$ collision in the future linear colliders. For the real photon ($`P^2=0`$) target, there exists only one spin-dependent struture function, $`g_1^\gamma (x,Q^2)`$, which is equivalent to the structure function $`W_4^\gamma (x,Q^2)`$ ($`g_1^\gamma 2W_4^\gamma `$) discussed some time ago in . The LO QCD corrections to $`g_1^\gamma `$ for the real photon target was first calculated by one of the authors and later in Refs., while the NLO QCD analysis was performed by Stratmann and Vogelsang . The first moment of the photon structure function $`g_1^\gamma `$ has recently attracted attention in the literature in connection with its relevance for the axial anomaly. More recently the present authors investigated the spin-dependent structure function $`g_1^\gamma (x,Q^2,P^2)`$ of the virtual photon up to NLO in QCD, where $`P^2`$ is in the above kinematical region (1.1). The analysis was made in the framework of the operator product expansion (OPE) supplemented by the renormalization group method and also in the framework of the QCD improved parton model using the DGLAP parton evolution equations. In the past few years, accuracy of the experimental data on the spin-dependent structure function $`g_1`$ of the nucleon has been significantly improved . Using these experimental data together with the already existing world data, several groups - have carried out the NLO QCD analyses on the polarized parton distributions in the nucleon. These parton distributions may be used for predicting the behaviors of other processes such as polarized Drell-Yan reactions and polarized semi-inclusive deep inelastic scatterings, and etc. However, parton distributions obtained from the NLO analyses are dependent on the factorization scheme employed. It is possible that parton distributions obtained in one scheme may be more appropriate to use than those in other schemes. In the case of nucleon target, however, it may be difficult to examine the features of each factorization scheme, since for the moment it is inevitable to resort to some assumptions in order to extract parton distributions from the experimental data. On the other hand, it is remarkable that, in the case of virtual photon target with a virtual mass $`P^2`$ being in the kinematical region of Eq.(1.1), not only the photon structure functions but also the parton distributions in the target can be predicted entirely up to NLO in QCD. Thus, comparing the parton distributions predicted by one scheme with those by other schemes, we can easily examine the features of each factorization scheme. In consequence, the virtual photon target may serve as an optimal place to study the behaviors of parton distributions and their factorization-scheme dependence. In this paper we examine in detail the polarized parton (i.e., quark and gluon) distributions in the virtual photon target. The polarized parton distributions are particularly interesting due to the fact that they have relevance to the axial anomaly . The interplay between the QCD axial anomaly and factorization schemes has already been discussed for the spin-dependent structure function $`g_1`$ of the nucleon -. It was explained there that the QCD axial anomaly effect is retained in the flavor-singlet quark distribution in the nucleon in the standard $`\overline{\mathrm{MS}}`$ scheme, but it is shifted to the gluon coefficient function in such a scheme called chirally-invariant (CI) factorization scheme. Now it should be pointed out that the polarized photon target is unique in the sense that not only QCD but also the OED axial anomaly takes place. The QED axial anomaly, which is U(1) anomaly, emerges when a quark has an electromagnetic charge. Thus the flavor-non-singlet quark distribution is also relevant, besides the flavor-singlet one. Depending upon factorization schemes, the QED axial anomaly effect resides in both the flavor-singlet and non-singlet polarized quark distributions in the virtual photon, or it is shifted to the photon coefficient function, in which case we arrive at an interesting result: the first moments of the polarized quark distributions in the virtual photon, both flavor singlet and non-singlet, vanish in NLO. Also we find that the large $`x`$-behaviors of polarized quark distributons dramatically vary from one factorization scheme to another. Indeed, for $`x1`$, the quark distributions positively diverge or negatively diverge or remain finite, depending on factorization schemes. We perform our analyses in six different factorization schemes, (i) $`\overline{\mathrm{MS}}`$, (ii) CI (chirally invariant) (it is also called as JET) , (iii) AB (Adler-Bardeen) , (iv) OS (off-shell) , (v) AR (Altarelli-Ross) , and finally (vi) DIS<sub>γ</sub> schemes , and see how the parton distributions change in each scheme. In particular, we study the axial anomaly effects on the first moments and the large-$`x`$ behaviors of parton distributions. Gluon distribution in the virtual photon is found to be the same up to NLO, at least among the factorization schemes considered in this paper. Furthermore, the first moment of gluon distribution turns out to be factorization-scheme independent up to NLO. Part of the result has been briefly reported elsewhere . In the next section we discuss the polarized parton distributions in the virtual photon. The explicit expressions for the flavor singlet-(non-singlet-)quark and gluon distributions predicted in QCD up to NLO are given in Appendix A. In Sec. 3, we derive the transformation rules for the relevant two-loop anomalous dimensions and one-loop photon matrix elements from the $`\overline{\mathrm{MS}}`$ scheme to other factorization schemes and then explain particular factorization schemes we consider in this paper. In Sec. 4, we examine the first moments of parton distributions with emphasis on the interplay between the QCD and QED axial anomalies and the facorization schemes. The behaviors of parton distributions near $`x=1`$ and their factorization-scheme dependence are discussed in Sec. 5. The numerical analyses of parton distributions predicted by different factorization schemes will be given in Sec. 6. The final section is devoted to the conclusion and discussion. ## 2 Polarized parton distributions in photon Let $`q_\pm ^i(x,Q^2,P^2)`$, $`G_\pm ^\gamma (x,Q^2,P^2)`$, $`\mathrm{\Gamma }_\pm ^\gamma (x,Q^2,P^2)`$ be quark with $`i`$-flavor, gluon, and photon distribution functions with $`\pm `$ helicities of the longitudinally polarized virtual photon with mass $`P^2`$. Then the spin-dependent parton distributions are defined as $`\mathrm{\Delta }q^iq_+^i+\overline{q}_+^iq_{}^i\overline{q}_{}^i`$ , $`\mathrm{\Delta }G^\gamma G_+^\gamma G_{}^\gamma `$ , and $`\mathrm{\Delta }\mathrm{\Gamma }^\gamma \mathrm{\Gamma }_+^\gamma \mathrm{\Gamma }_{}^\gamma `$ . In the leading order of the electromagnetic coupling constant, $`\alpha =e^2/4\pi `$, $`\mathrm{\Delta }\mathrm{\Gamma }^\gamma `$ does not evolve with $`Q^2`$ and is set to be $`\mathrm{\Delta }\mathrm{\Gamma }^\gamma (x,Q^2,P^2)=\delta (1x)`$. For later convenience we use, instead of $`\mathrm{\Delta }q^i`$, the flavor singlet and non-singlet combinations of spin-dependent quark distributions as follows: $`\mathrm{\Delta }q_S^\gamma `$ $``$ $`{\displaystyle \underset{i}{}}\mathrm{\Delta }q^i,`$ $`\mathrm{\Delta }q_{NS}^\gamma `$ $``$ $`{\displaystyle \underset{i}{}}e_i^2\left(\mathrm{\Delta }q^i{\displaystyle \frac{\mathrm{\Delta }q_S^\gamma }{N_f}}\right).`$ (2.1) In terms of these parton distributions, the polarized virtual photon structure function $`g_1^\gamma (x,Q^2,P^2)`$ is expressed in the QCD improved parton model as $`g_1^\gamma (x,Q^2,P^2)`$ $`=`$ $`{\displaystyle _x^1}{\displaystyle \frac{dy}{y}}\{\mathrm{\Delta }q_S^\gamma (y,Q^2,P^2)\mathrm{\Delta }C_S^\gamma ({\displaystyle \frac{x}{y}},Q^2)+\mathrm{\Delta }G^\gamma (y,Q^2,P^2)\mathrm{\Delta }C_G^\gamma ({\displaystyle \frac{x}{y}},Q^2)`$ (2.2) $`+\mathrm{\Delta }q_{NS}^\gamma (y,Q^2,P^2)\mathrm{\Delta }C_{NS}^\gamma ({\displaystyle \frac{x}{y}},Q^2)\}+\mathrm{\Delta }C^\gamma _\gamma (x,Q^2).`$ where $`\mathrm{\Delta }C_S^\gamma `$($`\mathrm{\Delta }C_{NS}^\gamma `$), $`\mathrm{\Delta }C_G^\gamma `$, and $`\mathrm{\Delta }C_\gamma ^\gamma `$ are the coefficient functions corresponding to singlet(non-singlet)-quark, gluon, and photon, respectively, and are independent of $`P^2`$. The Mellin moments of $`g_1^\gamma `$ is written as $$g_1^\gamma (n,Q^2,P^2)=\mathrm{\Delta }\stackrel{~}{𝑪}^\gamma (n,Q^2)\mathrm{\Delta }\stackrel{~}{𝒒}^\gamma (n,Q^2,P^2),$$ (2.3) where $`\mathrm{\Delta }\stackrel{~}{𝑪}^\gamma (n,Q^2)`$ $`=`$ $`(\mathrm{\Delta }C_S^\gamma ,\mathrm{\Delta }C_G^\gamma ,\mathrm{\Delta }C_{NS}^\gamma ,\mathrm{\Delta }C_\gamma ^\gamma ),`$ $`\mathrm{\Delta }\stackrel{~}{𝒒}^\gamma (n,Q^2,P^2)`$ $`=`$ $`(\mathrm{\Delta }q_S^\gamma ,\mathrm{\Delta }G^\gamma ,\mathrm{\Delta }q_{NS}^\gamma ,\mathrm{\Delta }\mathrm{\Gamma }^\gamma ),`$ and the matrix notation is implicit. The explicit expressions of $`\mathrm{\Delta }q_S^\gamma `$, $`\mathrm{\Delta }G^\gamma `$, and $`\mathrm{\Delta }q_{NS}^\gamma `$ up to NLO can be derived from Eq.(4.46) of Ref., which are given in Appendix A. They are written in terms of one-(two-) loop hadronic anomalous dimensions $`\mathrm{\Delta }\gamma _{ij}^{0,n}`$ ($`\mathrm{\Delta }\gamma _{ij}^{(1),n}`$) ($`i,j=\psi ,G`$) and $`\mathrm{\Delta }\gamma _{NS}^{0,n}`$ ($`\gamma _{NS}^{(1),n}`$), one-(two-) loop anomalous dimensions $`\mathrm{\Delta }K_i^{0,n}`$ ($`\mathrm{\Delta }K_i^{(1),n}`$) ($`i=\psi ,G,NS`$) which represent the mixing between photon and three hadronic operators $`R_i^n`$ ($`i=\psi ,G,NS`$), and finally $`\mathrm{\Delta }A_i^n`$, the one-loop photon matrix elements of hadronic operators renormalized at $`\mu ^2=P^2(=p^2)`$, $$\gamma (p)R_i^n(\mu )\gamma (p)\|_{\mu ^2=P^2}=\frac{\alpha }{4\pi }\mathrm{\Delta }A_i^n(i=\psi ,G,NS).$$ (2.4) The photon matrix elements $`\mathrm{\Delta }A_i^n`$ are scheme-dependent. In one-loop order, they are given, in the $`\overline{\mathrm{MS}}`$ scheme, by $`\mathrm{\Delta }A_{\psi ,\overline{\mathrm{MS}}}^n`$ $`=`$ $`{\displaystyle \frac{e^2}{e^4e^2^2}}\mathrm{\Delta }A_{NS,\overline{\mathrm{MS}}}^n`$ (2.5) $`=`$ $`12e^2N_f\left\{{\displaystyle \frac{n1}{n(n+1)}}S_1(n)+{\displaystyle \frac{4}{(n+1)^2}}{\displaystyle \frac{1}{n^2}}{\displaystyle \frac{1}{n}}\right\},`$ $`\mathrm{\Delta }A_{G,\overline{\mathrm{MS}}}^n`$ $`=`$ $`0,`$ where $`S_1(n)=_{j=1}^n\frac{1}{j}`$ . ## 3 Factorization schemes ### 3.1 Transformation rules from $`\overline{\mathrm{MS}}`$ scheme to $`a`$-scheme Although $`g_1^\gamma `$ is a physical quantity and thus unique, there remains a freedom in the factorization of $`g_1^\gamma `$ into $`\mathrm{\Delta }\stackrel{~}{𝑪}^\gamma `$ and $`\mathrm{\Delta }\stackrel{~}{𝒒}^\gamma `$. Given the formula Eq.(2.3), we can always redefine $`\mathrm{\Delta }\stackrel{~}{𝑪}^\gamma `$ and $`\mathrm{\Delta }\stackrel{~}{𝒒}^\gamma `$ as follows : $`\mathrm{\Delta }\stackrel{~}{𝑪}^\gamma (n,Q^2)\mathrm{\Delta }\stackrel{~}{𝑪}^\gamma (n,Q^2)\|_a`$ $``$ $`\mathrm{\Delta }\stackrel{~}{𝑪}^\gamma (n,Q^2)Z_a^1(n,Q^2),`$ (3.1) $`\mathrm{\Delta }\stackrel{~}{𝒒}^\gamma (n,Q^2,P^2)\mathrm{\Delta }\stackrel{~}{𝒒}(n,Q^2,P^2)\|_a`$ $``$ $`Z_a(n,Q^2)\mathrm{\Delta }\stackrel{~}{𝒒}^\gamma (n,Q^2,P^2),`$ (3.2) where $`\mathrm{\Delta }\stackrel{~}{𝑪}^\gamma \|_a`$ and $`\mathrm{\Delta }\stackrel{~}{𝒒}\|_a`$ correspond to the quantities in a new factorization scheme-$`a`$. Note that the coefficient functions and anomalous dimensions are closely connected under factorization. We will study the factorization scheme dependence of parton distribution up to NLO, by which we mean that a scheme transformation for the coefficient functions is considered up to the one-loop order, since a NLO prediction for $`g_1^\gamma `$ is given by the one-loop coefficient functions and anomalous dimensions up to the two-loop order. The most general form of a transformation for the coefficient functions in one-loop order, from the $`\overline{\mathrm{MS}}`$ scheme to a new factorization scheme-$`a`$, is given by $`\mathrm{\Delta }C_{S,a}^{\gamma ,n}`$ $`=`$ $`\mathrm{\Delta }C_{S,\overline{\mathrm{MS}}}^{\gamma ,n}e^2{\displaystyle \frac{\alpha _s}{2\pi }}\mathrm{\Delta }w_S(n,a),`$ $`\mathrm{\Delta }C_{G,a}^{\gamma ,n}`$ $`=`$ $`\mathrm{\Delta }C_{G,\overline{\mathrm{MS}}}^{\gamma ,n}e^2{\displaystyle \frac{\alpha _s}{2\pi }}\mathrm{\Delta }z(n,a),`$ $`\mathrm{\Delta }C_{NS,a}^{\gamma ,n}`$ $`=`$ $`\mathrm{\Delta }C_{NS,\overline{\mathrm{MS}}}^{\gamma ,n}{\displaystyle \frac{\alpha _s}{2\pi }}\mathrm{\Delta }w_{NS}(n,a),`$ (3.3) $`\mathrm{\Delta }C_{\gamma ,a}^{\gamma ,n}`$ $`=`$ $`\mathrm{\Delta }C_{\gamma ,\overline{\mathrm{MS}}}^{\gamma ,n}{\displaystyle \frac{\alpha }{\pi }}3e^4\mathrm{\Delta }\widehat{z}(n,a),`$ where $`e^4=_ie_i^4/N_f`$. The flavor-singlet(nonsinglet) quark coefficient functions are expanded up to the one-loop order as $`\mathrm{\Delta }C_S^{\gamma ,n}`$ $`=`$ $`e^2\left(1+{\displaystyle \frac{\alpha _s}{4\pi }}\mathrm{\Delta }B_S^n+𝒪(\alpha _s^2)\right),`$ (3.4) $`\mathrm{\Delta }C_{NS}^{\gamma ,n}`$ $`=`$ $`1+{\displaystyle \frac{\alpha _s}{4\pi }}\mathrm{\Delta }B_{NS}^n+𝒪(\alpha _s^2),`$ with $`\mathrm{\Delta }B_S^n=\mathrm{\Delta }B_{NS}^n`$ . The $`\mathrm{\Delta }z(n,a)`$ ($`\mathrm{\Delta }\widehat{z}(n,a)`$) term tells how much of the QCD (QED) axial anomaly effect is transferred to the coefficient function in the new factorization scheme. The gluon and photon coefficient functions $`\mathrm{\Delta }C_G^{\gamma ,n}`$ and $`\mathrm{\Delta }C_\gamma ^{\gamma ,n}`$ start from the one-loop order (i.e., from the NLO): $`\mathrm{\Delta }C_G^{\gamma ,n}`$ $`=`$ $`e^2\left({\displaystyle \frac{\alpha _s}{4\pi }}\mathrm{\Delta }B_G^n+𝒪(\alpha _s^2)\right),`$ $`\mathrm{\Delta }C_\gamma ^{\gamma ,n}`$ $`=`$ $`{\displaystyle \frac{\alpha }{4\pi }}3N_fe^4\left(\mathrm{\Delta }B_\gamma ^n+𝒪(\alpha _s)\right).`$ (3.5) In the $`\overline{\mathrm{MS}}`$ scheme, $`\mathrm{\Delta }C_{\gamma ,\overline{\mathrm{MS}}}^{\gamma ,n}`$ has been obtained from $`\mathrm{\Delta }C_{G,\overline{\mathrm{MS}}}^{\gamma ,n}`$, with changes: $`\alpha _s/2\pi (2\alpha /\alpha _s)\times (\alpha _s/2\pi )`$ , $`e^23e^4`$, and $`3`$ is the number of colors. Thus we have $$\mathrm{\Delta }B_{\gamma ,\overline{\mathrm{MS}}}^n=\frac{2}{N_f}\mathrm{\Delta }B_{G,\overline{\mathrm{MS}}}^n.$$ (3.6) Since, in the leading order, coefficient functions are given by $$\mathrm{\Delta }𝑪_{\overline{\mathrm{MS}}}^\gamma \|_{\mathrm{LO}}=\mathrm{\Delta }𝑪_a^\gamma \|_{\mathrm{LO}}=(e^2,0,1,0),$$ (3.7) the relations (3.3) between the coefficient functions in the $`a`$-scheme and $`\overline{\mathrm{MS}}`$ scheme lead to $`Z_a^1(n,Q^2)`$, which is expressed as $`Z_a^1(n,Q^2)`$ $`=`$ $`I\left(\begin{array}{cccc}\frac{\alpha _s}{2\pi }\mathrm{\Delta }w_S(n,a)& \frac{\alpha _s}{2\pi }\mathrm{\Delta }z(n,a)& 0& \frac{\alpha }{\pi }3e^2\mathrm{\Delta }\widehat{z}(n,a)\\ 0& & 0& 0\\ 0& 0& \frac{\alpha _s}{2\pi }\mathrm{\Delta }w_{NS}(n,a)& \frac{\alpha }{\pi }3(e^4e^2^2)\mathrm{\Delta }\widehat{z}(n,a)\\ 0& 0& 0& 0\end{array}\right),`$ where $`I`$ is a $`4\times 4`$ unit matrix. Now we derive corresponding transformation rules from $`\overline{\mathrm{MS}}`$ scheme to $`a`$-scheme for the relevant two-loop anomalous dimensions. The parton distribution functions $`\mathrm{\Delta }\stackrel{~}{𝒒}^\gamma (n,Q^2,P^2)`$ satisfy the following evolution equation : $$\frac{d\mathrm{\Delta }\stackrel{~}{𝒒}^\gamma (n,Q^2,P^2)}{d\mathrm{ln}Q^2}=\mathrm{\Delta }\stackrel{~}{P}(n,Q^2)\mathrm{\Delta }\stackrel{~}{𝒒}^\gamma (n,Q^2,P^2),$$ (3.9) where $$\mathrm{\Delta }\stackrel{~}{P}(n,Q^2)=\left(\begin{array}{cccc}\mathrm{\Delta }P_{\psi \psi }(n,Q^2)& \mathrm{\Delta }P_{\psi G}(n,Q^2)& 0& \mathrm{\Delta }k_S(n,Q^2)\\ \mathrm{\Delta }P_{G\psi }(n,Q^2)& \mathrm{\Delta }P_{GG}(n,Q^2)& 0& \mathrm{\Delta }k_G(n,Q^2)\\ 0& 0& \mathrm{\Delta }P_{NS}(n,Q^2)& \mathrm{\Delta }k_{NS}(n,Q^2)\\ 0& 0& 0& 0\end{array}\right).$$ From Eq.(3.2), we obtain $`{\displaystyle \frac{d\mathrm{\Delta }\stackrel{~}{𝒒}^\gamma (n,Q^2,P^2)\|_a}{d\mathrm{ln}Q^2}}`$ $`=`$ $`{\displaystyle \frac{dZ_a(n,Q^2)}{d\mathrm{ln}Q^2}}\mathrm{\Delta }\stackrel{~}{𝒒}^\gamma (n,Q^2,P^2)\|_{\overline{\mathrm{MS}}}+Z_a(n,Q^2){\displaystyle \frac{d\mathrm{\Delta }\stackrel{~}{𝒒}^\gamma (n,Q^2,P^2)\|_{\overline{\mathrm{MS}}}}{d\mathrm{ln}Q^2}}`$ (3.10) $`=`$ $`\mathrm{\Delta }\stackrel{~}{P}(n,Q^2)\|_a\mathrm{\Delta }\stackrel{~}{𝒒}^\gamma (n,Q^2,P^2)\|_a,`$ with $$\mathrm{\Delta }\stackrel{~}{P}(n,Q^2)\|_a=\left[\frac{dZ_a(n,Q^2)}{d\mathrm{ln}Q^2}+Z_a(n,Q^2)\mathrm{\Delta }\stackrel{~}{P}(n,Q^2)\|_{\overline{\mathrm{MS}}}\right]Z_a^1(n,Q^2).$$ (3.11) The splitting functions $`\mathrm{\Delta }P_i(n,Q^2)`$ ($`i=\psi \psi ,\psi G,G\psi ,GG`$, and $`NS`$) and $`\mathrm{\Delta }k_j(n,Q^2)`$ ($`j=S,G,NS`$) are expanded as $`\mathrm{\Delta }P_i(n,Q^2)`$ $`=`$ $`{\displaystyle \frac{\alpha _s(Q^2)}{2\pi }}\mathrm{\Delta }P_i^{(0)}(n)+\left[{\displaystyle \frac{\alpha _s(Q^2)}{2\pi }}\right]^2\mathrm{\Delta }P_i^{(1)}(n)+\mathrm{},`$ (3.12) $`\mathrm{\Delta }k_j(n,Q^2)`$ $`=`$ $`{\displaystyle \frac{\alpha }{2\pi }}\mathrm{\Delta }k_j^{(0)}(n)+{\displaystyle \frac{\alpha \alpha _s(Q^2)}{(2\pi )^2}}\mathrm{\Delta }k_j^{(1)}(n)+\mathrm{},`$ (3.13) Since the QCD effective coupling constant $`\alpha _s(Q^2)`$ satisfies $$\frac{d\alpha _s(Q^2)}{d\mathrm{ln}Q^2}=\beta _0\frac{\alpha _s(Q^2)^2}{4\pi }+\mathrm{},$$ (3.14) where $`\beta _0=11\frac{2}{3}N_f`$ is the one-loop coefficient of the QCD beta function, and the $`n`$-th anomalous dimensions are defined as $`\mathrm{\Delta }P_i^{(0)}(n)`$ $`=`$ $`{\displaystyle \frac{1}{4}}\mathrm{\Delta }\gamma _i^{0,n},\mathrm{\Delta }P_i^{(1)}(n)={\displaystyle \frac{1}{8}}\mathrm{\Delta }\gamma _i^{(1),n},`$ (3.15) $`\mathrm{\Delta }k_j^{(0)}(n)`$ $`=`$ $`{\displaystyle \frac{1}{4}}\mathrm{\Delta }K_j^{0,n},\mathrm{\Delta }k_j^{(1)}(n)={\displaystyle \frac{1}{8}}\mathrm{\Delta }K_j^{(1),n},`$ (3.16) we find for one-loop $$\mathrm{\Delta }\gamma _{i,a}^{0,n}=\mathrm{\Delta }\gamma _{i,\overline{\mathrm{MS}}}^{0,n},\mathrm{\Delta }K_{j,a}^{0,n}=\mathrm{\Delta }K_{j,\overline{\mathrm{MS}}}^{0,n},$$ (3.17) and for two-loop $`\mathrm{\Delta }\gamma _{\psi \psi ,a}^{(1),n}`$ $`=`$ $`\mathrm{\Delta }\gamma _{\psi \psi ,\overline{\mathrm{MS}}}^{(1),n}+2\mathrm{\Delta }z(n,a)\mathrm{\Delta }\gamma _{G\psi }^{0,n}+4\beta _0\mathrm{\Delta }w_S(n,a),`$ $`\mathrm{\Delta }\gamma _{\psi G,a}^{(1),n}`$ $`=`$ $`\mathrm{\Delta }\gamma _{\psi G,\overline{\mathrm{MS}}}^{(1),n}+2\mathrm{\Delta }z(n,a)\left[\mathrm{\Delta }\gamma _{GG}^{0,n}\mathrm{\Delta }\gamma _{\psi \psi }^{0,n}+2\beta _0\right]`$ $`+2\mathrm{\Delta }w_S(n,a)\mathrm{\Delta }\gamma _{\psi G}^{0,n},`$ $`\mathrm{\Delta }\gamma _{G\psi ,a}^{(1),n}`$ $`=`$ $`\mathrm{\Delta }\gamma _{G\psi ,\overline{\mathrm{MS}}}^{(1),n}2\mathrm{\Delta }w_S(n,a)\mathrm{\Delta }\gamma _{G\psi }^{0,n},`$ $`\mathrm{\Delta }\gamma _{GG,a}^{(1),n}`$ $`=`$ $`\mathrm{\Delta }\gamma _{GG,\overline{\mathrm{MS}}}^{(1),n}2\mathrm{\Delta }z(n,a)\mathrm{\Delta }\gamma _{G\psi }^{0,n},`$ $`\mathrm{\Delta }\gamma _{NS,a}^{(1),n}`$ $`=`$ $`\mathrm{\Delta }\gamma _{NS,\overline{\mathrm{MS}}}^{(1),n}+4\beta _0\mathrm{\Delta }w_{NS}(n,a),`$ (3.18) $`\mathrm{\Delta }K_{S,a}^{(1),n}`$ $`=`$ $`\mathrm{\Delta }K_{S,\overline{\mathrm{MS}}}^{(1),n}+2\mathrm{\Delta }w_S(n,a)\mathrm{\Delta }K_S^{0,n}+4\mathrm{\Delta }\widehat{z}(n,a)3e^2\mathrm{\Delta }\gamma _{\psi \psi }^{0,n},`$ $`\mathrm{\Delta }K_{G,a}^{(1),n}`$ $`=`$ $`\mathrm{\Delta }K_{G,\overline{\mathrm{MS}}}^{(1),n}+4\mathrm{\Delta }\widehat{z}(n,a)3e^2\mathrm{\Delta }\gamma _{G\psi }^{0,n},`$ $`\mathrm{\Delta }K_{NS,a}^{(1),n}`$ $`=`$ $`\mathrm{\Delta }K_{NS,\overline{\mathrm{MS}}}^{(1),n}+2\mathrm{\Delta }w_{NS}(n,a)\mathrm{\Delta }K_{NS}^{0,n}`$ $`+4\mathrm{\Delta }\widehat{z}(n,a)3(e^4e^2^2)\mathrm{\Delta }\gamma _{NS}^{0,n}.`$ The one-loop photon matrix elements of the hadronic operators, $`\mathrm{\Delta }A_\psi ^n`$ and $`\mathrm{\Delta }A_{NS}^n`$ in Eq.(2.4), are related to each other as $$\mathrm{\Delta }A_{NS}^n=\mathrm{\Delta }A_\psi ^n\frac{e^4e^2^2}{e^2},$$ (3.19) and the sum $$(\mathrm{\Delta }C_\gamma ^{\gamma ,n}/\frac{\alpha }{4\pi }+e^2\mathrm{\Delta }A_\psi ^n+\mathrm{\Delta }A_{NS}^n)$$ (3.20) is factorization-scheme-independent in one-loop order . Thus we obtain from Eq.(3.3) $`\mathrm{\Delta }A_{\psi ,a}^n`$ $`=`$ $`\mathrm{\Delta }A_{\psi ,\overline{\mathrm{MS}}}^n+12e^2\mathrm{\Delta }\widehat{z}(n,a),`$ $`\mathrm{\Delta }A_{G,a}^n`$ $`=`$ $`\mathrm{\Delta }A_{G,\overline{\mathrm{MS}}}^n=0,`$ (3.21) $`\mathrm{\Delta }A_{NS,a}^n`$ $`=`$ $`\mathrm{\Delta }A_{NS,\overline{\mathrm{MS}}}^n+12(e^4e^2^2)\mathrm{\Delta }\widehat{z}(n,a).`$ Note that $`\mathrm{\Delta }A_G^n=0`$ in one-loop order. It is possible to choose $`\mathrm{\Delta }z(n,a)`$ and $`\mathrm{\Delta }\widehat{z}(n,a)`$ arbitrarily. In the following, we take $`\mathrm{\Delta }\widehat{z}(n,a)=\mathrm{\Delta }z(n,a)`$ in the CI-like schemes and $`\mathrm{\Delta }\widehat{z}(n,\mathrm{DIS}_\gamma )\mathrm{\Delta }z(n,\mathrm{DIS}_\gamma )=0`$ in the $`\mathrm{DIS}_\gamma `$ scheme. In one-loop order we have $`\mathrm{\Delta }w_S(n,a)=\mathrm{\Delta }w_{NS}(n,a)`$. Thus from now on, we set $`\mathrm{\Delta }w_S(n,a)=\mathrm{\Delta }w_{NS}(n,a)\mathrm{\Delta }w(n,a)`$. Let us now discuss the features of several factorization schemes. ### 3.2 The $`\overline{\mathrm{MS}}`$ scheme This is the only scheme in which both relevant one-loop coefficient functions and two-loop anomalous dimensions were actually calculated . In fact there still remain ambuguities in the $`\overline{\mathrm{MS}}`$ scheme, depending on how to handle $`\gamma _5`$ in $`n`$ dimensions. The $`\overline{\mathrm{MS}}`$ scheme we call here is the one due to Mertig and van Neerven and Vogelsang , in which the first moment of the non-singlet quark operator vanishes, corresponding to the conservation of the non-singlet axial current. Indeed we have $`\mathrm{\Delta }\gamma _{NS,\overline{\mathrm{MS}}}^{(1),n=1}=0`$. Explicit expressions of the relevant one-loop coefficient functions and two-loop anomalous dimensions can be found, for example, in Appendix of Ref. . It is noted that, in the $`\overline{\mathrm{MS}}`$ scheme, both the QCD and QED axial anomalies reside in the quark distributions and not in the gluon and photon coefficient functions. In fact we observe $`\mathrm{\Delta }\gamma _{\psi \psi ,\overline{\mathrm{MS}}}^{(1),n=1}`$ $`=`$ $`24C_FT_f0,`$ (3.22) $`\mathrm{\Delta }B_{G,\overline{\mathrm{MS}}}^{n=1}`$ $`=`$ $`\mathrm{\Delta }B_{\gamma ,\overline{\mathrm{MS}}}^{n=1}=0.`$ (3.23) where $`C_F=\frac{4}{3}`$ and $`T_f=\frac{N_f}{2}`$. Also we find from Eq.(2.5) that the first moments of the one-loop photon matrix elements of quark operators gain the non-zero values, i.e., $$\mathrm{\Delta }A_{\psi ,\overline{\mathrm{MS}}}^{n=1}=\frac{e^2}{e^4e^2^2}\mathrm{\Delta }A_{NS,\overline{\mathrm{MS}}}^{n=1}=12e^2N_f,$$ (3.24) which is due to the QED axial anomaly. ### 3.3 The CI-like schemes The EMC measurement of the first moment of the proton spin structure function $`g_1^p(x,Q^2)`$ presented us with an issue called “proton spin crisis”. Since then many ideas have been proposed as solutions. One simple and plausible explanation was that there exists an anomalous gluon contribution to the first moment - originating from the QCD axial anomaly. This explanation was later supported with a notion of the factorization-scheme dependence. There is a set of the factorization schemes in which we obtain $$\mathrm{\Delta }B_G^{n=1}=2N_f,\mathrm{\Delta }\gamma _{\psi \psi }^{(1),n=1}=0.$$ (3.25) Let us call them CI-like schemes. In this paper we consider four CI-like schemes, in which we take $`\mathrm{\Delta }z(n,a)=\mathrm{\Delta }\widehat{z}(n,a)`$, since both QCD and QED anomalies originate from the similar triangle diagrams. With this choice, the relation between the one-loop gluon and photon coefficient functions, which holds in the $`\overline{\mathrm{MS}}`$ scheme, also holds in the CI-like schemes, $$\mathrm{\Delta }B_{\gamma ,\mathrm{CI}\mathrm{like}}^n=\frac{2}{N_f}\mathrm{\Delta }B_{G,\mathrm{CI}\mathrm{like}}^n.$$ (3.26) Thus, in addition to the relations in Eq.(3.25), we obtain in the CI-like schemes $$\mathrm{\Delta }B_{\gamma ,\mathrm{CI}\mathrm{like}}^{n=1}=4,\mathrm{\Delta }A_{\psi ,\mathrm{CI}\mathrm{like}}^{n=1}=\mathrm{\Delta }A_{NS,\mathrm{CI}\mathrm{like}}^{n=1}=0.$$ (3.27) (i) \[The chirally invariant (CI) scheme\] In this scheme the factorization of the photon-gluon (photon-photon) cross section into the hard and soft parts is made so that chiral symmetry is respected and the QCD and QED anomaly effects are absorbed into the gluon and photon coefficient functions. Thus the spin-dependent quark distributions in the CI scheme are anomaly-free. The transformation from the $`\overline{\mathrm{MS}}`$ scheme to the CI scheme is achieved by $$\mathrm{\Delta }w(n,a=\mathrm{CI})=0,\mathrm{\Delta }z(n,a=\mathrm{CI})=\mathrm{\Delta }\widehat{z}(n,a=\mathrm{CI})=2N_f\frac{1}{n(n+1)}.$$ (3.28) It has been argued by Cheng and Müller and Teryaev that the $`x`$-dependence of the axial-anomaly effect is uniquely fixed and that its $`x`$-behavior leads to the transformation rule (3.28) and thus to the CI scheme. (ii) \[The Adler-Bardeen (AB) scheme\] Ball, Forte and Ridolfi proposed several CI-like schemes for the analysis of the nucleon spin structure function $`g_1(x,Q^2)`$. One of them is the Adler-Bardeen (AB) scheme which was introduced by requiring that the change from the $`\overline{\mathrm{MS}}`$ scheme to this scheme be independent of $`x`$, so that the large and small $`x`$ behavior of the gluon coefficient function is unchanged. In our case, we have in moment space $$\mathrm{\Delta }w(n,a=\mathrm{AB})=0,\mathrm{\Delta }z(n,a=\mathrm{AB})=\mathrm{\Delta }\widehat{z}(n,a=\mathrm{AB})=N_f\frac{1}{n}.$$ (3.29) (iii) \[The off-shell (OS) scheme\] In this scheme we renormalize operators while keeping the incoming particle off-shell, $`p^20`$, so that at renormalization (factorization) point $`\mu ^2=p^2`$, the finite terms vanish. This is exactly the same as “the momentum subtraction scheme” which was used some time ago to calculate, for instance, the polarized quark and gluon coefficient functions . The CI-relations in Eqs.(3.25) and (3.27) also hold in the OS scheme , since the axial anomaly appears as a finite term in the calculation of the triangle graph for $`j_5^\mu `$ between external gluons (photons) and the finite term is thrown away in this scheme. The transformation from the $`\overline{\mathrm{MS}}`$ scheme to the OS scheme is made by choosing $`\mathrm{\Delta }w(n,a=\mathrm{OS})`$ $`=`$ $`C_F\{\left[S_1(n)\right]^2+3S_2(n)S_1(n)({\displaystyle \frac{1}{n}}{\displaystyle \frac{1}{(n+1)}})`$ $`{\displaystyle \frac{7}{2}}+{\displaystyle \frac{2}{n}}{\displaystyle \frac{3}{n+1}}{\displaystyle \frac{1}{n^2}}+{\displaystyle \frac{2}{(n+1)^2}}\},`$ $`\mathrm{\Delta }z(n,a=\mathrm{OS})`$ $`=`$ $`\mathrm{\Delta }\widehat{z}(n,a=\mathrm{OS})`$ (3.30) $`=`$ $`N_f\left\{{\displaystyle \frac{n1}{n(n+1)}}S_1(n)+{\displaystyle \frac{1}{n}}+{\displaystyle \frac{1}{n^2}}{\displaystyle \frac{4}{(n+1)^2}}\right\}.`$ It is noted that in the OS scheme we have $`\mathrm{\Delta }A_{\psi ,\mathrm{OS}}^n=\mathrm{\Delta }A_{NS,\mathrm{OS}}^n=0`$ not only for $`n=1`$ but also for all $`n`$. (iv) \[The Altarelli-Ross (AR) scheme\] Using massive quark as a regulator for collinear divergence, Altarelli and Ross derived the same one-loop gluon coefficient function $`\mathrm{\Delta }C_G^\gamma `$ as in the case of CI scheme. In order to obtain the one-loop quark coefficient function in this scheme, however, we need to do an extra subtraction so that the conservation of the nonsinglet axial currents is secured . The transformation rule is $`\mathrm{\Delta }w_S(n,a=\mathrm{AR})`$ $`=`$ $`C_F\{2\left[S_1(n)\right]^2+2S_2(n)S_1(n)({\displaystyle \frac{2}{n}}{\displaystyle \frac{2}{n+1}}+2)`$ (3.31) $`2+{\displaystyle \frac{1}{n}}{\displaystyle \frac{1}{n+1}}+{\displaystyle \frac{2}{(n+1)^2}}\},`$ $`\mathrm{\Delta }z(n,a=\mathrm{AR})`$ $`=`$ $`\mathrm{\Delta }\widehat{z}(n,a=\mathrm{AR})=2N_f{\displaystyle \frac{1}{n(n+1)}}.`$ (3.32) ### 3.4 The $`\mathrm{DIS}_\gamma `$ scheme An interesting factorization scheme, which is called $`\mathrm{DIS}_\gamma `$, was introduced some time ago into the NLO analysis of the unpolarized real photon structure function $`F_2^\gamma (x,Q^2)`$. Glück, Reya and Vogt observed that, in the $`\overline{\mathrm{MS}}`$ scheme, the $`\mathrm{ln}(1x)`$ term in the photonic coefficient function $`C_2^\gamma (x)`$ for $`F_2^\gamma `$, which becomes negative and divergent for $`x1`$, drives the ‘pointlike’ part of $`F_2^\gamma `$ to large negative values as $`x1`$, leading to a strong difference between the LO and the NLO results for $`F_{2,\mathrm{pointlike}}^\gamma `$ in the large-$`x`$ region. They introduced the $`\mathrm{DIS}_\gamma `$ scheme in which the photonic coefficient function $`C_2^\gamma `$, i.e., the direct-photon contribution to $`F_2^\gamma `$, is absorbed into the photonic quark distributions. It is noted that, for the real photon target, the structure function $`F_2^\gamma `$ is decomposed into a ‘pointlike’ and a ‘hadronic’ part, the former being perturbatively calculable but not the latter. And beyond the LO both the ‘pointlike’ and the ‘hadronic’ parts depend on the factorization scheme employed. A similar situation occurs in the polarized case, and the $`\mathrm{DIS}_\gamma `$ scheme was applied to the NLO analysis for the spin-dependent structure function $`g_1^\gamma (x,Q^2)`$ of the real photon target by Stratmann and Vogelsang . In the polarized version of $`\mathrm{DIS}_\gamma `$ scheme we take $`\mathrm{\Delta }w_S(n,\mathrm{DIS}_\gamma )`$ $`=`$ $`\mathrm{\Delta }w_{NS}(n,\mathrm{DIS}_\gamma )=\mathrm{\Delta }z(n,\mathrm{DIS}_\gamma )=0,`$ (3.33) $`\mathrm{\Delta }\widehat{z}(n,\mathrm{DIS}_\gamma )`$ $`=`$ $`{\displaystyle \frac{N_f}{4}}\mathrm{\Delta }B_{\gamma ,\overline{\mathrm{MS}}}^n`$ (3.34) $`=`$ $`N_f\left\{{\displaystyle \frac{n1}{n(n+1)}}S_1(n)+{\displaystyle \frac{3}{n}}{\displaystyle \frac{4}{n+1}}{\displaystyle \frac{1}{n^2}}\right\},`$ so that $`\mathrm{\Delta }B_{\gamma ,\mathrm{DIS}_\gamma }^n`$ $`=`$ $`\mathrm{\Delta }B_{\gamma ,\overline{\mathrm{MS}}}^n{\displaystyle \frac{4}{N_f}}\mathrm{\Delta }\widehat{z}(n,\mathrm{DIS}_\gamma )`$ (3.35) $`=`$ $`0.`$ Note that the relation $`\stackrel{`}{a}`$ la Eqs.(3.6) and (3.26) in the $`\overline{\mathrm{MS}}`$ and CI-like factorization schemes does not hold anymore in this scheme, i.e., $$\mathrm{\Delta }B_{\gamma ,\mathrm{DIS}_\gamma }^n\frac{2}{N_f}\mathrm{\Delta }B_{G,\mathrm{DIS}_\gamma }^n(=\frac{2}{N_f}\mathrm{\Delta }B_{G,\overline{\mathrm{MS}}}^n).$$ (3.36) For $`n=1`$, we have $$\mathrm{\Delta }\widehat{z}(n=1,\mathrm{DIS}_\gamma )=0,$$ (3.37) and thus, together with Eq.(3.33), we observe that as far as the first moments are concerned, $`\mathrm{DIS}_\gamma `$ scheme gives the same results with $`\overline{\mathrm{MS}}`$. In other words, in the $`\mathrm{DIS}_\gamma `$ scheme, both the QCD and QED axial anomaly effects are retained in the quark distributions. With these preparations, we now examine the factoraization scheme dependence of the polarized parton distributions in the virtual photon. The two-loop anomalous dimensions of the spin-dependent operators and one-loop photon matrix elements of the hadronic operators in the $`\overline{\mathrm{MS}}`$ scheme are already known. Corresponding quantities in a particular scheme are obtained through the transformation rules given in Eq.(3.18). Inserting these quantities into the formulas given in Appendix A, we get the NLO predictions for the moments of polarized parton distributions in a particular factorization scheme. ### 3.5 Gluon distribution in the virtual photon Let us start with the gluon distribution. We find that all the factorization schemes which we consider in this paper predict the same behavior for the gluon distribution up to NLO: $$\mathrm{\Delta }G^\gamma (n,Q^2,P^2)\|_a=\mathrm{\Delta }G^\gamma (n,Q^2,P^2)\|_{\overline{\mathrm{MS}}},$$ (3.38) where $`a`$ means factorization schemes of CI, AB, OS, AR and $`\mathrm{DIS}_\gamma `$. This can be seen from the direct calculation or from the notion that, up to NLO, $`\mathrm{\Delta }G^\gamma \|_a`$ satisfies the same evolution equation as $`\mathrm{\Delta }G^\gamma \|_{\overline{\mathrm{MS}}}`$ with the same initial condition at $`Q^2=P^2`$, namely, $`\mathrm{\Delta }G^\gamma (n,P^2,P^2)\|_a=\mathrm{\Delta }G^\gamma (n,P^2,P^2)\|_{\overline{\mathrm{MS}}}=0`$. If we consider a more general factorization scheme in which the hadronic part of $`Z_a^1(n,Q^2)`$ in Eq.(LABEL:Zinverse) is replaced with a new one as follows, $$\left(\begin{array}{cc}1\frac{\alpha _s}{2\pi }\mathrm{\Delta }w_S& \frac{\alpha _s}{2\pi }\mathrm{\Delta }z\\ 0& 1\end{array}\right)\left(\begin{array}{cc}1\frac{\alpha _s}{2\pi }\mathrm{\Delta }w_S& \frac{\alpha _s}{2\pi }\mathrm{\Delta }z\\ \frac{\alpha _s}{2\pi }\mathrm{\Delta }u& 1\frac{\alpha _s}{2\pi }\mathrm{\Delta }v\end{array}\right),$$ (3.39) then, in this new factorization scheme, the predicted gluon distribution is not the same with $`\mathrm{\Delta }G^\gamma (n,Q^2,P^2)\|_{\overline{\mathrm{MS}}}`$ in NLO. However, the first moment is found to be still the same. In other words, the first moment of the gluon distribution in the virtual photon, $`\mathrm{\Delta }G^\gamma (n=1,Q^2,P^2)`$, is factorization-scheme independent up to NLO. This is due to the fact that the new terms, which appear by the inclusion of $`\mathrm{\Delta }u`$ and $`\mathrm{\Delta }v`$, will be proportional to $`\mathrm{\Delta }K_\psi ^{0,n}`$ and that $`\mathrm{\Delta }K_\psi ^{0,n=1}=0`$. Also inclusion of $`\mathrm{\Delta }u`$ and $`\mathrm{\Delta }v`$ terms in $`Z_a^1`$ does not modify the photon structure function $`g_1^\gamma (x,Q^2,P^2)`$ itself up to NLO, since the gluon coefficient function starts in the order $`\alpha _s`$. Moreover, the quark distributions in the virtual photon do not change by the inclusion of $`\mathrm{\Delta }u`$ and $`\mathrm{\Delta }v`$ terms. ## 4 The $`n=1`$ moments of parton distributions The first moments of polarized parton distributions in the virtual photon are particularly interesting since they have relevance to the QCD and QED axial anomalies. The explicit expressions for the moments of $`\mathrm{\Delta }q_S^\gamma `$, $`\mathrm{\Delta }G^\gamma `$, and $`\mathrm{\Delta }q_{NS}^\gamma `$ up to NLO are given in Appendix A. We take the $`n1`$ limit in these expressions. Useful $`n=1`$ moments of one- and two-loop anomalous dimensions, photon matrix elements, and coefficient functions both in the $`\overline{\mathrm{MS}}`$ and CI-like schemes are enumerated in Appendix B. As far as the first moments are concerned, $`\mathrm{DIS}_\gamma `$ scheme gives the same results with $`\overline{\mathrm{MS}}`$. Note that we have $$\lambda _+^{n=1}=0,\lambda _{}^{n=1}=2\beta _0,\lambda _{NS}^{n=1}(=\mathrm{\Delta }\gamma _{NS}^{0,n=1})=0,$$ (4.1) where $`\lambda _\pm ^{n=1}`$ are eigenvalues of the one-loop hadronic anomalous dimension matrix $`\mathrm{\Delta }\gamma _{ij}^{0,n=1}`$. The zero eigenvalues $`\lambda _+^{n=1}=\lambda _{NS}^{n=1}=0`$ correspond to the conservation of the axial-vector current at one-loop order. ### 4.1 The $`n=1`$ moment of gluon distribution The expressions for the moments of gluon distribution are given in Appendix A.2. In these expressions the factors $$\frac{1}{\lambda _+^n},\frac{1}{2\beta _0+\lambda _{}^n},\frac{1}{2\beta _0+\lambda _{}^n\lambda _+^n}$$ (4.2) may develop singularities at $`n=1`$ and so we need a little care when we deal with them. Taking the limit of $`n`$ going to 1, we find $`\widehat{L}_G^{+n}`$ $``$ $`0,\widehat{L}_G^n\mathrm{finite},`$ $`\widehat{A}_G^{+n}`$ $``$ $`\mathrm{finite},\widehat{B}_G^{+n}0,\widehat{B}_G^n\mathrm{finite},`$ (4.3) $`\widehat{A}_G^n`$ $``$ $`72e^2N_fC_F.`$ The terms proportional to $`\widehat{L}_G^n`$, $`\widehat{B}_G^n`$, and $`\widehat{A}_G^{+n}`$ all vanish in the $`n=1`$ limit, since they are multiplied by the following vanishing factors: $$\left\{1\left[\frac{\alpha _s(Q^2)}{\alpha _s(P^2)}\right]^{\lambda _{}^n/2\beta _0+1}\right\},\left\{1\left[\frac{\alpha _s(Q^2)}{\alpha _s(P^2)}\right]^{\lambda _+^n/2\beta _0}\right\}.$$ (4.4) Only exception is the term proportional to $`\widehat{A}_G^n`$. We find for $`n1`$ $$\widehat{A}_G^n\left\{1\left[\frac{\alpha _s(Q^2)}{\alpha _s(P^2)}\right]^{\lambda _{}^n/2\beta _0}\right\}72e^2N_fC_F\frac{\alpha _s(Q^2)\alpha _s(P^2)}{\alpha _s(Q^2)}.$$ (4.5) Thus we obtain $$\mathrm{\Delta }G^\gamma (n=1,Q^2,P^2)=\frac{12\alpha }{\pi \beta _0}e^2N_f\frac{\alpha _s(Q^2)\alpha _s(P^2)}{\alpha _s(Q^2)}.$$ (4.6) for the first moment of the gluon distribution in the virtual photon. It should be emphasized that the result is factorization-scheme independent. ### 4.2 The $`n=1`$ moment of quark distributions The expressions for the moments of quark distributions are given in Appendix A.1 and A.3. In all the factorization schemes under study, i.e., $`\overline{\mathrm{MS}}`$, $`\mathrm{DIS}_\gamma `$, and CI-like schemes, we find for $`n1`$, $`\widehat{L}_S^{+n}`$ $``$ $`0,\widehat{L}_S^n0,\widehat{L}_{NS}^n0,`$ $`\widehat{A}_S^{+n}`$ $``$ $`\mathrm{finite},\widehat{A}_S^n0,\widehat{A}_{NS}^n\mathrm{finite}`$ (4.7) $`\widehat{B}_S^{+n}`$ $``$ $`0,\widehat{B}_S^n\mathrm{finite},\widehat{B}_{NS}^n0`$ The terms proportional to $`\widehat{A}_S^{+n}`$ and $`\widehat{B}_S^n`$ are multiplied by the vanishing factors in Eq.(4.4), and the $`\widehat{A}_{NS}^n`$ term multiplied by $$\left\{1\left[\frac{\alpha _s(Q^2)}{\alpha _s(P^2)}\right]^{\lambda _{NS}^n/2\beta _0}\right\},$$ (4.8) and, therefore, the $`\widehat{L}_S^{+n}`$, $`\widehat{L}_S^n`$, $`\widehat{L}_{NS}^n`$, $`\widehat{A}_S^{+n}`$, $`\widehat{A}_S^n`$, $`\widehat{A}_{NS}^n`$, $`\widehat{B}_S^{+n}`$, $`\widehat{B}_S^n`$, and $`\widehat{B}_{NS}^n`$ terms in Eqs.(A.1) and (A.17) all vanish in the $`n=1`$ limit. Then the first moments of quark distributions are given by $`\mathrm{\Delta }q_S^\gamma (n=1,Q^2,P^2)`$ $`=`$ $`{\displaystyle \frac{\alpha }{8\pi \beta _0}}\widehat{C}_S^{n=1}={\displaystyle \frac{\alpha }{4\pi }}\mathrm{\Delta }A_\psi ^{n=1},`$ (4.9) $`\mathrm{\Delta }q_{NS}^\gamma (n=1,Q^2,P^2)`$ $`=`$ $`{\displaystyle \frac{\alpha }{8\pi \beta _0}}\widehat{C}_{NS}^{n=1}={\displaystyle \frac{\alpha }{4\pi }}\mathrm{\Delta }A_{NS}^{n=1}.`$ (4.10) We now see that scheme dependence for the first moments of quark distributions is coming from the photon matrix elements $`\mathrm{\Delta }A_\psi ^n`$ and $`\mathrm{\Delta }A_{NS}^n`$. In the case of CI-like factorization schemes, $`a=\mathrm{CI},\mathrm{AB},\mathrm{OS},\mathrm{AR}`$, we have $`\mathrm{\Delta }w(n=1,a)`$ $`=`$ $`0,`$ $`\mathrm{\Delta }z(n=1,a)`$ $`=`$ $`\mathrm{\Delta }\widehat{z}(n=1,a)=N_f\mathrm{for}a=\mathrm{CI},\mathrm{AB},\mathrm{OS},\mathrm{AR}`$ (4.11) We find from Eqs.(3.21) and (3.24) that these schemes give $$\mathrm{\Delta }A_{\psi ,a}^{n=1}=\mathrm{\Delta }A_{NS,a}^{n=1}=0.$$ (4.12) This leads to an interesting result: The first moment of spin-dependent quark distributions in the virtual photon vanish in NLO for $`a=\mathrm{CI},\mathrm{AB},\mathrm{OS},\mathrm{AR}`$. $$\mathrm{\Delta }q_S^\gamma (n=1,Q^2,P^2)\|_a=\mathrm{\Delta }q_{NS}^\gamma (n=1,Q^2,P^2)\|_a=0$$ (4.13) The vanishing first moments imply that the axial anomaly effects do not reside in the quark distributions. In these CI-like schemes, the QCD and QED axial anomalies are transfered to the gluon and photon coefficient functions, respectively, and their first moments do not vanish. Indeed we obtain from Eqs.(3.1) and (LABEL:Zinverse) $`\mathrm{\Delta }C_{G,a}^{\gamma ,n=1}`$ $`=`$ $`e^2{\displaystyle \frac{\alpha _s(Q^2)}{2\pi }}N_f`$ (4.14) $`\mathrm{\Delta }C_{\gamma ,a}^{\gamma ,n=1}`$ $`=`$ $`{\displaystyle \frac{3\alpha }{\pi }}e^4N_f\left(1{\displaystyle \frac{\alpha _s(Q^2)}{\pi }}\right)\mathrm{for}a=\mathrm{CI},\mathrm{AB},\mathrm{OS},\mathrm{AR}`$ (4.15) where we have used the fact $`{\displaystyle \frac{1}{<e^2>}}\mathrm{\Delta }C_S^{\gamma ,n=1}\|_{\overline{\mathrm{MS}}}`$ $`=`$ $`\mathrm{\Delta }C_{NS}^{\gamma ,n=1}\|_{\overline{\mathrm{MS}}}=1{\displaystyle \frac{\alpha _s}{\pi }}+𝒪(\alpha _s^2)`$ (4.16) $`\mathrm{\Delta }C_\gamma ^{\gamma ,n=1}\|_{\overline{\mathrm{MS}}}`$ $`=`$ $`0+𝒪(\alpha \alpha _s^2)`$ (4.17) On the other hand, in the $`\overline{\mathrm{MS}}`$ (and also in $`\mathrm{DIS}_\gamma `$) we obtain from Eq.(2.5) $`\mathrm{\Delta }A_{\psi ,\overline{\mathrm{MS}}}^{n=1}`$ $`=`$ $`12e^2N_f`$ (4.18) $`\mathrm{\Delta }A_{NS,\overline{\mathrm{MS}}}^{n=1}`$ $`=`$ $`12(e^4e^2^2)N_f`$ (4.19) and thus $`\mathrm{\Delta }q_S^{\gamma ,n=1}\|_{\overline{\mathrm{MS}}}`$ and $`\mathrm{\Delta }q_{NS}^{\gamma ,n=1}\|_{\overline{\mathrm{MS}}}`$ are non-zero constant. Actually we can go one step further to the order of $`\alpha _s`$ QCD corrections. This is due to the fact that, in the $`\overline{\mathrm{MS}}`$ scheme, the parton distribution $`\mathrm{\Delta }𝒒^\gamma (n=1)\|_{\overline{\mathrm{MS}}}=(\mathrm{\Delta }q_S^\gamma ,\mathrm{\Delta }G^\gamma ,\mathrm{\Delta }q_{NS}^\gamma )\|_{\overline{\mathrm{MS}}}`$ satisfies a homogeneous differential equation without inhomogeneous LO and NLO $`\mathrm{\Delta }K`$ terms. Indeed we find $`\mathrm{\Delta }q_S^\gamma (n=1,Q^2,P^2)\|_{\overline{\mathrm{MS}}}`$ $`=`$ $`\left[{\displaystyle \frac{\alpha }{\pi }}3e^2N_f\right]\left\{1{\displaystyle \frac{2}{\beta _0}}{\displaystyle \frac{\alpha _s(P^2)\alpha _s(Q^2)}{\pi }}N_f\right\}`$ (4.20) $`\mathrm{\Delta }q_{NS}^\gamma (n=1,Q^2,P^2)\|_{\overline{\mathrm{MS}}}`$ $`=`$ $`\left[{\displaystyle \frac{\alpha }{\pi }}3\left(e^4e^2^2\right)N_f\right]\left\{1+𝒪(\alpha _s^2)\right\},`$ (4.21) the derivation of which is shown in Appendix C. In the $`\overline{\mathrm{MS}}`$ scheme, the axial anomaly effects are retained in the quark distributions. The factors $`\left[\frac{\alpha }{\pi }3e^2N_f\right]`$ and $`\left[\frac{\alpha }{\pi }3\left(e^4e^2^2\right)N_f\right]`$ are related to the QED axial anomaly and a term $`\frac{2}{\beta _0}\frac{\alpha _s(P^2)\alpha _s(Q^2)}{\pi }N_f`$ in $`\mathrm{\Delta }q_S^{\gamma ,n=1}\|_{\overline{\mathrm{MS}}}`$ is coming from the QCD axial anomaly. ### 4.3 The $`n=1`$ moment of $`g_1^\gamma (x,Q^2,P^2)`$ The polarized structure function $`g_1^\gamma (x,Q^2,P^2)`$ of the virtual photon satisfies the following sum rule: $`{\displaystyle _0^1}𝑑xg_1^\gamma (x,Q^2,P^2)`$ $`=`$ $`{\displaystyle \frac{3\alpha }{\pi }}e^4N_f\left(1{\displaystyle \frac{\alpha _s(Q^2)}{\pi }}\right)`$ (4.22) $`+`$ $`{\displaystyle \frac{6\alpha }{\pi \beta _0}}\left[e^2N_f\right]^2{\displaystyle \frac{\alpha _s(P^2)\alpha _s(Q^2)}{\pi }}+𝒪(\alpha _s^2).`$ This sum rule is of course the factorization-scheme independent. Now we examine how the scheme-dependent parton distributions contribute to this sum rule. In the CI-like schemes ($`a=\mathrm{CI},\mathrm{AB},\mathrm{OS},\mathrm{AR}`$), the first moment of the quark distributions vanish in NLO, and thus the contribution to the sum rule comes from the gluon and photon distributions. Equations (4.6) and (4.15) show that $$\mathrm{\Delta }C_{G,a}^{\gamma ,n=1}\mathrm{\Delta }G^\gamma (n=1,Q^2,P^2)\|_a+\mathrm{\Delta }C_{\gamma ,a}^{\gamma ,n=1}$$ (4.23) leads to the result (4.22). On the other hand, in the $`\overline{\mathrm{MS}}`$ scheme (and also in $`\mathrm{DIS}_\gamma `$), the one-loop gluon and photon coefficient functions vanish, $`\mathrm{\Delta }B_{G,\overline{\mathrm{MS}}}^{n=1}=\mathrm{\Delta }B_{\gamma ,\overline{\mathrm{MS}}}^{n=1}=0`$ and, therefore, the sum rule is derived from the quark contributions. Indeed we find from Eqs.(4.16), (4.20-4.21) $$\mathrm{\Delta }C_{S,\overline{\mathrm{MS}}}^{\gamma ,n=1}\mathrm{\Delta }q_S^\gamma (n=1,Q^2,P^2)\|_{\overline{\mathrm{MS}}}+\mathrm{\Delta }C_{NS,\overline{\mathrm{MS}}}^{\gamma ,n=1}\mathrm{\Delta }q_{NS}^\gamma (n=1,Q^2,P^2)\|_{\overline{\mathrm{MS}}}$$ (4.24) leads to the same result. It is interesting to note that the sum rule (4.22) is the consequence of the QCD and QED axial anomalies and that in the CI-like schemes the anomaly effect resides in the gluon contribution while, in $`\overline{\mathrm{MS}}`$, in the quark contributions. Furthermore, the first term of the sum rule (4.22) is coming from the QED axial anomaly and the second is from the QCD axial anomaly<sup>1</sup><sup>1</sup>1This notion was first pointed out by the authors of Ref.. ## 5 Behaviors of parton distributions near $`x=1`$ The behaviors of parton distributions near $`x=1`$ are governed by the large-$`n`$ limit of those moments. In the leading order, parton distributions are factorization-scheme independent. For large $`n`$$`\mathrm{\Delta }q_S^\gamma (n,Q^2,P^2)\|_{\mathrm{LO}}`$ and $`\mathrm{\Delta }q_{NS}^\gamma (n,Q^2,P^2)\|_{\mathrm{LO}}`$ behave as $`1/(n\mathrm{ln}n)`$, while $`\mathrm{\Delta }G^\gamma (n,Q^2,P^2)\|_{\mathrm{LO}}1/(n\mathrm{ln}n)^2`$. Thus in $`x`$ space, the parton distributions vanish for $`x1`$. In fact, we find $`\mathrm{\Delta }q_S^\gamma (x,Q^2,P^2)\|_{\mathrm{LO}}`$ $``$ $`{\displaystyle \frac{\alpha }{4\pi }}{\displaystyle \frac{4\pi }{\alpha _s(Q^2)}}N_fe^2{\displaystyle \frac{9}{4}}{\displaystyle \frac{1}{\mathrm{ln}(1x)}},`$ (5.1) $`\mathrm{\Delta }G^\gamma (x,Q^2,P^2)\|_{\mathrm{LO}}`$ $``$ $`{\displaystyle \frac{\alpha }{4\pi }}{\displaystyle \frac{4\pi }{\alpha _s(Q^2)}}N_fe^2{\displaystyle \frac{1}{2}}{\displaystyle \frac{\mathrm{ln}x}{\mathrm{ln}^2(1x)}}.`$ (5.2) The behaviors of $`\mathrm{\Delta }q_{NS}^\gamma (x,Q^2,P^2)`$ for $`x1`$, both in LO and NLO, are always given by the corresponding expressions for $`\mathrm{\Delta }q_S^\gamma (x,Q^2,P^2)`$ with replacement of the charge factor $`e^2`$ with $`(e^4e^2^2)`$. In the $`\overline{\mathrm{MS}}`$ scheme, the moments of the NLO parton distributions are written in large $`n`$ limit as $`\mathrm{\Delta }q_S^\gamma (n,Q^2,P^2)\|_{\mathrm{NLO},\overline{\mathrm{MS}}}`$ $``$ $`{\displaystyle \frac{\alpha }{4\pi }}N_fe^26{\displaystyle \frac{\mathrm{ln}n}{n}},`$ (5.3) $`\mathrm{\Delta }G^\gamma (n,Q^2,P^2)\|_{\mathrm{NLO},\overline{\mathrm{MS}}}`$ $``$ $`{\displaystyle \frac{\alpha }{4\pi }}N_fe^23{\displaystyle \frac{1}{n^2}}.`$ (5.4) So we have near $`x=1`$ $`\mathrm{\Delta }q_S^\gamma (x,Q^2,P^2)\|_{\mathrm{NLO},\overline{\mathrm{MS}}}`$ $``$ $`{\displaystyle \frac{\alpha }{4\pi }}N_fe^26\left[\mathrm{ln}(1x)\right],`$ (5.5) $`\mathrm{\Delta }G^\gamma (x,Q^2,P^2)\|_{\mathrm{NLO},\overline{\mathrm{MS}}}`$ $``$ $`{\displaystyle \frac{\alpha }{4\pi }}N_fe^23\left[\mathrm{ln}x\right].`$ (5.6) It is remarkable that, in the $`\overline{\mathrm{MS}}`$ scheme, quark parton distributions, $`\mathrm{\Delta }q_S^\gamma (x)\|_{\mathrm{NLO},\overline{\mathrm{MS}}}`$ and $`\mathrm{\Delta }q_{NS}^\gamma (x)\|_{\mathrm{NLO},\overline{\mathrm{MS}}}`$ positively diverge as $`[\mathrm{ln}(1x)]`$ for $`x1`$. Recall that $`\mathrm{\Delta }G^\gamma (x,Q^2,P^2)\|_{\mathrm{NLO}}`$ is the same among the schemes which we consider in this paper. The NLO quark distributions in the CI, AB, AR and $`\mathrm{DIS}_\gamma `$ schemes also diverge as $`x1`$, since their moments behave as $`\mathrm{ln}n/n`$ in the large $`n`$-limit. We find for large $`x`$, $`\mathrm{\Delta }q_S^\gamma (x,Q^2,P^2)\|_{\mathrm{NLO},\mathrm{CI}}`$ $``$ $`{\displaystyle \frac{\alpha }{4\pi }}N_fe^26\left[\mathrm{ln}(1x)\right],`$ (5.7) $`\mathrm{\Delta }q_S^\gamma (x,Q^2,P^2)\|_{\mathrm{NLO},\mathrm{AB}}`$ $``$ $`{\displaystyle \frac{\alpha }{4\pi }}N_fe^26\left[\mathrm{ln}(1x)+2\right],`$ (5.8) $`\mathrm{\Delta }q_S^\gamma (x,Q^2,P^2)\|_{\mathrm{NLO},\mathrm{AR}}`$ $``$ $`{\displaystyle \frac{\alpha }{4\pi }}N_fe^218\left[\mathrm{ln}(1x)\right],`$ (5.9) $`\mathrm{\Delta }q_S^\gamma (x,Q^2,P^2)\|_{\mathrm{NLO},\mathrm{DIS}_\gamma }`$ $``$ $`{\displaystyle \frac{\alpha }{4\pi }}N_fe^26\mathrm{ln}(1x).`$ (5.10) It is noted that $`\mathrm{\Delta }q_S^\gamma (x,Q^2,P^2)\|_{\mathrm{NLO},\mathrm{DIS}_\gamma }`$ negatively diverges as $`x1`$. This is due to the fact that the photonic coefficient function $`\mathrm{\Delta }C_\gamma ^\gamma (x)`$, which in $`\overline{\mathrm{MS}}`$ becomes negative and divergent for $`x1`$ , is absorbed into the quark distributions in the $`\mathrm{DIS}_\gamma `$ scheme. On the other hand, the OS scheme gives quite different behaviors near $`x=1`$ for the quark distributions. Since the typical two-loop anomalous dimensions in the OS scheme behave in the large $`n`$-limit as $$\mathrm{\Delta }\gamma _{NS,\mathrm{OS}}^{(1),n}\mathrm{\Delta }\gamma _{qq,\mathrm{OS}}^{(1),n}\mathrm{ln}^2n,\mathrm{\Delta }K_{S,\mathrm{OS}}^{(1),n}\frac{\mathrm{ln}n}{n},$$ (5.11) while in the $`\overline{\mathrm{MS}}`$ scheme $$\mathrm{\Delta }\gamma _{NS,\overline{\mathrm{MS}}}^{(1),n}\mathrm{\Delta }\gamma _{qq,\overline{\mathrm{MS}}}^{(1),n}\mathrm{ln}n,\mathrm{\Delta }K_{S,\overline{\mathrm{MS}}}^{(1),n}\frac{\mathrm{ln}^2n}{n},$$ (5.12) we find that the moment of $`\mathrm{\Delta }q_S^\gamma (n,Q^2,P^2)\|_{\mathrm{NLO}}`$ in the OS scheme is expressed in the large $`n`$-limit as $$\mathrm{\Delta }q_S^\gamma (n,Q^2,P^2)\|_{\mathrm{NLO},\mathrm{OS}}\frac{\alpha }{4\pi }N_fe^2\left[\frac{69}{8}+\frac{3}{4}N_f\right]\frac{1}{n}$$ (5.13) In $`x`$ space, $`\mathrm{\Delta }q_S^\gamma (x,Q^2,P^2)\|_{\mathrm{NLO},\mathrm{OS}}`$ does not diverge for $`x1`$ but approaches a constant value: $$\mathrm{\Delta }q_S^\gamma (x,Q^2,P^2)\|_{\mathrm{NLO},\mathrm{OS}}\frac{\alpha }{4\pi }N_fe^2\left[\frac{69}{8}+\frac{3}{4}N_f\right].$$ (5.14) Therefore, as far as the large $`x`$-behaviors of quark distributions, and gluon and photon coefficient functions (see Eqs.(LABEL:xCoeffiGluon\- LABEL:xCoeffiGamma) below) are concerned, the OS scheme is more appropriate than other schemes in the sense that they remain finite. Also the quark coefficient function in the OS scheme has a milder divergence for $`x1`$ than those predicted in other schemes (see Eq.(5.17)). Before ending this section, we now show that, as $`x1`$, the polarized virtual photon structure function $`g_1^\gamma (x,Q^2,P^2)`$ approaches a constant value $$\kappa =\frac{\alpha }{4\pi }N_fe^4\left[\frac{51}{8}+\frac{3}{4}N_f\right],$$ (5.15) in NLO. The result is, of course, factorization-scheme independent. It is interesting to note that the constant value $`\kappa `$ coincides exactly with the one given in Eq.(4.39) of Ref., which was derived as the large $`n`$ limit of the moment of the NLO term $`b_2(x)`$ for the unpolarized structure function $`F_2^\gamma `$ . In the leading order, Eq.(5.1) tells us that $`g_1^\gamma (x,Q^2,P^2)\|_{\mathrm{LO}}`$ $`=`$ $`e^2\mathrm{\Delta }q_S^\gamma (x,Q^2,P^2)\|_{\mathrm{LO}}+\mathrm{\Delta }q_{NS}^\gamma (x,Q^2,P^2)\|_{\mathrm{LO}}`$ (5.16) $``$ $`{\displaystyle \frac{\alpha }{4\pi }}{\displaystyle \frac{4\pi }{\alpha _s(Q^2)}}N_fe^4{\displaystyle \frac{9}{4}}{\displaystyle \frac{1}{\mathrm{ln}(1x)}},`$ and thus $`g_1^\gamma (x,Q^2,P^2)\|_{\mathrm{LO}}`$ vanishes as $`x1`$. In order to analyze the large $`x`$-bahavior of the next-leading order $`g_1^\gamma (x,Q^2,P^2)\|_{\mathrm{NLO}}`$, we need information on the coefficient functions. Note that $`\mathrm{\Delta }B_{NS}^n\|_a=\mathrm{\Delta }B_S^n\|_a`$. They behave, as $`x1`$, $`\mathrm{\Delta }B_S(x)\|_a`$ $``$ $`\{\begin{array}{cc}2C_F\left[\frac{2\mathrm{ln}(1x)}{1x}\right]_+\hfill & \text{for }a=\overline{\mathrm{MS}},\mathrm{CI},\mathrm{AB},\mathrm{DIS}_\gamma \text{ ,}\hfill \\ & \\ 2C_F\left[\frac{2\mathrm{ln}(1x)}{1x}\right]_+\hfill & \text{for }a=\mathrm{AR}\text{ ,}\hfill \\ & \\ 3C_F\frac{1}{(1x)_+}\hfill & \text{for }a=\mathrm{OS}\text{ ,}\hfill \end{array}`$ (5.17) $`\mathrm{\Delta }B_G(x)\|_a`$ $``$ $`\{\begin{array}{cc}2N_f\mathrm{ln}(1x)\hfill & \text{for }a=\overline{\mathrm{MS}},\mathrm{CI},\mathrm{AB},\mathrm{AR},\mathrm{DIS}_\gamma \text{ ,}\hfill \\ & \\ 4N_f\hfill & \text{for }a=\mathrm{OS}\text{ ,}\hfill \end{array}`$ (5.18) $`\mathrm{\Delta }C(x)_\gamma ^\gamma \|_a`$ $``$ $`\{\begin{array}{cc}\frac{\alpha }{4\pi }e^412N_f\mathrm{ln}(1x)\hfill & \text{for }a=\overline{\mathrm{MS}},\mathrm{CI},\mathrm{AB},\mathrm{AR}\text{ ,}\hfill \\ & \\ \frac{\alpha }{4\pi }e^424N_f\hfill & \text{for }a=\mathrm{OS}\text{ ,}\hfill \\ & \\ 0\hfill & \text{for }a=\mathrm{DIS}_\gamma \text{ .}\hfill \end{array}`$ (5.19) The coefficient functions $`\mathrm{\Delta }B_S(x)`$ and $`\mathrm{\Delta }B_G(x)`$ are the same that appear in the polarized nucleon structure function $`g_1(x,Q^2)`$. The $`\mathrm{\Delta }B_S(x)`$ in all schemes considered here diverges as $`x1`$, but OS scheme gives a milder divergence for $`\mathrm{\Delta }B_S(x)`$ than other schemes. Also note that $`\mathrm{\Delta }B_G(x)\|_{OS}`$ remains finite as $`x1`$, but $`\mathrm{\Delta }B_G(x)`$ in other schemes negatively diverge. Let us write $`g_1^\gamma (x,Q^2,P^2)\|_{\mathrm{NLO}}`$ in terms of partonic contributions as follows: $$g_1^\gamma (x,Q^2,P^2)\|_{\mathrm{NLO}}=g_1^\gamma (x)\|_{\mathrm{NLO}}^{\mathrm{quark}}+g_1^\gamma (x)\|_{\mathrm{NLO}}^{\mathrm{gluon}}+\mathrm{\Delta }C_\gamma ^\gamma (x),$$ (5.20) where $`g_1^\gamma (x)\|_{\mathrm{NLO}}^{\mathrm{quark}}`$ $``$ $`e^2\mathrm{\Delta }q_S^\gamma (x,Q^2,P^2)\|_{\mathrm{NLO}}+\mathrm{\Delta }q_{NS}^\gamma (x,Q^2,P^2)\|_{\mathrm{NLO}}`$ (5.21) $`+e^2{\displaystyle \frac{\alpha _s(Q^2)}{4\pi }}\mathrm{\Delta }B_S(x)\mathrm{\Delta }q_S^\gamma (x,Q^2,P^2)\|_{\mathrm{LO}}`$ $`+{\displaystyle \frac{\alpha _s(Q^2)}{4\pi }}\mathrm{\Delta }B_{NS}(x)\mathrm{\Delta }q_{NS}^\gamma (x,Q^2,P^2)\|_{\mathrm{LO}},`$ $`g_1^\gamma (x)\|_{\mathrm{NLO}}^{\mathrm{gluon}}`$ $``$ $`e^2{\displaystyle \frac{\alpha _s(Q^2)}{4\pi }}\mathrm{\Delta }B_G(x)\mathrm{\Delta }G^\gamma (x,Q^2,P^2)\|_{\mathrm{LO}}.`$ (5.22) Then we find, for $`x1`$, $`g_1^\gamma (x)\|_{\mathrm{NLO}}^{\mathrm{quark}}`$ $``$ $`\{\begin{array}{cc}\frac{\alpha }{4\pi }e^412N_f\mathrm{ln}(1x)\hfill & \text{for }a=\overline{\mathrm{MS}},\mathrm{CI},\mathrm{AB},\mathrm{AR}\text{ ,}\hfill \\ & \\ \frac{\alpha }{4\pi }e^4N_f\left[\frac{141}{8}+\frac{3}{4}N_f\right]\hfill & \text{for }a=\mathrm{OS}\text{ ,}\hfill \\ & \\ \frac{\alpha }{4\pi }e^4N_f\left[\frac{51}{8}+\frac{3}{4}N_f\right]\hfill & \text{for }a=\mathrm{DIS}_\gamma \text{ ,}\hfill \end{array}`$ (5.23) The NLO gluon contribution $`g_1^\gamma (x,Q^2,P^2)\|_{\mathrm{NLO}}^{\mathrm{gluon}}`$ vanishes faster than $`(\mathrm{ln}x)^2`$ in any scheme under consideration. As for the NLO quark contribution, $`g_1^\gamma (x,Q^2,P^2)\|_{\mathrm{NLO}}^{\mathrm{quark}}`$ in $`\overline{\mathrm{MS}}`$, CI, AB, AR schemes, diverges as $`[\mathrm{ln}(1x)]`$ for $`x1`$. However, Eq.(LABEL:xCoeffiGamma) shows that the one-loop photon coefficient function $`\mathrm{\Delta }C_\gamma ^\gamma (x)`$ in these schemes also diverges as $`[\mathrm{ln}(1x)]`$ with the opposite sign and the sum becomes finite. On the other hand, in the OS scheme, we observe from Eqs.(LABEL:QuarkContriNLO) and (LABEL:xCoeffiGamma) that both the quark contribution and photon coefficient function remain finite as $`x1`$, and it is easily seen that the sum $$g_1^\gamma (x)\|_{\mathrm{NLO},\mathrm{OS}}^{\mathrm{quark}}+\mathrm{\Delta }C(x)_{\gamma ,\mathrm{OS}}^\gamma $$ (5.24) approaches the constant value $`\kappa `$ given in Eq.(5.15). In the $`\mathrm{DIS}_\gamma `$ scheme, the NLO quark contribution $`g_1^\gamma (x)\|_{\mathrm{NLO},\mathrm{DIS}_\gamma }^{\mathrm{quark}}`$ reaches the finite value $`\kappa `$ as $`x1`$, since $`\mathrm{\Delta }C(x)_{\gamma ,\mathrm{DIS}_\gamma }^\gamma 0`$. In fact, as we see from Eq.(5.21), $`g_1^\gamma (x)\|_{\mathrm{NLO}}^{\mathrm{quark}}`$ is made up of two parts, the one from $`\mathrm{\Delta }q_S^\gamma (x,Q^2,P^2)\|_{\mathrm{NLO}}`$ and the other from $`\mathrm{\Delta }B_S(x)\mathrm{\Delta }q_S^\gamma (x,Q^2,P^2)\|_{\mathrm{LO}}`$, plus their non-singlet quark counterparts. In $`\mathrm{DIS}_\gamma `$, both contributions diverge as $`x1`$, but with the opposite sign, and the sum remains finite. The constant value $`\kappa `$ in Eq.(5.15) is negative unless $`N_f9`$. Consequently, it seems superficially that QCD with 8 flavors or less predicts that the structure function $`g_1^\gamma (x,Q^2,P^2)`$ turns out to be negative for $`x`$ very close to 1, since the leading term $`g_1^\gamma (x,Q^2,P^2)\|_{\mathrm{LO}}`$ vanishes as $`x1`$. But the fact is that $`x`$ cannot reach exactly one. The constraint $`(p+q)^20`$ gives $$xx_{\mathrm{max}}=\frac{Q^2}{Q^2+P^2},$$ (5.25) and we find $$g_1^\gamma (x=x_{\mathrm{max}},Q^2,P^2)\|_{\mathrm{LO}}>\frac{\alpha }{4\pi }N_fe^4\frac{3}{C_F}\beta _0$$ (5.26) and the sum $`g_1^\gamma (x=x_{\mathrm{max}},Q^2,P^2)\|_{\mathrm{LO}+\mathrm{NLO}}`$ is indeed positive. ## 6 Numerical analysis The parton distribution functions are recovered from the moments by the inverse Mellin transformation. In Fig. 2 we plot the factorization scheme dependence of the singlet quark distribution $`\mathrm{\Delta }q_S^\gamma (x,Q^2,P^2)`$ beyond the LO in units of $`(3N_fe^2\alpha /\pi )\mathrm{ln}(Q^2/P^2)`$. We have taken $`N_f=3`$, $`Q^2=30\mathrm{GeV}^2`$, $`P^2=1\mathrm{GeV}^2`$, and the QCD scale parameter $`\mathrm{\Lambda }=0.2\mathrm{GeV}`$. All four CI-like (i.e., CI, AB, OS and AR) curves cross the $`x`$-axis nearly at the same point, just below $`x=0.5`$, while the $`\overline{\mathrm{MS}}`$ curve crosses at above $`x=0.5`$. This is understandable since we saw from Eqs.(4.13, 4.20) that the first moment of $`\mathrm{\Delta }q_S^\gamma `$ vanishes in the CI-like schemes while it is negative in the $`\overline{\mathrm{MS}}`$ scheme. The $`\mathrm{DIS}_\gamma `$ curve crosses the $`x`$-axis below $`x=0.5`$, though the first moment of $`\mathrm{\Delta }q_S^\gamma \|_{\mathrm{DIS}_\gamma }`$ is negative, taking the same value with the one in the $`\overline{\mathrm{MS}}`$ scheme. Comparing the $`\mathrm{DIS}_\gamma `$ curve at large $`x`$ with the $`\overline{\mathrm{MS}}`$ one, we will see that rapid dropping of the $`\mathrm{DIS}_\gamma `$ curve as $`x1`$ drives the crossing point below $`x=0.5`$. As $`x1`$, we observe that the $`\overline{\mathrm{MS}}`$, CI, AB, and AR curves continue to increase. In fact we see that the $`\overline{\mathrm{MS}}`$ and CI curves tend to merge, the AB curve comes above those two curves and the AR curve diverges more rapidly than the other three. On the other hand, the OS and $`\mathrm{DIS}_\gamma `$ curves start to drop at large $`x`$. The OS curve continues to increase till near $`x=1`$, and then starts to drop to reach a finite positive value. The $`\mathrm{DIS}_\gamma `$ curve reaches maximum at $`x0.8`$ and drops to negative values. These behaviors are inferred from Eqs.(5.5, 5.7-5.10, 5.14). Concerning the non-singlet quark distribution $`\mathrm{\Delta }q_{NS}^\gamma (x,Q^2,P^2)`$, we find that when we take into account the charge factors, it falls on the singlet quark distribution in almost all $`x`$ region; namely two “normalized” distributions $`\mathrm{\Delta }\stackrel{~}{q}_S^\gamma \mathrm{\Delta }q_S^\gamma /e^2`$ and $`\mathrm{\Delta }\stackrel{~}{q}_{NS}^\gamma \mathrm{\Delta }q_{NS}^\gamma /(e^4e^2^2)`$ mostly overlap except at very small $`x`$ region. The situation is the same in all factorization schemes we have studied in this paper. This is attributable to the fact that once the charge factors are taken into account, the evolution equations for both $`\mathrm{\Delta }\stackrel{~}{q}_S^\gamma `$ and $`\mathrm{\Delta }\stackrel{~}{q}_{NS}^\gamma `$ have the same inhomogeneous LO and NLO $`\mathrm{\Delta }K`$ terms and the same initial conditions at $`Q^2=P^2`$ (see Eq.(3.19)). In Fig. 3 we plot again the OS and $`\mathrm{DIS}_\gamma `$ predictions for $`\mathrm{\Delta }q_S^\gamma (x,Q^2,P^2)`$ together with the LO result. The motivation of having introduced $`\mathrm{DIS}_\gamma `$ scheme into the analysis of the unpolaorized (polarized) real photon structure function $`F_2^\gamma `$ ($`g_1^\gamma `$) was to reduce the discrepancies at large-$`x`$ region between the LO and the NLO results for the ‘pointlike’ part of $`F_2^\gamma `$ ($`g_1^\gamma `$). When applied to the polarized virtual photon case, it is seen from Fig. 2 and 3 that $`\mathrm{DIS}_\gamma `$ scheme gives a better behavior for $`\mathrm{\Delta }q_S^\gamma (x,Q^2,P^2)`$ at large $`x`$ than $`\overline{\mathrm{MS}}`$ in the sense that $`\mathrm{DIS}_\gamma `$ curve is closer to the LO result. However, we observe that absorbing the photonic coefficient function $`\mathrm{\Delta }C_\gamma ^\gamma `$ into the quark distributions in the $`\mathrm{DIS}_\gamma `$ scheme has too much effect on their large-$`x`$ behaviors: The $`\mathrm{DIS}_\gamma `$ curve for $`\mathrm{\Delta }q_S^\gamma (x,Q^2,P^2)`$ goes under the LO one at $`x0.6`$ and the difference between the two grows as $`x1`$. In fact the $`\mathrm{DIS}_\gamma `$ curve drops to negative values near at $`x=1`$. From the viewpoint of ‘perturbative stabilities’ we find that the OS curve shows more appropriate behavior than the others. We see from Fig. 3 that the differences between the OS and LO curves are very small for the range $`0.05<x<0.7`$. And the OS curve comes above the LO for $`x>0.7`$. Fig. 4 shows the $`Q^2`$-dependence of $`\mathrm{\Delta }q_S^\gamma (x,Q^2,P^2)`$ in the OS scheme in units of $`(3N_fe^2\alpha /\pi )\mathrm{ln}(Q^2/P^2)`$. Three curves with $`Q^2=30,50`$ and $`100\mathrm{GeV}^2`$ almost overlap in whole $`x`$ region except in the vicinity of $`x=1`$. We see from Fig. 4 that, in the OS scheme, $`\mathrm{\Delta }q_S^\gamma `$ beyond the LO behaves approximately as the one obtained from the box (tree) diagram calculation, $$\mathrm{\Delta }q_S^{\gamma (\mathrm{Box})}(x,Q^2,P^2)=(2x1)3N_fe^2\frac{\alpha }{\pi }\mathrm{ln}\frac{Q^2}{P^2}.$$ (6.1) The gluon distribution $`\mathrm{\Delta }G^\gamma (x,Q^2,P^2)`$ beyond the LO is shown in Fig. 5 in units of $`(3N_fe^2\alpha /\pi )\mathrm{ln}(Q^2/P^2)`$, with three different $`Q^2`$ values. Recall that every scheme considered in this paper predicts the same behavior for the gluon distribution up to NLO. We do not see much difference in three curves with different $`Q^2`$. This means the $`\mathrm{\Delta }G^\gamma `$ is approximately proportional to $`\mathrm{ln}(Q^2/P^2)`$ . But, compared with quark distributions, $`\mathrm{\Delta }G^\gamma `$ is very much small in absolute value except at the small $`x`$ region. In Fig. 6 we plot the virtual photon structure function $`g_1^\gamma (x,Q^2,P^2)`$ in the NLO for $`N_f=3`$, $`Q^2=30`$ GeV<sup>2</sup> and $`P^2=1`$ GeV<sup>2</sup> and the QCD scale parameter $`\mathrm{\Lambda }=0.2`$ GeV. The vertical axis corresponds to $$g_1^\gamma (x,Q^2,P^2)/\frac{3\alpha }{\pi }N_f<e^4>\mathrm{ln}\frac{Q^2}{P^2}.$$ (6.2) Also shown are the LO result, the Box (tree) diagram contribution, $$g_1^{\gamma (\mathrm{Box})}(x,Q^2,P^2)=(2x1)\frac{3\alpha }{\pi }N_f<e^4>\mathrm{ln}\frac{Q^2}{P^2},$$ (6.3) and the Box diagram contribution including non-leading (NL) correction with mass being ignored $$g_1^{\gamma (\mathrm{Box}(\mathrm{NL}))}(x,Q^2,P^2)=\frac{3\alpha }{\pi }N_f<e^4>\left[(2x1)\mathrm{ln}\frac{Q^2}{P^2}2(2x1)(\mathrm{ln}x+1)\right].$$ (6.4) In our previous paper , there was an error in the program for numerical evaluation of the NLO $`g_1^\gamma (x,Q^2,P^2)`$. The corrected graph (NLO curve) here is different from the corresponding one in Fig.2 of Ref.. The new NLO curve appears lower than the previous one for $`x<0.7`$ and rather enhanced above $`x=0.7`$. We observe that the corrected NLO curve remains below the LO one, and that the NLO QCD corrections are significant at large $`x`$ as well as at low $`x`$. For the case of the real photon target, $`P^2=0`$, the structure function can be decomposed as $$g_1^\gamma (x,Q^2)=g_1^\gamma (x,Q^2)\|_{\mathrm{pert}.}+g_1^\gamma (x,Q^2)\|_{\mathrm{non}\mathrm{pert}.}.$$ (6.5) The first term, the point-like piece, can be calculated in a perturbative method. Actually, it can be obtained by setting $`P^2=\mathrm{\Lambda }^2`$ in the expressions of parton distributions in Eq.(2.2) or (2.3). The second term can only be computed by some non-perturbative methods. In Fig.7, we plot the point-like piece of the real photon $`g_1^\gamma (x,Q^2)`$ in the NLO, together with the LO result and the Box (tree) diagram contribution. The NLO curve, which is calculated by the corrected computer program, is different from the previous one in Fig. 6 in ref.. The new NLO curve appears lower than the previous one for $`x<0.6`$ and enhanced above $`x=0.6`$. Also it remains below the LO curve. The NLO result qualitatively consistent with the analysis by Stratmann and Vogelsang . In the unpolarized case, the moment of $`F_2^\gamma `$ has a singularity at $`n=2`$ which leads to the negative structure function at low $`x`$. Thus we need some regularization prescription to recover positive structure function as discussed in Refs. . Note that we do not have such complication at $`n=1`$ for the polarized case. Finally, in our numerical analyis, we took $`P^2=1\mathrm{G}\mathrm{e}\mathrm{V}^2`$, which may not be necessarily large enough for the non-perturbative effects to be dying away. For our normalized parton distributions, however, the larger values of $`P^2`$ would not give any sizable change in shape and magnitude. ## 7 Conclusion In the present paper, we have studied in detail the spin-dependent parton distributions inside the virtual photon, which can be predicted entirely up to NLO in perturbative QCD. The virtual photon target provides a good testing ground for examining the factorization scheme dependence of the quark and gluon distributions. We have investigated the polarized parton distributions in several different factorization schemes. We derived the explicit transformation rules from one scheme to another for the coefficient functions, the finite photon matrix elements and the two-loop anomalous dimensions or parton splitting functions. In particular, we studied the QCD and QED axial anomaly effects on the first moments of quark distributions to see the interplay between the axial anomalies and factorization schemes. We find that, in the CI-like schemes, the first moments of polarized quark distributions, both flavor singlet and non-singlet, vanish in NLO while the standard $`\overline{\mathrm{MS}}`$ scheme gives the non-zero value. Also we find that the large $`x`$-behaviors of polarized quark distributons dramatically vary from one factorization scheme to another. Indeed, for $`x1`$, the quark distributions positively diverge or negatively diverge or remain finite, depending on factorization schemes. The numerical analyses performed for the parton distributions reassures the above observations. From the viewpoint of ‘perturbative stabilities’ the OS scheme gives more appropriate behaviors for the quark distributions than the others. The gluon distribution turns out to be the same up to NLO among the six factorization schemes examined. Furthermore, its first moment is found to be factorization-scheme independent up to NLO. The same analysis on the factorization scheme dependence of the unpolarized parton distributions of the virtual photon can be carried out and will be discussed elsewhere. Acknowledgement We thank J. Blümlein, S. J. Brodsky, J. Kodaira, M. Stratmann, O. V. Teryaev and W.L. van Neerven for valuable discussions. One of the authors (K.S.) would like to thank J. Blümlein and T. Riemann for the hospitality extended to him at the 5th Zeuthen Workshop “Loops and Legs 2000”. Some of the results in this paper have been reported at the workshop. We also thank E. Reya and C. Sieg for the communication regarding the numerical evaluation of the structure function. This work is partially supported by the Monbusho Grant-in-Aid for Scientific Research NO.(C)(2)-12640266. Appendix ## Appendix A NLO expressions for polarized parton distributions in the virtual photon We give the explicit expressions of $`\mathrm{\Delta }q_S^\gamma `$, $`\mathrm{\Delta }G^\gamma `$, and $`\mathrm{\Delta }q_{NS}^\gamma `$ up to NLO. They are written in terms of one-(two-) loop anomalous dimensions $`\mathrm{\Delta }\gamma _{ij}^{0,n}`$ ($`\mathrm{\Delta }\gamma _{ij}^{(1),n}`$) ($`i,j=\psi ,G`$), $`\mathrm{\Delta }\gamma _{NS}^{0,n}`$ ($`\gamma _{NS}^{(1),n}`$), $`\mathrm{\Delta }K_l^{0,n}`$ ($`\mathrm{\Delta }K_l^{(1),n}`$) ($`l=\psi ,G,NS`$), and the one-loop photon matrix elements of hadronic operators, $`\mathrm{\Delta }A_l^n`$. The expressions of one-loop and $`\overline{\mathrm{MS}}`$ scheme-two-loop anomalous dimensions are found, for example, in Appendix of Ref.. ### A.1 Singlet quark distribution $`\mathrm{\Delta }q_S^\gamma (n,Q^2,P^2)/{\displaystyle \frac{\alpha }{8\pi \beta _0}}`$ $`={\displaystyle \frac{4\pi }{\alpha _s(Q^2)}}\widehat{L}_S^{+n}\left\{1\left[{\displaystyle \frac{\alpha _s(Q^2)}{\alpha _s(P^2)}}\right]^{\lambda _+^n/2\beta _0+1}\right\}+{\displaystyle \frac{4\pi }{\alpha _s(Q^2)}}\widehat{L}_S^n\left\{1\left[{\displaystyle \frac{\alpha _s(Q^2)}{\alpha _s(P^2)}}\right]^{\lambda _{}^n/2\beta _0+1}\right\}`$ $`+\widehat{A}_S^{+n}\left\{1\left[{\displaystyle \frac{\alpha _s(Q^2)}{\alpha _s(P^2)}}\right]^{\lambda _+^n/2\beta _0}\right\}+\widehat{A}_S^n\left\{1\left[{\displaystyle \frac{\alpha _s(Q^2)}{\alpha _s(P^2)}}\right]^{\lambda _{}^n/2\beta _0}\right\}`$ $`+\widehat{B}_S^{+n}\left\{1\left[{\displaystyle \frac{\alpha _s(Q^2)}{\alpha _s(P^2)}}\right]^{\lambda _+^n/2\beta _0+1}\right\}+\widehat{B}_S^n\left\{1\left[{\displaystyle \frac{\alpha _s(Q^2)}{\alpha _s(P^2)}}\right]^{\lambda _{}^n/2\beta _0+1}\right\}`$ $`+\widehat{C}_S^n`$ (A.1) where $`\widehat{L}_S^{+n}`$ $`=`$ $`\mathrm{\Delta }K_\psi ^{0,n}{\displaystyle \frac{\mathrm{\Delta }\gamma _{\psi \psi }^{0,n}\lambda _{}^n}{\lambda _+^n\lambda _{}^n}}{\displaystyle \frac{1}{1+\lambda _+^n/2\beta _0}}`$ (A.2) $`\widehat{L}_S^n`$ $`=`$ $`\mathrm{\Delta }K_\psi ^{0,n}{\displaystyle \frac{\mathrm{\Delta }\gamma _{\psi \psi }^{0,n}\lambda _+^n}{\lambda _{}^n\lambda _+^n}}{\displaystyle \frac{1}{1+\lambda _{}^n/2\beta _0}}`$ (A.3) $`\mathrm{with}`$ $`\lambda _\pm ^n`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left\{\mathrm{\Delta }\gamma _{\psi \psi }^{0,n}+\mathrm{\Delta }\gamma _{GG}^{0,n}\pm \left[(\mathrm{\Delta }\gamma _{\psi \psi }^{0,n}\mathrm{\Delta }\gamma _{GG}^{0,n})^2+4\mathrm{\Delta }\gamma _{\psi G}^{0,n}\mathrm{\Delta }\gamma _{G\psi }^{0,n}\right]^{1/2}\right\}`$ (A.4) $`\beta _0`$ $`=`$ $`112N_f/3,\beta _1=10238N_f/3.`$ (A.5) and $`\widehat{A}_S^{+n}`$ $`=`$ $`{\displaystyle \frac{1}{\lambda _+^n(\lambda _+^n\lambda _{}^n)(2\beta _0+\lambda _{}^n\lambda _+^n)}}`$ (A.6) $`\times [\mathrm{\Delta }K_\psi ^{0,n}\{(\mathrm{\Delta }\gamma _{\psi \psi }^{0,n}2\beta _0\lambda _{}^n)\mathrm{\Delta }\gamma _{\psi \psi }^{(1),n}+\mathrm{\Delta }\gamma _{G\psi }^{0,n}\mathrm{\Delta }\gamma _{\psi G}^{(1),n}\}(\mathrm{\Delta }\gamma _{\psi \psi }^{0,n}\lambda _{}^n)`$ $`+\mathrm{\Delta }K_\psi ^{0,n}\left\{(\mathrm{\Delta }\gamma _{\psi \psi }^{0,n}2\beta _0\lambda _{}^n)\mathrm{\Delta }\gamma _{G\psi }^{(1),n}+\mathrm{\Delta }\gamma _{G\psi }^{0,n}\mathrm{\Delta }\gamma _{GG}^{(1),n}\right\}\mathrm{\Delta }\gamma _{\psi G}^{0,n}`$ $`+2\beta _0(2\beta _0+\lambda _{}^n\lambda _+^n)\left\{\mathrm{\Delta }K_\psi ^{(1),n}(\mathrm{\Delta }\gamma _{\psi \psi }^{0,n}\lambda _{}^n)+\mathrm{\Delta }K_G^{(1),n}\mathrm{\Delta }\gamma _{\psi G}^{0,n}\right\}`$ $`2\beta _0(2\beta _0+\lambda _{}^n\lambda _+^n)\lambda _+^n\mathrm{\Delta }A_\psi ^n(\mathrm{\Delta }\gamma _{\psi \psi }^{0,n}\lambda _{}^n)`$ $`{\displaystyle \frac{\beta _1}{\beta _0}}\mathrm{\Delta }K_\psi ^{0,n}(2\beta _0+\lambda _{}^n\lambda _+^n)(2\beta _0\lambda _+^n)(\mathrm{\Delta }\gamma _{\psi \psi }^{0,n}\lambda _{}^n)]`$ $`\widehat{A}_S^n`$ $`=`$ $`{\displaystyle \frac{1}{\lambda _{}^n(\lambda _{}^n\lambda _+^n)(2\beta _0+\lambda _+^n\lambda _{}^n)}}`$ (A.7) $`\times [\mathrm{\Delta }K_\psi ^{0,n}\{(\mathrm{\Delta }\gamma _{\psi \psi }^{0,n}2\beta _0\lambda _+^n)\mathrm{\Delta }\gamma _{\psi \psi }^{(1),n}+\mathrm{\Delta }\gamma _{G\psi }^{0,n}\mathrm{\Delta }\gamma _{\psi G}^{(1),n}\}(\mathrm{\Delta }\gamma _{\psi \psi }^{0,n}\lambda _+^n)`$ $`+\mathrm{\Delta }K_\psi ^{0,n}\left\{(\mathrm{\Delta }\gamma _{\psi \psi }^{0,n}2\beta _0\lambda _+^n)\mathrm{\Delta }\gamma _{G\psi }^{(1),n}+\mathrm{\Delta }\gamma _{G\psi }^{0,n}\mathrm{\Delta }\gamma _{GG}^{(1),n}\right\}\mathrm{\Delta }\gamma _{\psi G}^{0,n}`$ $`+2\beta _0(2\beta _0+\lambda _+^n\lambda _{}^n)\left\{\mathrm{\Delta }K_\psi ^{(1),n}(\mathrm{\Delta }\gamma _{\psi \psi }^{0,n}\lambda _+^n)+\mathrm{\Delta }K_G^{(1),n}\mathrm{\Delta }\gamma _{\psi G}^{0,n}\right\}`$ $`2\beta _0(2\beta _0+\lambda _+^n\lambda _{}^n)\lambda _{}^n\mathrm{\Delta }A_\psi ^n(\mathrm{\Delta }\gamma _{\psi \psi }^{0,n}\lambda _+^n)`$ $`{\displaystyle \frac{\beta _1}{\beta _0}}\mathrm{\Delta }K_\psi ^{0,n}(2\beta _0+\lambda _+^n\lambda _{}^n)(2\beta _0\lambda _{}^n)(\mathrm{\Delta }\gamma _{\psi \psi }^{0,n}\lambda _+^n)]`$ $`\widehat{B}_S^{+n}`$ $`=`$ $`\mathrm{\Delta }K_\psi ^{0,n}{\displaystyle \frac{1}{(2\beta _0+\lambda _+^n)(\lambda _+^n\lambda _{}^n)(2\beta _0+\lambda _+^n\lambda _{}^n)}}`$ (A.8) $`\times [\{(\mathrm{\Delta }\gamma _{\psi \psi }^{0,n}\lambda _{}^n)\mathrm{\Delta }\gamma _{\psi \psi }^{(1),n}+\mathrm{\Delta }\gamma _{G\psi }^{0,n}\mathrm{\Delta }\gamma _{\psi G}^{(1),n}\}(2\beta _0+\mathrm{\Delta }\gamma _{\psi \psi }^{0,n}\lambda _{}^n)`$ $`+\left\{(\mathrm{\Delta }\gamma _{\psi \psi }^{0,n}\lambda _{}^n)\mathrm{\Delta }\gamma _{G\psi }^{(1),n}+\mathrm{\Delta }\gamma _{G\psi }^{0,n}\mathrm{\Delta }\gamma _{GG}^{(1),n}\right\}\mathrm{\Delta }\gamma _{\psi G}^{0,n}`$ $`{\displaystyle \frac{\beta _1}{\beta _0}}(2\beta _0+\lambda _+^n\lambda _{}^n)\lambda _+^n(\mathrm{\Delta }\gamma _{\psi \psi }^{0,n}\lambda _{}^n)]`$ $`\widehat{B}_S^n`$ $`=`$ $`\mathrm{\Delta }K_\psi ^{0,n}{\displaystyle \frac{1}{(2\beta _0+\lambda _{}^n)(\lambda _{}^n\lambda _+^n)(2\beta _0+\lambda _{}^n\lambda _+^n)}}`$ $`\times [\{(\mathrm{\Delta }\gamma _{\psi \psi }^{0,n}\lambda _+^n)\mathrm{\Delta }\gamma _{\psi \psi }^{(1),n}+\mathrm{\Delta }\gamma _{G\psi }^{0,n}\mathrm{\Delta }\gamma _{\psi G}^{(1),n}\}(2\beta _0+\mathrm{\Delta }\gamma _{\psi \psi }^{0,n}\lambda _+^n)`$ $`+\left\{(\mathrm{\Delta }\gamma _{\psi \psi }^{0,n}\lambda _+^n)\mathrm{\Delta }\gamma _{G\psi }^{(1),n}+\mathrm{\Delta }\gamma _{G\psi }^{0,n}\mathrm{\Delta }\gamma _{GG}^{(1),n}\right\}\mathrm{\Delta }\gamma _{\psi G}^{0,n}`$ $`{\displaystyle \frac{\beta _1}{\beta _0}}(2\beta _0+\lambda _{}^n\lambda _+^n)\lambda _{}^n(\mathrm{\Delta }\gamma _{\psi \psi }^{0,n}\lambda _+^n)]`$ $`\widehat{C}_S^n`$ $`=`$ $`2\beta _0\mathrm{\Delta }A_\psi ^n`$ (A.9) ### A.2 Gluon distribution $`\mathrm{\Delta }G^\gamma (n,Q^2,P^2)/{\displaystyle \frac{\alpha }{8\pi \beta _0}}`$ $`={\displaystyle \frac{4\pi }{\alpha _s(Q^2)}}\widehat{L}_G^{+n}\left\{1\left[{\displaystyle \frac{\alpha _s(Q^2)}{\alpha _s(P^2)}}\right]^{\lambda _+^n/2\beta _0+1}\right\}+{\displaystyle \frac{4\pi }{\alpha _s(Q^2)}}\widehat{L}_G^n\left\{1\left[{\displaystyle \frac{\alpha _s(Q^2)}{\alpha _s(P^2)}}\right]^{\lambda _{}^n/2\beta _0+1}\right\}`$ $`+\widehat{A}_G^{+n}\left\{1\left[{\displaystyle \frac{\alpha _s(Q^2)}{\alpha _s(P^2)}}\right]^{\lambda _+^n/2\beta _0}\right\}+\widehat{A}_G^n\left\{1\left[{\displaystyle \frac{\alpha _s(Q^2)}{\alpha _s(P^2)}}\right]^{\lambda _{}^n/2\beta _0}\right\}`$ $`+\widehat{B}_G^{+n}\left\{1\left[{\displaystyle \frac{\alpha _s(Q^2)}{\alpha _s(P^2)}}\right]^{\lambda _+^n/2\beta _0+1}\right\}+\widehat{B}_G^n\left\{1\left[{\displaystyle \frac{\alpha _s(Q^2)}{\alpha _s(P^2)}}\right]^{\lambda _{}^n/2\beta _0+1}\right\}`$ (A.10) where $`\widehat{L}_G^{+n}={\displaystyle \frac{\mathrm{\Delta }K_\psi ^{0,n}\mathrm{\Delta }\gamma _{G\psi }^{0,n}}{\lambda _+^n\lambda _{}^n}}{\displaystyle \frac{1}{1+\lambda _+^n/2\beta _0}}`$ (A.11) $`\widehat{L}_G^n={\displaystyle \frac{\mathrm{\Delta }K_\psi ^{0,n}\mathrm{\Delta }\gamma _{G\psi }^{0,n}}{\lambda _{}^n\lambda _+^n}}{\displaystyle \frac{1}{1+\lambda _{}^n/2\beta _0}}`$ (A.12) and $`\widehat{A}_G^{+n}`$ $`=`$ $`{\displaystyle \frac{1}{\lambda _+^n(\lambda _+^n\lambda _{}^n)(2\beta _0+\lambda _{}^n\lambda _+^n)}}`$ (A.13) $`\times [\mathrm{\Delta }K_\psi ^{0,n}\{(\mathrm{\Delta }\gamma _{\psi \psi }^{0,n}2\beta _0\lambda _{}^n)\mathrm{\Delta }\gamma _{\psi \psi }^{(1),n}+\mathrm{\Delta }\gamma _{G\psi }^{0,n}\mathrm{\Delta }\gamma _{\psi G}^{(1),n}\}\mathrm{\Delta }\gamma _{G\psi }^{0,n}`$ $`+\mathrm{\Delta }K_\psi ^{0,n}\left\{(\mathrm{\Delta }\gamma _{\psi \psi }^{0,n}2\beta _0\lambda _{}^n)\mathrm{\Delta }\gamma _{G\psi }^{(1),n}+\mathrm{\Delta }\gamma _{G\psi }^{0,n}\mathrm{\Delta }\gamma _{GG}^{(1),n}\right\}(\mathrm{\Delta }\gamma _{GG}^{0,n}\lambda _{}^n)`$ $`+2\beta _0(2\beta _0+\lambda _{}^n\lambda _+^n)\left\{\mathrm{\Delta }K_\psi ^{(1),n}\mathrm{\Delta }\gamma _{G\psi }^{0,n}+\mathrm{\Delta }K_G^{(1),n}(\mathrm{\Delta }\gamma _{GG}^{0,n}\lambda _{}^n)\right\}`$ $`2\beta _0(2\beta _0+\lambda _{}^n\lambda _+^n)\lambda _+^n\mathrm{\Delta }A_\psi ^n\mathrm{\Delta }\gamma _{G\psi }^{0,n}`$ $`{\displaystyle \frac{\beta _1}{\beta _0}}\mathrm{\Delta }K_\psi ^{0,n}(2\beta _0+\lambda _{}^n\lambda _+^n)(2\beta _0\lambda _+^n)\mathrm{\Delta }\gamma _{G\psi }^{0,n}]`$ $`\widehat{A}_G^n`$ $`=`$ $`{\displaystyle \frac{1}{\lambda _{}^n(\lambda _{}^n\lambda _+^n)(2\beta _0+\lambda _+^n\lambda _{}^n)}}`$ (A.14) $`\times [\mathrm{\Delta }K_\psi ^{0,n}\{(\mathrm{\Delta }\gamma _{\psi \psi }^{0,n}2\beta _0\lambda _+^n)\mathrm{\Delta }\gamma _{\psi \psi }^{(1),n}+\mathrm{\Delta }\gamma _{G\psi }^{0,n}\mathrm{\Delta }\gamma _{\psi G}^{(1),n}\}\mathrm{\Delta }\gamma _{G\psi }^{0,n}`$ $`+\mathrm{\Delta }K_\psi ^{0,n}\left\{(\mathrm{\Delta }\gamma _{\psi \psi }^{0,n}2\beta _0\lambda _+^n)\mathrm{\Delta }\gamma _{G\psi }^{(1),n}+\mathrm{\Delta }\gamma _{G\psi }^{0,n}\mathrm{\Delta }\gamma _{GG}^{(1),n}\right\}(\mathrm{\Delta }\gamma _{GG}^{0,n}\lambda _+^n)`$ $`+2\beta _0(2\beta _0+\lambda _+^n\lambda _{}^n)\left\{\mathrm{\Delta }K_\psi ^{(1),n}\mathrm{\Delta }\gamma _{G\psi }^{0,n}+\mathrm{\Delta }K_G^{(1),n}(\mathrm{\Delta }\gamma _{GG}^{0,n}\lambda _+^n)\right\}`$ $`2\beta _0(2\beta _0+\lambda _+^n\lambda _{}^n)\lambda _{}^n\mathrm{\Delta }A_\psi ^n\mathrm{\Delta }\gamma _{G\psi }^{0,n}`$ $`{\displaystyle \frac{\beta _1}{\beta _0}}\mathrm{\Delta }K_\psi ^{0,n}(2\beta _0+\lambda _+^n\lambda _{}^n)(2\beta _0\lambda _{}^n)\mathrm{\Delta }\gamma _{G\psi }^{0,n}]`$ $`\widehat{B}_G^{+n}`$ $`=`$ $`\mathrm{\Delta }K_\psi ^{0,n}{\displaystyle \frac{1}{(2\beta _0+\lambda _+^n)(\lambda _+^n\lambda _{}^n)(2\beta _0+\lambda _+^n\lambda _{}^n)}}`$ (A.15) $`\times [\{(\mathrm{\Delta }\gamma _{\psi \psi }^{0,n}\lambda _{}^n)\mathrm{\Delta }\gamma _{G\psi }^{(1),n}+\mathrm{\Delta }\gamma _{G\psi }^{0,n}\mathrm{\Delta }\gamma _{GG}^{(1),n}\}(2\beta _0+\mathrm{\Delta }\gamma _{GG}^{0,n}\lambda _{}^n)`$ $`+\left\{(\mathrm{\Delta }\gamma _{\psi \psi }^{0,n}\lambda _{}^n)\mathrm{\Delta }\gamma _{\psi \psi }^{(1),n}+\mathrm{\Delta }\gamma _{G\psi }^{0,n}\mathrm{\Delta }\gamma _{\psi G}^{(1),n}\right\}\mathrm{\Delta }\gamma _{G\psi }^{0,n}`$ $`{\displaystyle \frac{\beta _1}{\beta _0}}(2\beta _0+\lambda _+^n\lambda _{}^n)\lambda _+^n\mathrm{\Delta }\gamma _{G\psi }^{0,n}]`$ $`\widehat{B}_G^n`$ $`=`$ $`\mathrm{\Delta }K_\psi ^{0,n}{\displaystyle \frac{1}{(2\beta _0+\lambda _{}^n)(\lambda _{}^n\lambda _+^n)(2\beta _0+\lambda _{}^n\lambda _+^n)}}`$ (A.16) $`\times [\{(\mathrm{\Delta }\gamma _{\psi \psi }^{0,n}\lambda _+^n)\mathrm{\Delta }\gamma _{G\psi }^{(1),n}+\mathrm{\Delta }\gamma _{G\psi }^{0,n}\mathrm{\Delta }\gamma _{GG}^{(1),n}\}(2\beta _0+\mathrm{\Delta }\gamma _{GG}^{0,n}\lambda _+^n)`$ $`+\left\{(\mathrm{\Delta }\gamma _{\psi \psi }^{0,n}\lambda _+^n)\mathrm{\Delta }\gamma _{\psi \psi }^{(1),n}+\mathrm{\Delta }\gamma _{G\psi }^{0,n}\mathrm{\Delta }\gamma _{\psi G}^{(1),n}\right\}\mathrm{\Delta }\gamma _{G\psi }^{0,n}`$ $`{\displaystyle \frac{\beta _1}{\beta _0}}(2\beta _0+\lambda _{}^n\lambda _+^n)\lambda _{}^n\mathrm{\Delta }\gamma _{G\psi }^{0,n}]`$ ### A.3 Non-singlet quark $`\mathrm{\Delta }q_{NS}^\gamma (n,Q^2,P^2)/{\displaystyle \frac{\alpha }{8\pi \beta _0}}`$ $`=`$ $`{\displaystyle \frac{4\pi }{\alpha _s(Q^2)}}\widehat{L}_{NS}^n\left\{1\left[{\displaystyle \frac{\alpha _s(Q^2)}{\alpha _s(P^2)}}\right]^{\lambda _{NS}^n/2\beta _0+1}\right\}`$ (A.17) $`+\widehat{A}_{NS}^n\left\{1\left[{\displaystyle \frac{\alpha _s(Q^2)}{\alpha _s(P^2)}}\right]^{\lambda _{NS}^n/2\beta _0}\right\}`$ $`+\widehat{B}_{NS}^n\left\{1\left[{\displaystyle \frac{\alpha _s(Q^2)}{\alpha _s(P^2)}}\right]^{\lambda _{NS}^n/2\beta _0+1}\right\}`$ $`+\widehat{C}_{NS}^n`$ where $`\widehat{L}_{NS}^n`$ $`=`$ $`\mathrm{\Delta }K_{NS}^{0,n}{\displaystyle \frac{1}{1+\lambda _{NS}^n/2\beta _0}}`$ (A.18) $`\widehat{A}_{NS}^n`$ $`=`$ $`{\displaystyle \frac{1}{\lambda _{NS}^n}}\{\mathrm{\Delta }K_{NS}^{0,n}\mathrm{\Delta }\gamma _{NS}^{(1),n}+2\beta _0\mathrm{\Delta }K_{NS}^{(1),n}2\beta _0\lambda _{NS}^n\mathrm{\Delta }A_{NS}^n`$ (A.19) $`{\displaystyle \frac{\beta _1}{\beta _0}}\mathrm{\Delta }K_{NS}^{0,n}(2\beta _0\lambda _{NS}^n)\}`$ $`\widehat{B}_{NS}^n`$ $`=`$ $`\mathrm{\Delta }K_{NS}^{0,n}{\displaystyle \frac{1}{2\beta _0+\lambda _{NS}^n}}\left(\mathrm{\Delta }\gamma _{NS}^{(1),n}{\displaystyle \frac{\beta _1}{\beta _0}}\lambda _{NS}^n\right)`$ (A.20) $`\widehat{C}_{NS}^n`$ $`=`$ $`2\beta _0\mathrm{\Delta }A_{NS}^n`$ (A.21) $`\mathrm{with}`$ $`\lambda _{NS}^n`$ $`=`$ $`\mathrm{\Delta }\gamma _{NS}^{0,n}`$ (A.22) ## Appendix B The first moments ### B.1 One-loop order $`\mathrm{\Delta }\gamma _{NS}^{0,n=1}`$ $`=`$ $`\mathrm{\Delta }\gamma _{\psi \psi }^{0,n=1}=0`$ (B.1) $`\mathrm{\Delta }\gamma _{\psi G}^{0,n=1}`$ $`=`$ $`0,\mathrm{\Delta }\gamma _{G\psi }^{0,n=1}=6C_F`$ (B.2) $`\mathrm{\Delta }\gamma _{GG}^{0,n=1}`$ $`=`$ $`{\displaystyle \frac{22}{3}}C_A+{\displaystyle \frac{8}{3}}T_f=2\beta _0`$ (B.3) $`\lambda _+^{n=1}`$ $`=`$ $`0,\lambda _{}^{n=1}=2\beta _0`$ (B.4) $`\mathrm{\Delta }K_{NS}^{0,n=1}`$ $`=`$ $`\mathrm{\Delta }K_\psi ^{0,n=1}=0`$ (B.5) where $$C_A=3,C_F=\frac{4}{3},T_f=\frac{N_f}{2}$$ (B.6) with $`N_f`$ being the number of flavors. ### B.2 $`\overline{\mathrm{MS}}`$ scheme $`\mathrm{\Delta }\gamma _{NS,\overline{\mathrm{MS}}}^{(1),n=1}`$ $`=`$ $`0`$ (B.7) $`\mathrm{\Delta }\gamma _{\psi \psi ,\overline{\mathrm{MS}}}^{(1),n=1}`$ $`=`$ $`24C_FT_f`$ (B.8) $`\mathrm{\Delta }\gamma _{\psi G,\overline{\mathrm{MS}}}^{(1),n=1}`$ $`=`$ $`0`$ (B.9) $`\mathrm{\Delta }\gamma _{G\psi ,\overline{\mathrm{MS}}}^{(1),n=1}`$ $`=`$ $`18C_F^2{\displaystyle \frac{142}{3}}C_AC_F+{\displaystyle \frac{8}{3}}C_FT_f`$ (B.10) $`\mathrm{\Delta }\gamma _{GG,\overline{\mathrm{MS}}}^{(1),n=1}`$ $`=`$ $`8C_FT_f+{\displaystyle \frac{40}{3}}C_AT_f{\displaystyle \frac{68}{3}}C_A^2=2\beta _1`$ (B.11) $`\mathrm{\Delta }K_{\psi ,\overline{\mathrm{MS}}}^{(1),n=1}`$ $`=`$ $`\mathrm{\Delta }K_{G,\overline{\mathrm{MS}}}^{(1),n=1}=\mathrm{\Delta }K_{NS,\overline{\mathrm{MS}}}^{(1),n=1}=0`$ (B.12) $`\mathrm{\Delta }A_{\psi ,\overline{\mathrm{MS}}}^{n=1}`$ $`=`$ $`12e^2N_f`$ (B.13) $`\mathrm{\Delta }A_{G,\overline{\mathrm{MS}}}^{n=1}`$ $`=`$ $`0`$ (B.14) $`\mathrm{\Delta }A_{NS,\overline{\mathrm{MS}}}^{n=1}`$ $`=`$ $`12(e^4e^2^2)N_f`$ (B.15) $`\mathrm{\Delta }B_{\psi ,\overline{\mathrm{MS}}}^{n=1}`$ $`=`$ $`\mathrm{\Delta }B_{NS,\overline{\mathrm{MS}}}^{n=1}=3C_F`$ (B.16) $`\mathrm{\Delta }B_{G,\overline{\mathrm{MS}}}^{n=1}`$ $`=`$ $`{\displaystyle \frac{N_f}{2}}\mathrm{\Delta }B_{\gamma ,\overline{\mathrm{MS}}}^{n=1}=0`$ (B.17) ### B.3 CI-like schemes (CI, AB, OS, AR) $`\mathrm{\Delta }\gamma _{NS,a}^{(1),n=1}`$ $`=`$ $`0`$ (B.19) $`\mathrm{\Delta }\gamma _{\psi \psi ,a}^{(1),n=1}`$ $`=`$ $`0`$ (B.20) $`\mathrm{\Delta }\gamma _{\psi G,a}^{(1),n=1}`$ $`=`$ $`\mathrm{\Delta }\gamma _{\psi G,\overline{\mathrm{MS}}}^{(1),n=1}=0`$ (B.21) $`\mathrm{\Delta }\gamma _{G\psi ,a}^{(1),n=1}`$ $`=`$ $`\mathrm{\Delta }\gamma _{G\psi ,\overline{\mathrm{MS}}}^{(1),n=1}=18C_F^2{\displaystyle \frac{142}{3}}C_AC_F+{\displaystyle \frac{8}{3}}C_FT_f`$ (B.22) $`\mathrm{\Delta }\gamma _{GG,a}^{(1),n=1}`$ $`=`$ $`32C_FT_f+{\displaystyle \frac{40}{3}}C_AT_f{\displaystyle \frac{68}{3}}C_A^2=2\beta _1+12N_fC_F`$ (B.23) $`\mathrm{\Delta }K_{\psi ,a}^{(1),n=1}`$ $`=`$ $`\mathrm{\Delta }K_{NS,a}^{(1),n=1}=0`$ (B.24) $`\mathrm{\Delta }K_{G,a}^{(1),n=1}`$ $`=`$ $`72e^2N_fC_F`$ (B.25) $`\mathrm{\Delta }A_{\psi ,a}^{n=1}`$ $`=`$ $`\mathrm{\Delta }A_{G,a}^{n=1}=\mathrm{\Delta }A_{NS,a}^{n=1}=0`$ (B.26) $`\mathrm{\Delta }B_{\psi ,a}^{n=1}`$ $`=`$ $`\mathrm{\Delta }B_{NS,a}^{n=1}=3C_F`$ (B.27) $`\mathrm{\Delta }B_{G,a}^{n=1}`$ $`=`$ $`{\displaystyle \frac{N_f}{2}}\mathrm{\Delta }B_{\gamma ,a}^{n=1}=2N_f`$ (B.28) ## Appendix C Derivation of Eqs.(4.20) and (4.21) We observe that, in the $`\overline{\mathrm{MS}}`$ scheme, we have $`\mathrm{\Delta }𝑲^{0,n=1}=\mathrm{\Delta }𝑲^{(1),n=1}=0`$, where $`\mathrm{\Delta }𝑲^n=(\mathrm{\Delta }K_\psi ^n,\mathrm{\Delta }K_G^n,\mathrm{\Delta }K_{NS}^n)`$. (Note $`\mathrm{\Delta }K_{G,\mathrm{CI}\mathrm{like}}^{(1),n=1}0`$. See Eq.(B.25). ) Then, up to NLO, the parton distributions $`\mathrm{\Delta }𝒒^\gamma (n=1)\|_{\overline{\mathrm{MS}}}=(\mathrm{\Delta }q_S^\gamma ,\mathrm{\Delta }G^\gamma ,\mathrm{\Delta }q_{NS}^\gamma )\|_{\overline{\mathrm{MS}}}`$ satisfy a homogeneous differential equation instead of an inhomogenious one: $$\frac{d\mathrm{\Delta }𝒒^\gamma (n=1,Q^2,P^2)\|_{\overline{\mathrm{MS}}}}{d\mathrm{ln}Q^2}=\mathrm{\Delta }𝒒^\gamma (n=1,Q^2,P^2)\|_{\overline{\mathrm{MS}}}\mathrm{\Delta }P(n=1,Q^2)\|_{\overline{\mathrm{MS}}}$$ (C.1) where the $`3\times 3`$ splitting function matrix $`\mathrm{\Delta }P`$ is the hadronic part of $`\mathrm{\Delta }\stackrel{~}{P}`$ given in Eq.(3.1). Expanding $`\mathrm{\Delta }P(n=1,Q^2)\|_{\overline{\mathrm{MS}}}`$ as $$\mathrm{\Delta }P(n=1,Q^2)\|_{\overline{\mathrm{MS}}}=\frac{\alpha _s(Q^2)}{2\pi }\mathrm{\Delta }P_{n=1}^{(0)}+\left[\frac{\alpha _s(Q^2)}{2\pi }\right]^2\mathrm{\Delta }P_{n=1}^{(1)}\|_{\overline{\mathrm{MS}}}+\mathrm{},$$ (C.2) and introducing $`t`$ instead of $`Q^2`$ as the evolution variable $$t\frac{2}{\beta _0}\mathrm{ln}\frac{\alpha _s(P^2)}{\alpha _s(Q^2)},$$ (C.3) we find that Eq.(C.1) is rewritten as $$\frac{d\mathrm{\Delta }𝒒_{n=1}^\gamma (t)\|_{\overline{\mathrm{MS}}}}{dt}=\mathrm{\Delta }𝒒_{n=1}^\gamma (t)\|_{\overline{\mathrm{MS}}}\left\{\mathrm{\Delta }P_{n=1}^{(0)}+\frac{\alpha _s(t)}{2\pi }\left[\mathrm{\Delta }P_{n=1}^{(1)}\|_{\overline{\mathrm{MS}}}\frac{\beta _1}{2\beta _0}\mathrm{\Delta }P_{n=1}^{(0)}\right]+𝒪(\alpha _s^2)\right\}.$$ (C.4) We look for the solution in the following form: $$\mathrm{\Delta }𝒒_{n=1}^\gamma (t)\|_{\overline{\mathrm{MS}}}=\mathrm{\Delta }𝒒_{n=1}^{\gamma (0)}(t)+\mathrm{\Delta }𝒒_{n=1}^{\gamma (1)}(t)\|_{\overline{\mathrm{MS}}}$$ (C.5) with the initial condition (see Eq.(2.5)), $`\mathrm{\Delta }𝒒_{n=1}^{\gamma (0)}(0)`$ $`=`$ $`0`$ (C.6) $`\mathrm{\Delta }𝒒_{n=1}^{\gamma (1)}(0)\|_{\overline{\mathrm{MS}}}`$ $`=`$ $`{\displaystyle \frac{\alpha }{4\pi }}\mathrm{\Delta }𝑨_{n=1}\|_{\overline{\mathrm{MS}}}`$ (C.7) $`=`$ $`{\displaystyle \frac{3\alpha }{\pi }}N_f(e^2,0,e^4e^2^2)`$ In the LO, we easily find that $`\mathrm{\Delta }𝒒_{n=1}^{\gamma (0)}(t)=0`$ due to the initial condition (C.6). The evolution equation in the NLO is written as $$\frac{d\mathrm{\Delta }𝒒_{n=1}^{\gamma (1)}(t)\|_{\overline{\mathrm{MS}}}}{dt}=\mathrm{\Delta }𝒒_{n=1}^{\gamma (1)}(t)\|_{\overline{\mathrm{MS}}}\left\{\mathrm{\Delta }P_{n=1}^{(0)}+\frac{\alpha _s(t)}{2\pi }\left[\mathrm{\Delta }P_{n=1}^{(1)}\|_{\overline{\mathrm{MS}}}\frac{\beta _1}{2\beta _0}\mathrm{\Delta }P_{n=1}^{(0)}\right]\right\},$$ (C.8) and we obtain for the solution $$\mathrm{\Delta }𝒒_{n=1}^{\gamma (1)}(t)\|_{\overline{\mathrm{MS}}}=\mathrm{\Delta }𝒒_{n=1}^{\gamma (1)}(0)\|_{\overline{\mathrm{MS}}}\mathrm{exp}\left(M\right),$$ (C.9) where $$M=\mathrm{\Delta }P_{n=1}^{(0)}t+\frac{1}{\beta _0}\left[\frac{\alpha _s(0)}{\pi }\frac{\alpha _s(t)}{\pi }\right]\left[\mathrm{\Delta }P_{n=1}^{(1)}\|_{\overline{\mathrm{MS}}}\frac{\beta _1}{2\beta _0}\mathrm{\Delta }P_{n=1}^{(0)}\right]$$ (C.10) Since $$\mathrm{\Delta }P_{n=1}^{(0)}=\frac{1}{4}\mathrm{\Delta }\widehat{\gamma }_{n=1}^0,\mathrm{\Delta }P_{n=1}^{(1)}\|_{\overline{\mathrm{MS}}}=\frac{1}{8}\mathrm{\Delta }\widehat{\gamma }_{n=1}^{(1)}\|_{\overline{\mathrm{MS}}},$$ (C.11) and using the information on the first moments of anomalous dimensions which are listed in Appendices B.1 and B.2, we find that $`M`$ turns out to be a triangular matrix in the following form: $$M=\left(\begin{array}{ccc}a& b& 0\\ 0& c& 0\\ 0& 0& d\end{array}\right)$$ (C.12) with $`a`$ $`=`$ $`{\displaystyle \frac{1}{\beta _0}}\left[{\displaystyle \frac{\alpha _s(0)}{\pi }}{\displaystyle \frac{\alpha _s(t)}{\pi }}\right]\left({\displaystyle \frac{1}{8}}\gamma _{\psi \psi ,\overline{\mathrm{MS}}}^{(1),n=1}\right)`$ (C.13) $`b`$ $`=`$ $`{\displaystyle \frac{3}{2}}C_Ft+{\displaystyle \frac{1}{\beta _0}}\left[{\displaystyle \frac{\alpha _s(0)}{\pi }}{\displaystyle \frac{\alpha _s(t)}{\pi }}\right]\left({\displaystyle \frac{1}{8}}\gamma _{G\psi ,\overline{\mathrm{MS}}}^{(1),n=1}{\displaystyle \frac{3\beta _1}{4\beta _0}}C_F\right)`$ (C.14) $`c`$ $`=`$ $`{\displaystyle \frac{1}{2}}\beta _0t`$ (C.15) $`d`$ $`=`$ $`{\displaystyle \frac{1}{\beta _0}}\left[{\displaystyle \frac{\alpha _s(0)}{\pi }}{\displaystyle \frac{\alpha _s(t)}{\pi }}\right]\left({\displaystyle \frac{1}{8}}\gamma _{NS,\overline{\mathrm{MS}}}^{(1),n=1}\right)`$ (C.16) The matrix $`\mathrm{exp}(M)`$ is, therefore, written in the form $$\mathrm{exp}(M)=\left(\begin{array}{ccc}e^a& B& 0\\ 0& e^c& 0\\ 0& 0& e^d\end{array}\right)$$ (C.17) and thus we obtain from Eqs.(C.7) and (C.9), $`\mathrm{\Delta }q_S^\gamma (n=1,Q^2,P^2)\|_{\overline{\mathrm{MS}}}`$ $`=`$ $`{\displaystyle \frac{3\alpha }{\pi }}N_f<e^2>\mathrm{exp}\left\{{\displaystyle \frac{1}{8\beta _0}}\left[{\displaystyle \frac{\alpha _s(0)}{\pi }}{\displaystyle \frac{\alpha _s(t)}{\pi }}\right]\mathrm{\Delta }\gamma _{\psi \psi ,\overline{\mathrm{MS}}}^{(1),n=1}\right\}`$ $``$ $`{\displaystyle \frac{3\alpha }{\pi }}N_f<e^2>\left\{1{\displaystyle \frac{2}{\beta _0}}\left[{\displaystyle \frac{\alpha _s(P^2)}{\pi }}{\displaystyle \frac{\alpha _s(Q^2)}{\pi }}\right]N_f\right\}`$ $`\mathrm{\Delta }q_{NS}^\gamma (n=1,Q^2,P^2)\|_{\overline{\mathrm{MS}}}`$ $`=`$ $`{\displaystyle \frac{3\alpha }{\pi }}N_f(<e^4><e>^2)`$ (C.19) $`\times \mathrm{exp}\left\{{\displaystyle \frac{1}{8\beta _0}}\left[{\displaystyle \frac{\alpha _s(0)}{\pi }}{\displaystyle \frac{\alpha _s(t)}{\pi }}\right]\mathrm{\Delta }\gamma _{NS,\overline{\mathrm{MS}}}^{(1),n=1}\right\}`$ $`=`$ $`{\displaystyle \frac{3\alpha }{\pi }}N_f(<e^4><e>^2)`$ where in the last line we use the fact $`\mathrm{\Delta }\gamma _{NS,\overline{\mathrm{MS}}}^{(1),n=1}=0`$. Incidentally, under the following approximtion, $$b\frac{3}{2}C_Ft,a+cc,$$ (C.20) $`B`$ is evaluated as $`B`$ $``$ $`b\left\{1+{\displaystyle \frac{1}{2}}c+{\displaystyle \frac{1}{3!}}c^2+{\displaystyle \frac{1}{4!}}c^3+\mathrm{}\right\}={\displaystyle \frac{b}{c}}\left[e^c1\right]`$ (C.21) $``$ $`{\displaystyle \frac{3C_F}{\beta _0}}\left[{\displaystyle \frac{\alpha _s(P^2)}{\alpha _s(Q^2)}}1\right].`$ This leads to the expression for the first moment of gluon distribution $`\mathrm{\Delta }G^\gamma (n=1,Q^2,P^2)\|_{\overline{\mathrm{MS}}}`$ given in Eq.(4.6). Figure Captions 1. Deep inelastic scattering on a polarized virtual photon in polarized $`e^+e^{}`$ collision, $`e^+e^{}e^+e^{}+`$ hadrons (quarks and gluons). The arrows indicate the polarizations of the $`e^+`$, $`e^{}`$ and virtual photons. The mass squared of the ‘probe’ (‘target’) photon is $`Q^2(P^2)`$ ($`\mathrm{\Lambda }^2P^2Q^2`$). 2. Factorization scheme dependence of the polarized singlet quark distribution $`\mathrm{\Delta }q_S^\gamma (x,Q^2,P^2)`$ up to NLO in units of $`(3N_fe^2\alpha /\pi )\mathrm{ln}(Q^2/P^2)`$ with $`N_f=3`$, $`Q^2=30\mathrm{GeV}^2`$, $`P^2=1\mathrm{GeV}^2`$, and the QCD scale parameter $`\mathrm{\Lambda }=0.2\mathrm{GeV}`$, for $`\overline{\mathrm{MS}}`$ (dash-dotted line), CI (solid line), AB (short-dashed line), OS (long-dashed line) , AR (dashed line) and $`\mathrm{DIS}_\gamma `$ (dash-2dotted line) schemes. 3. The polarized singlet quark distribution $`\mathrm{\Delta }q_S^\gamma (x,Q^2,P^2)`$ up to NLO predicted by the OS and $`\mathrm{DIS}_\gamma `$ schemes in units of $`(3N_fe^2\alpha /\pi )\mathrm{ln}(Q^2/P^2)`$ for $`N_f=3`$, $`Q^2=30\mathrm{GeV}^2`$, $`P^2=1\mathrm{GeV}^2`$, and $`\mathrm{\Lambda }=0.2\mathrm{GeV}`$, together with the LO result. 4. The polarized singlet quark distribution $`\mathrm{\Delta }q_S^\gamma (x,Q^2,P^2)`$ up to NLO in the OS scheme in units of $`(3N_fe^2\alpha /\pi )\mathrm{ln}(Q^2/P^2)`$ with three different $`Q^2`$ values, for $`N_f=3`$, $`P^2=1\mathrm{GeV}^2`$, and $`\mathrm{\Lambda }=0.2\mathrm{GeV}`$. 5. The polarized gluon distribution $`\mathrm{\Delta }G^\gamma (x,Q^2,P^2)`$ beyond the LO in units of $`(3N_fe^2\alpha /\pi )\mathrm{ln}(Q^2/P^2)`$ with three different $`Q^2`$ values, for $`N_f=3`$, $`P^2=1\mathrm{GeV}^2`$, and $`\mathrm{\Lambda }=0.2\mathrm{GeV}`$. 6. Polarized virtual photon structure function $`g_1^\gamma (x,Q^2,P^2)`$ up to NLO in units of $`(3N_f\alpha e^4/\pi )\mathrm{ln}(Q^2/P^2)`$ for $`Q^2=30`$ GeV<sup>2</sup>, and $`P^2=1`$ GeV<sup>2</sup> and the QCD scale parameter $`\mathrm{\Lambda }=0.2`$ GeV with $`N_f=3`$ (solid line). We also plot the LO result (long-dashed line), the Box (tree) diagram (2dash-dotted line) and the Box including non-leading contribution, Box (NL) (short-dashed line). 7. Point-like piece of the real photon structure function $`g_1^\gamma (x,Q^2)`$ in NLO in units of $`(3N_f\alpha e^4/\pi )\mathrm{ln}(Q^2/\mathrm{\Lambda }^2)`$ for $`Q^2=30`$ GeV<sup>2</sup> with $`\mathrm{\Lambda }=0.2`$ GeV, $`N_f=3`$ (solid line). Also plotted are the LO result (long-dashed line) and the Box (tree) diagram contribution (short-dashed line).
warning/0007/astro-ph0007120.html	ar5iv	text	# SHALLOW CORES IN THE DARK MATTER HALOS: SELF-INTERACTION IN ACTION? ## 0.1 Introduction During several decades, dynamical studies of galaxies, and group and cluster of galaxies have pointed out to the existence of massive dark matter halos. On the other hand, according to current models of structure formation in the universe, luminous galaxies should form from the gas trapped within the deep gravitational wells of dark matter (DM) structures emerged from primordial density fluctuations. In these models, non-dissipative, cold, collisionless particles (cold dark matter, CDM) were required. The CDM structure formation scenario succesfully accounted for a wide range of observations, in particular on large scales. However, on small scales, compared with observations this scenario seems to predict too centrally concentrated halos and too much substructure in Milky Way-size halos. These discrepancies have induced to introduce some modifications to the CDM scenario, in particular, to the nature of the DM. From the point of view of particle physics, a large list of candidate DM particles has been proposed but unfortunately, none of the particles that might constitute the universe’s missing mass have been detected nowadays. Nevertheless, it is possible that astronomical observations may help us to constrict some of the properties of these particles. For example, as was mentioned above, the existence of soft cores in the dark halos appears to be not compatible with collisionless CDM particles. Therefore, astronomical studies about the halo properties — in particular of their cores — are crucial for understanding the nature of the dark particles and the structure formation in the universe, as was emphasized in a pioneering paper on this subject by J. Kormendy (1990; see also Kormendy 1988). Here, we summarize the halo core scaling relationships we have inferred from observations from dwarf galaxies through cluster of galaxies (Firmani et al. 2000a,b), and we dicuss some of the implications of our results on the nature of the DM particles and the formation of halos. ## 0.2 Halo core scaling laws from observations Analysis of the original virialized halo mass distribution for most of galaxies is uncertain due to the ambiguities in the estimate of the stellar mass-to-light ratios $`M/L`$ and the gravitational pull the collapsing gas exerted over the inner parts of the halo. This is why we limited our study to only galaxies (i) strongly dominated by DM and (ii) with accurately measured rotation curves. The sample taken from the literature consists of six dwarf galaxies, nine LSB galaxies, and two late-type low luminosity galaxies (Firmani et al. 2000b). In all these cases, the galaxies are DM dominated. Even so, we have substracted from the observed rotation curve the small disk contribution. The halo components were fitted to a non-singular isothemal model; thus, for each galaxy characterized by its maximum circular velocity V<sub>max</sub> we estimate its central density $`\rho _c`$ and core radius $`r_c`$. For the less DM dominated LSB and low-luminosity galaxies, we have roughly calculated the factor by which the halo component was “deformed” due to the disk pull over the DM using the adiabatic invariance approximation (see details in Firmani et al. 2000b). On galaxy cluster scale, we have used the surface mass distribution for the cluster CL0024+1654 derived from unprecedent high-resolution strong lensing mass maps (Tyson, Kochanski, & Dell’Antonio 1998), and for the cluster CL0016+16 derived from weak lensing studies (Smail et al. 1995). In both cases there is no evidence of a massive cD galaxy and the inner mass distribution is soft. In Figs. 1 we show the dependence of $`r_c`$ on V<sub>max</sub> we have found from the observational data. Although with a big scatter, within a large range in V<sub>max</sub> we estimate that: $$\rho _c(r)0.02\mathrm{M}_{}\mathrm{pc}^3\mathrm{and}r_c5.5\left[\frac{V_{\mathrm{max}}}{100\mathrm{k}\mathrm{m}\mathrm{s}^1}\right]^{0.95}\mathrm{kpc}.$$ (1) Similar results were found for an uniform sample of high and LSB galaxies of the Coma Ursa Mayor cluster (Verheijen 1997; §6). In this case, the rotation curve decompositions were made assuming $`M/L_K`$ constant for all galaxies, and the halo component was fitted to a pseudo-isothermal model. In contrast, from a sample of Sc-Im and dwarf galaxies, Kormendy (1988,1990) inferred that $`\rho _c`$ decreases with the galaxy luminosity (or V<sub>max</sub>). Certainly, more efforts should be done in the future in order to increase the sample of objects and to reduce the uncertainties in the rotation curve decomposition techniques. We remark the importance of strong gravitational lensing studies in order to directly probe the inner regions of the cluster of galaxies. ## 0.3 Implications of the inferred halo core scaling laws The existence of soft halo cores and even more, the scaling laws obtained for DM dominated systems \[eq.(1)\], are in complete disagreement with the predictions of CDM models. Warm dark matter (WDM) has been proposed in order to solve the other conflict of the CDM scenario —the overlying number of guest (satellite) halos in a Milky Way-size halo. Cosmological N-body simulations have shown that the latter problem is indeed solved for a filtering scale in the power spectrum of $`0.1`$ Mpc which corresponds to a warm particle of $`1`$ KeV (Colín, Avila-Reese, & Valenzuela 2000). These authors have also shown that the density profiles of halos with masses much larger than that corresponding to the filtering scale ($`10^9`$ h<sup>-1</sup>M) are very similar to those of the CDM models (see also Moore et al. 1999). Thus, even if the halos with masses near or smaller than the filtering mass would have a core, the more massive WDM halos will not obey the scaling laws inferred from observations in $`\mathrm{\S }2`$ (see also Avila-Reese, Firmani, & Hernández 1998). Spergel & Steinhardt (2000) suggested other modification to the nature of the DM particles: the introduction of self-interaction. In Firmani et al. 2000a, the gravothermal expansion was proposed as the mechanism able to produce soft cores in self-interacting CDM halos. The inner velocity dispersion profile of these halos raises with radius. Therefore, if particles are collisional, heat transfers inwards, the core expands and cools, exacerbating even more the temperature gradient. This process is similar to the postcollapse gravothermal oscillations in globular clusters (Bettwieser & Sugimoto 1984; Goodman 1987). For globular clusters, the core expansion halts when the inner dispersion velocity profile flattens; this occurs because there is also an outwards heat flux from the maximum of the velocity profile. In the case of DM halos, we propose a collisional cross-section $`\sigma `$ such that self-interaction is efficient only in the more dense halo regions. Besides, as the soft core grows, the core density decreases and at some moment, self-interaction should become inefficient even in the inner regions. On the other hand, it is important to beer in mind that the CDM halo does not form instantaneously, but by a hierarchical mass aggregation process which establishes a cuspy inner structure with a positive velocity dispersion gradient. Recent numerical simulations for a halo with Hernquist density profile and with relatively small cross-sections per unit of the particle mass $`m_X`$ ($`\sigma _{}=\sigma /m_X`$) have shown that the gravothermal processes act in time scales that depend on the value of $`\sigma _{}`$ (Burkert 2000; Kochanek & White 2000; see also Quinlan 1996). An important constriction is that the halo lifetime should be in between the core expansion time and the core collapse time; otherwise either the shallow core still have not been formed or the core is already in its collapse phase. In Firmani et al. (2000a,b), using the average observed $`\rho _c`$ and supossing that the collision time $`t_{\mathrm{col}}`$ in the core is close to the Hubble time, a lower limit for $`\sigma _{}`$ was estimated<sup>1</sup><sup>1</sup>1Here we assume that V$`{}_{\mathrm{max}}{}^{}v_{\mathrm{rms},\mathrm{max}}`$; in fact, for CDM halos V<sub>max</sub> is roughly 1.3-1.7 times larger than $`v_{\mathrm{rms},\mathrm{max}}`$: $`\sigma _{}\mathrm{4\; 10}^{25}V_{\mathrm{max},100}^1`$ cm<sup>2</sup>/GeV, where V<sub>max,100</sub> is V<sub>max</sub> in units of 100 km/s. An important point to be noted is that $`\sigma _{}`$ depends on V<sub>max</sub> or the maximum velocity dispersion, i.e. the cross-section is a function of the particle energy as in other classical physical interactions. For velocity dispersions corresponding to galaxy clusters, this value is close to the limit estimated by Miralda-Escudé (2000) from the observationally inferred ellipticity of the cluster MS2137-23. One may think that the evolution of the collisionless DM halo occurs in scales of dynamical times, $`t_{\mathrm{dyn}}`$, while those central regions of the halo affected by the gravothermal processes, evolve in relaxation time scales, $`t_{\mathrm{rel}}`$. The final halo density profile is the result of both dynamical processes. The simulations carried out by Burkert (2000) and Kochanek & White (2000; KW00) are for a halo already virialized. Therefore, these simulations do not describe the cosmological process of halo collapse and virialization. Kochanek & White find that the gravothermal core collapse occurs in scale times less than $`5`$ times the core formation time $`t_c`$ independent of the value of $`\sigma _{}`$. On the other hand, $`t_c1/\sigma _{}`$. Thus, for $`\sigma _{}`$ small enough the halos may still be in their core expansion phase. Besides, if $`\sigma _{}`$ depends on V<sub>max</sub> as we have inferred from observations, then larger halos should be today in earlier stages of gravothermal expansion than smaller halos, i.e. their central densities have not decreased too much. This, combined with the fact that in the hierarchical scenario smaller halos are intrinsically more concentrated than larger ones, could produce the invariance of $`\rho _c`$ with the halo scale. The simulation of KW00 are for a Hernquist halo, and they express $`\sigma _{}`$ in unities of $`r_H^2/M_h`$, where $`M_h`$ is the halo mass and $`r_H`$ is the scale radius of the Hernquist profile. Fitting the Hernquist profile to halos obtained in a N-body CDM simulation, one finds that $`r_H^2/M_h`$ is roughly constant. In order to obtain more quantitative estimates, we have used results for a $`\mathrm{\Lambda }`$CDM<sub>0.3</sub>, h=0.7 model (Avila-Reese et al. 1999). We calculate $`r_H`$ as $`r_H=r_v/c_H`$, where the virial radius $`r_v`$ is defined as the radius where the average halo density is $`\mathrm{\Delta }_c`$ times the background density (for our cosmology, $`\mathrm{\Delta }_c=340`$), and $`c_H`$ is the concentration parameter which ultimately depends on the halo mass or V<sub>max</sub> and is the only free parameter in the cosmological halo density profiles. From the results of the simulation, we find on the average c$`{}_{H}{}^{}=37.5/(`$V$`{}_{\mathrm{max}}{}^{}/\mathrm{kms}^1)^{0.36}`$. The virial radius is proportional to V<sub>v</sub>, the circular velocity at this radius, and for the Hernquist profile, V$`{}_{v}{}^{}=2`$V<sub>max</sub>c$`{}_{H}{}^{1/2}/(1+c_H)`$. We find that $`r_H^2/M_h\mathrm{7\; 10}^{24}`$ cm<sup>2</sup>/GeV. Thus, in KW00 $`\sigma _{}`$ would be $`\widehat{\sigma }\times \mathrm{7\; 10}^{24}`$ cm<sup>2</sup>/GeV. For $`\widehat{\sigma }=1`$, KW00 find that $`t_c1.7`$ dynamical times; after this the halo suffers the gravothermal core collapse. For a value of $`\sigma _{}`$ as that we have inferred from observations, for a V$`{}_{\mathrm{max}}{}^{}100`$ km/s halo for example ($`\sigma _{}\mathrm{4\; 10}^{25}`$ cm<sup>2</sup>/GeV), $`\widehat{\sigma }0.05`$. The dynamical time (as deffined in KW00) for a V$`{}_{\mathrm{max}}{}^{}100`$ km/s halo is $`\mathrm{5\; 10}^8`$ years. Therefore, the core formation time would be of the order of a Hubble time. Halos larger than V$`{}_{\mathrm{max}}{}^{}100`$ km/s would have even larger core formation times. ## 0.4 Conclusions $``$ The halo core scaling laws inferred from observations of dwarf galaxies to galaxy clusters show that $`\rho _c`$ does not depend on the halo mass or V<sub>max</sub> and the core radius is roughly proportional to V<sub>max</sub>. $``$ If the dark particles are self-interacting with not very large cross sections, then gravothermal processes may produce a soft core in the DM halos. Using the observational data, we estimated the value of $`\sigma _{}`$ and found that is roughly proportional to V$`{}_{\mathrm{max}}{}^{1}v_{\mathrm{rms},\mathrm{max}}^1`$. $``$ Results from numerical simulations of already virialized halos with self-interaction, show that if $`\sigma _{}`$ is of the order we inferred from observations, then $`t_c`$ for small halos is close to the Hubble time, while for larger halos, $`t_c`$ is probably even larger, i.e. these halos are still in early stages of gravothermal expansion. Numerical simulations and theoretical studies of collapsing and virializing DM halos where self-interacion is efficient only in the more dense inner regions are necessary in order to attain more quantitative conclusions.
warning/0007/hep-ph0007294.html	ar5iv	text	# TTP00-13 NIKHEF-2000-018 Some higher moments of deep inelastic structure functions at next-to-next-to leading order of perturbative QCD ## 1 Introduction The determination of the next-to-next-to-leading (NNL) order QCD approximation for the structure functions of deep inelastic scattering has become important for the understanding of perturbative QCD and necessary for an accurate comparison of perturbative QCD with the increasing precision of experiments. Such calculations however are rather complicated and hence a complete NNL result does not exist as of yet. The one-loop anomalous dimensions were calculated in . In (see also the references therein) the complete one-loop coefficient functions were obtained. Anomalous dimensions at 2-loop order were obtained in . The 2-loop coefficient functions were calculated in , . Analytical results of the 3-loop anomalous dimensions and coefficient functions of the moments of $`F_2`$ and $`F_L`$ are only known for the moments $`N`$=2,4,6,8 in both the singlet and non-singlet case and additionally for $`N`$=10 in the non-singlet case from and . In addition the Gross-Llevellyn Smith sum rule, which corresponds to the first moment of $`F_3^{\nu \mathrm{p}+\overline{\nu }\mathrm{p}}`$ has been calculated at this order . For a complete reconstruction of the $`x`$-dependence of the structure functions via an inverse Mellin-transformation one would need the moments for all $`N`$ (that is either all even or all odd integer values). Additionally one needs them both for $`F_2`$ and $`F_3`$ in order to untangle the various quark and gluon contributions. The determination of the NNL approximation for generic $`N`$ is work in progress , but probably will not be finished in the near future. The available moments of $`F_2`$ have been used by a number of authors to make a reconstruction of the complete structure functions at NNL by a variety of means . Additionally they can be used to obtain a better value of $`\alpha _S`$ It should be clear that it is important to have as large a number of moments as possible. First, these results can be immediately used to increase the precision of phenomenological investigations of deep inelastic scattering . Second, the moments will be a very important check for the new methods and programs needed for the determination of the 3-loop results for arbitrary $`N`$. Unfortunately it is not very easy to increase the number of moments, because each new moment requires roughly five times the computer resources that its predecessor needs. With the advent of better computers this means that by now it has been possible to obtain two more moments for the singlet case and one additional moment for the non singlet $`F_2`$ case. This should allow for instance a somewhat better determination of $`\alpha _S`$. More important however is the determination of the first seven odd moments of $`F_3^{\nu \mathrm{p}+\overline{\nu }\mathrm{p}}`$ to three loops. To this end we used the the same programs as in , state of the art computers and a new version of the symbolic manipulation program FORM that supports now 64-bit architectures and to some extend parallel computers (see also ). We could push the limit in these calculations to include two new moments ($`N`$=10,12) in the calculation of the flavour singlet structure functions $`F_L`$ and $`F_2`$, and two new moments ($`N`$ = 12,14) for the flavour nonsinglet structure functions $`F_L`$ and $`F_2`$. Additionally we have computed the moments $`N`$=3,5,7,9,11,13 of the structure function $`F_3`$. We do not expect more moments to become available before the complete results for all $`N`$ will be presented. ## 2 The formalism This calculation follows the one presented in (see also ) in every detail, so we only will give a very short review on the methods used. We need to calculate the hadronic part of the amplitude for unpolarized deep inelastic scattering which is given by the hadronic tensor $`W_{\mu \nu }(x,Q^2)`$ $`=`$ $`{\displaystyle \frac{1}{4\pi }}{\displaystyle \mathrm{d}^4z\mathrm{e}^{\mathrm{i}qz}p,\mathrm{nucl}\|J_\mu (z)J_\nu (0)\|\mathrm{nucl},p}`$ (1) $`=`$ $`(g_{\mu \nu }q_\mu q_\nu ){\displaystyle \frac{1}{2x}}F_L(x,Q^2)`$ $`+\left(g_{\mu \nu }p_\mu p_\nu {\displaystyle \frac{4x^2}{q^2}}(p_\mu q_\nu +p_\nu q_\mu ){\displaystyle \frac{2x}{q^2}}\right){\displaystyle \frac{1}{2x}}F_2(x,Q^2)`$ $`+\mathrm{i}ϵ_{\mu \nu \rho \sigma }{\displaystyle \frac{p^\rho q^\sigma }{pq}}F_3(x,Q^2)`$ where the $`J^\mu `$ are either electromagnetic or weak hadronic currents and $`x=Q^2/(2pq)`$ is the Bjorken scaling variable with $`0<x1`$. $`Q^2=q^2`$ is the transfered momentum and $`\|p,\mathrm{nucl}`$ is the nucleon state with momentum $`p`$. In these equations spin averaging is assumed. The longitudinal structure function $`F_L`$ is related to the structure function $`F_1`$ by $`F_L=F_22xF_1`$. For electron-nucleon scattering $`J^\mu `$ is the electromagnetic quark current and $`F_3`$ vanishes. For neutrino-nucleon scattering $`J^\mu `$ is an electroweak quark current which has an axial vector contribution and $`F_3`$, which describes parity violating effects that arise from vector and axial-vector interference, will not vanish. Using the dispersion relation technique one can relate the hadronic tensor to the following 4-point Green functions: $$W^{\mu \nu }(p,q)=\frac{1}{2\pi }\mathrm{Im}T^{\mu \nu }(p,q),T_{\mu \nu }(p,q)=\mathrm{i}\mathrm{d}^4z\mathrm{e}^{\mathrm{i}qz}p,\mathrm{nucl}\|T\left[J_\mu (z)J_\nu (0)\right]\|\mathrm{nucl},p.$$ Applying a formal operator product expansion in terms of local operators to the time-ordered product of two quark currents leads to: $$\begin{array}{c}\hfill \mathrm{i}\mathrm{d}^4z\mathrm{e}^{\mathrm{i}qz}T\left[J_{\nu _1}(z)J_{\nu _2}(0)\right]=\underset{N,j}{}\left(\frac{1}{Q^2}\right)^N[(g_{\nu _1\nu _2}\frac{q_{\nu _1}q_{\nu _2}}{q^2})q_{\mu _1}q_{\mu _2}C_{L,N}^j(\frac{Q^2}{\mu ^2},a_s)\\ \hfill \left(g_{\nu _1\mu _1}g_{\mu _2\nu _2}q^2g_{\nu _1\mu _1}q_{\nu _2}q_{\mu _2}g_{\nu _2\mu _2}q_{\nu _1}q_{\mu _1}+g_{\nu _1\nu _2}q_{\mu _1}q_{\mu _2}\right)C_{2,N}^j(\frac{Q^2}{\mu ^2},a_s)\\ \hfill +\mathrm{i}ϵ_{\nu _1\nu _2\mu _1\nu _3}g^{\nu _3\nu _4}q_{\nu _4}q_{\mu _2}C_{3,N}^j(\frac{Q^2}{\mu ^2},a_s)]\times q_{\mu _3}\mathrm{}q_{\mu _N}O^{j,\{\mu _1,\mathrm{}\mu _N\}}(0)+\mathrm{higher}\mathrm{twists},\\ \hfill j=\alpha ,\psi ,G\end{array}$$ (2) Here we have introduced the notation $`a_s=\alpha _s/(4\pi )=g^2/(4\pi )^2`$ and everything is assumed to be renormalized. The sum over $`N`$ runs over the standard set of the spin-$`N`$ twist-2 irreducible flavour non-singlet quark operators and the singlet quark and gluon operators: $`O^{\alpha ,\{\mu _1,\mathrm{},\mu _N\}}`$ $`=\overline{\psi }\lambda ^\alpha \gamma ^{\{\mu _1}D^{\mu _2}\mathrm{}D^{\mu _N\}}\psi ,\alpha =1,2,\mathrm{},(n_f^21)`$ $`O^{\mathrm{\Psi },\{\mu _1,\mathrm{},\mu _N\}}`$ $`=\overline{\psi }\gamma ^{\{\mu _1}D^{\mu _2}\mathrm{}D^{\mu _N\}}\psi `$ $`O^{G,\{\mu _1,\mathrm{},\mu _N\}}`$ $`=G^{\{\mu \mu _1}D^{\mu _2}\mathrm{}D^{\mu _{N1}}G^{\mu _N\mu \}}`$ Application of this OPE to Eq. (1) leads to an expansion for the unphysical values $`x\mathrm{}`$. From the proper analytical continuation to the physical region $`0<x1`$ one finds for the moments of the structure functions $`F_2`$, $`F_L`$ and $`F_3`$: $`M_{k,N2}`$ $`={\displaystyle \underset{0}{\overset{1}{}}}dxx^{N2}F_k(x,Q^2)={\displaystyle \underset{i=\alpha ,\psi ,G}{}}C_{k,N}^i({\displaystyle \frac{Q^2}{\mu ^2}},a_s)A_{\mathrm{nucl},N}^i,k=2,L`$ (3) $`M_{3,N1}`$ $`={\displaystyle \underset{0}{\overset{1}{}}}dxx^{N1}F_3(x,Q^2)={\displaystyle \underset{i=\alpha }{}}C^{i,N}({\displaystyle \frac{Q^2}{\mu ^2}},a_s)A_{\mathrm{nucl},N}^i`$ (4) with the spin averaged matrix elements $$p,\mathrm{nucl}\|O^{j,\{\mu _1\mathrm{}\mu _N\}}\|p,\mathrm{nucl}=p^{\{\mu _1}\mathrm{}p^{\mu _N\}}A_{\mathrm{nucl},N}^j(\frac{p^2}{\mu ^2})$$ (5) In the derivation of (3) one needs the symmetry properties of $`T_{\mu \nu }`$ under $`xx`$. This is why one can only find either even or odd moments from these equations, dependent on the process under consideration. For $`F_3`$ we will only consider the flavor non-singlet contributions, due to the properties of the operators $`O^\psi `$ and $`O^G`$ under charge conjugation there should not be a singlet contribution (see e.g. ). The scale-dependence of the coefficient functions is then covered by the renormalization group equations: $`\left[\mu ^2{\displaystyle \frac{}{\mu ^2}}+\beta (a_s){\displaystyle \frac{}{a_s}}\gamma _N^{\mathrm{ns}}\right]C_{i,N}^{\mathrm{ns}}({\displaystyle \frac{Q^2}{\mu ^2}},a_s)`$ $`=0,i=2,3,L`$ (6) $`{\displaystyle \underset{k=\psi ,G}{}}\left[\left\{\mu ^2{\displaystyle \frac{}{\mu ^2}}+\beta (a_s){\displaystyle \frac{}{a_s}}\right\}\delta ^{jk}\gamma _N^{jk}\right]C_{i,N}^k({\displaystyle \frac{Q^2}{\mu ^2}},a_s)`$ $`=0,i=2,L;j=\psi ,G`$ (7) The non-singlet coefficient functions and anomalous dimensions don’t depend on the index $`\alpha `$, and we have adopted the conventional collective denotation “ns” for them. ## 3 The even moments of $`F_2`$ and $`F_L`$ Equation (2) is a relation between operators and does not depend on the hadronic states to which the OPE is applied. Using the method of projectors one can find both the coefficient functions and the anomalous dimensions for the even moments of $`F_2`$ and $`F_L`$ as defined in Eq. (6) from the following 4-point Green functions: $`T_{\mu \nu }^{q\gamma q\gamma }`$ $`=\mathrm{i}{\displaystyle \mathrm{d}^4z\mathrm{e}^{\mathrm{i}qz}p,\mathrm{quark}\|T\left[J_\mu (z)J_\nu (0)\right]\|\mathrm{quark},p}`$ (8) $`T_{\mu \nu }^{g\gamma g\gamma }`$ $`=\mathrm{i}{\displaystyle \mathrm{d}^4z\mathrm{e}^{\mathrm{i}qz}p,\mathrm{gluon}\|T\left[J_\mu (z)J_\nu (0)\right]\|\mathrm{gluon},p}`$ (9) Applying to Eqs. (8) the projectors $$P_N\left[\frac{q^{\{\mu _1}\mathrm{}q^{\mu _N\}}}{N!}\frac{^N}{p^{\mu _1}\mathrm{}p^{\mu _N}}\right]\|_{p=0}$$ (10) and, to project out the different Lorentz projections (these projectors are valid in $`D=42ϵ`$ and for the leading twist approximation): $`P_L`$ $`={\displaystyle \frac{q^2}{(pq)^2}}p^\mu p^\nu `$ $`P_2`$ $`=\left({\displaystyle \frac{32ϵ}{22ϵ}}{\displaystyle \frac{q^2}{(pq)^2}}p^\mu p^\nu +{\displaystyle \frac{1}{22ϵ}}g^{\mu \nu }\right),`$ as well as the corresponding flavour projections results in the equations: $`T_{k,N}^{q\gamma q\gamma ,\mathrm{s}}({\displaystyle \frac{Q^2}{\mu ^2}},a_s,ϵ)`$ $`=\left(C_{k,N}^\psi ({\displaystyle \frac{Q^2}{\mu ^2}},a_s,ϵ)Z_N^{\psi \psi }(a_s,{\displaystyle \frac{1}{ϵ}})+C_{k,N}^G({\displaystyle \frac{Q^2}{\mu ^2}},a_s,ϵ)Z_N^{G\psi }(a_s,{\displaystyle \frac{1}{ϵ}})\right)`$ $`A_{\mathrm{quark},N}^{\psi ,\mathrm{tree}}(ϵ)`$ $`T_{k,N}^{g\gamma g\gamma ,\mathrm{s}}({\displaystyle \frac{Q^2}{\mu ^2}},a_s,ϵ)`$ $`=\left(C_{k,N}^\psi ({\displaystyle \frac{Q^2}{\mu ^2}},a_s,ϵ)Z_N^{\psi G}(a_s,{\displaystyle \frac{1}{ϵ}})+C_{k,N}^G({\displaystyle \frac{Q^2}{\mu ^2}},a_s,ϵ)Z_N^{GG}(a_s,{\displaystyle \frac{1}{ϵ}})\right)`$ $`A_{\mathrm{gluon},N}^{G,\mathrm{tree}}(ϵ)`$ $`T_{k,N}^{q\gamma q\gamma ,\mathrm{ns}}({\displaystyle \frac{Q^2}{\mu ^2}},a_s,ϵ)`$ $`=C_{k,N}^{\mathrm{ns}}({\displaystyle \frac{Q^2}{\mu ^2}},a_s,ϵ)Z_N^{\mathrm{ns}}(a_s,{\displaystyle \frac{1}{ϵ}})A_{\mathrm{quark},N}^{\mathrm{ns},\mathrm{tree}}(ϵ)`$ $`k=2,L`$ (11) From these equations the coefficient functions and from the $`Z_N^{ij}`$ the anomalous dimensions can be calculated in the usual way. It should be mentioned that in the Eqs. (3) on the left hand side after applying the projectors (10) we are left with only diagrams of the massless propagator type, a problem solved at 3-loop order long ago and implemented in an efficient way in the FORM package MINCER . On the right hand side only the tree level diagrams contributing to the Matrix elements survive. It turns out that much computing time can be saved when calculating additionally Green functions with external ghosts to get rid of the unphysical polarization states of the external gluons instead of using the very complicated projection onto physical states. Also, from Eqs. (3) one can determine the $`Z_N^{GG}`$ and $`Z_N^{G\psi }`$ only to order $`\alpha _s^2`$. To obtain the $`\alpha _s^3`$-contributions one can calculate additionally Greens functions with external scalar fields $`\varphi `$ that couple to gluons only at tree level. Altogether, to obtain the coefficient functions and anomalous dimensions for the even moments with $`N`$=10,12 of $`F_2`$ and $`F_L`$ the following diagrams had to be calculated (q=quark, g=gluon, $`\gamma `$=photon, h=ghost, $`\varphi `$=scalar field): \| \| tree \| 1-loop \| 2-loops \| 3-loops \| Lorentz projections \| \| --- \| --- \| --- \| --- \| --- \| --- \| \| $`q\gamma q\gamma `$ \| 1 \| 3 \| 27 \| 413 \| 2 \| \| $`q\varphi q\varphi `$ \| \| 1 \| 24 \| 697 \| 1 \| \| $`g\gamma g\gamma `$ \| \| 2 \| 20 \| 366 \| 2 \| \| $`h\gamma h\gamma `$ \| \| \| 2 \| 53 \| 2 \| \| $`g\varphi g\varphi `$ \| 1 \| 11 \| 241 \| 1266 \| 1 \| \| $`h\varphi h\varphi `$ \| \| 11 \| 241 \| 1266 \| 1 \| \| Total \| 3 \| 23 \| 399 \| 10846 \| \| The $`q^{\{\mu _1}\mathrm{}q^{\mu _N\}}`$ in Eq. (10) are the harmonic (i.e. symmetrical and traceless) part of the tensor $`q^{\mu _1}\mathrm{}q^{\mu _N}`$. The number of terms in these harmonic tensors explodes as $`N`$ increases and this is the real limitation in these calculations considering the computing time as well as disk-space usage. In spite of a very efficient implementation of these tensors (see ) for $`N=12`$, singlet and $`N=14`$, nonsinglet, individual diagrams had a disk space usage up to and over 100 GB. Altogether the calculation of all the above diagrams for $`N=10,12`$ took approximately 5 weeks on a Compaq Server with 8 Alpha 21264 processors running at 700 MHz, 4 GB of RAM and 12$`\times `$17 GB of disk-space. The $`N=14`$ nonsinglet calculation took comparable resources. ## 4 The odd moments of $`F_3`$ The coefficient functions and anomalous dimensions of the odd moments of the structure function $`F_3`$ can be obtained in the same way as the non-singlet part of $`F_2`$ and $`F_L`$ but now considering the time-ordered product of one vector current $`V_\mu `$ and one axial vector current $`A_\nu `$. The axial current introduces the appearance of a $`\gamma _5`$ and some care has to be taken to treat it correctly within the framework of dimensional regularization. We adopt the definition used in : $$\gamma _\mu \gamma _5=\mathrm{i}\frac{1}{6}ϵ_{\mu \nu \rho \sigma }\gamma ^\mu \gamma ^\nu \gamma ^\rho \gamma ^\sigma $$ Projecting out the flavour non-singlet part and the corresponding Lorentz structure with: $$P_3=\mathrm{i}\frac{1}{(12ϵ)(22ϵ)}ϵ^{\mu \nu \alpha \beta }\frac{p_\alpha q_\beta }{pq}$$ one finds products of metric tensors which have to be considered as $`D`$-dimensional objects. Since this definition of $`\gamma _5`$ in $`D`$ dimensions violates the axial Ward identity one needs to renormalize $`A_\mu `$ with a renormalization constant $`Z_A`$ and additionally apply a finite renormalization with $`Z_5`$, both of these constants are given to 3-loop order in . Combining all this finally leads to $$Z_A(a_s,\frac{1}{ϵ})Z_5(a_s,ϵ)T_{3,N}^{\mathrm{ns}}(\frac{Q^2}{\mu ^2},a_s,ϵ)=C_{3,N}(\frac{Q^2}{\mu ^2},a_s,ϵ)Z_N^1(a_s,\frac{1}{ϵ})A_{N,\mathrm{tree}}^{\mathrm{ns}}(ϵ)$$ Due to the $`\gamma _5`$ insertion at one of the vertices, some of the symmetries that were used to minimize the number of diagrams could not be applied in this case and to determine the $`T_{3,N}^{\mathrm{ns}}`$ 1076 (= 1 + 4 + 55 + 1016) diagrams had to be evaluated, which took about 6 weeks for the moments $`N`$=1,3,5,7,9,11,13. ## 5 Results Using the strategies sketched in the previous section we find the following results for the coefficient functions and anomalous dimensions. Again following Ref. , we present the combined singlet and non-singlet results for $`F_2`$ and $`F_L`$ in terms of flavour factors which are defined in the following table for $`n_f`$ number of flavours: \| \| $`fl_2`$ \| $`fl_{11}`$ \| $`fl_{02}`$ \| $`fl_2^g`$ \| $`fl_{11}^g`$ \| \| --- \| --- \| --- \| --- \| --- \| --- \| \| non-singlet \| 1 \| $`\frac{3}{n_f}_{f=1}^{n_f}e_f`$ \| 0 \| - \| - \| \| singlet \| 1 \| $`\frac{1}{n_f}\frac{\left(_{f=1}^{n_f}e_f\right)^2}{_{f=1}^{n_f}e_f^2}`$ \| 1 \| 1 \| $`\frac{1}{n_f}\frac{\left(_{f=1}^{n_f}e_f\right)^2}{_{f=1}^{n_f}e_f^2}`$ \| The numerical values of the anomalous dimensions for the spin-even operators contributing to $`F_2`$ and $`F_L`$ that now are known to NNL order are: $`\gamma _2^{\psi \psi }`$ $`=3.555555556a_s+a_s^2\left(48.329218113.160493827n_f1.975308642\mathrm{fl}_{02}n_f\right)`$ $`+a_s^3(859.4478372133.4381617n_f1.229080933n_f^2`$ $`+\mathrm{fl}_{02}(42.21182429n_f3.445816187n_f^2))`$ $`\gamma _4^{\psi \psi }`$ $`=6.977777778a_s+a_s^2\left(86.286650216.553580247n_f0.1060246914\mathrm{fl}_{02}n_f\right)`$ $`+a_s^3(1515.562363244.728592n_f2.108515775n_f^2`$ $`+\mathrm{fl}_{02}(5.17013312n_f0.6789278464n_f^2))`$ $`\gamma _6^{\psi \psi }`$ $`=9.003174603a_s+a_s^2\left(108.01846978.62925674n_f0.02080408883\mathrm{fl}_{02}n_f\right)`$ $`+a_s^3(1891.827779307.4236889n_f2.570638992n_f^2`$ $`+\mathrm{fl}_{02}(2.526091192n_f0.2884556035n_f^2))`$ $`\gamma _8^{\psi \psi }`$ $`=10.45820106a_s+a_s^2\left(123.776452510.14583662n_f0.006586485074\mathrm{fl}_{02}n_f\right)`$ $`+a_s^3(2164.091836352.3116596n_f2.882493484n_f^2`$ $`+\mathrm{fl}_{02}(1.682156519n_f0.1620816452n_f^2))`$ $`\gamma _{10}^{\psi \psi }`$ $`=11.5969216a_s+a_s^2\left(136.274177511.34594534n_f0.00269944007\mathrm{fl}_{02}n_f\right)`$ $`+a_s^3(2379.919952387.6422968n_f3.115523145n_f^2`$ $`+\mathrm{fl}_{02}(1.256562245n_f0.1054295071n_f^2))`$ $`\gamma _{12}^{\psi \psi }`$ $`=12.53336293a_s+a_s^2\left(146.677196412.34063169n_f0.001302012432\mathrm{fl}_{02}n_f\right)`$ $`+a_s^3(2559.641948416.9091301n_f3.300349343n_f^2`$ $`+\mathrm{fl}_{02}(0.9957297081n_f0.07498848634n_f^2))`$ $`\gamma _{14}^{\psi \psi ,NS}`$ $`=13.32896733a_s+a_s^2\left(155.607196213.19066078n_f\right)`$ $`+a_s^3\left(2714.031720441.9494048n_f3.452881831n_f^2\right)`$ $`\gamma _2^{\psi G}`$ $`=0.6666666667a_sn_f7.543209877a_s^2n_f`$ $`+a_s^3\left(37.62337275n_f+12.11248285n_f^2\right)`$ $`\gamma _4^{\psi G}`$ $`=0.3666666667a_sn_f+1.290703704a_s^2n_f`$ $`+a_s^3\left(33.58149273n_f+6.06027262n_f^2\right)`$ $`\gamma _6^{\psi G}`$ $`=0.2619047619a_sn_f+2.761104812a_s^2n_f`$ $`+a_s^3\left(33.41602135n_f+3.537682102n_f^2\right)`$ $`\gamma _8^{\psi G}`$ $`=0.2055555556a_sn_f+3.243957223a_s^2n_f`$ $`+a_s^3\left(28.7612615n_f+2.225433112n_f^2\right)`$ $`\gamma _{10}^{\psi G}`$ $`=0.1696969697a_sn_f+3.407168695a_s^2n_f`$ $`+a_s^3\left(23.93704198n_f+1.449828678n_f^2\right)`$ $`\gamma _{12}^{\psi G}`$ $`=0.1446886447a_sn_f+3.438705999a_s^2n_f`$ $`+a_s^3\left(19.63230379n_f+0.9524545446n_f^2\right)`$ $`\gamma _2^{G\psi }`$ $`=3.555555556a_s+a_s^2\left(48.32921811+5.135802469n_f\right)`$ $`+a_s^3\left(859.4478372+175.649986n_f+4.674897119n_f^2\right)`$ $`\gamma _4^{G\psi }`$ $`=0.9777777778a_s+a_s^2\left(16.1752428+0.6182716049n_f\right)`$ $`+a_s^3\left(315.276255+39.82571027n_f+1.801843621n_f^2\right)`$ $`\gamma _6^{G\psi }`$ $`=0.5587301587a_s+a_s^2\left(9.496317796+0.08884857647n_f\right)`$ $`+a_s^3\left(188.9088124+19.67944546n_f+1.087843741n_f^2\right)`$ $`\gamma _8^{G\psi }`$ $`=0.3915343915a_s+a_s^2\left(6.7576035060.07061952353n_f\right)`$ $`+a_s^3\left(134.7055042+12.3754454n_f+0.7536013741n_f^2\right)`$ $`\gamma _{10}^{G\psi }`$ $`=0.3016835017a_s+a_s^2\left(5.2975769450.1348941718n_f\right)`$ $`+a_s^3\left(104.911278+8.796702078n_f+0.5579674847n_f^2\right)`$ $`\gamma _{12}^{G\psi }`$ $`=0.2455322455a_s+a_s^2\left(4.3986259170.1639529655n_f\right)`$ $`+a_s^3\left(86.18107998+6.735609285n_f+0.4293379075n_f^2\right)`$ $`\gamma _2^{GG}`$ $`=0.6666666667a_sn_f+7.543209877a_s^2n_f`$ $`+a_s^3\left(37.62337275n_f12.11248285n_f^2\right)`$ $`\gamma _4^{GG}`$ $`=a_s^2\left(128.17813.64948148n_f\right)+a_s\left(12.6+0.6666666667n_f\right)`$ $`+a_s^3\left(2066.19278401.3127939n_f10.43150645n_f^2\right)`$ $`\gamma _6^{GG}`$ $`=a_s^2\left(183.053814420.46668466n_f\right)+a_s\left(17.78571429+0.6666666667n_f\right)`$ $`+a_s^3\left(2987.042058566.6373298n_f10.78060861n_f^2\right)`$ $`\gamma _8^{GG}`$ $`=a_s^2\left(219.624098824.69926432n_f\right)+a_s\left(21.26666667+0.6666666667n_f\right)`$ $`+a_s^3\left(3609.35419673.9430658n_f11.20133837n_f^2\right)`$ $`\gamma _{10}^{GG}`$ $`=a_s^2\left(247.665548427.82178573n_f\right)+a_s\left(23.92337662+0.6666666667n_f\right)`$ $`+a_s^3\left(4089.236943755.1340541n_f11.57068198n_f^2\right)`$ $`\gamma _{12}^{GG}`$ $`=a_s^2\left(270.642889230.31377688n_f\right)+a_s\left(26.08168498+0.6666666667n_f\right)`$ $`+a_s^3\left(4483.563048821.1236576n_f11.88665683n_f^2\right)`$ The corresponding coefficient functions read: $`C_{2,2}^\psi `$ $`=1+0.4444444444a_s+a_s^2\left(17.693765895.333333333n_f2.189300412\mathrm{fl}_{02}n_f\right)`$ $`+a_s^3(442.7409693165.1971095n_f24.09201335\mathrm{fl}_{11}n_f+6.030272415n_f^2`$ $`+\mathrm{fl}_{02}(79.04486142n_f+3.325504478n_f^2))`$ $`C_{2,4}^\psi `$ $`=1+6.066666667a_s+a_s^2\left(142.343471916.98791358n_f+0.4858308642\mathrm{fl}_{02}n_f\right)`$ $`+a_s^3(4169.267888901.2351626n_f18.21884618\mathrm{fl}_{11}n_f+23.35503924n_f^2`$ $`+\mathrm{fl}_{02}(16.64834849n_f2.208630689n_f^2))`$ $`C_{2,6}^\psi `$ $`=1+11.17671958a_s+a_s^2\left(302.39873528.0130504n_f+0.4868787285\mathrm{fl}_{02}n_f\right)`$ $`+a_s^3(10069.630851816.322929n_f16.14271761\mathrm{fl}_{11}n_f+42.66273116n_f^2`$ $`+\mathrm{fl}_{02}(24.11778813n_f1.525489143n_f^2))`$ $`C_{2,8}^\psi `$ $`=1+15.52989418a_s+a_s^2\left(470.80741937.9248228n_f+0.3859585393\mathrm{fl}_{02}n_f\right)`$ $`+a_s^3(17162.372452787.297692n_f15.09203827\mathrm{fl}_{11}n_f+61.91177997n_f^2`$ $`+\mathrm{fl}_{02}(22.33201938n_f1.036308122n_f^2))`$ $`C_{2,10}^\psi `$ $`=1+19.30061568a_s+a_s^2\left(639.21066346.86131842n_f+0.3045901308\mathrm{fl}_{02}n_f\right)`$ $`+a_s^3(24953.134973770.10212n_f14.45874451\mathrm{fl}_{11}n_f+80.52097973n_f^2`$ $`+\mathrm{fl}_{02}(19.53359559n_f0.7372434464n_f^2))`$ $`C_{2,12}^\psi `$ $`=1+22.62841097a_s+a_s^2\left(804.585432154.99446579n_f+0.2451231747\mathrm{fl}_{02}n_f\right)`$ $`+a_s^3(33171.455014746.440949n_f14.03541028\mathrm{fl}_{11}n_f+98.3483124n_f^2`$ $`+\mathrm{fl}_{02}(16.98652635n_f0.5471547625n_f^2))`$ $`C_{2,14}^{\psi ,NS}`$ $`=1+2.561093284a_s+a_s^2\left(965.813256462.46549093n_f\right)`$ $`+a_s^3\left(41657.115685708.215623n_f13.73240102\mathrm{fl}_{11}n_f+115.3919490n_f^2\right)`$ $`C_{2,2}^G`$ $`=0.5a_sn_f8.918338961a_s^2n_f`$ $`+a_s^3\left(130.7340963n_f+29.37933515n_f^20.9007972776\mathrm{fl}_{11}^gn_f^2\right)`$ $`C_{2,4}^G`$ $`=0.7388888889a_sn_f14.27158692a_s^2n_f`$ $`+a_s^3\left(346.4612756n_f+46.52017564n_f^21.611816512\mathrm{fl}_{11}^gn_f^2\right)`$ $`C_{2,6}^G`$ $`=0.7051587302a_sn_f20.06849828a_s^2n_f`$ $`+a_s^3\left(715.0372438n_f+61.28545096n_f^21.496036938\mathrm{fl}_{11}^gn_f^2\right)`$ $`C_{2,8}^G`$ $`=0.6440873016a_sn_f23.17873524a_s^2n_f`$ $`+a_s^3\left(996.5038709n_f+68.66467304n_f^21.286400915\mathrm{fl}_{11}^gn_f^2\right)`$ $`C_{2,10}^G`$ $`=0.5861279461a_sn_f24.76678064a_s^2n_f`$ $`+a_s^3\left(1201.206903n_f+72.23614791n_f^21.094394334\mathrm{fl}_{11}^gn_f^2\right)`$ $`C_{2,12}^G`$ $`=0.5358430591a_sn_f25.51669345a_s^2n_f`$ $`+a_s^3\left(1351.047836n_f+73.7936445n_f^20.9344248731\mathrm{fl}_{11}^gn_f^2\right)`$ $`C_{L,2}^\psi `$ $`=1.777777778a_s+a_s^2\left(56.755301524.543209877n_f3.950617284\mathrm{fl}_{02}n_f\right)`$ $`+a_s^3(2544.598087421.6908885n_f7.736698288\mathrm{fl}_{11}n_f+11.8957476n_f^2`$ $`+\mathrm{fl}_{02}(213.9253076n_f+17.91326528n_f^2))`$ $`C_{L,4}^\psi `$ $`=1.066666667a_s+a_s^2\left(47.993989313.413333333n_f0.6945185185\mathrm{fl}_{02}n_f\right)`$ $`+a_s^3(2523.73902383.0520013n_f5.058869512\mathrm{fl}_{11}n_f+10.88895473n_f^2`$ $`+\mathrm{fl}_{02}(55.5530456n_f+2.348005487n_f^2))`$ $`C_{L,6}^\psi `$ $`=0.7619047619a_s+a_s^2\left(40.99619762.69569161n_f0.2524824533\mathrm{fl}_{02}n_f\right)`$ $`+a_s^3(2368.193775340.0691069n_f3.705612526\mathrm{fl}_{11}n_f+9.472190428n_f^2`$ $`+\mathrm{fl}_{02}(24.01322539n_f+0.7652692585n_f^2))`$ $`C_{L,8}^\psi `$ $`=0.5925925926a_s+a_s^2\left(35.876644062.231471683n_f0.1217397796\mathrm{fl}_{02}n_f\right)`$ $`+a_s^3(2215.210875305.4730329n_f2.913702563\mathrm{fl}_{11}n_f+8.337149534n_f^2`$ $`+\mathrm{fl}_{02}(12.97185267n_f+0.344362391n_f^2))`$ $`C_{L,10}^\psi `$ $`=0.4848484848a_s+a_s^2\left(32.017659471.908598248n_f0.06856377261\mathrm{fl}_{02}n_f\right)`$ $`+a_s^3(2081.213222278.0172177n_f2.397641695\mathrm{fl}_{11}n_f+7.452505612n_f^2`$ $`+\mathrm{fl}_{02}(7.947555677n_f+0.1841598535n_f^2))`$ $`C_{L,12}^\psi `$ $`=0.4102564103a_s+a_s^2\left(29.00580651.671007435n_f0.04265241396\mathrm{fl}_{02}n_f\right)`$ $`+a_s^3(1965.791047255.8431044n_f2.035689631\mathrm{fl}_{11}n_f+6.751061503n_f^2`$ $`+\mathrm{fl}_{02}(5.284573837n_f+0.1098463658n_f^2))`$ $`C_{L,14}^{\psi ,NS}`$ $`=0.3555555556a_s+a_s^2\left(26.58488441.488624298n_f\right)`$ $`+a_s^3\left(1866.009187237.5642566n_f1.768138102\mathrm{fl}_{11}n_f+6.182499654n_f^2\right)`$ $`C_{L,2}^G`$ $`=0.6666666667a_sn_f+12.94776709a_s^2n_f`$ $`+a_s^3\left(407.280632n_f20.23959748n_f^20.388939664\mathrm{fl}_{11}^gn_f^2\right)`$ $`C_{L,4}^G`$ $`=0.2666666667a_sn_f+13.81659259a_s^2n_f`$ $`+a_s^3\left(767.7125421n_f36.78419232n_f^20.3984298404\mathrm{fl}_{11}^gn_f^2\right)`$ $`C_{L,6}^G`$ $`=0.1428571429a_sn_f+10.26095364a_s^2n_f`$ $`+a_s^3\left(694.5092121n_f27.9895081n_f^20.3055276256\mathrm{fl}_{11}^gn_f^2\right)`$ $`C_{L,8}^G`$ $`=0.08888888889a_sn_f+7.733545104a_s^2n_f`$ $`+a_s^3\left(592.3307972n_f21.30333681n_f^20.2322211886\mathrm{fl}_{11}^gn_f^2\right)`$ $`C_{L,10}^G`$ $`=0.06060606061a_sn_f+6.023053074a_s^2n_f`$ $`+a_s^3\left(504.5424832n_f16.70544537n_f^20.1805482229\mathrm{fl}_{11}^gn_f^2\right)`$ $`C_{L,12}^G`$ $`=0.04395604396a_sn_f+4.830676204a_s^2n_f`$ $`+a_s^3\left(433.9050534n_f13.47418408n_f^20.1439099345\mathrm{fl}_{11}^gn_f^2\right)`$ The numerical values of the anomalous dimensions for the odd moments of $`F_3`$ read: $`\gamma _1^{\mathrm{ns}}`$ $`=0`$ $`\gamma _3^{\mathrm{ns}}`$ $`=5.555555556a_s+a_s^2\left(70.884773665.12345679n_f\right)`$ $`+a_s^3\left(1244.913602196.4738081n_f1.762002743n_f^2\right)`$ $`\gamma _5^{\mathrm{ns}}`$ $`=8.088888889a_s+a_s^2\left(98.199407417.68691358n_f\right)`$ $`+a_s^3\left(1720.942172278.1581739n_f2.366211248n_f^2\right)`$ $`\gamma _7^{\mathrm{ns}}`$ $`=9.780952381a_s+a_s^2\left(116.41589039.437457798n_f\right)`$ $`+a_s^3\left(2036.492478330.8816595n_f2.739358023n_f^2\right)`$ $`\gamma _9^{\mathrm{ns}}`$ $`=11.05820106a_s+a_s^2\left(130.341404510.77682428n_f\right)`$ $`+a_s^3\left(2277.19805370.5905277n_f3.006446616n_f^2\right)`$ $`\gamma _{11}^{\mathrm{ns}}`$ $`=12.08581049a_s+a_s^2\left(141.690790111.86441897n_f\right)`$ $`+a_s^3\left(2473.311857402.691565n_f3.212753995n_f^2\right)`$ $`\gamma _{13}^{\mathrm{ns}}`$ $`=12.94606135a_s+a_s^2\left(151.298904412.78102552n_f\right)`$ $`+a_s^3\left(2639.409887429.7314605n_f3.37996738n_f^2\right)`$ and the coefficient functions are: $`C_{3,1}^{\mathrm{ns}}`$ $`=14a_s+a_s^2\left(73.33333333+5.333333333n_f\right)`$ $`+a_s^3\left(2652.154437+513.3100408n_f11.35802469n_f^2\right)`$ $`C_{3,3}^{\mathrm{ns}}`$ $`=1+1.666666667a_s+a_s^2\left(14.254040156.742283951n_f\right)`$ $`+a_s^3\left(839.763871745.09953407n_f+1.747689309n_f^2\right)`$ $`C_{3,5}^{\mathrm{ns}}`$ $`=1+7.748148148a_s+a_s^2\left(173.00062919.39801646n_f\right)`$ $`+a_s^3\left(4341.081057961.2756356n_f+22.24125078n_f^2\right)`$ $`C_{3,7}^{\mathrm{ns}}`$ $`=1+12.72248677a_s+a_s^2\left(345.991077730.52332666n_f\right)`$ $`+a_s^3\left(11119.000531960.237096n_f+43.10377964n_f^2\right)`$ $`C_{3,9}^{\mathrm{ns}}`$ $`=1+16.9152381a_s+a_s^2\left(520.005961540.35464229n_f\right)`$ $`+a_s^3\left(18771.996422975.924131n_f+63.17127673n_f^2\right)`$ $`C_{3,11}^{\mathrm{ns}}`$ $`=1+20.54831329a_s+a_s^2\left(690.871966649.17096968n_f\right)`$ $`+a_s^3\left(26941.479873984.411605n_f+82.24581704n_f^2\right)`$ $`C_{3,13}^{\mathrm{ns}}`$ $`=1+23.76237745a_s+a_s^2\left(857.177881757.18099124n_f\right)`$ $`+a_s^3\left(35426.828684976.080869n_f+100.3509187n_f^2\right)`$ ## 6 Acknowledgements A.R. would like to thank Sven Moch, Timo van Ritbergen and Thomas Gehrmann for many useful discussions and Denny Fliegner for technical support. J.V. would like to thank the University of Karlsruhe and the DFG for repeated hospitality. This work was supported by the DFG under Contract Ku 502/8-1 (DFG-Forschergruppe “Quantenfeldtheorie, Computeralgebra und Monte-Carlo-Simulationen” ) and the Graduiertenkolleg “Elementarteilchenphysik an Beschleunigern” at the Univerität Karlsruhe. ## Appendix A Conventions Here we give the complete expressions for the newly computed moments and coefficient functions. The notation of the color factors is as usual: The Casimir operators of the fundamental and adjoint representation are denoted by $`C_F`$ and $`C_A`$ and their values for the color group SU$`(3)`$ are $`\frac{4}{3}`$ and $`3`$, respectively. Additionally we are using the symmetric structure constants of SU$`(3)`$ for which $`d^{abc}d^{abc}=\frac{40}{3}`$. For the trace normalization of the fundamental representation we have inserted $`T_F=\frac{1}{2}`$. For a generic SU$`(n)`$ group the number of generators equals $`N_A=n^21`$. The values of the Riemann $`\zeta `$ function are written as $`\zeta _n=\zeta (n)`$. The only check for the newly computed 3-loop results is from a large $`n_f`$-expansion of where the $`n_f^2`$ terms are calculated. These terms are in agreement with ours. To test the new parallel version of FORM we have also recalculated the lower moments for $`F_2`$ and $`F_L`$ and found complete agreement with except for one missing term. In $`C_{2,2}^G`$ there should be the extra term $`a_s^3n_fC_F^2(\frac{160}{3}\zeta _5)`$. The numerical values given in this reference are correct. It should be noted that the terms in $`d^{abc}d^{abc}`$ enter for the first time at the three loop level and help in the determination of $`P_{qq}^SP_{q\overline{q}}^S`$. ## Appendix B Results for the moments of $`F_2`$ and $`F_L`$ $`\gamma _{10}^{GG}`$ $`=`$ $`a_sC_A\left[+{\displaystyle \frac{18421}{2310}}\right]+a_sn_f\left[+{\displaystyle \frac{2}{3}}\right]`$ $`+a_s^2C_A^2\left[+{\displaystyle \frac{339202487377}{12326391000}}\right]+a_s^2C_Fn_f\left[+{\displaystyle \frac{9284182}{4492125}}\right]+a_s^2C_An_f\left[{\displaystyle \frac{17481908}{1715175}}\right]`$ $`+a_s^3C_A^3\left[+{\displaystyle \frac{239083526238286750523}{1578596520362400000}}\right]`$ $`+a_s^3C_FC_An_f\left[{\displaystyle \frac{3374081335517123191}{36243287457300000}}+{\displaystyle \frac{17831164}{190575}}\zeta _3\right]`$ $`+a_s^3C_F^2n_f\left[{\displaystyle \frac{3009386129483453}{3883209370425000}}{\displaystyle \frac{1344}{3025}}\zeta _3\right]`$ $`+a_s^3C_A^2n_f\left[+{\displaystyle \frac{43228502203851731}{2196562876200000}}{\displaystyle \frac{17746492}{190575}}\zeta _3\right]`$ $`+a_s^3C_Fn_f^2\left[{\displaystyle \frac{453912946493}{420260754375}}\right]+a_s^3C_An_f^2\left[{\displaystyle \frac{2752314359}{815051160}}\right]`$ $`\gamma _{12}^{GG}`$ $`=`$ $`a_sC_A\left[+{\displaystyle \frac{71203}{8190}}\right]+a_sn_f\left[+{\displaystyle \frac{2}{3}}\right]+a_s^2C_A^2\left[+{\displaystyle \frac{16519839244157}{549353259000}}\right]`$ $`+a_s^2C_An_f\left[{\displaystyle \frac{23220103}{2108106}}\right]+a_s^2C_Fn_f\left[+{\displaystyle \frac{74823503}{36540504}}\right]`$ $`+a_s^3C_A^3\left[+{\displaystyle \frac{6993119873800651614841}{42112541869725600000}}\right]`$ $`+a_s^3C_FC_An_f\left[{\displaystyle \frac{49693541388602890695713}{498238162032042432000}}+{\displaystyle \frac{58085396}{585585}}\zeta _3\right]`$ $`+a_s^3C_F^2n_f\left[{\displaystyle \frac{4699124115250144376149}{5480619782352466752000}}{\displaystyle \frac{158}{507}}\zeta _3\right]`$ $`+a_s^3C_A^2n_f\left[+{\displaystyle \frac{74338654569222233539}{3871314390303360000}}{\displaystyle \frac{57902906}{585585}}\zeta _3\right]`$ $`+a_s^3C_Fn_f^2\left[{\displaystyle \frac{50627726543561953}{45626205314289600}}\right]`$ $`+a_s^3C_An_f^2\left[{\displaystyle \frac{2635361358193}{759677078160}}\right]`$ $`\gamma _{10}^{G\psi }`$ $`=`$ $`a_sC_F\left[{\displaystyle \frac{112}{495}}\right]`$ $`+a_s^2C_FC_A\left[{\displaystyle \frac{133349533}{80858250}}\right]+a_s^2C_F^2\left[+{\displaystyle \frac{88631998}{121287375}}\right]+a_s^2C_Fn_f\left[{\displaystyle \frac{74368}{735075}}\right]`$ $`+a_s^3C_FC_A^2\left[{\displaystyle \frac{165894431274725803}{48324383276400000}}{\displaystyle \frac{645584}{81675}}\zeta _3\right]`$ $`+a_s^3C_F^2C_A\left[{\displaystyle \frac{4224421791031474951}{217459724743800000}}+{\displaystyle \frac{645584}{27225}}\zeta _3\right]`$ $`+a_s^3C_F^3\left[+{\displaystyle \frac{6457897459084371893}{326189587115700000}}{\displaystyle \frac{1291168}{81675}}\zeta _3\right]`$ $`+a_s^3C_FC_An_f\left[{\displaystyle \frac{18846629176433}{47069204490000}}+{\displaystyle \frac{1792}{495}}\zeta _3\right]`$ $`+a_s^3C_F^2n_f\left[+{\displaystyle \frac{529979902254031}{1294403123475000}}{\displaystyle \frac{1792}{495}}\zeta _3\right]+a_s^3C_Fn_f^2\left[+{\displaystyle \frac{152267426}{363862125}}\right]`$ $`\gamma _{12}^{G\psi }`$ $`=`$ $`a_sC_F\left[{\displaystyle \frac{79}{429}}\right]+a_s^2C_FC_A\left[{\displaystyle \frac{70863259553}{50645138544}}\right]+a_s^2C_F^2\left[+{\displaystyle \frac{9387059226553}{13927413099600}}\right]`$ $`+a_s^2C_Fn_f\left[{\displaystyle \frac{14257247}{115945830}}\right]`$ $`+a_s^3C_FC_A^2\left[{\displaystyle \frac{24751969909767240119153}{14947144860961272960000}}{\displaystyle \frac{46179137}{6441435}}\zeta _3\right]`$ $`+a_s^3C_F^2C_A\left[{\displaystyle \frac{115986599183122809974023}{5872092623949071520000}}+{\displaystyle \frac{46179137}{2147145}}\zeta _3\right]`$ $`+a_s^3C_F^3\left[+{\displaystyle \frac{2165927305724992215894703}{113037783011019626760000}}{\displaystyle \frac{92358274}{6441435}}\zeta _3\right]`$ $`+a_s^3C_FC_An_f\left[{\displaystyle \frac{64190493078139789}{195540879918384000}}+{\displaystyle \frac{1264}{429}}\zeta _3\right]`$ $`+a_s^3C_F^2n_f\left[+{\displaystyle \frac{1401404001326440151}{13981172914164456000}}{\displaystyle \frac{1264}{429}}\zeta _3\right]`$ $`+a_s^3C_Fn_f^2\left[+{\displaystyle \frac{13454024393417}{41782239298800}}\right]`$ $`\gamma _{10}^{\psi G}`$ $`=`$ $`a_sn_f\left[{\displaystyle \frac{28}{165}}\right]+a_s^2C_Fn_f\left[{\displaystyle \frac{379479917}{125779500}}\right]+a_s^2C_An_f\left[+{\displaystyle \frac{373810079}{150935400}}\right]`$ $`+a_s^3C_FC_An_f\left[+{\displaystyle \frac{926990216580622991}{24162191638200000}}{\displaystyle \frac{643396}{27225}}\zeta _3\right]`$ $`+a_s^3C_F^2n_f\left[{\displaystyle \frac{1091980048536213833}{54364931185950000}}+{\displaystyle \frac{17712}{3025}}\zeta _3\right]`$ $`+a_s^3C_A^2n_f\left[{\displaystyle \frac{21025430857658971}{1022738270400000}}+{\displaystyle \frac{483988}{27225}}\zeta _3\right]`$ $`+a_s^3C_Fn_f^2\left[+{\displaystyle \frac{1584713325754369}{2588806246950000}}\right]+a_s^3C_An_f^2\left[+{\displaystyle \frac{1669885489}{7906140000}}\right]`$ $`\gamma _{12}^{\psi G}`$ $`=`$ $`a_sn_f\left[{\displaystyle \frac{79}{546}}\right]`$ $`+a_s^2C_Fn_f\left[{\displaystyle \frac{9256843807}{3197294100}}\right]+a_s^2C_An_f\left[+{\displaystyle \frac{653436358741}{268572704400}}\right]`$ $`+a_s^3C_FC_An_f\left[+{\displaystyle \frac{4046032530021008148641959}{104630014026728910720000}}{\displaystyle \frac{171207527}{8198190}}\zeta _3\right]`$ $`+a_s^3C_F^2n_f\left[{\displaystyle \frac{2960118366121154186145047}{143866269286752252240000}}+{\displaystyle \frac{2563}{507}}\zeta _3\right]`$ $`+a_s^3C_A^2n_f\left[{\displaystyle \frac{10876559659107463644949}{543532540398591744000}}+{\displaystyle \frac{129763817}{8198190}}\zeta _3\right]`$ $`+a_s^3C_Fn_f^2\left[+{\displaystyle \frac{149081947693135635881}{249119081016021216000}}\right]`$ $`+a_s^3C_An_f^2\left[+{\displaystyle \frac{226617401255197}{4399220898072000}}\right]`$ $`\gamma _{10}^{\psi \psi }`$ $`=`$ $`a_sC_F\left[+{\displaystyle \frac{12055}{1386}}\right]`$ $`+a_s^2C_FC_A\left[+{\displaystyle \frac{19524247733}{523908000}}\right]+a_s^2C_F^2\left[{\displaystyle \frac{9579051036701}{1331250228000}}\right]`$ $`+a_s^2C_Fn_f\left[{\displaystyle \frac{2451995507}{288149400}}\right]+a_s^2\mathrm{fl}_{02}C_Fn_f\left[{\displaystyle \frac{27284}{13476375}}\right]`$ $`+a_s^3C_FC_A^2\left[+{\displaystyle \frac{94091568579766453}{435681892800000}}+{\displaystyle \frac{151796299}{8004150}}\zeta _3\right]`$ $`+a_s^3C_F^2C_A\left[{\displaystyle \frac{16389982059548833}{465937579800000}}{\displaystyle \frac{151796299}{2668050}}\zeta _3\right]`$ $`+a_s^3C_F^3\left[{\displaystyle \frac{2207711300808736405687}{127866318149354400000}}+{\displaystyle \frac{151796299}{4002075}}\zeta _3\right]`$ $`+a_s^3C_FC_An_f\left[{\displaystyle \frac{9007773127403}{389001690000}}{\displaystyle \frac{48220}{693}}\zeta _3\right]`$ $`+a_s^3C_F^2n_f\left[{\displaystyle \frac{75522073210471127}{1230075210672000}}+{\displaystyle \frac{48220}{693}}\zeta _3\right]`$ $`+a_s^3C_Fn_f^2\left[{\displaystyle \frac{27995901056887}{11981252052000}}\right]`$ $`+a_s^3\mathrm{fl}_{02}C_FC_An_f\left[{\displaystyle \frac{1028766412107043}{5177612493900000}}{\displaystyle \frac{12544}{27225}}\zeta _3\right]`$ $`+a_s^3\mathrm{fl}_{02}C_F^2n_f\left[+{\displaystyle \frac{209966063746798}{485401171303125}}+{\displaystyle \frac{12544}{27225}}\zeta _3\right]`$ $`+a_s^3\mathrm{fl}_{02}C_Fn_f^2\left[{\displaystyle \frac{33230913134}{420260754375}}\right]`$ $`\gamma _{12}^{\psi \psi }`$ $`=`$ $`a_sC_F\left[+{\displaystyle \frac{423424}{45045}}\right]`$ $`+a_s^2C_FC_A\left[+{\displaystyle \frac{19487270392267}{486972486000}}\right]+a_s^2C_F^2\left[{\displaystyle \frac{5507868301548461}{731189187729000}}\right]`$ $`+a_s^2C_Fn_f\left[{\displaystyle \frac{90143221429}{9739449720}}\right]+a_s^2\mathrm{fl}_{02}C_Fn_f\left[{\displaystyle \frac{249775}{255783528}}\right]`$ $`+a_s^3C_FC_A^2\left[+{\displaystyle \frac{1395004186652448755863}{6016642459027200000}}+{\displaystyle \frac{25648239313}{1352701350}}\zeta _3\right]`$ $`+a_s^3C_F^2C_A\left[{\displaystyle \frac{98204412073910020058227}{2634913356900224400000}}{\displaystyle \frac{25648239313}{450900450}}\zeta _3\right]`$ $`+a_s^3C_F^3\left[{\displaystyle \frac{81630141011772791446330057}{4747586886462824323920000}}+{\displaystyle \frac{25648239313}{676350675}}\zeta _3\right]`$ $`+a_s^3C_FC_An_f\left[{\displaystyle \frac{25478252190337435009}{1052912430329760000}}{\displaystyle \frac{3387392}{45045}}\zeta _3\right]`$ $`+a_s^3C_F^2n_f\left[{\displaystyle \frac{35346062280941906036867}{526982671380044880000}}+{\displaystyle \frac{3387392}{45045}}\zeta _3\right]`$ $`+a_s^3C_Fn_f^2\left[{\displaystyle \frac{65155853387858071}{26322810758244000}}\right]`$ $`+a_s^3\mathrm{fl}_{02}C_FC_An_f\left[{\displaystyle \frac{69697489543846494691}{332158774688028288000}}{\displaystyle \frac{12482}{39039}}\zeta _3\right]`$ $`+a_s^3\mathrm{fl}_{02}C_F^2n_f\left[+{\displaystyle \frac{86033255402443256197}{219224791294098670080}}+{\displaystyle \frac{12482}{39039}}\zeta _3\right]`$ $`+a_s^3\mathrm{fl}_{02}C_Fn_f^2\left[{\displaystyle \frac{2566080055386457}{45626205314289600}}\right]`$ $`\gamma _{14}^{\mathrm{ns}}`$ $`=`$ $`a_sC_F\left[+{\displaystyle \frac{180121}{18018}}\right]`$ $`+a_s^2C_FC_A\left[+{\displaystyle \frac{288858136265399}{6817614804000}}\right]`$ $`+a_s^2C_Fn_f\left[{\displaystyle \frac{481761665447}{48697248600}}\right]`$ $`+a_s^2C_F^2\left[{\displaystyle \frac{22819142381313407}{2924756750916000}}\right]`$ $`+a_s^3C_FC_An_f\left[{\displaystyle \frac{92531316363319241549}{3685193506154160000}}{\displaystyle \frac{720484}{9009}}\zeta _3\right]`$ $`+a_s^3C_FC_A^2\left[+{\displaystyle \frac{126653245164236390142889}{515927090861582400000}}+{\displaystyle \frac{3663695353}{193243050}}\zeta _3\right]`$ $`+a_s^3C_Fn_f^2\left[{\displaystyle \frac{68167166257767019}{26322810758244000}}\right]`$ $`+a_s^3C_F^2C_A\left[{\displaystyle \frac{96001333716903621488387}{2459252466440209440000}}{\displaystyle \frac{3663695353}{64414350}}\zeta _3\right]`$ $`+a_s^3C_F^2n_f\left[{\displaystyle \frac{37908544797975614512733}{526982671380044880000}}+{\displaystyle \frac{720484}{9009}}\zeta _3\right]`$ $`+a_s^3C_F^3\left[{\displaystyle \frac{40552395746064871242211709}{2373793443231412161960000}}+{\displaystyle \frac{3663695353}{96621525}}\zeta _3\right]`$ $`C_{2,10}^\psi `$ $`=`$ $`1+a_sC_F\left[+{\displaystyle \frac{2006299}{138600}}\right]+a_s^2C_FC_A\left[+{\displaystyle \frac{6124093193824187}{29045459520000}}{\displaystyle \frac{104674}{1155}}\zeta _3\right]`$ $`+a_s^2C_Fn_f\left[{\displaystyle \frac{561457267429757}{15975002736000}}\right]+a_s^2C_F^2\left[+{\displaystyle \frac{558708799987324013}{14760902528064000}}+{\displaystyle \frac{88798}{1155}}\zeta _3\right]`$ $`+a_s^2\mathrm{fl}_{02}C_Fn_f\left[+{\displaystyle \frac{3584203788491}{15689734830000}}\right]`$ $`+a_s^3C_FC_An_f\left[{\displaystyle \frac{21664244926039357214987}{23550349033411200000}}+{\displaystyle \frac{10519793104}{42567525}}\zeta _3{\displaystyle \frac{24110}{693}}\zeta _4\right]`$ $`+a_s^3C_FC_A^2[+{\displaystyle \frac{709221119965457939095237}{235503490334112000000}}{\displaystyle \frac{14713925739913}{6243237000}}\zeta _3`$ $`+{\displaystyle \frac{151796299}{16008300}}\zeta _4+{\displaystyle \frac{190858}{231}}\zeta _5]`$ $`+a_s^3C_Fn_f^2\left[+{\displaystyle \frac{57084428047851551911}{996360920644320000}}+{\displaystyle \frac{48220}{18711}}\zeta _3\right]`$ $`+a_s^3C_F^2C_A[+{\displaystyle \frac{16350009304926933389608829}{8369431733412288000000}}+{\displaystyle \frac{1430215936081}{6163195500}}\zeta _3`$ $`{\displaystyle \frac{151796299}{5336100}}\zeta _4{\displaystyle \frac{22658}{99}}\zeta _5]`$ $`+a_s^3C_F^2n_f\left[{\displaystyle \frac{1521387460036994061010049}{2720065313358993600000}}{\displaystyle \frac{3997754476}{42567525}}\zeta _3+{\displaystyle \frac{24110}{693}}\zeta _4\right]`$ $`+a_s^3C_F^3[{\displaystyle \frac{3247779532370920623770610131}{92155812816602703168000000}}+{\displaystyle \frac{2182208825245282}{1622461215375}}\zeta _3`$ $`+{\displaystyle \frac{151796299}{8004150}}\zeta _4{\displaystyle \frac{75212}{99}}\zeta _5]`$ $`+a_s^3\mathrm{fl}_{02}C_FC_An_f\left[+{\displaystyle \frac{3566946294536415188593}{861140509985448000000}}{\displaystyle \frac{79527463}{56600775}}\zeta _3{\displaystyle \frac{6272}{27225}}\zeta _4\right]`$ $`+a_s^3\mathrm{fl}_{02}C_Fn_f^2\left[{\displaystyle \frac{148475806971656561}{244642190336775000}}+{\displaystyle \frac{33008}{735075}}\zeta _3\right]`$ $`+a_s^3\mathrm{fl}_{02}C_F^2n_f\left[+{\displaystyle \frac{422577250875954453771617}{58772839806506826000000}}{\displaystyle \frac{12949105012}{11037151125}}\zeta _3+{\displaystyle \frac{6272}{27225}}\zeta _4\right]`$ $`+a_s^3\mathrm{fl}_{11}n_f{\displaystyle \frac{d^{abc}d^{abc}}{n}}\left[+{\displaystyle \frac{3753913187503}{352066176000}}+{\displaystyle \frac{81388}{606375}}\zeta _3{\displaystyle \frac{448}{33}}\zeta _5\right]`$ $`C_{2,12}^\psi `$ $`=`$ $`1+a_sC_F\left[+{\displaystyle \frac{183473419}{10810800}}\right]`$ $`+a_s^2C_FC_A\left[+{\displaystyle \frac{21388499873332252399}{87742702527480000}}{\displaystyle \frac{1477711}{15015}}\zeta _3\right]`$ $`+a_s^2C_F^2\left[+{\displaystyle \frac{9127110915702407798941}{131745667845011220000}}+{\displaystyle \frac{1261726}{15015}}\zeta _3\right]`$ $`+a_s^2C_Fn_f\left[{\displaystyle \frac{57904356630607013}{1403883240439680}}\right]+a_s^2\mathrm{fl}_{02}C_Fn_f\left[+{\displaystyle \frac{103555928663269}{563286485361600}}\right]`$ $`+a_s^3C_FC_An_f[{\displaystyle \frac{83712626229337204073275967}{75885504678726462720000}}+{\displaystyle \frac{36842282041}{133783650}}\zeta _3`$ $`{\displaystyle \frac{1693696}{45045}}\zeta _4]+a_s^3C_FC_A^2[+{\displaystyle \frac{163181620367687907864404054279}{45531302807235877632000000}}`$ $`{\displaystyle \frac{641004330821357}{243486243000}}\zeta _3+{\displaystyle \frac{25648239313}{2705402700}}\zeta _4+{\displaystyle \frac{2642336}{3003}}\zeta _5]`$ $`+a_s^3C_Fn_f^2\left[+{\displaystyle \frac{2003755100099657438423887}{28457064254522423520000}}+{\displaystyle \frac{3387392}{1216215}}\zeta _3\right]`$ $`+a_s^3C_F^2C_A[+{\displaystyle \frac{4843191986526849157300333804109}{1709131279126616756611200000}}+{\displaystyle \frac{4948562279669051}{97491891697200}}\zeta _3`$ $`{\displaystyle \frac{25648239313}{901800900}}\zeta _4{\displaystyle \frac{241094}{1287}}\zeta _5]`$ $`+a_s^3C_F^2n_f[{\displaystyle \frac{848389810670975600831798557}{1130004151488672235776000}}{\displaystyle \frac{1326399435709}{12174312150}}\zeta _3`$ $`+{\displaystyle \frac{1693696}{45045}}\zeta _4]+a_s^3C_F^3[+{\displaystyle \frac{131770367393773533536363889790633}{3421680820811486746735622400000}}`$ $`+{\displaystyle \frac{42634681331415644}{24926904127125}}\zeta _3+{\displaystyle \frac{25648239313}{1352701350}}\zeta _4{\displaystyle \frac{88004}{99}}\zeta _5]`$ $`+a_s^3\mathrm{fl}_{02}C_FC_An_f[+{\displaystyle \frac{86565903158314138806418902631}{26931765610480021619328000000}}{\displaystyle \frac{3109476222803}{3165321159000}}\zeta _3`$ $`{\displaystyle \frac{6241}{39039}}\zeta _4]+a_s^3\mathrm{fl}_{02}C_Fn_f^2[{\displaystyle \frac{1960867603733060624851}{4404069467961803640000}}+{\displaystyle \frac{2780}{95823}}\zeta _3]`$ $`+a_s^3\mathrm{fl}_{02}C_F^2n_f[+{\displaystyle \frac{33528586559068805780843402423}{5290168244915718532368000000}}{\displaystyle \frac{751723347541}{791330289750}}\zeta _3`$ $`+{\displaystyle \frac{6241}{39039}}\zeta _4]`$ $`+a_s^3\mathrm{fl}_{11}n_f{\displaystyle \frac{d^{abc}d^{abc}}{n}}\left[+{\displaystyle \frac{809917806143013559}{67966375276800000}}+{\displaystyle \frac{622064791}{851350500}}\zeta _3{\displaystyle \frac{200}{13}}\zeta _5\right]`$ $`C_{2,14}^{\mathrm{ns}}`$ $`=`$ $`1+a_sC_F\left[+{\displaystyle \frac{90849502}{4729725}}\right]`$ $`+a_s^2C_FC_A\left[+{\displaystyle \frac{1345455725874078602801}{4913591341538880000}}{\displaystyle \frac{315626}{3003}}\zeta _3\right]`$ $`+a_s^2C_Fn_f\left[{\displaystyle \frac{1644267296654871017}{35097081010992000}}\right]`$ $`+a_s^2C_F^2\left[+{\displaystyle \frac{31002322638187643268973}{301132955074311360000}}+{\displaystyle \frac{271010}{3003}}\zeta _3\right]`$ $`+a_s^3C_FC_An_f[{\displaystyle \frac{28812973254576289068812626927}{22575937641921122659200000}}+{\displaystyle \frac{31112773830559}{103481653275}}\zeta _3`$ $`{\displaystyle \frac{360242}{9009}}\zeta _4]`$ $`+a_s^3C_FC_A^2[+{\displaystyle \frac{7823621156350047175125815627621}{1896378761921374303372800000}}{\displaystyle \frac{83168919211026563}{28974862917000}}\zeta _3`$ $`+{\displaystyle \frac{3663695353}{386486100}}\zeta _4+{\displaystyle \frac{2738146}{3003}}\zeta _5]`$ $`+a_s^3C_Fn_f^2\left[+{\displaystyle \frac{2361466163828853440218087}{28457064254522423520000}}+{\displaystyle \frac{720484}{243243}}\zeta _3\right]`$ $`+a_s^3C_F^2C_A[+{\displaystyle \frac{64116556842577537005311631547547}{17091312791266167566112000000}}{\displaystyle \frac{13682992796062061}{81243243081000}}\zeta _3`$ $`{\displaystyle \frac{3663695353}{128828700}}\zeta _4{\displaystyle \frac{865562}{9009}}\zeta _5]`$ $`+a_s^3C_F^2n_f\left[{\displaystyle \frac{358451845381207175240043774137}{377340672014967335875200000}}{\displaystyle \frac{12775152582499}{103481653275}}\zeta _3+{\displaystyle \frac{360242}{9009}}\zeta _4\right]`$ $`+a_s^3C_F^3[+{\displaystyle \frac{12924407779985031268428280066600921}{72710717442244093368131976000000}}+{\displaystyle \frac{3556746663996971701}{1695029480644500}}\zeta _3`$ $`+{\displaystyle \frac{3663695353}{193243050}}\zeta _4{\displaystyle \frac{9512108}{9009}}\zeta _5]`$ $`+a_s^3fl_{11}n_f{\displaystyle \frac{d^{abc}d^{abc}}{n}}\left[+{\displaystyle \frac{637395762233410021}{49587712574400000}}+{\displaystyle \frac{78376866703}{65553988500}}\zeta _3{\displaystyle \frac{352}{21}}\zeta _5\right]`$ $`C_{L,10}^\psi `$ $`=`$ $`a_sC_F\left[+{\displaystyle \frac{4}{11}}\right]`$ $`+a_s^2C_FC_A\left[+{\displaystyle \frac{89670761}{8731800}}{\displaystyle \frac{48}{11}}\zeta _3\right]`$ $`+a_s^2C_Fn_f\left[{\displaystyle \frac{163679}{114345}}\right]`$ $`+a_s^2C_F^2\left[{\displaystyle \frac{1999510607}{528273900}}+{\displaystyle \frac{96}{11}}\zeta _3\right]`$ $`+a_s^2\mathrm{fl}_{02}C_Fn_f\left[{\displaystyle \frac{415796}{8085825}}\right]`$ $`+a_s^3C_FC_An_f\left[{\displaystyle \frac{176183576988227323}{1699159381920000}}+{\displaystyle \frac{55485434}{1216215}}\zeta _3\right]`$ $`+a_s^3C_FC_A^2\left[+{\displaystyle \frac{2366034921481985137}{6796637527680000}}{\displaystyle \frac{95022195887}{187297110}}\zeta _3+{\displaystyle \frac{3760}{11}}\zeta _5\right]`$ $`+a_s^3C_Fn_f^2\left[+{\displaystyle \frac{63272639}{11320155}}\right]`$ $`+a_s^3C_F^2C_A\left[{\displaystyle \frac{323139848004267269}{3354750574560000}}+{\displaystyle \frac{22904191}{17325}}\zeta _3{\displaystyle \frac{14240}{11}}\zeta _5\right]`$ $`+a_s^3C_F^2n_f\left[+{\displaystyle \frac{9048874326307637}{190368782604000}}{\displaystyle \frac{1174256}{15015}}\zeta _3\right]`$ $`+a_s^3C_F^3\left[{\displaystyle \frac{887562386698645967383}{3166213592269728000}}{\displaystyle \frac{357031607224}{468242775}}\zeta _3+{\displaystyle \frac{13440}{11}}\zeta _5\right]`$ $`+a_s^3\mathrm{fl}_{02}C_FC_An_f\left[{\displaystyle \frac{68379915239899511}{43491944948760000}}{\displaystyle \frac{36224}{147015}}\zeta _3\right]`$ $`+a_s^3\mathrm{fl}_{02}C_Fn_f^2\left[+{\displaystyle \frac{21670644503}{156897348300}}\right]`$ $`+a_s^3\mathrm{fl}_{02}C_F^2n_f\left[{\displaystyle \frac{319520059852805113}{282697642166940000}}+{\displaystyle \frac{1373248}{1911195}}\zeta _3\right]`$ $`+a_s^3\mathrm{fl}_{11}n_f{\displaystyle \frac{d^{abc}d^{abc}}{n}}\left[{\displaystyle \frac{5073093424963}{528099264000}}{\displaystyle \frac{1820773}{363825}}\zeta _3+{\displaystyle \frac{160}{11}}\zeta _5\right]`$ $`C_{L,12}^\psi `$ $`=`$ $`a_sC_F\left[+{\displaystyle \frac{4}{13}}\right]+a_s^2C_FC_A\left[+{\displaystyle \frac{7126442885209}{811620810000}}{\displaystyle \frac{48}{13}}\zeta _3\right]`$ $`+a_s^2C_Fn_f\left[{\displaystyle \frac{2201663}{1756755}}\right]+a_s^2C_F^2\left[{\displaystyle \frac{12296141077867}{5275535265000}}+{\displaystyle \frac{96}{13}}\zeta _3\right]`$ $`+a_s^2\mathrm{fl}_{02}C_Fn_f\left[{\displaystyle \frac{16072451}{502431930}}\right]+a_s^3C_Fn_f^2\left[+{\displaystyle \frac{5203557911}{1027701675}}\right]`$ $`+a_s^3C_FC_An_f\left[{\displaystyle \frac{61335054761825219657}{658070268956100000}}+{\displaystyle \frac{2151978544}{52026975}}\zeta _3\right]`$ $`+a_s^3C_FC_A^2\left[+{\displaystyle \frac{14897818968827338964411}{52645621516488000000}}{\displaystyle \frac{5467652405147}{10145260125}}\zeta _3+{\displaystyle \frac{5600}{13}}\zeta _5\right]`$ $`+a_s^3C_F^2C_A\left[+{\displaystyle \frac{1508964059584735294818791}{26349133569002244000000}}+{\displaystyle \frac{3065115509662}{2029052025}}\zeta _3{\displaystyle \frac{21600}{13}}\zeta _5\right]`$ $`+a_s^3C_F^2n_f\left[+{\displaystyle \frac{12143703744427185053}{316848648015900000}}{\displaystyle \frac{47440688}{675675}}\zeta _3\right]`$ $`+a_s^3C_F^3\left[{\displaystyle \frac{16091344678479668687707841}{42817342049628646500000}}{\displaystyle \frac{10180155542576}{10145260125}}\zeta _3+1600\zeta _5\right]`$ $`+a_s^3\mathrm{fl}_{02}C_FC_An_f\left[{\displaystyle \frac{592278098533386728681}{576664539388938000000}}{\displaystyle \frac{561855512}{3381753375}}\zeta _3\right]`$ $`+a_s^3\mathrm{fl}_{02}C_Fn_f^2\left[+{\displaystyle \frac{435058406339681}{5280810800265000}}\right]`$ $`+a_s^3\mathrm{fl}_{02}C_F^2n_f\left[{\displaystyle \frac{65983065928499265747263}{85634684099257293000000}}+{\displaystyle \frac{1570446544}{3381753375}}\zeta _3\right]`$ $`+a_s^3\mathrm{fl}_{11}n_f{\displaystyle \frac{d^{abc}d^{abc}}{n}}\left[{\displaystyle \frac{503821438649257451}{62302510670400000}}{\displaystyle \frac{302982523}{70945875}}\zeta _3+{\displaystyle \frac{160}{13}}\zeta _5\right]`$ $`C_{L,14}^{\mathrm{ns}}`$ $`=`$ $`a_sC_F\left[+{\displaystyle \frac{4}{15}}\right]`$ $`+a_s^2C_FC_A\left[+{\displaystyle \frac{3736751546509}{486972486000}}{\displaystyle \frac{16}{5}}\zeta _3\right]`$ $`+a_s^2C_Fn_f\left[{\displaystyle \frac{2263109}{2027025}}\right]`$ $`+a_s^2C_F^2\left[{\displaystyle \frac{6706232197}{4969107000}}+{\displaystyle \frac{32}{5}}\zeta _3\right]`$ $`+a_s^3C_FC_An_f\left[{\displaystyle \frac{14479744644122508496943}{170858971648965600000}}+{\displaystyle \frac{16947716309}{447972525}}\zeta _3\right]`$ $`+a_s^3C_FC_A^2\left[+{\displaystyle \frac{1704131316277822389591041}{7517794752554486400000}}{\displaystyle \frac{47503400282843}{82785322620}}\zeta _3+{\displaystyle \frac{1552}{3}}\zeta _5\right]`$ $`+a_s^3C_Fn_f^2\left[+{\displaystyle \frac{422957746}{91216125}}\right]`$ $`+a_s^3C_F^2C_A\left[+{\displaystyle \frac{1016513978878471248683819}{5269826713800448800000}}+{\displaystyle \frac{51799344112951}{30435780375}}\zeta _32016\zeta _5\right]`$ $`+a_s^3C_F^2n_f\left[+{\displaystyle \frac{58345189864914275683}{1864532428708950000}}{\displaystyle \frac{438969448}{6891885}}\zeta _3\right]`$ $`+a_s^3C_F^3\left[{\displaystyle \frac{8328506007291432016890749}{17917410826921525920000}}{\displaystyle \frac{640509641943719}{517408266375}}\zeta _3+{\displaystyle \frac{5888}{3}}\zeta _5\right]`$ $`+a_s^3fl_{11}n_f{\displaystyle \frac{d^{abc}d^{abc}}{n}}\left[{\displaystyle \frac{110339419075223314037}{15793686454946400000}}{\displaystyle \frac{17420356529}{4682427750}}\zeta _3+{\displaystyle \frac{32}{3}}\zeta _5\right]`$ $`C_{2,10}^G`$ $`=`$ $`a_sn_f\left[{\displaystyle \frac{4352}{7425}}\right]`$ $`+a_s^2C_An_f\left[{\displaystyle \frac{651112454591}{185952412800}}{\displaystyle \frac{54}{55}}\zeta _3\right]`$ $`+a_s^2C_Fn_f\left[{\displaystyle \frac{72533010722807}{6973215480000}}+{\displaystyle \frac{108}{55}}\zeta _3\right]`$ $`+a_s^3C_An_f^2\left[+{\displaystyle \frac{30367858943250477461}{1739677797950400000}}{\displaystyle \frac{18965}{137214}}\zeta _3\right]`$ $`+a_s^3C_A^2n_f[{\displaystyle \frac{114563089982050063118447}{669775952210904000000}}+{\displaystyle \frac{427785744377}{7924108500}}\zeta _3`$ $`+{\displaystyle \frac{241994}{27225}}\zeta _4+{\displaystyle \frac{8}{33}}\zeta _5]`$ $`+a_s^3C_FC_An_f[{\displaystyle \frac{110850413975318446061177}{2487739251069072000000}}+{\displaystyle \frac{7772983651}{7358100750}}\zeta _3`$ $`{\displaystyle \frac{321698}{27225}}\zeta _4+{\displaystyle \frac{244}{11}}\zeta _5]`$ $`+a_s^3C_Fn_f^2\left[+{\displaystyle \frac{25164738348656825457229}{1399353328726353000000}}{\displaystyle \frac{2762978063}{1226350125}}\zeta _3\right]`$ $`+a_s^3C_F^2n_f[{\displaystyle \frac{26284376777719892724358177}{235091359226027304000000}}+{\displaystyle \frac{1782724946402}{77260057875}}\zeta _3`$ $`+{\displaystyle \frac{8856}{3025}}\zeta _4{\displaystyle \frac{1048}{33}}\zeta _5]`$ $`+a_s^3\mathrm{fl}_{11}^gn_f^2{\displaystyle \frac{d^{abc}d^{abc}}{N_A}}\left[{\displaystyle \frac{11661390042871}{128024064000}}{\displaystyle \frac{29154339}{107800}}\zeta _3+{\displaystyle \frac{4408}{11}}\zeta _5\right]`$ $`C_{2,12}^G`$ $`=`$ $`a_sn_f\left[{\displaystyle \frac{8110049}{15135120}}\right]`$ $`+a_s^2C_An_f\left[{\displaystyle \frac{39540735563386331}{10558130155372800}}{\displaystyle \frac{11}{13}}\zeta _3\right]`$ $`+a_s^2C_Fn_f\left[{\displaystyle \frac{6258789011950819}{598533455520000}}+{\displaystyle \frac{22}{13}}\zeta _3\right]`$ $`+a_s^3C_An_f^2\left[+{\displaystyle \frac{1102976289541525512630679}{62778008416037346432000}}{\displaystyle \frac{1574273}{8049132}}\zeta _3\right]`$ $`+a_s^3C_A^2n_f[{\displaystyle \frac{18295361799349570193991309242431}{102830377785469173455616000000}}+{\displaystyle \frac{6726830671058291}{132943488678000}}\zeta _3`$ $`+{\displaystyle \frac{129763817}{16396380}}\zeta _4+{\displaystyle \frac{20}{21}}\zeta _5]`$ $`+a_s^3C_FC_An_f[{\displaystyle \frac{9699486762665487150445457223019}{188522359273360151335296000000}}+{\displaystyle \frac{789241683301607}{132943488678000}}\zeta _3`$ $`{\displaystyle \frac{171207527}{16396380}}\zeta _4+{\displaystyle \frac{4390}{273}}\zeta _5]`$ $`+a_s^3C_Fn_f^2\left[+{\displaystyle \frac{25016184592875203289560351501}{1346588280524001080966400000}}{\displaystyle \frac{13373036672}{7193911725}}\zeta _3\right]`$ $`+a_s^3C_F^2n_f[{\displaystyle \frac{194382513719914205037949320835327}{1555309464005221248516192000000}}+{\displaystyle \frac{542538591728921}{33235872169500}}\zeta _3`$ $`+{\displaystyle \frac{2563}{1014}}\zeta _4{\displaystyle \frac{2220}{91}}\zeta _5]`$ $`+a_s^3\mathrm{fl}_{11}^gn_f^2{\displaystyle \frac{d^{abc}d^{abc}}{N_A}}\left[{\displaystyle \frac{1656600471440498533}{10297935648000000}}{\displaystyle \frac{82799792129}{180589500}}\zeta _3+{\displaystyle \frac{62436}{91}}\zeta _5\right]`$ $`C_{L,10}^G`$ $`=`$ $`a_sn_f\left[+{\displaystyle \frac{2}{33}}\right]`$ $`+a_s^2C_An_f\left[+{\displaystyle \frac{2460678191}{1056547800}}\right]`$ $`+a_s^2C_Fn_f\left[{\displaystyle \frac{509195549}{704365200}}\right]`$ $`+a_s^3C_An_f^2\left[{\displaystyle \frac{140853814103239}{21965628762000}}{\displaystyle \frac{4}{33}}\zeta _3\right]`$ $`+a_s^3C_A^2n_f\left[+{\displaystyle \frac{127219094296097749861}{1623699278087040000}}+{\displaystyle \frac{5906419}{24012450}}\zeta _3{\displaystyle \frac{80}{11}}\zeta _5\right]`$ $`+a_s^3C_FC_An_f\left[{\displaystyle \frac{148061845621707638477}{2638511326891440000}}{\displaystyle \frac{210786157}{31216185}}\zeta _3+{\displaystyle \frac{320}{11}}\zeta _5\right]`$ $`+a_s^3C_Fn_f^2\left[+{\displaystyle \frac{491491787586698683}{188465094777960000}}{\displaystyle \frac{136}{429}}\zeta _3\right]`$ $`+a_s^3C_F^2n_f\left[+{\displaystyle \frac{9053411269935853949}{376930189555920000}}+{\displaystyle \frac{860883476}{156080925}}\zeta _3{\displaystyle \frac{320}{11}}\zeta _5\right]`$ $`+a_s^3\mathrm{fl}_{11}^gn_f^2{\displaystyle \frac{d^{abc}d^{abc}}{N_A}}\left[+{\displaystyle \frac{112883693141257}{4526565120000}}+{\displaystyle \frac{347273501}{4365900}}\zeta _3{\displaystyle \frac{1280}{11}}\zeta _5\right]`$ $`C_{L,12}^G`$ $`=`$ $`a_sn_f\left[+{\displaystyle \frac{4}{91}}\right]+a_s^2C_An_f\left[+{\displaystyle \frac{4978992299}{2685727044}}\right]+a_s^2C_Fn_f\left[{\displaystyle \frac{98593150597}{179847793125}}\right]`$ $`+a_s^3C_An_f^2\left[{\displaystyle \frac{116919410865341069}{22683482755683750}}{\displaystyle \frac{8}{91}}\zeta _3\right]`$ $`+a_s^3C_A^2n_f\left[+{\displaystyle \frac{96912603479263207273229}{1467873372990024000000}}{\displaystyle \frac{16493287}{717341625}}\zeta _3{\displaystyle \frac{480}{91}}\zeta _5\right]`$ $`+a_s^3C_FC_An_f\left[{\displaystyle \frac{2084160567995424969918773}{46709827690503978000000}}{\displaystyle \frac{9006563627}{2152024875}}\zeta _3+{\displaystyle \frac{1920}{91}}\zeta _5\right]`$ $`+a_s^3C_Fn_f^2\left[+{\displaystyle \frac{49342300647723867029}{24328035255470821875}}{\displaystyle \frac{1696}{6825}}\zeta _3\right]`$ $`+a_s^3C_F^2n_f\left[+{\displaystyle \frac{281711596115081884853551}{15569942563501326000000}}+{\displaystyle \frac{2347741862}{717341625}}\zeta _3{\displaystyle \frac{1920}{91}}\zeta _5\right]`$ $`+a_s^3\mathrm{fl}_{11}^gn_f^2{\displaystyle \frac{d^{abc}d^{abc}}{N_A}}\left[+{\displaystyle \frac{2232852976776993919}{56638646064000000}}+{\displaystyle \frac{115891012697}{993242250}}\zeta _3{\displaystyle \frac{15776}{91}}\zeta _5\right]`$ ## Appendix C Results for the moments of $`F_3`$ $`\gamma _3^{\mathrm{ns}}`$ $`=`$ $`a_sC_F\left[+{\displaystyle \frac{25}{6}}\right]+a_s^2C_FC_A\left[+{\displaystyle \frac{535}{27}}\right]+a_s^2C_Fn_f\left[{\displaystyle \frac{415}{108}}\right]+a_s^2C_F^2\left[{\displaystyle \frac{2035}{432}}\right]`$ $`+a_s^3C_FC_An_f\left[{\displaystyle \frac{62249}{3888}}{\displaystyle \frac{100}{3}}\zeta _3\right]+a_s^3C_FC_A^2\left[+{\displaystyle \frac{889433}{7776}}+{\displaystyle \frac{55}{3}}\zeta _3\right]`$ $`+a_s^3C_Fn_f^2\left[{\displaystyle \frac{2569}{1944}}\right]+a_s^3C_F^2C_A\left[{\displaystyle \frac{311213}{15552}}55\zeta _3\right]`$ $`+a_s^3C_F^2n_f\left[{\displaystyle \frac{203627}{7776}}+{\displaystyle \frac{100}{3}}\zeta _3\right]+a_s^3C_F^3\left[{\displaystyle \frac{244505}{15552}}+{\displaystyle \frac{110}{3}}\zeta _3\right]`$ $`+a_s^3n_f{\displaystyle \frac{d^{abc}d^{abc}}{n}}\left[+{\displaystyle \frac{205}{288}}\right]`$ $`\gamma _5^{\mathrm{ns}}`$ $`=`$ $`a_sC_F\left[+{\displaystyle \frac{91}{15}}\right]+a_s^2C_FC_A\left[+{\displaystyle \frac{73223}{2700}}\right]+a_s^2C_Fn_f\left[{\displaystyle \frac{7783}{1350}}\right]+a_s^2C_F^2\left[{\displaystyle \frac{2891}{500}}\right]`$ $`+a_s^3C_FC_An_f\left[{\displaystyle \frac{38587}{2000}}{\displaystyle \frac{728}{15}}\zeta _3\right]+a_s^3C_FC_A^2\left[+{\displaystyle \frac{305342801}{1944000}}+{\displaystyle \frac{1414}{75}}\zeta _3\right]`$ $`+a_s^3C_Fn_f^2\left[{\displaystyle \frac{215621}{121500}}\right]+a_s^3C_F^2C_A\left[{\displaystyle \frac{63892213}{2430000}}{\displaystyle \frac{1414}{25}}\zeta _3\right]`$ $`+a_s^3C_F^2n_f\left[{\displaystyle \frac{5494973}{135000}}+{\displaystyle \frac{728}{15}}\zeta _3\right]+a_s^3C_F^3\left[{\displaystyle \frac{51831073}{3037500}}+{\displaystyle \frac{2828}{75}}\zeta _3\right]`$ $`+a_s^3n_f{\displaystyle \frac{d^{abc}d^{abc}}{n}}\left[+{\displaystyle \frac{931}{4050}}\right]`$ $`\gamma _7^{\mathrm{ns}}`$ $`=`$ $`a_sC_F\left[+{\displaystyle \frac{1027}{140}}\right]+a_s^2C_FC_A\left[+{\displaystyle \frac{67710257}{2116800}}\right]+a_s^2C_Fn_f\left[{\displaystyle \frac{3745727}{529200}}\right]`$ $`+a_s^2C_F^2\left[{\displaystyle \frac{106801937}{16464000}}\right]+a_s^3C_FC_An_f\left[{\displaystyle \frac{2257057261}{106686720}}{\displaystyle \frac{2054}{35}}\zeta _3\right]`$ $`+a_s^3C_FC_A^2\left[+{\displaystyle \frac{219582793861}{1185408000}}+{\displaystyle \frac{92741}{4900}}\zeta _3\right]+a_s^3C_Fn_f^2\left[{\displaystyle \frac{1369936511}{666792000}}\right]`$ $`+a_s^3C_F^2C_A\left[{\displaystyle \frac{629922436973}{20744640000}}{\displaystyle \frac{278223}{4900}}\zeta _3\right]+a_s^3C_F^2n_f\left[{\displaystyle \frac{3150205788689}{62233920000}}+{\displaystyle \frac{2054}{35}}\zeta _3\right]`$ $`+a_s^3C_F^3\left[{\displaystyle \frac{151689902637457}{8712748800000}}+{\displaystyle \frac{92741}{2450}}\zeta _3\right]+a_s^3n_f{\displaystyle \frac{d^{abc}d^{abc}}{n}}\left[+{\displaystyle \frac{56527729}{508032000}}\right]`$ $`\gamma _9^{\mathrm{ns}}`$ $`=`$ $`a_sC_F\left[+{\displaystyle \frac{1045}{126}}\right]+a_s^2C_FC_A\left[+{\displaystyle \frac{242855129}{6804000}}\right]+a_s^2C_Fn_f\left[{\displaystyle \frac{19247947}{2381400}}\right]`$ $`+a_s^2C_F^2\left[{\displaystyle \frac{6993510271}{1000188000}}\right]+a_s^3C_FC_An_f\left[{\displaystyle \frac{63405201707}{2813028750}}{\displaystyle \frac{4180}{63}}\zeta _3\right]`$ $`+a_s^3C_FC_A^2\left[+{\displaystyle \frac{744184331602493}{3600676800000}}+{\displaystyle \frac{1253219}{66150}}\zeta _3\right]+a_s^3C_Fn_f^2\left[{\displaystyle \frac{20297329837}{9001692000}}\right]`$ $`+a_s^3C_F^2C_A\left[{\displaystyle \frac{5022136344247}{150028200000}}{\displaystyle \frac{1253219}{22050}}\zeta _3\right]+a_s^3C_F^2n_f\left[{\displaystyle \frac{1630263834317}{28005264000}}+{\displaystyle \frac{4180}{63}}\zeta _3\right]`$ $`+a_s^3C_F^3\left[{\displaystyle \frac{13829238556849837}{793949234400000}}+{\displaystyle \frac{1253219}{33075}}\zeta _3\right]+a_s^3n_f{\displaystyle \frac{d^{abc}d^{abc}}{n}}\left[+{\displaystyle \frac{13172190779}{200037600000}}\right]`$ $`\gamma _{11}^{\mathrm{ns}}`$ $`=`$ $`a_sC_F\left[+{\displaystyle \frac{31408}{3465}}\right]+a_s^2C_FC_A\left[+{\displaystyle \frac{1448235599}{37422000}}\right]+a_s^2C_Fn_f\left[{\displaystyle \frac{512808781}{57629880}}\right]`$ $`+a_s^2C_F^2\left[{\displaystyle \frac{2454220717793}{332812557000}}\right]+a_s^3C_FC_An_f\left[{\displaystyle \frac{1031510572686647}{43568189280000}}{\displaystyle \frac{251264}{3465}}\zeta _3\right]`$ $`+a_s^3C_FC_A^2\left[+{\displaystyle \frac{390549244457621303}{1742727571200000}}+{\displaystyle \frac{151689577}{8004150}}\zeta _3\right]+a_s^3C_Fn_f^2\left[{\displaystyle \frac{28869611542843}{11981252052000}}\right]`$ $`+a_s^3C_F^2C_A\left[{\displaystyle \frac{43074459020106191}{1198125205200000}}{\displaystyle \frac{151689577}{2668050}}\zeta _3\right]`$ $`+a_s^3C_F^2n_f\left[{\displaystyle \frac{1188145134622636787}{18451128160080000}}+{\displaystyle \frac{251264}{3465}}\zeta _3\right]`$ $`+a_s^3C_F^3\left[{\displaystyle \frac{221430869576690083141}{12786631814935440000}}+{\displaystyle \frac{151689577}{4002075}}\zeta _3\right]`$ $`+a_s^3n_f{\displaystyle \frac{d^{abc}d^{abc}}{n}}\left[+{\displaystyle \frac{5083969985783}{116181838080000}}\right]`$ $`\gamma _{13}^{\mathrm{ns}}`$ $`=`$ $`a_sC_F\left[+{\displaystyle \frac{874733}{90090}}\right]+a_s^2C_FC_A\left[+{\displaystyle \frac{31236494566127}{757512756000}}\right]+a_s^2C_Fn_f\left[{\displaystyle \frac{93360116539}{9739449720}}\right]`$ $`+a_s^2C_F^2\left[{\displaystyle \frac{22445960639039759}{2924756750916000}}\right]`$ $`+a_s^3C_FC_An_f\left[{\displaystyle \frac{90849626920977361109}{3685193506154160000}}{\displaystyle \frac{3498932}{45045}}\zeta _3\right]`$ $`+a_s^3C_FC_A^2\left[+{\displaystyle \frac{41070753377638233401027}{171975696953860800000}}+{\displaystyle \frac{3662719609}{193243050}}\zeta _3\right]`$ $`+a_s^3C_Fn_f^2\left[{\displaystyle \frac{66727681292862571}{26322810758244000}}\right]`$ $`+a_s^3C_F^2C_A\left[{\displaystyle \frac{1400874681707762602284653}{36888786996603141600000}}{\displaystyle \frac{3662719609}{64414350}}\zeta _3\right]`$ $`+a_s^3C_F^2n_f\left[{\displaystyle \frac{36688336888519925613757}{526982671380044880000}}+{\displaystyle \frac{3498932}{45045}}\zeta _3\right]`$ $`+a_s^3C_F^3\left[{\displaystyle \frac{40795447722180713788820819}{2373793443231412161960000}}+{\displaystyle \frac{3662719609}{96621525}}\zeta _3\right]`$ $`+a_s^3n_f{\displaystyle \frac{d^{abc}d^{abc}}{n}}\left[+{\displaystyle \frac{62160363128061559}{1984334964852240000}}\right]`$ $`C_1^{\mathrm{ns}}`$ $`=`$ $`1+a_sC_F\left[3\right]`$ $`+a_s^2C_FC_A\left[23\right]+a_s^2C_Fn_f\left[+4\right]+a_s^2C_F^2\left[+{\displaystyle \frac{21}{2}}\right]`$ $`+a_s^3C_FC_An_f\left[+{\displaystyle \frac{3535}{27}}+24\zeta _3{\displaystyle \frac{80}{3}}\zeta _5\right]+a_s^3C_FC_A^2\left[{\displaystyle \frac{10874}{27}}+{\displaystyle \frac{440}{3}}\zeta _5\right]`$ $`+a_s^3C_Fn_f^2\left[{\displaystyle \frac{230}{27}}\right]+a_s^3C_F^2C_A\left[+{\displaystyle \frac{1241}{9}}{\displaystyle \frac{176}{3}}\zeta _3\right]+a_s^3C_F^2n_f\left[{\displaystyle \frac{133}{18}}{\displaystyle \frac{40}{3}}\zeta _3\right]`$ $`+a_s^3C_F^3\left[{\displaystyle \frac{3}{2}}\right]+a_s^3n_f{\displaystyle \frac{d^{abc}d^{abc}}{n}}\left[{\displaystyle \frac{11}{3}}+8\zeta _3\right]`$ $`C_3^{\mathrm{ns}}`$ $`=`$ $`1+a_sC_F\left[+{\displaystyle \frac{5}{4}}\right]`$ $`+a_s^2C_FC_A\left[+{\displaystyle \frac{5209}{144}}33\zeta _3\right]+a_s^2C_Fn_f\left[{\displaystyle \frac{4369}{864}}\right]+a_s^2C_F^2\left[{\displaystyle \frac{34763}{10368}}+16\zeta _3\right]`$ $`+a_s^3C_FC_An_f\left[{\displaystyle \frac{24877649}{699840}}+{\displaystyle \frac{18539}{405}}\zeta _3{\displaystyle \frac{50}{3}}\zeta _4\right]`$ $`+a_s^3C_FC_A^2\left[+{\displaystyle \frac{92517547}{699840}}{\displaystyle \frac{225337}{405}}\zeta _3+{\displaystyle \frac{55}{6}}\zeta _4+390\zeta _5\right]`$ $`+a_s^3C_Fn_f^2\left[{\displaystyle \frac{12125}{69984}}+{\displaystyle \frac{100}{81}}\zeta _3\right]+a_s^3C_F^2C_A\left[+{\displaystyle \frac{222157399}{559872}}+{\displaystyle \frac{206}{9}}\zeta _3{\displaystyle \frac{55}{2}}\zeta _4260\zeta _5\right]`$ $`+a_s^3C_F^2n_f\left[{\displaystyle \frac{4292227}{34992}}+{\displaystyle \frac{527}{9}}\zeta _3+{\displaystyle \frac{50}{3}}\zeta _4\right]`$ $`+a_s^3C_F^3\left[{\displaystyle \frac{1041473}{373248}}+{\displaystyle \frac{2183}{324}}\zeta _3+{\displaystyle \frac{55}{3}}\zeta _440\zeta _5\right]+a_s^3n_f{\displaystyle \frac{d^{abc}d^{abc}}{n}}\left[{\displaystyle \frac{19477}{5184}}+5\zeta _3\right]`$ $`C_5^{\mathrm{ns}}`$ $`=`$ $`1+a_sC_F\left[+{\displaystyle \frac{523}{90}}\right]`$ $`+a_s^2C_FC_A\left[+{\displaystyle \frac{30312449}{324000}}{\displaystyle \frac{276}{5}}\zeta _3\right]+a_s^2C_Fn_f\left[{\displaystyle \frac{2356859}{162000}}\right]`$ $`+a_s^2C_F^2\left[{\displaystyle \frac{14729261}{1620000}}+{\displaystyle \frac{188}{5}}\zeta _3\right]`$ $`+a_s^3C_FC_An_f\left[{\displaystyle \frac{33317596477}{122472000}}+{\displaystyle \frac{1413442}{14175}}\zeta _3{\displaystyle \frac{364}{15}}\zeta _4\right]`$ $`+a_s^3C_FC_A^2\left[+{\displaystyle \frac{142639414763}{163296000}}{\displaystyle \frac{30424087}{28350}}\zeta _3+{\displaystyle \frac{707}{75}}\zeta _4+548\zeta _5\right]`$ $`+a_s^3C_Fn_f^2\left[+{\displaystyle \frac{784091}{54000}}+{\displaystyle \frac{728}{405}}\zeta _3\right]`$ $`+a_s^3C_F^2C_A\left[+{\displaystyle \frac{179060877287}{255150000}}{\displaystyle \frac{512039}{7875}}\zeta _3{\displaystyle \frac{707}{25}}\zeta _4180\zeta _5\right]`$ $`+a_s^3C_F^2n_f\left[{\displaystyle \frac{28068155797}{127575000}}+{\displaystyle \frac{205712}{4725}}\zeta _3+{\displaystyle \frac{364}{15}}\zeta _4\right]`$ $`+a_s^3C_F^3\left[{\displaystyle \frac{29707273013}{2187000000}}+{\displaystyle \frac{4045486}{10125}}\zeta _3+{\displaystyle \frac{1414}{75}}\zeta _4376\zeta _5\right]`$ $`+a_s^3n_f{\displaystyle \frac{d^{abc}d^{abc}}{n}}\left[{\displaystyle \frac{24801551}{17496000}}+{\displaystyle \frac{86}{45}}\zeta _3\right]`$ $`C_7^{\mathrm{ns}}`$ $`=`$ $`1+a_sC_F\left[+{\displaystyle \frac{48091}{5040}}\right]`$ $`+a_s^2C_FC_A\left[+{\displaystyle \frac{502042084559}{3556224000}}{\displaystyle \frac{4917}{70}}\zeta _3\right]+a_s^2C_Fn_f\left[{\displaystyle \frac{20352710029}{889056000}}\right]`$ $`+a_s^2C_F^2\left[+{\displaystyle \frac{972223395949}{248935680000}}+{\displaystyle \frac{1836}{35}}\zeta _3\right]`$ $`+a_s^3C_FC_An_f\left[{\displaystyle \frac{3373161364214023}{6721263360000}}+{\displaystyle \frac{4838581}{33075}}\zeta _3{\displaystyle \frac{1027}{35}}\zeta _4\right]`$ $`+a_s^3C_FC_A^2\left[+{\displaystyle \frac{5386502616041647}{3360631680000}}{\displaystyle \frac{134407981}{88200}}\zeta _3+{\displaystyle \frac{92741}{9800}}\zeta _4+{\displaystyle \frac{14108}{21}}\zeta _5\right]`$ $`+a_s^3C_Fn_f^2\left[+{\displaystyle \frac{9986153999891}{336063168000}}+{\displaystyle \frac{2054}{945}}\zeta _3\right]`$ $`+a_s^3C_F^2C_A\left[+{\displaystyle \frac{2316127855865945089}{1881953740800000}}{\displaystyle \frac{450138923}{4630500}}\zeta _3{\displaystyle \frac{278223}{9800}}\zeta _4{\displaystyle \frac{4118}{21}}\zeta _5\right]`$ $`+a_s^3C_F^2n_f\left[{\displaystyle \frac{3455733053370931}{9801842400000}}+{\displaystyle \frac{483863}{26460}}\zeta _3+{\displaystyle \frac{1027}{35}}\zeta _4\right]`$ $`+a_s^3C_F^3\left[{\displaystyle \frac{906919428644261623}{14637417984000000}}+{\displaystyle \frac{27358836137}{37044000}}\zeta _3+{\displaystyle \frac{92741}{4900}}\zeta _4{\displaystyle \frac{11224}{21}}\zeta _5\right]`$ $`+a_s^3n_f{\displaystyle \frac{d^{abc}d^{abc}}{n}}\left[{\displaystyle \frac{1521158314519}{1920360960000}}+{\displaystyle \frac{529}{504}}\zeta _3\right]`$ $`C_9^{\mathrm{ns}}`$ $`=`$ $`1+a_sC_F\left[+{\displaystyle \frac{17761}{1400}}\right]`$ $`+a_s^2C_FC_A\left[+{\displaystyle \frac{6247950073621}{34292160000}}{\displaystyle \frac{2858}{35}}\zeta _3\right]+a_s^2C_Fn_f\left[{\displaystyle \frac{363260060687}{12002256000}}\right]`$ $`+a_s^2C_F^2\left[+{\displaystyle \frac{26950029280781}{1008189504000}}+{\displaystyle \frac{6698}{105}}\zeta _3\right]`$ $`+a_s^3C_FC_An_f\left[{\displaystyle \frac{1532454569897512927}{2138802019200000}}+{\displaystyle \frac{121711969}{654885}}\zeta _3{\displaystyle \frac{2090}{63}}\zeta _4\right]`$ $`+a_s^3C_FC_A^2\left[+{\displaystyle \frac{15496996833630287207}{6805279152000000}}{\displaystyle \frac{7620848351}{3969000}}\zeta _3+{\displaystyle \frac{1253219}{132300}}\zeta _4+{\displaystyle \frac{49634}{63}}\zeta _5\right]`$ $`+a_s^3C_Fn_f^2\left[+{\displaystyle \frac{3023214594598871}{68052791520000}}+{\displaystyle \frac{4180}{1701}}\zeta _3\right]`$ $`+a_s^3C_F^2C_A\left[+{\displaystyle \frac{1111314791170240118909}{571643448768000000}}{\displaystyle \frac{4449188411}{41674500}}\zeta _3{\displaystyle \frac{1253219}{44100}}\zeta _4{\displaystyle \frac{2678}{9}}\zeta _5\right]`$ $`+a_s^3C_F^2n_f\left[{\displaystyle \frac{80521228766527611043}{157201948411200000}}{\displaystyle \frac{22156972}{3274425}}\zeta _3+{\displaystyle \frac{2090}{63}}\zeta _4\right]`$ $`+a_s^3C_F^3\left[{\displaystyle \frac{388781612300063841847}{4001504141376000000}}+{\displaystyle \frac{95897878894}{93767625}}\zeta _3+{\displaystyle \frac{1253219}{66150}}\zeta _4{\displaystyle \frac{5092}{9}}\zeta _5\right]`$ $`+a_s^3n_f{\displaystyle \frac{d^{abc}d^{abc}}{n}}\left[{\displaystyle \frac{390642677252603}{756142128000000}}+{\displaystyle \frac{95839}{141750}}\zeta _3\right]`$ $`C_{11}^{\mathrm{ns}}`$ $`=`$ $`1+a_sC_F\left[+{\displaystyle \frac{12815983}{831600}}\right]`$ $`+a_s^2C_FC_A\left[+{\displaystyle \frac{56653343996617}{259334460000}}{\displaystyle \frac{104947}{1155}}\zeta _3\right]+a_s^2C_Fn_f\left[{\displaystyle \frac{117825956269669}{3195000547200}}\right]`$ $`+a_s^2C_F^2\left[+{\displaystyle \frac{25436786950269217}{461278204002000}}+{\displaystyle \frac{84262}{1155}}\zeta _3\right]`$ $`+a_s^3C_FC_An_f\left[{\displaystyle \frac{5339200983355026311}{5825689309440000}}+{\displaystyle \frac{206111663893}{936485550}}\zeta _3{\displaystyle \frac{125632}{3465}}\zeta _4\right]`$ $`+a_s^3C_FC_A^2[+{\displaystyle \frac{5458605391363920639632689}{1884027922672896000000}}{\displaystyle \frac{8531884151173}{3745942200}}\zeta _3+{\displaystyle \frac{151689577}{16008300}}\zeta _4`$ $`+{\displaystyle \frac{209072}{231}}\zeta _5]`$ $`+a_s^3C_Fn_f^2\left[+{\displaystyle \frac{58243229827221303967}{996360920644320000}}+{\displaystyle \frac{251264}{93555}}\zeta _3\right]`$ $`+a_s^3C_F^2C_A[+{\displaystyle \frac{37894668775960020571538737}{13600326566794968000000}}{\displaystyle \frac{52724308943557}{576875098800}}\zeta _3{\displaystyle \frac{151689577}{5336100}}\zeta _4`$ $`{\displaystyle \frac{46958}{99}}\zeta _5]`$ $`+a_s^3C_F^2n_f\left[{\displaystyle \frac{7504061852485890170936687}{10880261253435974400000}}{\displaystyle \frac{2560572011}{85135050}}\zeta _3+{\displaystyle \frac{125632}{3465}}\zeta _4\right]`$ $`+a_s^3C_F^3[{\displaystyle \frac{358252596866786733349510147}{3544454339100103968000000}}+{\displaystyle \frac{156213323360401}{124804708875}}\zeta _3+{\displaystyle \frac{151689577}{8004150}}\zeta _4`$ $`{\displaystyle \frac{49124}{99}}\zeta _5]`$ $`+a_s^3n_f{\displaystyle \frac{d^{abc}d^{abc}}{n}}\left[{\displaystyle \frac{1182009270543683429}{3220560551577600000}}+{\displaystyle \frac{6758}{14175}}\zeta _3\right]`$ $`C_{13}^{\mathrm{ns}}`$ $`=`$ $`1+a_sC_F\left[+{\displaystyle \frac{84292133}{4729725}}\right]`$ $`+a_s^2C_FC_A\left[+{\displaystyle \frac{137082775015804598489}{545954593504320000}}{\displaystyle \frac{1480186}{15015}}\zeta _3\right]`$ $`+a_s^2C_Fn_f\left[{\displaystyle \frac{301032882286150933}{7019416202198400}}\right]`$ $`+a_s^2C_F^2\left[+{\displaystyle \frac{5233449307353261454457}{60226591014862272000}}+{\displaystyle \frac{1210906}{15015}}\zeta _3\right]`$ $`+a_s^3C_FC_An_f\left[{\displaystyle \frac{879328426004243254681619899}{796797799126627858560000}}+{\displaystyle \frac{3046809425749}{12174312150}}\zeta _3{\displaystyle \frac{1749466}{45045}}\zeta _4\right]`$ $`+a_s^3C_FC_A^2[+{\displaystyle \frac{645026845788679520269283647601}{185919486462879833664000000}}{\displaystyle \frac{37039408217872}{14203364175}}\zeta _3`$ $`+{\displaystyle \frac{3662719609}{386486100}}\zeta _4+{\displaystyle \frac{3086186}{3003}}\zeta _5]`$ $`+a_s^3C_Fn_f^2\left[+{\displaystyle \frac{2043359170316436833750887}{28457064254522423520000}}+{\displaystyle \frac{3498932}{1216215}}\zeta _3\right]`$ $`+a_s^3C_F^2C_A[+{\displaystyle \frac{9081431864860519767694223551709}{2441616113038023938016000000}}{\displaystyle \frac{2292102879924479}{48745945848600}}\zeta _3`$ $`{\displaystyle \frac{3662719609}{128828700}}\zeta _4{\displaystyle \frac{6487178}{9009}}\zeta _5]`$ $`+a_s^3C_F^2n_f\left[{\displaystyle \frac{19506028401623230293412698913}{22196510118527490345600000}}{\displaystyle \frac{313559848919}{6087156075}}\zeta _3+{\displaystyle \frac{1749466}{45045}}\zeta _4\right]`$ $`+a_s^3C_F^3[{\displaystyle \frac{298680730292949187763497559489677}{4277101026014358433419528000000}}+{\displaystyle \frac{142217314466274683}{99707616508500}}\zeta _3`$ $`+{\displaystyle \frac{3662719609}{193243050}}\zeta _4{\displaystyle \frac{235868}{693}}\zeta _5]`$ $`+a_s^3n_f{\displaystyle \frac{d^{abc}d^{abc}}{n}}\left[{\displaystyle \frac{65696095155620706347863}{238358315978051068800000}}+{\displaystyle \frac{11802932}{33108075}}\zeta _3\right]`$
warning/0007/hep-ph0007304.html	ar5iv	text	# BUHEP-00-15hep-ph/0007304 TECHNICOLOR 2000 ## 1. The Motivation for Technicolor and Extended Technicolor The elements of the standard model of elementary particles have been in place for more than 25 years now. These include the $`SU(3)SU(2)U(1)`$ gauge model of strong and electroweak interactions . And, they include the Higgs mechanism used to break spontaneously electroweak $`SU(2)U(1)`$ down to the $`U(1)`$ of electromagnetism . In the standard model, couplings of the elementary Higgs scalar bosons also break explicitly quark and lepton chiral–flavor symmetries, giving them hard (Lagrangian) masses. In this quarter century, the standard model has stood up to the most stringent experimental tests . The only indications we have of physics beyond this framework are the existence of neutrino mixing and, presumably, masses (though some would say this physics is accommodated within the standard model); the enormous range of masses, about $`10^{12}`$, between the neutrinos and the top quark; the need for a new source of CP–violation to account for the baryon asymmetry of the universe; the likely presence of cold dark matter; and, possibly, the cosmological constant. These hints are powerful. But they are also obscure, and they do not point unambiguously to any particular extension of the standard model. In addition to these experimental facts, considerable theoretical discomfort and dissatisfaction with the standard model have dogged it from the beginning. All of it concerns the elementary Higgs boson picture of electroweak and flavor symmetry breaking—the cornerstone of the standard model. In particular: 1. Elementary Higgs models provide no dynamical explanation for electroweak symmetry breaking. 2. Elementary Higgs models are unnatural, requiring fine tuning of parameters to enormous precision. 3. Elementary Higgs models with grand unification have a “hierarchy” problem of widely different energy scales. 4. Elementary Higgs models are “trivial”. 5. Elementary Higgs models provide no insight to flavor physics. In nonsupersymmetric Higgs models, there is no explanation why electroweak symmetry breaking occurs and why it has the energy scale of 1 TeV. The Higgs doublet self–interaction potential is $`V(\varphi )=\lambda (\varphi ^{}\varphi v^2)^2`$, where $`v`$ is the vacuum expectation of the Higgs field $`\varphi `$ when $`v^20`$. Its experimental value is $`v=2^{1/4}G_F^{1/2}=246\mathrm{GeV}`$. But what dynamics makes $`v^2>0`$? What dynamics sets its magnitude? In supersymmetric Higgs models, the large top–quark Yukawa coupling drives $`v^2`$ positive, but this just replaces one problem with another or, to be generous, replaces two with one. Elementary Higgs boson models are unnatural. The Higgs boson’s mass, $`M_H^2=2\lambda v^2`$ is quadratically unstable against radiative corrections . Thus, there is no natural reason why $`M_H`$ and $`v`$ should be much less than the energy scale at which the essential physics of the model changes, e.g., a unification scale or the Planck scale of $`10^{16}\mathrm{TeV}`$. To make $`M_H`$ very much less that $`M_P`$, say 1 TeV, the bare Higgs mass must be balanced against its radiative corrections to the fantastic precision of a part in $`M_P^2/M_H^210^{32}`$. In grand–unified Higgs boson models, supersymmetric or not, there are two very different scales of gauge symmetry breaking, the GUT scale of about $`10^{16}\mathrm{GeV}`$ and the electroweak scale of a few 100 GeV. This hierarchy is put in by hand, and must be maintained by unnaturally–fine tuning in ordinary Higgs models, or by the “set it and forget it” feature of supersymmetry. Taken at face value, elementary Higgs boson models are free field theories . To a good approximation, the self–coupling $`\lambda (\mu )`$ of the minimal one–doublet Higgs boson at an energy scale $`\mu `$ is given by $$\lambda (\mu )\frac{\lambda (\mathrm{\Lambda })}{1+(24/16\pi ^2)\lambda (\mathrm{\Lambda })\mathrm{log}(\mathrm{\Lambda }/\mu )}.$$ (1) This coupling vanishes for all $`\mu `$ as the cutoff $`\mathrm{\Lambda }`$ is taken to infinity, hence the description “trivial”. This feature persists in a general class of two–Higgs doublet models and it is probably true of all Higgs models. Triviality really means that elementary–Higgs Lagrangians are meaningful only for scales $`\mu `$ below some cutoff $`\mathrm{\Lambda }_{\mathrm{}}`$ at which new physics sets in. The larger the Higgs couplings are, the lower the scale $`\mathrm{\Lambda }_{\mathrm{}}`$. This relationship translates into the so–called triviality bounds on Higgs masses. For the minimal model, the connection between $`M_H`$ and $`\mathrm{\Lambda }_{\mathrm{}}`$ is $$M_H(\mathrm{\Lambda }_{\mathrm{}})\sqrt{2\lambda (M_H)}v=\frac{2\pi v}{\sqrt{3\mathrm{log}(\mathrm{\Lambda }_{\mathrm{}}/M_H)}}.$$ (2) Clearly, the cutoff has to be greater than the Higgs mass for the effective theory to have some range of validity. From lattice–based arguments , $`\mathrm{\Lambda }_{\mathrm{}}>2\pi M_H`$. Since $`v`$ is fixed at 246 GeV in the minimal model, this implies the triviality bound $`M_H<700\mathrm{GeV}`$<sup>1</sup><sup>1</sup>1Precision electroweak measurements suggesting that $`M_H<200\mathrm{GeV}`$ do not take into account additional interactions that occur if the Higgs is heavy and the scale $`\mathrm{\Lambda }`$ relatively low. Chivukula and Evans have argued that these interactions allow $`M_H=400`$$`500\mathrm{GeV}`$ to be consistent with the precision measurements . If the standard Higgs boson were to be found with a mass this large or larger, we would know for sure that additional new physics is lurking in the range of a few TeV. If the Higgs boson is light, less than 200–300 GeV, as it is expected to be in supersymmetric models, this transition to a more fundamental theory may be postponed until very high energy, but what lies up there worries us nonetheless. Finally, in all elementary Higgs models, supersymmetric or not, every aspect of flavor is completely mysterious, from the primordial symmetry defining the number of quark and lepton generations to the bewildering patterns of flavor breaking. The presence of Higgs bosons has no connection to the existence of multiple identical fermion generations. The flavor–symmetry breaking Yukawa couplings of Higgs bosons to fermions are arbitrary free parameters, put in by hand. As far as we know, it is a logically consistent state of affairs that we may not understand flavor until we understand the physics of the Planck scale. I do not believe this. And, I cannot see how this problem, more pressing and immmediate than any other save electroweak symmetry break itself, can be so cavalierly set aside by those pursuing the “theory of everything”. <sup>2</sup><sup>2</sup>2This is not quite fair. In the early days of the second string revolution, in the mid 1980s, there was a great deal of hope and even expectation that string theory would provide the spectrum—quantum numbers and masses—of the quarks and leptons. Those string pioneers and their descendants have learned how hard the flavor problem is. The dynamical approach to electroweak and flavor symmetry breaking known as technicolor (TC) and extended technicolor (ETC), emerged in the late 1970s in response to these shortcomings of the standard model. This picture was motivated first of all by the premise that every fundamental energy scale should have a dynamical origin. Thus, the weak scale embodied in the Higgs vacuum expectation value $`v=246\mathrm{GeV}`$ should reflect the characteristic energy of a new strong interaction—technicolor—just as the pion decay constant $`f_\pi =93\mathrm{MeV}`$ reflects QCD’s scale $`\mathrm{\Lambda }_{QCD}200\mathrm{MeV}`$. For this reason, I write $`F_\pi =2^{1/4}G_F^{1/2}=246\mathrm{GeV}`$ to emphasize that this quantity has a dynamical origin. Technicolor, a gauge theory of fermions with no elementary scalars, is modeled on the precedent of QCD: The electroweak assignments of quarks to left–handed doublets and right–handed singlets prevent their bare mass terms. Thus, if there are no elementary Higgses to couple to, quarks have a large chiral symmetry, $`SU(6)_LSU(6)_R`$ for three generations. This symmetry is spontaneously broken to the diagonal (vectorial) $`SU(6)`$ subgroup when the QCD gauge coupling grows strong near $`\mathrm{\Lambda }_{QCD}`$. This produces 35 massless Goldstone bosons, the “pions”. According to the Higgs mechanism—whose operation requires no elementary scalar bosons —this yields weak boson masses of $`M_W=M_Z\mathrm{cos}\theta _W=\frac{1}{2}\sqrt{3}gf_\pi 50\mathrm{MeV}`$ . These masses are 1600 times too small, but they do have the right ratio. Suppose, then, that there are technifermions belonging to a complex representation of a technicolor gauge group (taken to be $`SU(N_{TC})`$) whose coupling $`\alpha _{TC}`$ becomes strong at $`\mathrm{\Lambda }_{TC}=100`$s of GeV. If, like quarks, technifermions form left–handed doublets and right–handed singlets under $`SU(2)U(1)`$, then they have no bare masses. When $`\alpha _{TC}`$ becomes strong, the technifermions’ chiral symmetry is spontaneously broken, Goldstone bosons appear, three of them become the longitudinal components of $`W^\pm `$ and $`Z^0`$, and the masses become $`M_W=M_Z\mathrm{cos}\theta _W=\frac{1}{2}gF_\pi `$. Here, $`F_\pi \mathrm{\Lambda }_{TC}`$ is the decay constant of the linear combination of the absorbed “technipions”. Thus, technicolor provides a dynamical basis for electroweak symmetry breaking, one that is based on the familiar and well–understood precedent of QCD. Technicolor, like QCD, is asymptotically free. This solves in one stroke the naturalness, hierarchy, and triviality problems. The mass of all ground–state technihadrons, including Higgs–like scalars (though that language is neither accurate nor useful in technicolor) is of order $`\mathrm{\Lambda }_{TC}`$ or less. There are no large renormalizations of bound state masses, hence no fine–tuning of parameters. If the technicolor gauge symmetry is embedded at a very high energy $`\mathrm{\Lambda }`$ in some grand unified gauge group with a relatively weak coupling, then the characteristic energy scale $`\mathrm{\Lambda }_{TC}`$—where the coupling $`\alpha _{TC}`$ becomes strong enough to trigger chiral symmetry breaking—is naturally exponentially smaller than $`\mathrm{\Lambda }`$. Finally, asymptotically free field theories are nontrivial. A minus sign in the denominator of the analog of Eq. (1) for $`\alpha _{TC}(\mu )`$ prevents one from concluding that it tends to zero for all $`\mu `$ as the cutoff is taken to infinity. No other scenario for the physics of the TeV scale solves these problems so neatly. Period. Technicolor alone does not address the flavor problem. It does not tell us why there are multiple generations and it does not provide explicit breaking of quark and lepton chiral symmetries. Something must play the role of Higgs bosons to communicate electroweak symmetry breaking to quarks and leptons. Furthermore, in all but the minimal TC model with just one doublet of technifermions, there are Goldstone bosons, technipions $`\pi _T`$, in addition to $`W_L^\pm `$ and $`Z_L^0`$. These must be given mass and their masses must be more than 50–100 GeV for them to have escaped detection. Extended technicolor (ETC) was invented to address all these aspects of flavor physics . It was also motivated by the desire to make flavor understandable at energies well below the GUT scale solely in terms of gauge dynamics of the kind that worked so neatly for electroweak symmetry breaking, namely, technicolor. Let me repeat: the ETC approach is based on the gauge dynamics of fermions only. There can be no elementary scalar fields to lead us into the difficulties technicolor itself was invented to escape. ## 2. Dynamical Basics In extended technicolor, ordinary $`SU(3)`$ color, $`SU(N_{TC})`$ technicolor, and flavor symmetries are unified into the ETC gauge group, $`G_{ETC}`$. Thus, we understand flavor, color, and technicolor as subsets of the quantum numbers of extended technicolor. Technicolor and color are exact gauge symmetries. Flavor gauge symmetries are broken at a one or more high energy scales $`\mathrm{\Lambda }_{ETC}M_{ETC}/g_{ETC}`$ where $`M_{ETC}`$ is a typical flavor gauge boson mass. In these lectures, I assume that $`G_{ETC}`$ commutes with electroweak $`SU(2)`$. In this case, it must not commute with electroweak $`U(1)`$, i.e., some part of that $`U(1)`$ must be contained in $`G_{ETC}`$. Otherwise, there will be very light pseudoGoldstone bosons which behave like classical axions and are ruled out experimentally . More generally, all fermions—technifermions, quarks, and leptons—must form no more than four irreducible ETC representations: two equivalent ones for left–handed up and down–type fermions and two inequivalent ones for right–handed up and down fermions (so that up and down mass matrices are not identical). In other words, ETC interactions explicitly break all global flavor symmetries so that there are no very light pseudoGoldstone bosons or fermions. <sup>3</sup><sup>3</sup>3I leave neutrinos out of this discussion. Their very light masses are not yet understood in the ETC framework. The energy scale of ETC gauge symmetry breaking is high, well above the TC scale of 0.1–1.0 TeV, into $`SU(3)SU(N_{TC})`$. The broken gauge interactions, mediated by massive ETC boson exchange, give mass to quarks and leptons by connecting them to technifermions (Fig. 1a). They give mass to technipions by connecting technifermions to each other (Fig. 1b). The graphs in Figs. 1 are convergent: The changes in chirality imply insertions on the technifermion lines of the momentum–dependent dynamical mass, $`\mathrm{\Sigma }(p)`$. This function falls off as $`1/p^2(\mathrm{log}(p/\mathrm{\Lambda }_{TC}))^c`$ in an asymptotically free theory at weak coupling and, in any case, at least as fast as $`1/p`$ . For such a power law, the dominant momentum running around the loop is $`M_{ETC}`$. Then, the operator product expansion tells us that the generic quark or lepton mass and technipion mass are given by the expressions $`m_q(M_{ETC})m_{\mathrm{}}(M_{ETC}){\displaystyle \frac{g_{ETC}^2}{M_{ETC}^2}}\overline{T}_LT_R_{ETC};`$ (3) $`F_T^2M_{\pi _T}^22{\displaystyle \frac{g_{ETC}^2}{M_{ETC}^2}}\overline{T}_LT_R\overline{T}_RT_L_{ETC}.`$ (4) Here, $`m_q(M_{ETC})`$ is the quark mass renormalized at $`M_{ETC}`$. It is a hard mass in that it scales like one for energies below $`M_{ETC}`$. Above that, it falls off more rapidly, like $`\mathrm{\Sigma }(p)`$. The technipion decay constant $`F_T=F_\pi /\sqrt{N}`$ in TC models containing $`N`$ identical electroweak doublets of color–singlet technifermions. The vacuum expectation values $`\overline{T}_LT_R_{ETC}`$ and $`\overline{T}_LT_R\overline{T}_RT_L_{ETC}`$ are the bilinear and quadrilinear technifermion condensates renormalized at $`M_{ETC}`$. The bilinear condensate is related to the one renormalized at $`\mathrm{\Lambda }_{TC}`$, expected by scaling from QCD to be $$\overline{T}_LT_R_{TC}=\frac{1}{2}\overline{T}T_{TC}2\pi F_T^3,$$ (5) by the equation $$\overline{T}T_{ETC}=\overline{T}T_{TC}\mathrm{exp}\left(_{\mathrm{\Lambda }_{TC}}^{M_{ETC}}\frac{d\mu }{\mu }\gamma _m(\mu )\right).$$ (6) The anomalous dimension $`\gamma _m`$ of the operator $`\overline{T}T`$ is given in perturbation theory by $$\gamma _m(\mu )=\frac{3C_2(R)}{2\pi }\alpha _{TC}(\mu )+O(\alpha _{TC}^2),$$ (7) where $`C_2(R)`$ is the quadratic Casimir of the technifermion $`SU(N_{TC})`$–representation $`R`$. For the fundamental representation of $`SU(N_{TC})`$, $`C_2(N_{TC})=(N_{TC}^21)/2N_{TC}`$. Finally, in the large–$`N_{TC}`$ approximation (which will be questionable in the walking technicolor theories we discuss later, but which we adopt anyway for rough estimates) $$\overline{T}_LT_R\overline{T}_RT_L_{ETC}\overline{T}_LT_R_{ETC}\overline{T}_RT_L_{ETC}=\frac{1}{4}\overline{T}T_{ETC}^2.$$ (8) We can obtain an estimate of $`M_{ETC}`$ if we assume that technicolor is QCD–like. In that case, its asymptotic freedom sets in quickly (or “precociously”) at energies above $`\mathrm{\Lambda }_{TC}`$ and $`\gamma _m(\mu )1`$ for $`\mu `$ greater than a few times $`\mathrm{\Lambda }_{TC}`$. Then Eq. (5) applies to $`\overline{T}T_{ETC}`$. For $`N`$ technidoublets, the ETC scale required to generate $`m_q(M_{ETC})1\mathrm{GeV}`$ is $$\mathrm{\Lambda }_{ETC}\frac{M_{ETC}}{g_{ETC}}\sqrt{\frac{2\pi F_\pi ^3}{m_qN^{3/2}}}\frac{10\mathrm{TeV}}{N^{3/4}}.$$ (9) This is pretty low, but the estimate is rough. The typical technipion mass implied by this ETC scale is $$M_{\pi _T}\frac{\overline{T}T_{TC}}{\sqrt{2}\mathrm{\Lambda }_{ETC}F_T}\frac{55\mathrm{GeV}}{N^{1/4}}.$$ (10) Finally, some phenomenological basics: In any model of technicolor, one expects bound technihadrons with a spectrum of mesons paralleling what we see in QCD. The principal targets of collider experiments are the spin–zero technipions and spin–one isovector technirhos and isoscalar techniomegas. In the minimal one–technidoublet model ($`T=(T_U,T_D)`$), the three technipions are the longitudinal compononents $`W_L`$ of the massive weak gauge bosons. Susskind pointed out that the analog of the QCD decay $`\rho \pi \pi `$ is $`\rho _TW_LW_L`$. In the limit that $`M_{\rho _T}M_{W,Z}`$, the equivalence theorem states that the amplitude for $`\rho _TW_LW_L`$ has the same form as the one for $`\rho \pi \pi `$. If we scale technicolor from QCD and use large–$`N_{TC}`$ arguments, it is easy to estimate the strength of this amplitude and the $`\rho _T`$ mass and decay rate : $`M_{\rho _T}=\sqrt{{\displaystyle \frac{3}{N_{TC}}}}{\displaystyle \frac{F_\pi }{f_\pi }}M_\rho 2\sqrt{{\displaystyle \frac{3}{N_{TC}}}}\mathrm{TeV},`$ $`\mathrm{\Gamma }(\rho _TW_LW_L)={\displaystyle \frac{2\alpha _{\rho _T}p_W^3}{3M_{\rho _T}^2}}500\left({\displaystyle \frac{3}{N_{TC}}}\right)^{3/2}\mathrm{GeV}.`$ (11) Here, the naive scaling argument gives $`\alpha _{\rho _T}=(3/N_{TC})\alpha _\rho `$ where $`\alpha _\rho =2.91`$. In the minimal model, a very high energy collider, such as the ill–fated Superconducting Super Collider (SSC) or a $`2\mathrm{TeV}`$ linear collider, is needed to discover the lightest technihadrons.<sup>4</sup><sup>4</sup>4It is possible that, like the attention paid to discovering the minimal standard model Higgs boson, this emphasis on the $`W_LW_L`$ decay mode of the $`\rho _T`$ is somewhat misguided . Since the minimal $`\rho _T`$ is so much heavier than $`2M_W`$, this mode may be suppressed by the high $`W`$–momentum in its decay form factor. Then, $`\rho _T`$ decays to four or more weak bosons may be competitive or even dominate. This means that the minimal $`\rho _T`$ may be wider than indicated in Eq. (2.) and, in any case, that its decays are much more complicated than previously thought. Furthermore, walking technicolor , discussed below, implies that the spectrum of technihadrons cannot be exactly QCD–like. Rather, there must be something like a tower of technirhos extending almost up to $`M_{ETC}>`$ several 100 TeV. Whether or not these would appear as discernible resonances is an open question . All these remarks apply as well to the isoscalar $`\omega _T`$ and its excitations. In nonminimal models, where $`N2`$, the signatures of technicolor ought to be accessible at the Large Hadron Collider (LHC) and at a comparable lepton collider. We shall argue later that technicolor signatures are even likely to be within reach of the Tevatron Collider in Run II! <sup>5</sup><sup>5</sup>5Run II of the Tevatron Collider begins in Spring 2001. The first stage, Run IIa, is intended to collect $`2\mathrm{fb}^1`$ of data with significantly enhanced CDF and DØ detectors featuring new silicon tracking systems. It is planned that, after a brief shutdown to replace damaged silicon, Run IIb will bring the total data sets for each detector to $`15\mathrm{fb}^1`$ or more before the LHC is in full swing in 2006 or so. Before we can do that, however, we must face the obstacles to technicolor dynamics and see how they are overcome. ## 3. Dynamical Perils Technicolor and extended technicolor are challenged by a number of phenomenological hurdles, but the most widely cited causes of the “death of technicolor” are flavor–changing neutral current interactions (FCNC) , precision measurements of electroweak quantities (STU) , and the large mass of the top quark. We discuss these in turn. <sup>6</sup><sup>6</sup>6Much of the discussion here on FCNC and STU is a slightly updated version of that appearing in my 1993 TASI lectures . ### 3.1 Flavor–Changing Neutral Currents Extended technicolor interactions are expected to have flavor–changing neutral currents involving quarks and leptons. The reason for this is simple: Realistic quark mass matrices require ETC transitions between different flavors: $`qTq^{}`$. Thus, there must be ETC currents of the form $`\overline{q}_{L,R}^{}\gamma _\mu T_{L,R}`$ and $`\overline{T}_{L,R}\gamma _\mu q_{L,R}`$; their commutator algebra includes the ETC currents $`\overline{q}_{L,R}^{}\gamma _\mu q_{L,R}`$, and ETC interactions necessarily produce $`\overline{q}q\overline{q}q`$ operators at low energy. Similarly, there will be $`\overline{q}q\overline{\mathrm{}}\mathrm{}`$ and $`\overline{\mathrm{}}\mathrm{}\overline{\mathrm{}}\mathrm{}`$ operators. Even if these interactions are electroweak–eigenstate conserving (or generation–conserving), they will induce FCNC four–fermion operators after diagonalization of mass matrices and transformation to the mass–eigenstate basis. No satisfactory GIM mechanism has ever been found that eliminates these FCNC interactions . The most stringent constraint on ETC comes from $`\|\mathrm{\Delta }S\|=2`$ interactions. Such an interaction has the generic form $$_{\|\mathrm{\Delta }S\|=2}^{}=\frac{g_{ETC}^2V_{ds}^2}{M_{ETC}^2}\overline{d}\mathrm{\Gamma }^\mu s\overline{d}\mathrm{\Gamma }_\mu ^{}s+\mathrm{h}.\mathrm{c}.$$ (12) Here, $`V_{ds}`$ is a mixing–angle factor; it may be complex and seems unlikely to be much smaller in magnitude than the Cabibbo angle, say $`0.1<\|V_{ds}\|<1`$. The matrices $`\mathrm{\Gamma }_\mu `$ and $`\mathrm{\Gamma }_\mu ^{}`$ are left– and/or right–chirality Dirac matrices. I shall put $`\mathrm{\Gamma }_\mu ,\mathrm{\Gamma }_\mu ^{}=\frac{1}{2}\gamma _\mu (1\gamma _5)`$ and count the interaction twice to allow for different chirality terms in $`_{\|\mathrm{\Delta }S\|=2}^{}`$. The contribution of this interaction to the $`K_LK_S`$ mass difference is then estimated to be $`(\mathrm{\Delta }M_K)_{ETC}`$ $``$ $`2\mathrm{R}\mathrm{e}(M_{12})_{ETC}={\displaystyle \frac{4g_{ETC}^2\mathrm{Re}(V_{ds}^2)}{8M_KM_{ETC}^2}}K^0\|\overline{d}\gamma ^\mu (1\gamma _5)s\overline{d}\gamma _\mu (1\gamma _5)s\|\overline{K}^0`$ (13) $``$ $`{\displaystyle \frac{g_{ETC}^2\mathrm{Re}(V_{ds}^2)}{M_{ETC}^2}}f_K^2M_K,`$ where I used the vacuum insertion approximation with $`\mathrm{\Omega }\|\overline{d}\gamma _\mu \gamma _5s\|\overline{K}^0(p)=i\sqrt{2}f_Kp_\mu `$ with $`f_K110\mathrm{MeV}`$. This ETC contribution must be less than the measured mass difference, $`\mathrm{\Delta }M_K=3.5\times 10^{18}\mathrm{TeV}`$. This gives the limit $$\frac{M_{ETC}}{g_{ETC}\sqrt{\mathrm{Re}(V_{ds}^2)}}>1300\mathrm{TeV}.$$ (14) If $`V_{ds}`$ is complex, $`_{\|\mathrm{\Delta }S\|=2}^{}`$ contributes to the imaginary part of the $`K^0\overline{K}^0`$ mass matrix. Using $`\mathrm{Im}(M_{12})=\sqrt{2}\mathrm{\Delta }M_K\|ϵ\|1.15\times 10^{20}\mathrm{TeV}`$, the limit is $$\frac{M_{ETC}}{g_{ETC}\sqrt{\mathrm{Im}(V_{ds}^2)}}>16000\mathrm{TeV}.$$ (15) If we use these large ETC masses and scale the technifermion condensates in Eqs. (3,4) from QCD, i.e., assume the anomalous dimension $`\gamma _m`$ is small so that $`\overline{T}T\overline{T}T_{ETC}\overline{T}T_{ETC}^2\overline{T}T_{TC}^2(4\pi F_T^3)^2`$, we obtain quark and lepton and technipion masses that are 10–1000 times too small, depending on the size of $`V_{ds}`$. This is the FCNC problem. It is remedied by the non–QCD–like dynamics of technicolor with a slowly running gauge coupling, walking technicolor, which will be described in the next section. ### 3.2 Precision Electroweak Measurements Precision electroweak measurements actually challenge technicolor, not extended technicolor. The basic parameters of the standard $`SU(2)U(1)`$ model—$`\alpha (M_Z)`$, $`M_Z`$, $`\mathrm{sin}^2\theta _W`$—are measured so precisely that they may be used to limit new physics at energy scales above 100 GeV . The quantities most sensitive to new physics are defined in terms of correlation functions of the electroweak currents: $$d^4xe^{iqx}\mathrm{\Omega }\|T\left(j_i^\mu (x)j_j^\nu (0)\right)\|\mathrm{\Omega }=ig^{\mu \nu }\mathrm{\Pi }_{ij}(q^2)+q^\mu q^\nu \mathrm{terms}.$$ (16) Once one has accounted for the contributions from standard model physics, including a single Higgs boson (whose mass $`M_H`$ must be assumed), new high–mass physics affects the $`\mathrm{\Pi }_{ij}`$ functions. Assuming that the scale of this physics is well above $`M_{W,Z}`$, it enters the “oblique” correction factors $`S`$, $`T`$, $`U`$ defined by $`S=16\pi {\displaystyle \frac{d}{dq^2}}\left[\mathrm{\Pi }_{33}(q^2)\mathrm{\Pi }_{3Q}(q^2)\right]_{q^2=0}\mathrm{\hspace{0.17em}16}\pi \left[\mathrm{\Pi }_{33}^{^{}}(0)\mathrm{\Pi }_{3Q}^{^{}}(0)\right],`$ $`T={\displaystyle \frac{4\pi }{M_Z^2\mathrm{cos}^2\theta _W\mathrm{sin}^2\theta _W}}\left[\mathrm{\Pi }_{11}(0)\mathrm{\Pi }_{33}(0)\right],`$ $`U=16\pi \left[\mathrm{\Pi }_{11}^{^{}}(0)\mathrm{\Pi }_{33}^{^{}}(0)\right].`$ (17) The parameter $`S`$ is a measure of the splitting between $`M_W`$ and $`M_Z`$ induced by weak–isospin conserving effects; the $`\rho `$–parameter is given by $`\rho M_W^2/M_Z^2\mathrm{cos}^2\theta _W=1+\alpha T`$; the $`U`$–parameter measures weak–isospin breaking in the $`W`$ and $`Z`$ mass splitting. The experimental limits on $`S,T,U`$ are $`S=0.07\pm 0.11(0.09),`$ $`T=0.10\pm 0.14(+0.09),`$ $`U=+0.11\pm 0.15(+0.01).`$ (18) The central values assume $`M_H=100\mathrm{GeV}`$, and the parentheses contain the change for $`M_H=300\mathrm{GeV}`$. The $`S`$ and $`T`$–parameters and $`M_H`$ cannot be obtained simultaneously from data because the Higgs loops resemble oblique effects. The $`S`$–parameter is the one most touted as a show–stopper for technicolor . The value obtained in technicolor by scaling from QCD is $`𝒪(1)`$. For example, for $`N`$ color–singlet technidoublets, Peskin and Takeuchi found the positive result $$S=4\pi \left(1+\frac{M_{\rho _T}^2}{M_{a_{1T}}^2}\right)\frac{F_\pi ^2}{M_{\rho _T}^2}0.25N\frac{N_{TC}}{3}.$$ (19) The resolution to this problem may also be found in walking technicolor. One thing is sure: naive scaling of $`S`$ from QCD is unjustified and probably incorrect in walking gauge theories. No reliable estimate exists because no data on walking gauge theories are available to put into the calculation of $`S`$. ### 3.3 The Top Quark Mass The ETC scale required to produce $`m_t=175\mathrm{GeV}`$ in Eq. (3) is $`0.75\mathrm{TeV}/N^{3/4}`$ for $`N`$ technidoublets. This is uncomfortably close to the TC scale itself. In effect, TC becomes strong and ETC is broken at the same energy; the representation of broken ETC interactions as contact operators is wrong; and all our mass estimates are questionable. It is possible to raise the ETC scale so that it is considerably greater than $`m_t`$, but then one runs into the problem of fine–tuning the ETC coupling $`g_{ETC}`$ (just as in the Nambu–Jona-Lasinio model, where requiring the dynamical fermion mass to be much less than the four–fermion mass scale $`\mathrm{\Lambda }`$ requires fine–tuning the NJL coupling very close to $`4\pi `$. This flouts our cherished principle of naturalness, and we reject it. Another, more direct, problem with ETC generation of the top mass is that there must be large weak isospin violation to raise it so high above the bottom mass. This adversely affects the $`\rho `$ parameter . The large effective ETC coupling to top quarks also makes a large, unwanted contribution to the $`Z\overline{b}b`$ decay rate, in conflict with experiment . In the end, there is no plausible way to understand the top quark’s large mass from ETC. Something more is needed. The best idea so far is topcolor–assisted technicolor , in which a new gauge interaction, topcolor , becomes strong near 1 TeV and generates a large $`\overline{t}t`$ condensate and top mass. This, too, will be described in the next section. ## 4. Dynamical Rescues The FCNC and STU difficulties of technicolor have a common cause: the assumption that technicolor is a just a scaled–up version of QCD. Let us focus on Eqs.(3,4,6), the key equations of extended technicolor. In a QCD–like technicolor theory, asymptotic freedom sets in quickly above $`\mathrm{\Lambda }_{TC}`$, the anomalous dimension $`\gamma _m1`$, and $`\overline{T}T_{ETC}\overline{T}T_{TC}`$. The conclusion that fermion and technipion masses are one or more orders of magnitude too small then followed from the FCNC requirement in Eqs. (14,15) that $`M_{ETC}/g_{ETC}\|V_{ds}\|>1000\mathrm{TeV}`$. Scaling from QCD also means that the technihadron spectrum is just a magnified image of the QCD–hadron spectrum, hence that $`S`$ is too large for all technicolor models except, possibly, the minimal one–doublet model with $`N_{TC}<4`$. The solution to these difficulties lies in technicolor gauge dynamics that are distinctly not QCD–like. The only plausible example is one in which the gauge coupling $`\alpha _{TC}(\mu )`$ evolves slowly, or “walks”, over the large range of energy $`\mathrm{\Lambda }_{TC}<\mu <M_{ETC}`$ . In the extreme walking limit in which $`\alpha _{TC}(\mu )`$ is constant, it is possible to obtain an approximate nonperturbative formula for the $`\overline{T}T`$ anomalous dimension $`\gamma _m`$, namely, $$\gamma _m(\mu )=1\sqrt{1\alpha _{TC}(\mu )/\alpha _C}\mathrm{where}\alpha _C=\frac{\pi }{3C_2(R)}.$$ (20) This reduces to the expression in Eq. (7) for small $`\alpha _{TC}`$. It has been argued that $`\gamma _m=1`$ is the signal for spontaneous chiral symmetry breaking , and, so, $`\alpha _C`$ is called the critical coupling for $`\chi `$SB, with $`\pi /3C_2(R)`$ its approximate value.<sup>7</sup><sup>7</sup>7An attempt to improve upon this approximation and study its accuracy is in Ref. . If we identify $`\mathrm{\Lambda }_{TC}`$ with the scale at which technifermions in the $`SU(N_{TC})`$ fundamental representation condense, then $`\alpha _{TC}(\mathrm{\Lambda }_{TC})=\alpha _C`$. In walking technicolor, $`\alpha _{TC}(\mu )`$ is presumed to remain close to its critical value from $`\mathrm{\Lambda }_{TC}`$ almost up to $`M_{ETC}`$. This implies $`\gamma _m(\mu )1`$, and by Eq. (6), the condensate $`\overline{T}T_{ETC}`$ is enhanced by a factor of 100 or more. This yields quark masses up to a few GeV and reasonably large technipion masses despite the very large ETC mass scale. This is still not enough to account for the top mass; more on that soon. Another consequence of the walking $`\alpha _{TC}`$ is that the spectrum of technihadrons, especially the $`I=0,1`$ vector and axial vector mesons, $`\rho _T`$, $`\omega _T`$, $`a_{1T}`$ and $`f_{1T}`$, cannot be QCD–like . In QCD, the lowest lying isovector $`\rho `$ and $`a_1`$ saturate the spectral functions appearing in Weinberg’s sum rules . Then, the relevant combination $`\rho _V\rho _A`$ of spectral functions falls off like $`1/p^6`$ for $`p>M_{\rho ,a_1}\mathrm{\Lambda }_{QCD}`$, and the spectral integrals converge very rapidly. This “vector meson dominance” of the spectral integrals is related to the precocious onset of asymptotic freedom in QCD. The $`1/p^6`$ momentum dependence is just what one would deduce from a naive, lowest–order calculation of $`\rho _V\rho _A`$ using the asymptotic $`1/p^2`$ behavior of the quark dynamical mass $`\mathrm{\Sigma }(p)`$ . In walking technicolor, the technifermion’s $`\mathrm{\Sigma }(p)`$ falls only like $`1/p^{(2\gamma _m)}1/p`$ for $`\mathrm{\Lambda }_{TC}<M_{ETC}`$, so that $`\rho _V\rho _A1/p^4`$ up to very high energies. To account for this in terms of spin–one technihadrons, there must be something like a tower of $`\rho _T`$ and $`\omega _T`$ extending up to $`M_{ETC}`$. Their mass spectrum, widths, and couplings to currents cannot be predicted. Thus, without experimental knowledge of these states, it is impossible to estimate $`S`$ reliably, any more than it would have been in QCD before the $`\rho `$ and $`a_1`$ were discovered and measured. Another issue that may affect $`S`$ is that it is usually defined assuming that the new physics appears at energies well above $`M_{W,Z}`$. We shall see below that, on the contrary, walking technicolor suggests that there are $`\pi _T`$ and $`\rho _T`$ starting near or not far above $`100\mathrm{GeV}`$. We have seen that extended technicolor cannot explain the top quark’s large mass. An alternative approach was developed in the early 90s based on a new interaction of the third generation quarks. This interaction, called topcolor, was invented as a minimal dynamical scheme to reproduce the simplicity of the one–doublet Higgs model and explain a very large top–quark mass . Here, a large top–quark condensate, $`\overline{t}t`$, is formed by strong interactions at the energy scale, $`\mathrm{\Lambda }_t`$ . To preserve electroweak $`SU(2)`$, topcolor must treat $`t_L`$ and $`b_L`$ the same. To prevent a large $`b`$–condensate and mass, it must violate weak isospin and treat $`t_R`$ and $`b_R`$ differently. In order that the resulting low–energy theory simulate the standard model, particularly its small violation of weak isospin, the topcolor scale must be very high—$`\mathrm{\Lambda }_t10^{15}\mathrm{GeV}m_t`$. Therefore, this original topcolor scenario is highly unnatural, requiring a fine–tuning of couplings of order one part in $`\mathrm{\Lambda }_t^2/m_t^210^{25}`$ (remember Nambu–Jona-Lasinio!). Technicolor is still the most natural mechanism for electroweak symmetry breaking, while topcolor dynamics most aptly explains the top mass. Hill proposed to combine the two into what he called topcolor–assisted technicolor (TC2) . In TC2, electroweak symmetry breaking is driven mainly by technicolor interactions strong near $`1\mathrm{TeV}`$. Light quark, lepton, and technipion masses are still generated by ETC. The topcolor interaction, whose scale is also near $`1\mathrm{TeV}`$, generate $`\overline{t}t`$ and the large top–quark mass.<sup>8</sup><sup>8</sup>8Three massless Goldstone “top–pions” arise from top-quark condensation. Thus, ETC interactions must contribute a few GeV to $`m_t`$ to give the top–pions a mass large enough that $`tb\pi _t^+`$ is not a major decay mode. The scale of ETC interactions still must be at least several $`100\mathrm{TeV}`$ to suppress flavor-changing neutral currents and, so, the technicolor coupling still must walk. Their marriage neatly removes the objections that topcolor is unnatural and that technicolor cannot generate a large top mass. In this scenario, the nonabelian part of topcolor is an ordinary asymptotically free gauge theory. Hill’s original TC2 scheme assumes separate color $`SU(3)`$ and weak hypercharge $`U(1)`$ gauge interactions for the third and for the first two generations of quarks and leptons. In the simplest example, the (electroweak eigenstate) third generation $`(t,b)_{L,R}`$ transform with the usual quantum numbers under the topcolor gauge group $`SU(3)_1U(1)_1`$ while $`(u,d)`$, $`(c,s)`$ transform under a separate group $`SU(3)_2U(1)_2`$. Leptons of the third and the first two generations transform in the obvious way to cancel gauge anomalies. At a scale of order $`1\mathrm{TeV}`$, $`SU(3)_1SU(3)_2U(1)_1U(1)_2`$ is dynamically broken to the diagonal subgroup of ordinary color and weak hypercharge, $`SU(3)_CU(1)_Y`$. At this energy, the $`SU(3)_1U(1)_1`$ couplings are strong while the $`SU(3)_2U(1)_2`$ couplings are weak. This breaking gives rise to massive gauge bosons—a color octet of “colorons” $`V_8`$ and a color singlet $`Z^{}`$. Top, but not bottom, condensation is driven by the fact that the $`SU(3)_1U(1)_1`$ interactions are supercritical for top quarks, but subcritical for bottom.<sup>9</sup><sup>9</sup>9A large bottom condensate is not generated by $`SU(3)_1`$ because it is broken and its coupling does not grow stronger as one descends to lower energies. The difference between top and bottom is caused by the $`U(1)_1`$ couplings of $`t_R`$ and $`b_R`$. If this TC2 scenario is to be natural, i.e., there is no fine–tuning of the $`SU(3)_1`$, the $`U(1)_1`$ couplings cannot be weak. To avoid large violations of weak isospin in this and all other TC2 models , right as well as left–handed members of individual technifermion doublets $`T_{L,R}=(T_U,T_D)_{L,R}`$ must carry the same $`U(1)_1`$ quantum quantum numbers, $`Y_{1L}`$ and $`Y_{1R}`$, respectively . Hill’s simplest TC2 model does not how explain how topcolor is broken. Since natural topcolor requires it to occur near 1 TeV, the most likely cause is technifermion condensation. In Ref. , it was argued that this can be done for $`SU(3)_1SU(3)_2SU(3)_C`$ by arranging that technifermion doublets $`T_1`$ and $`T_2`$ transforming under $`SU(N_{TC})SU(3)_1SU(3)_2`$ as $`(N_{TC},3,1)`$ and $`(N_{TC},1,3)`$ condense with each other as well as themselves, i.e., $$\overline{T}_{iL}T_{jR}=U_{ij}\mathrm{\Delta }_T(i,j=1,2),$$ (21) where $`U`$ is a nondiagonal unitary matrix and $`\mathrm{\Delta }_T`$ the technifermion condensate of $`𝒪(\mathrm{\Lambda }_{TC}^3)`$. The strongly coupled $`U(1)_1`$ plays a critical role in tilting $`U`$ away from the identity, which is the form of the condensate preferred by the color interactions. The breaking $`U(1)_1U(1)_2U(1)_Y`$ is trickier. In order that there is a well–defined $`U(1)_Y`$ boson with standard couplings to all quarks and leptons, this must occur at a somewhat higher scale, several TeV. Thus, the $`Z^{}`$ boson from this breaking has a mass of several TeV and is strongly coupled to technifermions, at least.<sup>10</sup><sup>10</sup>10In Ref. the fermions of the first two generations also need to couple to $`U(1)_1`$. The limits on these strong couplings and $`M_Z^{}`$ from precision electroweak measurements were studied by Chivukula and Terning . Another variant of TC2 has all three generations transforming in the same way under topcolor . This “flavor–universal topcolor” has certain phenomenological advantages (see the second paper of Ref. ), but the problems of the strong $`U(1)_1`$ coupling afflict it too. To employ technicolor in this $`U(1)`$ breaking too, technifermions $`\psi _{L,R}`$ belonging to a higher–dimensional $`SU(N_{TC})`$ representation are introduced. They condense at higher energy than the fundamentals $`T_{iL,R}`$ . The critical reader will note that this scenario also flirts with unnatural fine tuning because the multi–TeV $`Z^{}`$ plays a critical role in top and bottom quark condensation. Another pitfall is that the strong $`U(1)_1`$ coupling may blow up at a Landau singularity at a relatively low energy . To avoid this, unification of $`U(1)_1`$ with the nonabelian $`G_{ETC}`$ must occur at a lower energy still. This is not a very satisfactory state of affairs, but that is how things stand for now with TC2. There are many opportunities for improvement. A variant of topcolor models is called the “top seesaw” mechanism . Its motivation is to realize the original, supposedly more economical, top–condensate idea of the Higgs boson as a fermion–antifermion bound state . Apart from its fine tuning problem, that way failed because it implied a top mass of about 250 GeV. In top seesaw models, an electroweak singlet fermion $`F`$ acquires a dynamical mass of several TeV. Through mixing of $`F`$ with the top quark, it gives the latter a much smaller mass (the seesaw) and the scalar $`\overline{F}F`$ bound state acquires a component with an electroweak symmetry breaking vacuum expectation value. The latest twist on this variant is called the “topcolor jungle gym” . We’ll say no more about these approaches here as they are off our main line of technicolor and extended technicolor. The interested reader should consult the literature. ## 5. Technicolor Phenomenology The coupling $`\alpha _{TC}`$ in walking technicolor decreases slowly if the beta–function $`\beta (\alpha _{TC})=\mu d\alpha _{TC}/d\mu `$ is negative and near zero for a large range of energy $`\mu `$ above $`\mathrm{\Lambda }_{TC}`$. This small $`\beta `$–function may be achieved by having many technifermions in the fundamental representation of $`SU(N_{TC})`$, or a few in higher-dimensional representations, or both . For different reasons, models of topcolor–assisted technicolor also seem to require many technifermions . The technidoublets include $``$ 5 that are color singlets as well as the color triplets $`T_1(N_{TC},3,1)`$ and $`T_2(N_{TC},1,3)`$ mentioned above. The color singlets insure that all quarks and leptons get the appropriate ETC mass and that there is sufficient mixing between the third generation quarks and the two light ones (so that weak decays of the $`b`$–quark are allowed). These requirements suggest that the technicolor scale is much lower than previously thought. If the number $`N`$ of technidoublets is $`𝒪(10)`$ (including 3 for each color triplet), then $`\mathrm{\Lambda }_{TC}F_T=F_\pi /\sqrt{N}<100\mathrm{GeV}`$. This sets the mass scale for the lightest color–singlet technivector mesons, $`M_{\rho _T}M_{\omega _T}2\mathrm{\Lambda }_{TC}<200\mathrm{GeV}`$. These states are produced in hadron and lepton colliders. The mechanism is good old–fashioned vector meson dominance of the $`s`$–channel production of $`\gamma `$, $`Z^0`$, and $`W^\pm `$. The lightest color–octet $`\rho _{T8}`$, bound states of the color–triplet technifermions of TC2, will be heavier, starting, perhaps, at 400–500 GeV. Hadron colliders are needed to produce these states. Isosinglet $`\rho _{T8}`$ bosons are produced by their couplings to the QCD gluon and the $`V_8`$ colorons. In the limit that color interactions are weak compared to technicolor, the chiral symmetry of $`N`$ technidoublets is $`SU(2N)_LSU(2N)_R`$. When it is spontaneously broken, there result $`4N^24`$ technipions in addition to $`W_L^\pm `$ and $`Z_L^0`$, a large number of states if $`N`$ is large. In QCD–like technicolor, these technipions would be very light and the $`\rho _T`$ and $`\omega _T`$ would decay to two or more technipions, with $`\rho _{T8}`$ decaying to color–octet and color–triplet (leptoquark) pairs. Walking technicolor dramatically changes this expectation. In the extreme walking limit, $`\overline{T}T_{ETC}(M_{ETC}/\mathrm{\Lambda }_{TC})\overline{T}T_{TC}`$, so that technipions masses are of order $`\mathrm{\Lambda }_{TC}`$, and they are not pseudoGoldstone bosons at all. Though this extreme limit is theoretically problematic because it is exactly scale–invariant, it is clear that walking technicolor enhances $`\pi _T`$ masses significantly more than it does the $`\rho _T`$ and $`\omega _T`$ masses. Thus, it is likely that $`M_{\pi _T}>\frac{1}{2}M_{\rho _T,\omega _T}`$ and, so, the nominal isospin–conserving decay channels $`\rho _T\pi _T\pi _T`$ and $`\omega _T\pi _T\pi _T\pi _T`$ are closed . If the color–singlet $`\rho _T`$ start near $`200\mathrm{GeV}`$, we expect $`M_{\pi _T}>100\mathrm{GeV}`$. Of course, an explicit ETC model will be needed to make firm mass estimates. This “low–scale technicolor” may be within the reach of CDF and DØ in Run II of the Tevatron Collider <sup>11</sup><sup>11</sup>11Many of these signatures are now encoded in Pythia . It certainly will be accessible at the LHC. Color–singlet $`\rho _T`$ and $`\omega _T`$ may even be detected at LEP200. If a lepton collider with $`\sqrt{s}<500\mathrm{GeV}`$ is built, it will be able carry out precision studies of color–singlet technihadrons. The Very Large Hadron Collider or a multi–TeV lepton collider will be needed to explore more fully the strongly coupled region of walking technicolor. In the rest of this section, we describe a simple model, suitable for experimental studies, of our expectations for the low–lying states of low–scale technicolor—first for the the color–singlet sector, then for color–nonsinglets. ### 5.1 Theory and Experiment for Color–Singlet Technihadrons The flavor problem is hard whether it is attacked with extended technicolor or from any other direction. We theorists need experimental guidance. Experimentalists, in turn, need input from theorists to help design useful searches. Supersymmetry has its MSSM. What follows is a description of the corresponding thing for technicolor, in the sense that it defines a set of incisive experimental tests in terms of a limited number of adjustable parameters. I call this the “Technicolor Straw Man” model (or TCSM). In the TCSM, we assume that we can consider in isolation the lowest-lying bound states of the lightest technifermion doublet, $`(T_U,T_D)`$. If these technifermions belong to the fundamental representation of $`SU(N_{TC})`$, they probably are color singlets. In walking technicolor, ordinary color interactions contribute significantly to the hard mass of $`SU(3)`$ triplets . The lightest technidoublet’s electric charges are unknown; we denote them by $`Q_U`$ and $`Q_D=Q_U1`$. The bound states in question are vector and pseudoscalar mesons. The vectors include a spin–one triplet $`\rho _T^{\pm ,0}`$ and a singlet $`\omega _T`$. In topcolor–assisted technicolor, there is no need to invoke large isospin–violating extended technicolor interactions to explain the top–bottom splitting. Techni–isospin can be, and likely must be, a good approximate symmetry. Then, $`\rho _T`$ and $`\omega _T`$ will be mostly isovector and isoscalar, respectively, and they will be nearly degenerate. Their production in $`\overline{q}q`$ and $`e^+e^{}`$ annihilation is described using vector meson dominance, with propagator matrices that mix them with $`W^\pm `$ and $`\gamma `$, $`Z^0`$. The details are given in Ref. , called TCSM–1 below. I reiterate, mixing of these $`\rho _T`$ and $`\omega _T`$ with their excitations is ignored in the TCSM as is the production of the axial vector $`a_{1T}`$ and the like. The lightest pseudoscalar $`\overline{T}T`$ bound states, the technipions, also comprise an isotriplet $`\mathrm{\Pi }_T^{\pm ,0}`$ and an isosinglet $`\mathrm{\Pi }_T^0`$. However, these are not mass eigenstates; all color–singlet isovector technipions have a $`W_L`$ component. To limit the number of parameters in the TCSM, we make the simplifying assumption that the isotriplets are simple two–state mixtures of the $`W_L^\pm `$, $`Z_L^0`$ and the lightest mass eigenstate pseudo–Goldstone technipions $`\pi _T^\pm ,\pi _T^0`$: $$\|\mathrm{\Pi }_T=\mathrm{sin}\chi \|W_L+\mathrm{cos}\chi \|\pi _T.$$ (22) Here, $`\mathrm{sin}\chi =F_T/F_\pi =1/\sqrt{N}1`$ is an adjustable parameter. The isosinglet is also an admixture, $`\|\mathrm{\Pi }_T^0=\mathrm{cos}\chi ^{}\|\pi _T^0+\mathrm{}`$, where $`\chi ^{}`$ is another adjustable mixing angle and the ellipsis refers to other technipions needed to eliminate the two-technigluon anomaly from the $`\mathrm{\Pi }_T^0`$ chiral current. It is unclear whether, like $`\rho _T^0`$ and $`\omega _T`$, the neutral technipions $`\pi _T^0`$ and $`\pi _T^0`$ will be degenerate as we have previously supposed . On one hand, they both contain the lightest $`\overline{T}T`$ as constituents. On the other, $`\pi _T^0`$ must contain other, presumably heavier, technifermions as a consequence of anomaly cancellation. The calculations and searches presented here assume that $`\pi _T^0`$ and $`\pi _T^0`$ are nearly degenerate. If this is true, and if their widths are roughly equal, there will be appreciable $`\pi _T^0`$$`\pi _T^0`$ mixing. Then, the lightest neutral technipions will be ideally-mixed $`\overline{T}_UT_U`$ and $`\overline{T}_DT_D`$ bound states. In any case, these technipions are expected to couple most strongly to the heaviest fermion pairs that they can. The reason for this is that $`\pi _T`$ couple to ordinary fermions via extended technicolor, $`\pi _T\overline{T}T\overline{f}f`$. Figure 1 suggests that this coupling is proportional to $`m_f`$ (more precisely, the ETC contribution to $`m_f`$). In our studies we assume technipions to be lighter than $`m_t+m_b`$. Then, we expect them to decay as follows: $`\pi _T^+c\overline{b}`$ or $`c\overline{s}`$ or even $`\tau ^+\nu _\tau `$; $`\pi _T^0b\overline{b}`$ and, perhaps $`c\overline{c}`$, $`\tau ^+\tau ^{}`$; and $`\pi _T^0gg`$, $`b\overline{b}`$, $`c\overline{c}`$, $`\tau ^+\tau ^{}`$<sup>12</sup><sup>12</sup>12See Ref. for a discussion and estimate of $`\pi _T`$ decay rates. This puts a premium on heavy–flavor identification in collider experiments. However, this is only an educated guess. The reader is cautioned that the mass–eigenstate neutral $`\pi _T`$ may have a sizable branching ratio to gluon (or even light–quark) pairs. For vanishing electroweak couplings $`g`$ and $`g^{}`$, the $`\rho _T`$ and $`\omega _T`$ decay as $`\rho _T`$ $``$ $`\mathrm{\Pi }_T\mathrm{\Pi }_T=\mathrm{cos}^2\chi (\pi _T\pi _T)+2\mathrm{sin}\chi \mathrm{cos}\chi (W_L\pi _T)+\mathrm{sin}^2\chi (W_LW_L);`$ $`\omega _T`$ $``$ $`\mathrm{\Pi }_T\mathrm{\Pi }_T\mathrm{\Pi }_T=\mathrm{cos}^3\chi (\pi _T\pi _T\pi _T)+\mathrm{}.`$ (23) As noted above however, the all–$`\pi _T`$ modes are likely to be closed. Thus, major decay modes of the $`\rho _T`$ will be $`W_L\pi _T`$ or, if $`M_{\rho _T}<180\mathrm{GeV}`$ (a possibility we regard as unlikely, if not already eliminated by LEP data), $`W_LW_L`$. The $`W^\pm \pi _T^{,0}`$ and $`Z^0\pi _T^\pm `$ decays of $`\rho _T`$ have striking signatures in any collider. Only at LEP is it now possible to detect $`\rho _T^0W^+W^{}`$ above the standard model background. If $`M_{\omega _T}<250\mathrm{GeV}`$, all the $`\omega _T\mathrm{\Pi }_T\mathrm{\Pi }_T\mathrm{\Pi }_T`$ modes are closed. In all cases, the $`\rho _T`$ and $`\omega _T`$ are very narrow, $`\mathrm{\Gamma }(\omega _T)<\mathrm{\Gamma }(\rho _T)<1\mathrm{GeV}`$, because of the smallness of $`\mathrm{sin}\chi `$ and the limited phase space. Therefore, we must consider other decay modes. These are electroweak, suppressed by powers of $`\alpha `$, but not by phase space. The decays $`\rho _T,\omega _TG\pi _T`$, where $`G`$ is a transversely polarized electroweak gauge boson, and $`\rho _T,\omega _T\overline{f}f`$ were calculated in TCSM–1. The $`G\pi _T`$ modes have rates of $`𝒪(\alpha )`$, while the fermion mode $`\overline{f}f`$ rates are $`𝒪(\alpha ^2)`$. The $`\mathrm{\Gamma }(\rho _T,\omega _TG\pi _T)`$ are suppressed by $`1/M_V^2`$ or $`1/M_A^2`$, depending on whether the vector or axial vector part of the electroweak current is involved in the decay. Here, $`M_{V,A}`$ are masses of order $`\mathrm{\Lambda }_{TC}`$ occuring in the dimension–5 operators for these decays. We usually take them equal and vary them from 100 to 400 GeV. For the smaller values of $`M_{V,A}`$, these modes, especially $`\rho _T,\omega _T\gamma \pi _T`$, are as important as the $`W_L\pi _T`$ modes. For larger $`M_{V,A}`$ and $`\|Q_U+Q_D\|>1`$, the $`\overline{f}f`$ decay modes may become competitive. As an illustration, Table 1 lists the relative strengths of the decay amplitudes for the $`\rho _T,\omega _TG\pi _T`$ processes. Figure 2 gives a sense of the $`M_{V,A}`$ dependence of the total decay rates of $`\rho _T`$ and $`\omega _T`$ for $`M_{\rho _T}=210\mathrm{GeV}`$, $`M_{\omega _T}=200`$$`220\mathrm{GeV}`$, $`M_{\pi _T}=110\mathrm{GeV}`$, and $`Q_U=Q_D+1=4/3`$. Figure 3 shows the decay rates for $`Q_U=Q_D=1/2`$. Note how narrow the $`\rho _T`$ and $`\omega _T`$ are. These and all subsequent calculations assume that $`N_{TC}=4`$ and $`\mathrm{sin}\chi =\mathrm{sin}\chi ^{}=1/3`$. Experimental analyses quoted below use the same defaults and (usually) $`Q_U=Q_D+1=4/3`$. Figures 4 and 5 show the cross sections in $`\overline{p}p`$ collisions at $`\sqrt{s}=2\mathrm{TeV}`$ for production of $`\gamma \pi _T`$ and for $`W\pi _T`$, $`Z\pi _T`$ and $`\pi _T\pi _T`$ as a function of $`M_V=M_A`$. Figure 6 shows the $`e^+e^{}`$ rate for $`M_V=100\mathrm{GeV}`$. The production rates in these figures, all in the picobarn range, are typical for the Tevatron for $`M_{\rho _T,\omega _T}<250\mathrm{GeV}`$ and $`M_{\pi _T}<150\mathrm{GeV}`$. That is why we believe Run II will probe a significant portion of the parameter space of low–scale technicolor. Let us turn to the recent searches for color–singlet technihadrons. We begin with analyses by the L3 and DELPHI collaborations at LEP. Note that the LEP experiments can be sensitive to $`\rho _T`$ and $`\omega _T`$ masses significantly above the $`e^+e^{}`$ c.m. energy, $`\sqrt{s}`$. This is because the $`e^+e^{}`$ cross section on resonance is very large for the narrow $`\rho _T`$. Furthermore, masses below the nominal c.m. energy are scanned by the process of radiative return, $`e^+e^{}\rho _T^0/\omega _T+n\gamma `$. The L3 search is based on $`176\mathrm{pb}^1`$ of data taken at an average energy of $`189\mathrm{GeV}`$. The analysis used TCSM–1 for the channels $`e^+e^{}\rho _T^0W^+W^{}`$; $`W_L^\pm \pi _T^{}\mathrm{}\nu _{\mathrm{}}bc`$; $`\pi _T^+\pi _T^{}c\overline{b}b\overline{c}`$; and $`\gamma \pi _T^0\gamma b\overline{b}`$. The TC–scale masses were fixed at $`M_V=M_A=200\mathrm{GeV}`$ and the technifermion charges ranged over $`Q_U+Q_D=5/3,0,1`$. The resulting 95% confidence limits in the $`M_{\rho _T}`$$`M_{\pi _T}`$ plane are shown in Fig. 7. The DELPHI collaboration searched for $`\rho _T^0W^+W^{}`$; $`W^\pm \pi _T^{}`$; $`\pi _T^+\pi _T^{}`$; and $`\mu ^+\mu ^{}`$. Data was taken over a range of $`\sqrt{s}`$ between 161 and $`202\mathrm{GeV}`$ with a variety of integrated luminosities. The modes $`\rho _T^0,\omega _T\pi _T^0\gamma `$, $`\pi _T^0\gamma `$ were neglected. This is not a good assumption if $`M_V<200\mathrm{GeV}`$. The DELPHI exclusion plot is shown in Fig. 8. Since these LEP analyses were done, I have realized that the cross section formulae stated in TCSM–1 are inappropriate for $`\sqrt{s}`$ well below $`M_{\rho _T}`$. This is unimportant for the Tevatron and LHC, where the production rate comes mainly from integrating parton distributions over the resonance pole. However, it may have a significant effect on limits derived from $`e^+e^{}`$ annihilation. This is especially true for the $`W^+W^{}`$ channel, which has a large standard model amplitude interfering with the TCSM one.<sup>13</sup><sup>13</sup>13I thank F. Richard for drawing my attention to this shortcoming of TCSM–1. A correction will be issued soon. Another feature of these analyses not evident in the exclusion plots is that limits on $`M_{\pi _T}`$ approaching $`\sqrt{s}/2`$ should be derivable from $`e^+e^{}\pi _T^+\pi _T^{}`$. We look forward to the new LEP limits that will be announced in the summer of 2000. Both Tevatron collider collaborations have searched for signals of low–scale technicolor. In Run I, only CDF had a vertex detector to find the detached vertices of $`b`$–quark decays. The collaboration used this capability to search for processes signalled by a $`W`$ or photon plus two jets, one of which is $`b`$–tagged: $`\overline{q}qW^\pm ,\gamma ,Z^0`$ $``$ $`\rho _T^{\pm ,0}W_L^\pm \pi _T\mathrm{}^\pm \nu _{\mathrm{}}b+\mathrm{jet}`$ (24) $``$ $`\rho _T^{\pm ,0},\omega _T\gamma \pi _T\gamma b+\mathrm{jet}.`$ These analyses were carried out before the publication of TCSM–1, so they do not include the $`G\pi _T`$ and $`\overline{f}f`$ processes and corresponding branching ratios. They will be included in Run II data analyses. Figure 9 shows data for the $`W\pi _T`$ search on top of a background and signal expected for default parameters with $`M_{\rho _T}=180\mathrm{GeV}`$ and $`M_{\pi _T}=90\mathrm{GeV}`$. The topological cuts leading to the lower figure are described in the second paper of Ref. . The region excluded at 95% confidence level is shown in Fig. 10 . Figure 11 shows the invariant mass of the tagged and untagged jets and the invariant mass difference $`M(\gamma +b+\mathrm{jet})M(b+\mathrm{jet})`$ in a search for $`\omega _T,\rho _T\gamma \pi _T`$ . The good resolution in this mass difference is controlled mainly by that of the electromagnetic energy. The exclusion plot is shown in Fig. 12. It is amusing that the $`2\sigma `$ excesses in Figs. 9 and 12 are both consistent with expectations for a signal with $`M_{\rho _T,\omega _T}200\mathrm{GeV}`$ and $`M_{\pi _T}100\mathrm{GeV}`$. The expected reach of CDF in Run IIa for the $`\rho _TW^\pm \pi _T\mathrm{}^\pm \nu _{\mathrm{}}b\mathrm{jet}`$ processes is shown in Fig. 13 for $`M_V=M_A=200\mathrm{GeV}`$ . This study uses all the processes of TCSM–1. It also assumes the same selections and systematic uncertainty as in the published Run I data , but double the signal efficiency (1.38% vs. 0.69%). The $`5\sigma `$ discovery reach goes up to $`M_{\rho _T}=210\mathrm{GeV}`$ and $`M_{\pi _T}=110\mathrm{GeV}`$, larger than the 95% excluded region in Run I. The region that can be excluded in Run IIa extends up to $`M_{\rho _T}=250\mathrm{GeV}`$ and $`M_{\pi _T}=145\mathrm{GeV}`$. When $`M_V=400\mathrm{GeV}`$, the $`5\sigma `$ discovery and 95% exclusion regions are only slightly larger than this. The Run I DØ detector had superior calorimetry and hermiticity. The collaboration studied its Drell–Yan data to search for $`\rho _T,\omega _Te^+e^{}`$ . The data and the excluded region are shown in Fig. 14 for $`Q_U=Q_D+1=4/3`$, $`M_V=100`$$`400\mathrm{GeV}`$ and $`M_{\rho _T}M_{\pi _T}=100\mathrm{GeV}`$. Increasing $`M_V`$ and decreasing $`M_{\rho _T}M_{\pi _T}`$ both increase the branching ratio for the $`e^+e^{}`$ channel. For the parameters considered here, $`M_{\rho _T}=M_{\omega _T}<150`$$`200\mathrm{GeV}`$ is excluded at the 95% CL. The expected reach of DØ in Run IIa for $`\rho _T,\omega _Te^+e^{}`$ with $`M_V=100`$ and $`200\mathrm{GeV}`$ and other TCSM parameters (see above) is shown in Fig. 15 . As long as $`Q_U+Q_D=𝒪(1)`$, masses $`M_{\rho _T,\omega _T}`$ up to 450–500 GeV should be accessible in the $`e^+e^{}`$ channel. The ATLAS collaboration has studied its reach for $`\rho _TW^\pm Z,W^\pm \pi _T,Z\pi _T`$ and for $`\omega _T\gamma \pi _T`$ . Figure 16 shows $`\rho _T^\pm W^\pm Z\mathrm{}^\pm \nu _{\mathrm{}}\mathrm{}^+\mathrm{}^{}`$ for several $`\rho _T`$ and $`\pi _T`$ masses and a luminosity of $`10\mathrm{fb}^1`$. Detailed studies have not been published in which all the TCSM processes have been included and the parameters varied over a wide range. Still, it is clear from Fig. 16 that the higher energy and luminosity of the LHC ought to make it possible to completely exclude, or discover, low–scale technicolor for any reasonable parameterization. ### 5.2 Color–Nonsinglet Technihadrons The experimental searches so far in the color–nonsinglet sector of low–scale technicolor have been inspired by the phenomenology of a pre–TCSM, one–family TC model. This model contains a single doublet each of color–triplet techniquarks $`Q=(U,D)`$ and of color–singlet technileptons $`L=(N,E)`$ . We consider these searches first, commenting on the status of the TCSM for color nonsinglets at the end. Assuming that techni–isospin is conserved, production of color–nonsinglet states is assumed to proceed through the lightest isoscalar color–octet technirho, $`\rho _{T8}`$: $`\overline{q}q,ggg\rho _{T8}`$ $``$ $`\pi _{T8}\pi _{T8},\pi _{\overline{L}Q}\pi _{\overline{Q}L}`$ (25) $``$ $`\overline{q}q,gg\mathrm{dijets}.`$ Here, $`\pi _{T8}=\pi _{T8}^\pm ,\pi _{T8}^0,\pi _{T8}^0^{}\eta _T`$ are four color–octet technipions that are expected to decay to heavy $`\overline{q}q`$ pairs; $`\pi _{\overline{L}Q}`$ are four color–triplet leptoquarks expected to decay to heavy $`\overline{\mathrm{}}q`$ with the corresponding charges. If TC2 is invoked, the neutral $`\pi _{T8}`$ decay to $`\overline{b}b`$ and, possibly, $`gg`$ as readily as to $`\overline{t}t`$<sup>14</sup><sup>14</sup>14The ATLAS collaboration has studied $`gg\eta _T\overline{t}t`$ . Even if this is not the dominant decay mode of an $`\eta _T`$ in a TC2 model, other bosons, such as the colorons $`V_8`$, will have a sizable $`\overline{t}t`$ branching ratio and the ATLAS study serves as a promising prototype of a search for this process. The two $`\rho _{T8}\pi _T\pi _T`$ searches are by CDF for leptoquark technipions: $`\pi _{\overline{E}D}\tau ^+b`$ where the $`b`$ is not tagged and $`\pi _{\overline{N}D}\nu b`$, $`\nu c`$ with one or two tagged jets . These are based on $`110\mathrm{pb}^1`$ and $`88\mathrm{pb}^1`$ of Run I data, respectively. The exclusion plot for the $`\tau ^+\tau ^{}+`$dijet signal is shown in Fig. 17 as a function of the $`\pi _{T8}`$$`\pi _{\overline{L}Q}`$ mass difference. The theoretically likely case is that this mass difference is about 50 GeV, implying a 95% excluded region extending over $`200<M_{\rho _{T8}}2M_{\pi _{\overline{L}Q}}<500\mathrm{GeV}`$. Figure 18 shows the reach for $`\rho _{T8}^0b\overline{b}\nu \overline{\nu }`$ with at least one $`b`$–jet tagged. Here the 95% limits extends over $`300<M_{\rho _{T8}}2M_{\pi _{\overline{L}Q}}<600\mathrm{GeV}`$. The search for $`\pi _{\overline{L}Q}c\nu `$ excludes a similar range. These limits are quite impressive. However, it is not clear how they will be affected by the complications of topcolor–assisted technicolor. Given the walking technicolor enhancement of $`\pi _T`$ masses, it is likely that the $`\rho _{T8}\pi _T\pi _T`$ channels are closed. In that case, one looks for $`\rho _{T8}`$ as a dijet resonance: $`\overline{q}q,ggg\rho _{T8}g\overline{q}q,gg`$<sup>15</sup><sup>15</sup>15The decays $`\rho _{T8}g\pi _{T8}`$ and $`g\pi _{T1}`$ may occur and deplete the $`\rho _{T8}\mathrm{dijets}`$ rate. These modes were not expected to be important in Ref. and they were not taken into account in the CDF analyses. They are included in the studies in Ref. (TCSM–2) discussed below. Searches have been made by CDF for both untagged and $`b`$–tagged dijet resonances. The latter mode has a better signal–to–background ratio, but the rates and $`b`$–identification efficiencies in Run I were not high enough to make this advantage significant; the best limits come from untagged–dijets. The results of such a search are shown in Fig. 19. The region $`260<M_{\rho _{T8}}<460\mathrm{GeV}`$ is excluded at the 95% confidence level. This is a stringent constraint, but its applicability to TC2 models is uncertain. Finally, and very briefly, we turn to the effect of topcolor–assisted technicolor on experimental studies of the color–nonsinglet sector . As I mentioned, the simplest implementation of TC2 models requires two color $`SU(3)`$ groups, one strongly–coupled at 1 TeV for the third generation quarks $`(t,b)`$ and one weakly–coupled for the two light generations. These two color groups must be broken down to the diagonal $`SU(3)`$ near 1 TeV, and this remaining symmetry is identified with ordinary color. The most economical way I know to achieve this is to have the two technifermion doublets $`T_1=(U_1,D_1)(N_{TC},3,1)`$ and $`T_2=(U_2,D_2)(N_{TC},1,3)`$ condense with each other to achieve the desired breaking to $`SU(3)_C`$ . The main phenomenological consequence of this scenario is that the $`SU(3)`$ gluons mix with the $`SU(3)_1`$ octet of massive colorons, $`V_8`$, and with four color–octet technirhos, $`\rho _{ij}\overline{T}_i\lambda _AT_j`$ ($`i,j=1,2`$. The colorons decay strongly to top and bottom quarks and weakly to the light quarks . Alternatively, in the flavor–universal variant of TC2 , the colorons decay with equal strength to all quark flavors. In Ref. , we assume for simplicity that all $`\rho _{ij}`$ are too light to decay to pairs of technipions.<sup>16</sup><sup>16</sup>16The colored technipion sector of a TC2 model is bound to be very rich. Thus, it is not clear how the limits on leptoquarks discussed above are to be interpreted. This is work for the future. Then, they decay (via gluon and coloron dominance) into $`\overline{q}q`$ and $`gg`$ dijets and into $`g\pi _{T8}`$ and $`g\pi _{T1}`$. Even this simplified minimal TC2 version of the TCSM has a much richer set of dijet spectra and other hadron collider signals than the one–family model discussed above . We are just beginning to study it. Some preliminary examples of dijet production based on the assumptions of TCSM–2 are shown in Figs. 20 and 21 for $`\sqrt{s}=2\mathrm{TeV}`$ at the Tevatron. In both figures the coloron mass is 1.2 TeV while the input $`\rho _{T8}`$ masses range from 350 to 500 GeV. <sup>17</sup><sup>17</sup>17The pole masses are shifted somewhat from these input values by mixing effects. Figure 20 shows $`\overline{b}b`$ production with a strong resonance at 300 GeV (i,e., below $`\overline{t}t`$ threshold). Figure 21 shows $`\overline{t}t`$ production with roughly a factor of two enhancement of the total cross section over that predicted in the standard model. Both signals are ruled out by Run I measurements of the $`\overline{b}b`$ and $`\overline{t}t`$ cross sections . Many more studies of both the color–singlet and nonsinglet sectors of the TCSM need to be carried out. The Fermilab Workshop on Strong Dynamics at Run II will begin these studies this autumn, in time to be of use when the run starts in Spring 2001. The CDF and DØ collaborations will carry out detector–specific simulations in the next year or two. More detailed and more incisive $`e^+e^{}\rho _T,\omega _T`$ studies will come from the LEP experiments this year. The ATLAS and CMS collaborations likewise ought to study a broad range of signals for strong dynamics before they begin their runs later in the decade. ## 6. Open Problems My main goal in these lectures has been to attract some bright young people to the dynamical approach to electroweak and flavor symmetry breaking. Many difficult problems remain open for study there. These lectures provide a basis for starting to tackle them. All that’s needed now are new ideas, new data, and good luck. Here are the problems that vex me: 1. First and foremost, we need a reasonably realistic model of extended technicolor, or any other natural, dynamical description of flavor. To repeat: This is the hardest problem we face in particle physics. It deserves much more effort. I think the difficulty of this problem and the lack of a “standard model” of flavor are what have led to ETC’s being in such disfavor. Experiments will be of great help, possibly inspiring the right new ideas. Certainly, experiments that will be done in this decade will rule out, or vindicate, the ideas outlined in these lectures. That is an exciting prospect! 2. More tractable is the problem of constructing a dynamical theory of the top–quark mass that is natural, i.e., requires no fine–tuning of parameters, and has no nearby Landau pole. Like topcolor–assisted technicolor and top–seesaw models, such a theory is bound to have testable consequences below 2–3 TeV. So hurry—before the experiments get done! 3. Neutrino masses are at least as difficult a problem as the top mass. In particular, it is a great puzzle how ETC interactions could produce $`m_\nu <10^7m_e`$. It seems unnatural to have to assume an extra large ETC mass scale just for the neutrinos. Practically no thought has been has been given to this problem. Is there some simple way to tinker with the basic ETC mass–generating mechanism, some way to implement a seesaw mechanism, or must the whole ETC idea be scrapped? The area is wide open. 4. My favorite problem is “vacuum alignment” and CP violation . The basic idea is this: Spontaneous chiral symmetry breaking implies the existence of infinitely many degenerate ground states. These are manifested by the presence of massless Goldstone bosons (technipions). The “correct” ground state, i.e., the one on which consistent chiral perturbation theory for the technipions is to be carried out, is the one which minimizes the vacuum expectation value of the explicit chiral symmetry breaking Hamiltonian $`^{}`$ generated by ETC. As Dashen discovered, it is possible that an $`^{}`$ that appears to conserve CP actually violates it in the correct ground state. This provides a beautiful dynamical mechanism for the CP violation we observe. Or it could lead to disaster—strong CP violation, with a neutron electric dipole moment ten orders of magnitude larger than its upper limit. This field of research is just beginning in earnest. If the strong–CP problem can be controlled (there is reason to hope that it can be!), there are bound to be new sources of CP violation that are accessible to experiment. ## Acknowledgements I am grateful to Giulia Pancheri for inviting me to lecture at the 2000 Frascati Spring School. I owe much of my pleasure to the School’s enthusiastic students and their openness to the subversive science I told them about. I give great thanks to the Frascati Spring School secretariat, M. Legramante and A. Mantella, for many large and small assistances. I thank Yogi Srivastava for his invitation to visit the beautiful town and University of Perugia and the opportunity to speak there. Lia’s and Yogi’s generous support and hospitality leave me in their debt. I thank my colleagues at BU, members of the Fermilab Run II Strong Dynamics Workshop, and others for their help and constructive comments. In particular, I thank Georges Azuelos, Sekhar Chivukula, Estia Eichten, Andre Kounine, Greg Landsberg, Richard Haas, Takanobu Handa, Robert Harris, Kaori Maeshima, Meenakshi Narain, Steve Mrenna, Stephen Parke, Tonguç Rador, Francois Richard, Elizabeth Simmons, and John Womersley. This research was supported in part by the U. S. Department of Energy under Grant No. DE–FG02–91ER40676.
warning/0007/hep-ph0007167.html	ar5iv	text	# 1 Introduction ## 1 Introduction Nearly everyone who has made a serious study of quantum field theory in the past half century is familiar with Schwinger’s treatment of quantum electrodynamics (QED) in the presence of a constant electromagnetic field . A complicating feature of this background is that particle production occurs when the electric field is nonzero, so no state can be stable even in the free theory. Schwinger was well aware of this feature and took pains to work around it by assuming, where ever necessary, that the electromagnetic field was purely magnetic. He first computed the electron propagator then used it to evaluate what would later be known as the in-out effective action at one loop. He computed the rate of particle production for the unstable case of an electric field by inferring the imaginary part of the effective action through an analytic continuation from the stable case of a magnetic field. Schwinger’s analysis is rightly regarded as one of the great achievements of quantum field theory. However, its generality gives rise to a question. Schwinger actually obtained a form for the electron propagator for any constant electromagnetic field, magnetic or electric. What would result from using his electron propagator to calculate the one loop expectation value of the current operator for the unstable case of a constant electric field? This is how one might begin computing the back-reaction of particle production on the electric field. The result turns out to be zero. Of course this does not mean that there is no particle creation at one loop! Rather it implies that, for the case of a non-zero electric field, Schwinger’s “propagator” is really only a Green’s function and not the expectation value of the time-ordered product of $`\psi (x)\overline{\psi }(x^{})`$ in the presence of any fixed state. Schwinger never said otherwise, and it was presumably to avoid this problem that he employed the circuitous analytic continuation procedure. The distinction between propagators and Green’s functions can be understood most clearly in the context of the one dimensional harmonic oscillator with mass $`m`$ and frequency $`\omega `$. We can use the Heisenberg operator equations to express the position operator $`q(t)`$ in terms of its initial values $`q_0`$ and $`\dot{q}_0=p_0/m`$, $$q(t)=q_0\mathrm{cos}(\omega t)+\frac{p_0}{m\omega }\mathrm{sin}(\omega t).$$ (1) The time-ordered product of two such operators can be expressed in terms of the $`\mathrm{C}\text{ }\text{t}`$-number commutator function and the operator anti-commutator, $`T\{q(t)q(t^{})\}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\mathrm{sgn}(\mathrm{\Delta }t)[q(t),q(t^{})]+{\displaystyle \frac{1}{2}}\{q(t),q(t^{})\},`$ (2) $`=`$ $`{\displaystyle \frac{i}{2m\omega }}\mathrm{sin}(m\|\mathrm{\Delta }t\|)+{\displaystyle \frac{1}{2}}\{q(t),q(t^{})\},`$ (3) where $`\mathrm{\Delta }ttt^{}`$. A general propagator is the expectation value of this in the presence of some normalized state, $`i\mathrm{\Delta }(t;t^{})`$ $``$ $`S\|T\{q(t)q(t^{})\}\|S,`$ (5) $`=`$ $`{\displaystyle \frac{i}{2m\omega }}\mathrm{sin}(\omega \|\mathrm{\Delta }t\|)+\alpha \mathrm{cos}(\omega t)\mathrm{cos}(\omega t^{})`$ $`+\beta \mathrm{sin}[\omega (t+t^{})]+\gamma \mathrm{sin}(\omega t)\mathrm{sin}(\omega t^{}).`$ We have expressed the result in terms of three real numbers $`\alpha `$, $`\beta `$ and $`\gamma `$ defined as follows: $$\alpha S\|q_0^2\|S,\beta \frac{1}{2m\omega }S\|q_0p_0+p_0q_0\|S,\gamma \frac{1}{m^2\omega ^2}S\|p_0^2\|S.$$ (6) The key point is that the uncertainty principle imposes an inequality on $`\alpha `$ and $`\gamma `$, $$\alpha \gamma \left(\frac{1}{2m\omega }\right)^2.$$ (7) But $`\alpha `$, $`\beta `$ and $`\gamma `$ can be any $`\mathrm{C}\text{ }\text{t}`$-numbers if one only requires the Green’s function equation, $$m\left(\frac{d^2}{dt^2}+\omega ^2\right)i\mathrm{\Delta }(t;t^{})=i\delta (tt^{}),$$ (8) and symmetry under interchange of $`t`$ and $`t^{}`$. For example, $`\alpha =\beta =\gamma =0`$ gives a Green’s function but not a propagator. This paper contains six sections of which this introduction is the first. In Section 2 we show that Schwinger’s propagator gives zero for the expectation value of the current operator at one loop. To see that Schwinger’s “propagator” is really only a Green’s function it is useful to first express it in a diagonal function and spinor basis. This is the work of Section 3. In Section 4 we exhibit an initial value solution for the Heisenberg field operators analogous to the one presented above for the harmonic oscillator. We then show that there is no fixed state which gives Schwinger’s result. The manner in which it fails also explains the zero current result, and incidentally provides what is probably the simplest picture we shall ever get of particle production. In Section 5 we work out a true propagator in the presence of a natural state. It is significant that we can actually do this for a class of backgrounds which is general enough to include the actual electric field as it evolves under the influence of quantum electrodynamic back-reaction. Our conclusions comprise Section 6. ## 2 Zero current with Schwinger’s propagator We begin by representing the propagator as Schwinger did, as the expectation value of a first quantized operator, $`iS(x;x^{})`$ $``$ $`x\left\|{\displaystyle \frac{i}{\text{P}/e\text{A}/(X)m+iϵ}}\right\|x^{},`$ (9) $`=`$ $`x\left\|[\text{P}/e\text{A}/+m]{\displaystyle _0^{\mathrm{}}}dse^{is[(PeA)^2\frac{1}{2}eF_{\mu \nu }\sigma ^{\mu \nu }m^2+iϵ]}\right\|x^{}.`$ (10) The position and momentum operators of the first quantized theory are $`X^\mu `$ and $`P^\nu `$, respectively. We assign the standard meaning to $`\sigma ^{\mu \nu }\frac{i}{2}[\gamma ^\mu ,\gamma ^\nu ]`$, and we will assume that the vector potential $`A_\mu (X)`$ is a liner function of $`X^\nu `$. We also follow Schwinger in using the proper time method to regulate expressions involving the propagator, $`iS(x;x^{})=\underset{s_00^+}{lim}{\displaystyle _{s_0}^{\mathrm{}}}𝑑se^{is(m^2iϵ)}`$ (11) $`\times x\left\|[\text{P}/e\text{A}/+m]e^{\frac{i}{2}eF_{\mu \nu }\sigma ^{\mu \nu }}e^{is[(PeA)^2}\right\|x^{}.`$ Note that if the electric field has magnitude $`E`$ then the exponential of the matrix term is, $$e^{\frac{i}{2}eF_{\mu \nu }\sigma ^{\mu \nu }}=\mathrm{cosh}(eEs)I+\frac{1}{2E}\mathrm{sinh}(eEs)F_{\mu \nu }\gamma ^\mu \gamma ^\nu .$$ (12) At one loop order the expectation value of the current operator is just the trace of $`e\gamma ^\mu `$ up against the coincident propagator, $`\mathrm{\Omega }\left\|J^\mu (x)\right\|\mathrm{\Omega }=\mathrm{Tr}\left[e\gamma ^\mu iS(x;x)\right],`$ (14) $`=`$ $`\underset{s_00^+}{lim}4e{\displaystyle _{s_0}^{\mathrm{}}}dse^{is(m^2iϵ)}\{\mathrm{cosh}(eEs)x\left\|[P^\mu eA^\mu ]e^{is(PeA)^2}\right\|x`$ $`{\displaystyle \frac{1}{E}}\mathrm{sinh}(eEs)F^{\mu \nu }x\left\|[P_\nu eA_\nu ]e^{is(PeA)^2}\right\|x\}.`$ Note that the coincidence limit is regulated by the factor of $`e^{is(PeA)^2}`$ as long as $`s0`$. Note also that we can get the factor of $`P^\mu eA^\mu `$ by commutation, $$P^\mu eA^\mu =\frac{i}{2}[X^\mu ,(PeA)^2].$$ (15) The commutator of $`X^\mu `$ with the exponential gives $`is`$ times a special operator ordering of the product of the exponential times with this factor. Because the first quantized bra and ket are the same the original ordering can be restored, $$x\left\|[X^\mu ,e^{is(PeA)^2}]\right\|x=2sx\left\|(P^\mu eA^\mu )e^{is(PeA)^2}\right\|x.$$ (16) But the commutator vanishes for the same reason. So the expectation value of the current operator computed using Schwinger’s propagator vanishes. ## 3 Going to lightcone momentum space For definiteness we assume the electric field has magnitude $`E`$ and is directed along the positive $`z`$ axis. If we define the separation 4-vector as $`\mathrm{\Delta }x^\mu x^\mu x^\mu `$ then Schwinger’s result for the electron propagator is $`iS(x;x^{})={\displaystyle \frac{i}{32\pi ^2}}\mathrm{exp}\left[ie{\displaystyle _x^{}^x}𝑑\xi ^\mu A_\mu (\xi )\right]{\displaystyle _0^{\mathrm{}}}𝑑se^{is(m^2iϵ)}`$ (17) $`\times \mathrm{exp}\left[{\displaystyle \frac{i}{4s}}\left(\mathrm{\Delta }x^2+\mathrm{\Delta }y^2+eEs\mathrm{coth}(eEs)[\mathrm{\Delta }z^2\mathrm{\Delta }t^2]\right)\right]`$ $`\times \{[{\displaystyle \frac{2eE}{s}}m+{\displaystyle \frac{eE}{s^2}}(\gamma ^1\mathrm{\Delta }x+\gamma ^2\mathrm{\Delta }y)][\mathrm{coth}(eEs)+\gamma ^0\gamma ^3]`$ $`+{\displaystyle \frac{e^2E^2}{s}}\mathrm{csch}^2(eEs)[\gamma ^3\mathrm{\Delta }z\gamma ^0\mathrm{\Delta }t]\}.`$ The special role assumed by the $`0`$ and $`3`$ directions strongly suggests that $`iS(x;x^{})`$ should be expressed in terms of lightcone coordinates. The fact that translation invariance is broken only by the exponential of the line integral of the vector potential also suggests that the propagator should be transformed to momentum space. We define the lightcone coordinates and gamma matrices as follows: $$x_\pm \frac{1}{\sqrt{2}}(x^0\pm x^3),\gamma _\pm \frac{1}{\sqrt{2}}(\gamma ^0\pm \gamma ^3).$$ (18) The other (“transverse”) components of $`x^\mu `$ and $`\gamma ^\mu `$ comprise the 2-vectors $`\stackrel{~}{x}`$ and $`\stackrel{~}{\gamma }`$, and the invariant contraction decomposes as follows, $$\gamma ^\mu x_\mu =\gamma ^0x^0\gamma ^3x^3\stackrel{~}{\gamma }\stackrel{~}{x}=\gamma _+x_{}+\gamma _{}x_+\stackrel{~}{\gamma }\stackrel{~}{x}.$$ (19) Note that $`(\gamma _\pm )^2=0`$. We follow Kogut and Soper in defining lightcone spinor projection operators, $$P_\pm \frac{1}{2}\left(I\pm \gamma ^0\gamma ^3\right)=\frac{1}{2}\gamma _{}\gamma _\pm .$$ (20) With these conventions the propagator takes the form, $`iS(x;x^{})`$ $`=`$ $`{\displaystyle \frac{ie^2E^2}{32\pi ^2}}\mathrm{exp}\left[ie{\displaystyle _x^{}^x}𝑑\xi ^\mu A_\mu (\xi )\right]{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{ds}{s}}\mathrm{csch}^2(eEs)e^{is(m^2iϵ)}`$ (21) $`\times \mathrm{exp}\left[{\displaystyle \frac{i}{4s}}\mathrm{\Delta }\stackrel{~}{x}\mathrm{\Delta }\stackrel{~}{x}{\displaystyle \frac{i}{2}}eE\mathrm{coth}(eEs)\mathrm{\Delta }x_+\mathrm{\Delta }x_{}\right]`$ $`\times \{[{\displaystyle \frac{m}{eE}}{\displaystyle \frac{\stackrel{~}{\gamma }\mathrm{\Delta }\stackrel{~}{x}}{2eEs}}]P_+(e^{2eEs}1)+\gamma _+\mathrm{\Delta }x_{}`$ $`+\gamma _{}\mathrm{\Delta }x_++[{\displaystyle \frac{m}{eE}}{\displaystyle \frac{\stackrel{~}{\gamma }\mathrm{\Delta }\stackrel{~}{x}}{2eEs}}]P_{}(1e^{2eEs})\}.`$ We define the lightcone components of the vector potential $`A_\mu `$ as $$A_\pm \frac{1}{\sqrt{2}}\left(A_0\pm A_3\right).$$ (22) Our gauge condition is $`A_+=0`$, and to get $`F^{30}=E`$ we take $`A_{}=Ex_+`$ with $`\stackrel{~}{A}=0`$. It is now possible to compute the initial phase. The path is $`\xi ^\mu (\tau )=x^\mu +\mathrm{\Delta }x^\mu \tau `$ so we get $`ie{\displaystyle _x^{}^x}𝑑\xi ^\mu A_\mu (\xi )`$ $`=`$ $`ieE\mathrm{\Delta }x_{}{\displaystyle _0^1}𝑑\tau \left(x_+^{}+\mathrm{\Delta }x_+\tau \right),`$ (23) $`=`$ $`{\displaystyle \frac{i}{2}}eE\mathrm{\Delta }x_{}\left(x_++x_+^{}\right).`$ (24) The propagator is invariant under translations of $`x_{}`$ and $`\stackrel{~}{x}`$ so those are the variables on which we shall Fourier transform, $`i\stackrel{~}{S}(x_+,x_+^{};k_+\stackrel{~}{k}){\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\mathrm{\Delta }x_{}e^{ik_+\mathrm{\Delta }x_{}}{\displaystyle d^2\mathrm{\Delta }\stackrel{~}{x}e^{i\stackrel{~}{k}\mathrm{\Delta }\stackrel{~}{x}}iS(x;x^{})},`$ (25) $`={\displaystyle \frac{e^2E^2}{4}}{\displaystyle _0^{\mathrm{}}}ds\mathrm{csch}^2(eEs)e^{is(\stackrel{~}{\omega }^2iϵ)}\{\left[{\displaystyle \frac{m\stackrel{~}{\gamma }\stackrel{~}{k}}{eE}}\right]P_+(e^{2eEs}1)`$ $`i\gamma _+{\displaystyle \frac{}{k_+}}+\gamma _{}\mathrm{\Delta }x_++\left[{\displaystyle \frac{m\stackrel{~}{\gamma }\stackrel{~}{k}}{eE}}\right]P_{}(1e^{2eEs})\}`$ $`\times \delta \left(k_++{\displaystyle \frac{1}{2}}eE(x_++x_+^{}){\displaystyle \frac{1}{2}}eE\mathrm{coth}(eEs)\mathrm{\Delta }x_+\right),`$ (26) where $`\stackrel{~}{\omega }^2=m^2+\stackrel{~}{k}\stackrel{~}{k}`$. The delta function can be recast to determine $`s`$, $$\delta \left(k_++\frac{1}{2}eE(x_++x_+^{})\frac{1}{2}eE\mathrm{coth}(eEs)\mathrm{\Delta }x_+\right)=\frac{\delta \left(s\frac{1}{2eE}\mathrm{ln}\left[\frac{k_++eEx_+}{k_++eEx_+^{}}\right]\right)}{\frac{1}{2}e^2E^2\mathrm{csch}^2(eEs)\|\mathrm{\Delta }x_+\|}.$$ (27) This brings the propagator to the form, $`i\stackrel{~}{S}(x_+,x_+^{};k_+\stackrel{~}{k})={\displaystyle \frac{1}{2\|\mathrm{\Delta }x_+\|}}\{\left[{\displaystyle \frac{m\stackrel{~}{\gamma }\stackrel{~}{k}}{eE}}\right]\left({\displaystyle \frac{eE\mathrm{\Delta }x_+}{k_++eEx_+^{}}}\right)P_+`$ $`\gamma _+{\displaystyle \frac{}{k_+}}+\gamma _{}\mathrm{\Delta }x_++\left[{\displaystyle \frac{m\stackrel{~}{\gamma }\stackrel{~}{k}}{eE}}\right]\left({\displaystyle \frac{eE\mathrm{\Delta }x_+}{k_++eEx_+}}\right)P_{}\}`$ $`\times {\displaystyle _0^{\mathrm{}}}dse^{i(\stackrel{~}{\omega }^2iϵ)}\delta (s{\displaystyle \frac{1}{2eE}}\mathrm{ln}\left[{\displaystyle \frac{k_++eEx_+}{k_++eEx_+^{}}}\right]).`$ (28) At this stage it becomes crucial to recall that the electron charge is negative so $`eE<0`$. The delta function can only become singular for $`s>0`$ if $`k_+>eEx_+`$ for $`\mathrm{\Delta }x_+>0`$, or if $`k_+>eEx_+`$ for $`\mathrm{\Delta }x_+<0`$. The integration over $`s`$ therefore gives, $$\left[\frac{k_++eEx_+}{k_++eEx_+^{}}\right]^{\frac{i(\stackrel{~}{\omega }^2iϵ)}{2eE}}\left\{\theta (\mathrm{\Delta }x_+)\theta \left(k_++eEx_+\right)+\theta (\mathrm{\Delta }x_+)\theta \left(k_+eEx_+\right)\right\}.$$ (29) Note that the $`iϵ`$ in the exponent means that the term in square brackets is raised to a power with a small positive real part. This has the important consequence that multiplying by $`\delta (k_++eEx_+)`$ gives zero, so we need not worry about the $`i/k_+`$ acting on the theta functions. Acting on the power it gives, $$i\frac{}{k_+}\left[\frac{k_++eEx_+}{k_++eEx_+^{}}\right]^{\frac{i\stackrel{~}{\omega }^2}{2eE}}=\frac{\frac{1}{2}\stackrel{~}{\omega }^2\mathrm{\Delta }x_+}{(k_++eEx_+)(k_++eEx_+^{})}\left[\frac{k_++eEx_+}{k_++eEx_+^{}}\right]^{\frac{i\stackrel{~}{\omega }^2}{2eE}}.$$ (30) The final result for the propagator is, $`i\stackrel{~}{S}(x_+,x_+^{};k_+,\stackrel{~}{k})=\mathrm{sgn}(\mathrm{\Delta }x_+)\theta \left[\mathrm{sgn}(\mathrm{\Delta }x_+)(k_++eEx_+)\right]`$ (31) $`\times \left[{\displaystyle \frac{k_++eEx_+}{k_++eEx_+^{}}}\right]^{\frac{i\stackrel{~}{\omega }^2}{2eE}}{\displaystyle \frac{1}{2}}\{\left({\displaystyle \frac{m\stackrel{~}{\gamma }\stackrel{~}{k}}{k_++eEx_+^{}}}\right)P_+`$ $`+{\displaystyle \frac{\frac{1}{2}\stackrel{~}{\omega }^2\gamma _+}{(k_++eEx_+)(k_++eEx_+^{})}}+\gamma _{}+\left({\displaystyle \frac{m\stackrel{~}{\gamma }\stackrel{~}{k}}{k_++eEx_+}}\right)P_{}\}.`$ It is more illuminating for the work of the next section to right-multiply this by $`\gamma ^0=(\gamma _++\gamma _{})/\sqrt{2}`$ and slightly re-arrange the order in which the four spinor matrices appear, $`i\stackrel{~}{S}(x_+,x_+^{};k_+,\stackrel{~}{k})\gamma ^0=\mathrm{sgn}(\mathrm{\Delta }x_+)\theta \left[\mathrm{sgn}(\mathrm{\Delta }x_+)(k_++eEx_+)\right]`$ (32) $`\times \left[{\displaystyle \frac{k_++eEx_+}{k_++eEx_+^{}}}\right]^{\frac{i\stackrel{~}{\omega }^2}{2eE}}\{{\displaystyle \frac{1}{\sqrt{2}}}P_++{\displaystyle \frac{1}{\sqrt{2}}}P_+{\displaystyle \frac{1}{2}}\gamma _{}\left({\displaystyle \frac{m+\stackrel{~}{\gamma }\stackrel{~}{k}}{k_++eEx_+^{}}}\right)`$ $`+\left({\displaystyle \frac{m\stackrel{~}{\gamma }\stackrel{~}{k}}{k_++eEx_+}}\right){\displaystyle \frac{1}{2}}\gamma _+{\displaystyle \frac{1}{\sqrt{2}}}P_++{\displaystyle \frac{\frac{1}{2}\stackrel{~}{\omega }^2\frac{1}{\sqrt{2}}P_{}}{(k_++eEx_+)(k_++eEx_+^{})}}\}.`$ ## 4 Inferring the state The point of this section is to understand Schwinger’s propagator in terms of operators and states. Let us start with notation for the Fourier transform (on $`x_{}`$ and $`\stackrel{~}{x}`$) of the electron field operator $`\psi (x_+,x_{},\stackrel{~}{x})`$, $$\mathrm{\Psi }(x_+,k_+,\stackrel{~}{k})𝑑x_{}e^{ik_+x_{}}d^2\stackrel{~}{x}e^{i\stackrel{~}{k}\stackrel{~}{x}}\psi (x_+,x_{},\stackrel{~}{x}).$$ (33) This section is about finding a state $`\|S`$ such that, $`i\stackrel{~}{S}(x_+,x_+^{};k_+,\stackrel{~}{k})\gamma ^0(2\pi )^3\delta (k_+q_+)\delta ^2(\stackrel{~}{k}\stackrel{~}{q})`$ (35) $`=\theta (\mathrm{\Delta }x_+)S\left\|\mathrm{\Psi }(x_+,k_+,\stackrel{~}{k})\mathrm{\Psi }^{}(x_+^{},q_+,\stackrel{~}{q})\right\|S`$ $`\theta (\mathrm{\Delta }x_+)S\left\|\mathrm{\Psi }^{}(x_+^{},q_+,\stackrel{~}{q})\mathrm{\Psi }(x_+,k_+,\stackrel{~}{k})\right\|S,`$ $`S\left\|X_+\left\{\mathrm{\Psi }(x_+,k_+,\stackrel{~}{k})\mathrm{\Psi }^{}(x_+^{},q_+,\stackrel{~}{q})\right\}\right\|S.`$ Note that the adjoint is taken after Fourier transforming. We define the $`\pm `$ components of the electron field (and its Fourier transform) using the $`P_\pm `$ projectors, $$\psi _\pm (x_+,x_{},\stackrel{~}{x})P_\pm \psi (x_+,x_{},\stackrel{~}{x}).$$ (36) By acting $`P_+`$ on the left and right of expression (32) it is easy to see that the state $`\|S`$ must obey $`S\left\|X_+\left\{\mathrm{\Psi }_+(x_+,k_+,\stackrel{~}{k})\mathrm{\Psi }_+^{}(x_+^{},q_+,\stackrel{~}{q})\right\}\right\|S=(2\pi )^3\delta (k_+q_+)\delta ^2(\stackrel{~}{k}\stackrel{~}{q})`$ (37) $`\times \mathrm{sgn}(\mathrm{\Delta }x_+)\theta \left[\mathrm{sgn}(\mathrm{\Delta }x_+)(k_++eEx_+)\right]\left[{\displaystyle \frac{k_++eEx_+}{k_++eEx_+^{}}}\right]^{\frac{i\stackrel{~}{\omega }^2}{2eE}}\times {\displaystyle \frac{1}{\sqrt{2}}}P_+.`$ The other components can be recognized similarly, $`S\left\|X_+\left\{\mathrm{\Psi }_+\mathrm{\Psi }_{}^{}\right\}\right\|S`$ $`=`$ $`\mathrm{same}\times {\displaystyle \frac{1}{\sqrt{2}}}P_+{\displaystyle \frac{1}{2}}\gamma _{}\left({\displaystyle \frac{m+\stackrel{~}{\gamma }\stackrel{~}{k}}{k_++eEx_+^{}}}\right),`$ (38) $`S\left\|X_+\left\{\mathrm{\Psi }_{}\mathrm{\Psi }_+^{}\right\}\right\|S`$ $`=`$ $`\mathrm{same}\times \left({\displaystyle \frac{m\stackrel{~}{\gamma }\stackrel{~}{k}}{k_++eEx_+^{}}}\right){\displaystyle \frac{1}{2}}\gamma _+{\displaystyle \frac{1}{\sqrt{2}}}P_+,`$ (39) $`S\left\|X_+\left\{\mathrm{\Psi }_{}\mathrm{\Psi }_{}^{}\right\}\right\|S`$ $`=`$ $`\mathrm{same}\times {\displaystyle \frac{\frac{1}{2}\stackrel{~}{\omega }^2\frac{1}{\sqrt{2}}P_{}}{(k_++eEx_+)(k_++eEx_+^{})}}.`$ (40) Of relations (37-40) only the first is really independent, the others follow from the equations of motion. To see this consider the Dirac equation for our vector potential, $$\left(\gamma ^\mu i_\mu \gamma ^\mu eA_\mu m\right)\psi =\left(\gamma _+i_++\gamma _{}(i_{}+eEx_+)+\stackrel{~}{\gamma }i\stackrel{~}{}m\right)\psi .$$ (41) Multiplication alternately with $`\gamma _{}`$ and $`\gamma _+`$ gives $`i_+\psi _+(x_+,x_{},\stackrel{~}{x})`$ $`=`$ $`\left(m+\stackrel{~}{\gamma }i\stackrel{~}{}\right){\displaystyle \frac{1}{2}}\gamma _{}\psi _{}(x_+,x_{},\stackrel{~}{x}),`$ (42) $`\left({\displaystyle \frac{}{}}i_{}+eEx_+\right)\psi _{}(x_+,x_{},\stackrel{~}{x})`$ $`=`$ $`\left(m+\stackrel{~}{\gamma }i\stackrel{~}{}\right){\displaystyle \frac{1}{2}}\gamma _+\psi _+(x_+,x_{},\stackrel{~}{x}).`$ (43) Fourier transforming (42) and multiplying by $`(m\stackrel{~}{\gamma }\stackrel{~}{k})\gamma _+/\stackrel{~}{\omega }^2`$ gives, $$\mathrm{\Psi }_{}(x_+,k_+,\stackrel{~}{k})=\left(\frac{m\stackrel{~}{\gamma }\stackrel{~}{k}}{\stackrel{~}{\omega }^2}\right)\gamma _+i_+\mathrm{\Psi }_+(x_+,k_+,\stackrel{~}{k}).$$ (44) It follows that $`\mathrm{\Psi }_{}`$ can be eliminated in favor of $`\mathrm{\Psi }_+`$, $`S\left\|X_+\left\{\mathrm{\Psi }_+\mathrm{\Psi }_{}^{}\right\}\right\|S`$ $`=`$ $`S\left\|X_+\left\{\mathrm{\Psi }_+(i_+^{})\mathrm{\Psi }_+^{}\right\}\right\|S\gamma _{}\left({\displaystyle \frac{m+\stackrel{~}{\gamma }\stackrel{~}{q}}{\stackrel{~}{\omega }^2}}\right),`$ (45) $`S\left\|X_+\left\{\mathrm{\Psi }_{}\mathrm{\Psi }_+^{}\right\}\right\|S`$ $`=`$ $`\left({\displaystyle \frac{m\stackrel{~}{\gamma }\stackrel{~}{k}}{\stackrel{~}{\omega }^2}}\right)\gamma _+S\left\|X_+\left\{i_+\mathrm{\Psi }_+\mathrm{\Psi }_+^{}\right\}\right\|S,`$ (46) $`S\left\|X_+\left\{\mathrm{\Psi }_{}\mathrm{\Psi }_{}^{}\right\}\right\|S`$ $`=`$ $`\left({\displaystyle \frac{m\stackrel{~}{\gamma }\stackrel{~}{k}}{\stackrel{~}{\omega }^2}}\right)\gamma _+S\left\|X_+\left\{_+\mathrm{\Psi }_+_+^{}\mathrm{\Psi }_+^{}\right\}\right\|S`$ (47) $`\times \gamma _{}\left({\displaystyle \frac{m+\stackrel{~}{\gamma }\stackrel{~}{q}}{\stackrel{~}{\omega }^2}}\right).`$ The derivatives with respect to $`x_+`$ and $`x_+^{}`$ can be pulled outside the $`x_+`$–ordering symbol if we agree that they do not act on $`\theta (\pm \mathrm{\Delta }x_+)`$. The delta functions in (37) set $`q_+=k_+`$ and $`\stackrel{~}{q}=\stackrel{~}{k}`$ and the derivatives give, $`i_+\left[{\displaystyle \frac{k_++eEx_+}{k_++eEx_+^{}}}\right]^{\frac{i\stackrel{~}{\omega }^2}{2eE}}`$ $`=`$ $`{\displaystyle \frac{\frac{1}{2}\stackrel{~}{\omega }^2}{k_++eEx_+}}\left[{\displaystyle \frac{k_++eEx_+}{k_++eEx_+^{}}}\right]^{\frac{i\stackrel{~}{\omega }^2}{2eE}},`$ (48) $`i_+^{}\left[{\displaystyle \frac{k_++eEx_+}{k_++eEx_+^{}}}\right]^{\frac{i\stackrel{~}{\omega }^2}{2eE}}`$ $`=`$ $`{\displaystyle \frac{\frac{1}{2}\stackrel{~}{\omega }^2}{k_++eEx_+^{}}}\left[{\displaystyle \frac{k_++eEx_+}{k_++eEx_+^{}}}\right]^{\frac{i\stackrel{~}{\omega }^2}{2eE}}.`$ (49) So relations (38-40) indeed follow from (37). To identify the state $`\|S`$ we must express the operators as functions of $`(x_+,k_+,\stackrel{~}{k})`$ and functionals of the initial value operators. These initial value operators are the only true degrees of freedom of any Heisenberg operator. On the lightcone the “initial value surface” can be taken as $`x_+=0`$ with $`x_{}>L`$ and $`x_+>0`$ with $`x_{}=L`$. Taking $`L`$ to $`\mathrm{}`$ gives the following result for QED in this background , $`\mathrm{\Psi }_+(x_+,k_+,\stackrel{~}{k})=\left[{\displaystyle \frac{k_++eEx_++iϵ}{k_++iϵ}}\right]^{i\lambda }\mathrm{\Xi }_0(k_+,\stackrel{~}{k})`$ (50) $`\theta (k_+)\theta (eEx_+k_+)\left[{\displaystyle \frac{k_++eEx_++iϵ}{iϵ}}\right]^{i\lambda }{\displaystyle \frac{\sqrt{2\pi \lambda }}{\mathrm{\Gamma }(1+i\lambda )}}\mathrm{\Xi }_{\mathrm{}}(k_+,\stackrel{~}{k}),`$ where $`\lambda \stackrel{~}{\omega }^2/(2eE)`$ and the initial value operators are, $`\mathrm{\Xi }_0(k_+,\stackrel{~}{k})`$ $``$ $`\mathrm{\Psi }_+(0,k_+,\stackrel{~}{k}),`$ (51) $`\mathrm{\Xi }_{\mathrm{}}(k_+,\stackrel{~}{k})`$ $``$ $`\sqrt{{\displaystyle \frac{2\pi }{\lambda }}}\left({\displaystyle \frac{m\stackrel{~}{\gamma }\stackrel{~}{k}}{eE}}\right){\displaystyle \frac{i}{2}}\gamma _{}`$ (52) $`\times {\displaystyle }d^2\stackrel{~}{x}e^{i\stackrel{~}{k}\stackrel{~}{x}}\underset{x_{}\mathrm{}}{lim}e^{ik_+x_{}}\psi _{}({\displaystyle \frac{k_+}{eE}},x_{},\stackrel{~}{x}).`$ The initial value operators anti-commute canonically with their adjoints, $`\{\mathrm{\Xi }_0(k_+,\stackrel{~}{k}),\mathrm{\Xi }_0^{}(q_+,\stackrel{~}{q})\}`$ $`=`$ $`(2\pi )^3\delta (k_+q_+)\delta ^2(\stackrel{~}{k}\stackrel{~}{q}){\displaystyle \frac{1}{\sqrt{2}}}P_+,`$ (53) $`\{\mathrm{\Xi }_{\mathrm{}}(k_+,\stackrel{~}{k}),\mathrm{\Xi }_{\mathrm{}}^{}(q_+,\stackrel{~}{q})\}`$ $`=`$ $`(2\pi )^3\delta (k_+q_+)\delta ^2(\stackrel{~}{k}\stackrel{~}{q}){\displaystyle \frac{1}{\sqrt{2}}}P_+.`$ (54) The “$`0`$“ operators anti-commute with the “$`\mathrm{}`$” ones by causality since the two surfaces are spacelike related for all finite $`L`$. It is worth pausing at this point to comment on the significant features of our operator solution. First, note that for $`k_+>0`$ the mode functions experience a characteristic drop in amplitude as they evolve through the singular point at $`x_+=k_+/eE`$, $$\left\|\left[\frac{k_++eEx_++iϵ}{k_++iϵ}\right]^{i\lambda }\right\|=\{\begin{array}{cc}1\hfill & x_+<k_+/eE\hfill \\ e^{\pi \lambda }\hfill & x_+>k_+/eE\hfill \end{array}.$$ (55) The amplitude lost by $`\mathrm{\Xi }_0(k_+,\stackrel{~}{k})`$ passes to $`\mathrm{\Xi }_{\mathrm{}}(k_+,\stackrel{~}{k})`$ by virtue of the identity , $$\left\|\left[\frac{k_++eEx_++iϵ}{iϵ}\right]^{i\lambda }\frac{\sqrt{2\pi \lambda }}{\mathrm{\Gamma }(1+i\lambda )}\right\|=\sqrt{1e^{2\pi \lambda }},$$ (56) for $`0<k_+<eEx_+`$. The physical interpretation of the amplitude drop is particle production. This is a discrete and instantaneous event on the lightcone. As each mode passes through its singularity at $`x_+=k_+/eE`$ the eigenvalue of $`i_+`$ changes sign and the operator coefficient switches interpretation from annihilator to creator. Since Heisenberg states are fixed a state which was originally empty in that mode seems, after singularity, to have acquired a particle with probability equal to the square of the post-singular amplitude (55). The second point worthy of mention about (50) is that we cannot follow the usual lightcone practice, for nonzero mass and/or more than two spacetime dimensions, of ignoring the “$`\mathrm{}`$” operators. These are always technically present in lightcone field theory, but they usually remain segregated to sector at $`k_+=0`$. In computing scattering amplitudes one can ignore this sector and recover the zero momentum limit instead by analytic continuation. We cannot get away with this for QED in a constant electric field. The physical momentum is the minimally coupled one, $`p_+k_++eEx_+`$, which reaches the far infrared near singularity. At this point operators from the surface at $`x_{}=\mathrm{}`$ can and do mix with the “$`0`$” operators in (50). We should end this digression by noting that part of our operator solution (50) has been obtained in a different context by Srinivasan and Padmanabhan . The mode solution they give corresponds to the term proportional to $`\mathrm{\Xi }_0(k_+,\stackrel{~}{k})`$, although without our $`iϵ`$ convention. They do not get the part proportional to $`\mathrm{\Xi }_{\mathrm{}}(k_+,\stackrel{~}{k})`$. They employed a WKB approach to evolve around the singularity at $`k_++eEx_+=0`$, and they claim this results in a second term proportional to $`\mathrm{\Xi }_0(k_+,\stackrel{~}{k})`$. We have been unable to understand why the WKB technique applies to a first order evolution equation, nor can we reproduce their results. Our operator solution (50) implies the following anti-commutation relation for $`x_+x_+^{}`$: $$\{\mathrm{\Psi }_+(x_+,k_+\stackrel{~}{k}),\mathrm{\Psi }_+^{}(x_+^{},q_+,\stackrel{~}{q})\}=(2\pi )^3\delta ^3(kq)\left\|\frac{k_++eEx_+}{k_++eEx_+^{}}\right\|^{i\lambda }\frac{1}{\sqrt{2}}P_+.$$ (57) Comparison with (37) reveals that $`\|S`$ must obey, $`S\left\|\mathrm{\Psi }_+(x_+,k_+,\stackrel{~}{k}),\mathrm{\Psi }_+^{}(x_+^{},q_+,\stackrel{~}{q})\right\|S`$ $`=`$ $`\theta (k_++eEx_+)\{\mathrm{\Psi }_+,\mathrm{\Psi }_+^{}\},`$ (58) $`S\left\|\mathrm{\Psi }_+^{}(x_+^{},q_+,\stackrel{~}{q}),\mathrm{\Psi }_+(x_+,k_+,\stackrel{~}{k})\right\|S`$ $`=`$ $`\theta (k_+eEx_+)\{\mathrm{\Psi }_+,\mathrm{\Psi }_+^{}\}.`$ (59) The state $`\|S`$ must therefore be annihilated by $`\mathrm{\Psi }_+`$ for $`k_++eEx_+>0`$ and by $`\mathrm{\Psi }_+^{}`$ for $`k_++eEx_+<0`$. In terms of the fundamental operators this translates to, $`\mathrm{\Xi }_0(k_++eEx_+,\stackrel{~}{k})\|S=`$ $`0`$ $`=\mathrm{\Xi }_0^{}(k_+eEx_+,\stackrel{~}{k})\|S,`$ (60) $`\mathrm{\Xi }_{\mathrm{}}(k_++eEx_+,\stackrel{~}{k})\|S=`$ $`0`$ $`=\mathrm{\Xi }_{\mathrm{}}^{}(k_+eEx_+,\stackrel{~}{k})\|S,`$ (61) for all $`k_+>0`$. It is immediately apparent that $`\|S`$ is not a proper Heisenberg state in the sense of remaining fixed. We must rather use a different state for each value of $`x_+`$ in order to recover Schwinger’s result. Hence it is only a Green’s function and not a true propagator. We can also understand the curious result of Section 2 that Schwinger’s “propagator” gives a null result for the expectation vlaue of the current operator at one loop. Recall from our operator solution that the eigenvalue of $`i_+`$ on $`\mathrm{\Psi }_+(x_+,k_+,\stackrel{~}{k})`$ is negative for $`k_++eEx_+>0`$ and positive for $`k_++eEx_+<0`$. This means that at fixed $`x_+`$ the electron annihilators are proportional to $`\mathrm{\Psi }_+(x_+,k_+,\stackrel{~}{k})`$, for all $`k_++eEx_+>0`$, and the positron annihilators are proportional to $`\mathrm{\Psi }_+^{}(x_+,k_+,\stackrel{~}{k})`$, again for $`k_++eEx_+>0`$. Conditions (60-61) guarantee that these operators annihilate $`\|S`$, which means that the state is empty at $`x_+`$. This is no contradiction with the fact that particle production really happens in this background because it is always possible to find a state which is empty at one particular instant. If the state were held fixed one would see a nonzero (actually infinite) current at later $`x_+`$ , but that is not what Schwinger’s “propagator” does. As $`x_+`$ changes the state also changes to the one which happens to be instantaneously empty at the new value of $`x_+`$. So one of course sees zero current. End of mystery. ## 5 A true propagator We have been able to solve the Heisenberg operator equations for a vector potential $`A_{}(x_+)`$ which depends arbitrarily upon $`x_+`$.<sup>1</sup><sup>1</sup>1At certain points we do assume that $`eA_{}(x_+)`$ is an increasing function of $`x_+`$. This could be avoided at the cost of more complicated expressions. Since there is no significant simplification arising from the assumption that $`A_{}(x_+)=Ex_+`$, we begin by stating the more general solution. The derivation can be found elsewhere . Modes still undergo particle production when the minimally coupled momentum $`k_+eA_{}(x_+)`$ vanishes. We define this value of $`x_+`$ as $`X(k_+)`$, $$k_+=A_{}\left(X(k_+)\right).$$ (62) Since the electric field is no longer necessarily constant we must generalize the definition of $`\lambda `$, $$\lambda (k_+,\stackrel{~}{k})\frac{\stackrel{~}{\omega }^2}{2eA_{}^{}(X(k_+))}.$$ (63) The mode functions require a similar generalization, $$[A_{}](y_+,x_+;k_+,\stackrel{~}{k})\mathrm{exp}\left[\frac{i}{2}\stackrel{~}{\omega }^2_{y_+}^{x_+}\frac{du}{k_+eA_{}(u)+iϵ}\right].$$ (64) With these conventions the operator solution can be written as follows:<sup>2</sup><sup>2</sup>2We have actually quoted a simplified form in which a distributional limit was taken assuming that $`k_+`$ is somewhat displaced from the singular points at $`k_+=0`$ and $`k_+=eA_{}(x_+)`$. The more accurate expression must be used in taking derivatives with respect to $`x_+`$ . $`\mathrm{\Psi }_+(x_+,k_+,\stackrel{~}{k})=[A_{}](0,x_+;k_+,\stackrel{~}{k})\mathrm{\Xi }_0(k_+,\stackrel{~}{k})`$ (65) $`\theta (k_+)\theta (eA_{}(x_+)k_+)(X(k_+),x_+;k_+,\stackrel{~}{k}){\displaystyle \frac{\sqrt{2\pi \lambda }}{\mathrm{\Gamma }(1+i\lambda )}}\mathrm{\Xi }_{\mathrm{}}(k_+,\stackrel{~}{k}),`$ where the initial value operators are, $`\mathrm{\Xi }_0(k_+,\stackrel{~}{k})`$ $``$ $`\mathrm{\Psi }_+(0,k_+,\stackrel{~}{k}),`$ (66) $`\mathrm{\Xi }_{\mathrm{}}(k_+,\stackrel{~}{k})`$ $``$ $`\sqrt{{\displaystyle \frac{2\pi }{\lambda (k_+,\stackrel{~}{k})}}}\left({\displaystyle \frac{m\stackrel{~}{\gamma }\stackrel{~}{k}}{eA_{}^{}(X(k_+))}}\right){\displaystyle \frac{i}{2}}\gamma _{}`$ (67) $`\times {\displaystyle }d^2\stackrel{~}{x}e^{i\stackrel{~}{k}\stackrel{~}{x}}\underset{x_{}\mathrm{}}{lim}e^{iek_+x_{}}\psi _{}({\displaystyle \frac{}{}}X(k_+),x_{},\stackrel{~}{x}).`$ The canonical anti-commutation relations (53-54) are unchanged. It is natural to study the state that is empty at $`x_+=0`$, which means $$\mathrm{\Xi }_0(k_+,\stackrel{~}{k})\|\mathrm{\Omega }=0=\mathrm{\Xi }_0^{}(k_+,\stackrel{~}{k})\|\mathrm{\Omega },$$ (68) for all $`k_+>0`$. It is also natural to forbid particles from entering via the surface at $`x_{}=\mathrm{}`$, which implies $$\mathrm{\Xi }_{\mathrm{}}^{}(k_+,\stackrel{~}{k})\|\mathrm{\Omega }=0,$$ (69) also for $`k_+>0`$. The two operator orderings of $`\mathrm{\Psi }_+`$ and $`\mathrm{\Psi }_+^{}`$ give, $`\mathrm{\Omega }\left\|\mathrm{\Psi }_+(x_+,k_+,\stackrel{~}{k})\mathrm{\Psi }_+^{}(x_+^{},q_+,\stackrel{~}{q})\right\|\mathrm{\Omega }=(2\pi )^3\delta (k_+q_+)\delta ^2\left(\stackrel{~}{k}\stackrel{~}{q}\right){\displaystyle \frac{P_+}{\sqrt{2}}}`$ $`\times \theta (k_+)(0,x_+;k_+,\stackrel{~}{k})^{}(0,x_+^{};k_+,\stackrel{~}{k}),`$ (70) $`\mathrm{\Omega }\left\|\mathrm{\Psi }_+^{}(x_+^{},q_+,\stackrel{~}{q})\mathrm{\Psi }_+(x_+,k_+,\stackrel{~}{k})\right\|\mathrm{\Omega }=(2\pi )^3\delta (k_+q_+)\delta ^2\left(\stackrel{~}{k}\stackrel{~}{q}\right){\displaystyle \frac{P_+}{\sqrt{2}}}`$ $`\times \{{\displaystyle \frac{}{}}\theta (k_+)\theta (x_+X)(X,x_+;k_+,\stackrel{~}{k})^{}(X,x_+^{};k_+,\stackrel{~}{k})(e^{\pi \lambda }e^{\pi \lambda })`$ $`{\displaystyle \frac{}{}}+\theta (k_+)(0,x_+;k_+,\stackrel{~}{k})^{}(0,x_+^{};k_+,\stackrel{~}{k})\}.`$ (71) Now note that the conjugated mode function can be re-expressed using the Dirac identity, $`^{}(0,x_+^{};k_+,\stackrel{~}{k})=\mathrm{exp}\left[{\displaystyle \frac{i}{2}}\stackrel{~}{\omega }^2{\displaystyle _0^{x_+^{}}}{\displaystyle \frac{du}{k_+eA_{}(u)iϵ}}\right],`$ (74) $`\mathrm{exp}\left[{\displaystyle \frac{i}{2}}\stackrel{~}{\omega }^2{\displaystyle _0^{x_+^{}}}𝑑u\left\{{\displaystyle \frac{1}{k_+eA_{}(u)+iϵ}}+2\pi i\delta \left(k_+eA_{}(u)\right)\right\}\right],`$ $`=\mathrm{exp}\left[{\displaystyle \frac{i}{2}}\stackrel{~}{\omega }^2{\displaystyle _0^{x_+^{}}}{\displaystyle \frac{du}{k_+eA_{}(u)+iϵ}}\right]e^{2\pi \lambda \theta (k_+)\theta (x_+^{}X)}.`$ We can therefore write the various $`\times ^{}`$ products as, $`(0,x_+;k_+,\stackrel{~}{k})^{}(0,x_+^{};ik_+,\stackrel{~}{k})`$ (75) $`(x_+^{},x_+;k_+\stackrel{~}{k})e^{2\pi \lambda \theta (k_+)\theta (x_+^{}X)},`$ $`(X,x_+;k_+,\stackrel{~}{k})^{}(X,x_+^{};k_+,\stackrel{~}{k})`$ (76) $`(x_+^{},x_+;k_+\stackrel{~}{k})e^{\pi \lambda \theta (k_+)\theta (x_+^{}X)}.`$ Assembling all the pieces gives the following result for the expectation value of the $`x_+`$–ordered product: $`\mathrm{\Omega }\left\|X_+\left\{\mathrm{\Psi }_+(x_+,k_+,\stackrel{~}{k})\mathrm{\Psi }_+^{}(x_+^{},q_+,\stackrel{~}{q})\right\}\right\|\mathrm{\Omega }=(2\pi )^3\delta ^3(kq){\displaystyle \frac{P_+}{\sqrt{2}}}`$ (77) $`\times (x_+^{},x_+;k_+\stackrel{~}{k})\{\theta (\mathrm{\Delta }x_+)\theta (k_+)e^{2\pi \lambda \theta (x_+^{}X)}\theta (\mathrm{\Delta }x_+)\theta (k_+)`$ $`\theta (\mathrm{\Delta }x_+)\theta (k_+)\theta (x_+X)(1e^{2\pi \lambda })\}.`$ Expanding the exponential which contains the theta function, $$e^{2\pi \lambda \theta (x_+^{}X)}=\theta (Xx_+^{})+\theta (x_+^{}X)e^{2\pi \lambda },$$ (78) and simplifying under the assumption that $`eA_{}(u)`$ is an increasing function gives the following more compact result, $`\mathrm{\Omega }\left\|X_+\left\{\mathrm{\Psi }_+(x_+,k_+,\stackrel{~}{k})\mathrm{\Psi }_+^{}(x_+^{},q_+,\stackrel{~}{q})\right\}\right\|\mathrm{\Omega }=(2\pi )^3\delta ^3(kq){\displaystyle \frac{P_+}{\sqrt{2}}}`$ (79) $`\times (x_+^{},x_+;k_+\stackrel{~}{k})\{\theta (\mathrm{\Delta }x_+)\theta (Xx_+^{})\theta (\mathrm{\Delta }x_+)\theta (Xx_+)`$ $`+\theta (x_+X)\theta (x_+^{}X)\theta (k_+)e^{2\pi \lambda }\}.`$ We can still obtain the minus components from the plus ones, $$\mathrm{\Psi }_{}(x_+,k_+,\stackrel{~}{k})=\left(\frac{m\stackrel{~}{\gamma }\stackrel{~}{k}}{\stackrel{~}{\omega }^2}\right)\gamma _+i_+\mathrm{\Psi }_+(x_+,k_+,\stackrel{~}{k}),$$ (80) so that relations (45-47) still determine the other components of the propagator from the $`++`$ ones. So the full 2-point function is, $`\mathrm{\Omega }\left\|X\left\{\mathrm{\Psi }(x_+,k_+,\stackrel{~}{k})\mathrm{\Psi }^{}(x_+^{},q_+,\stackrel{~}{q})\right\}\right\|\mathrm{\Omega }=(2\pi )^3\delta (k_+q_+)\delta ^2\left(\stackrel{~}{k}\stackrel{~}{q}\right)`$ (81) $`\times (x_+^{},x_+;k_+\stackrel{~}{k})\{\theta (\mathrm{\Delta }x_+)\theta (Xx_+^{})\theta (\mathrm{\Delta }x_+)\theta (Xx_+)`$ $`+\theta (x_+X)\theta (x_+^{}X)\theta (k_+)e^{2\pi \lambda }\left\}\right\{{\displaystyle \frac{1}{\sqrt{2}}}P_+`$ $`+{\displaystyle \frac{P_+\gamma _{}}{2\sqrt{2}}}\left({\displaystyle \frac{m+\stackrel{~}{\gamma }\stackrel{~}{k}}{k_+eA_{}(x_+^{})iϵ}}\right)+\left({\displaystyle \frac{m\stackrel{~}{\gamma }\stackrel{~}{k}}{k_+eA_{}(x_+)+iϵ}}\right){\displaystyle \frac{\gamma _+P_+}{2\sqrt{2}}}`$ $`+{\displaystyle \frac{\frac{1}{2}\stackrel{~}{\omega }^2\frac{1}{\sqrt{2}}P_{}}{(k_+eA_{}(x_+)+iϵ)(k_+eA_{}(x_+^{})iϵ)}}\}.`$ ## 6 Discussion We have shown that, if one computes the one loop expectation value of the current using Schwinger’s propagator, the result is zero. The reason for this emerges when one attempts to write the propagator as the expectation value of the time-ordered product of $`\psi (x)\overline{\psi }(x^{})`$ in the presence of some state. The only way to do so is with a state which changes as the time parameter ($`x_+`$) of the $`\psi `$ field does. At each different value of $`x_+`$ the appropriate state is the one which happens to be free of particles at that instant, so of course the current is zero. Because of the special roles played by time and by the direction of the electric field, lightcone coordinates are particularly well suited to this problem. At the level of operators and states this implies a profound departure from the usual picture in which one imagines specifying a state on a surface of constant $`x^0`$. For a state specified instead on a surface of constant $`x_+`$ the phenomenon of pair production is a discrete and instantaneous event. For each mode of fixed $`k_+`$ it occurs when the minimally coupled momentum $`p_+=k_+eA_(x_+)`$ vanishes. This can be understood quite simply by representing the lightcone system in the standard way as the infinite boost limit of a conventional system in which the states are defined on surfaces of constant time . Particles produced with finite $`p^3`$ in that frame must have $`p^3=\mathrm{}`$ in the lightcone frame, and $`p^3=\mathrm{}`$ implies $`p_+=0^+`$ for a particle which is on shell . A remarkable feature of our treatment is that explicit (and quite simple!) mode functions can be obtained for a background in which the electric field depends arbitrarily upon the lightcone coordinate $`x_+`$.<sup>3</sup><sup>3</sup>3The simplicity of the lightcone mode functions has been noted previously, for the special case of constant electric field, by Srinivasan and Padmanabhan . Note, however, that we do not agree with some aspects of their solution. We have in fact worked out the actual electron propagator for such a background in the presence of a state which is empty at $`x_+=0`$. Since neither the photon propagator nor the vertices of QED are affected by the background, expression (81) completes the Feynman rules with which one can compute the quantum electrodynamic back-reaction to any order. It is significant that we can get the propagator for a class of backgrounds general enough to include the actual solution as the electric field evolves under the impact of the current it generates. Acknowledgements We thank D. Boyanovsky, C. B. Thorn and T. N. Tomaras for discussions on related subjects. This work was partially supported by DOE contract DE-FG02-97ER41029, by the Greek General Secretariat of Research and Technology grant 97 E$`\mathrm{\Lambda }`$–120, and by the Institute for Fundamental Theory.
warning/0007/hep-ph0007275.html	ar5iv	text	# THE LIGHTEST STRANGE SCALAR MESONTalk presented by Stuart Cherry ## 1 Introduction The scalar mesons are one of the most controversial subjects in hadron physics. In the last edition of the PDG, there were four scalar isoscalars listed below 1.5 GeV (the $`f_0(4001200)`$, $`f_0(980)`$, $`f_0(1370)`$ and $`f_0(1500)`$, with a possible fifth at 1.7 GeV, the $`f_J(1710)`$. There were also two isovectors listed, the $`a_0(980)`$ and the $`a_0(1450)`$. This is obviously too many for a standard $`q\overline{q}`$ nonet. However, many QCD-motivated models predict the the existence of non-$`q\overline{q}`$ mesons, such as $`qq\overline{qq}`$ states, $`K\overline{K}`$ molecules and glueballs. It is precisely in the scalar-isoscalar that these unconventional mesons are most like to be found. This excess of isoscalars and isovectors has led to the suggestion that there are in fact two scalar nonets: the conventional one centred around 1.4 GeV and an unconventional, possibly $`qq\overline{qq}`$, one centred around 1 GeV. However, the PDG lists only one pair of strange scalar mesons in this region, the $`K_0^{}(1430)`$ and so some authors have postulated a light strange meson, known as the $`\kappa (900)`$ Evidence for and against the $`\kappa (900)`$ has been claimed within various models. It should be kept in mind that the existence of a resonance is not defined by the ability to fit a particular formula to the experimental data. A resonance is entirely defined by the presence of a pole of the $`S`$matrix on the nearby unphysical sheet. I would like to present to you the results of a model-independent calculation to determine whether the $`\kappa (900)`$ is indeed present in the experimental data by determining the number and positions of the poles of the $`S`$matrix. The method, due to Nagová et al, does not require the artificial separation of data into background and resonance components and I will begin by briefly outlining it. ## 2 Mapping To determine the position of a pole in the complex plane from data along the real axis we must perform some form of analytic continuation. In order to maximise the region of validity of our analytic continuation we begin by conformally mapping the unphysical sheet of the $`\pi K`$ partial wave amplitude (see Fig. 1) into the unit disc. The mapping is designed so that the cuts in the $`s`$plane are mapped onto the circumference of the circle in the $`z`$plane. The mapping is accomplished in two steps. First $$y(s)=\left(\frac{ss_c}{s+s_c}\right)^2\mathrm{and}\mathrm{then}z(s)=\frac{i\beta \sqrt{y(s)}\sqrt{y(s)y(s_{th})}}{i\beta \sqrt{y(s)}+\sqrt{y(s)y(s_{th})}},$$ (1) where $`\beta `$ is a real number chosen to minimise the distance the continuation must cover. From Fig. 2 we notice that physical data could never cover the whole circle. We also see that the proportion of the circle covered by each increment in $`s`$ falls very sharply as $`s`$ increases. Hence we can neglect radial excitations from our analysis. ## 3 Continuation Let $`Y(z)`$ and $`ϵ(z)`$ be continuous functions describing the experimental data and errors around the entire circle. If $`F(z)`$ is a square-integrable function around the circle, then we can test how well it describes the data through a $`\chi ^2`$ defined by $$\chi ^2=\frac{1}{2\pi }_C\left\|\frac{F(z)Y(z)}{ϵ(z)}\right\|^2\|\mathrm{d}z\|.$$ (2) We introduce a weight function $`g(z)`$ which we define to be real analytic, non-zero throughout the disc and constrained by $`\|g(z)\|=ϵ(z)`$ around the circle. Expanding our data and trial functions as Laurent series about the origin, i.e. $$y(z)=\frac{Y(z)}{g(z)}=\underset{k=\mathrm{}}{\overset{\mathrm{}}{}}y_kz^k\mathrm{and}f(z)=\frac{F(z)}{g(z)}=\underset{k=\mathrm{}}{\overset{\mathrm{}}{}}a_kz^k,$$ (3) gives $$\chi ^2=\underset{k=1}{\overset{\mathrm{}}{}}\left(a_ky_k\right)^2+\underset{k=0}{\overset{\mathrm{}}{}}\left(a_ky_k\right)^2$$ (4) The pole structure of our partial wave amplitude, $`Y(z)`$, can then be determined by finding the test function $`F(z)`$ which minimises the first summation in Eq. (4). Since partial wave amplitudes are real analytic, the coefficients $`a_k`$ and $`y_k`$ will be real. To test if the amplitude is free of resonances we set $`a_k0k>0`$. Then the quantity $`\chi _0^2=_{k=1}^{\mathrm{}}y_k^2`$ will have a $`\chi ^2`$distribution. To test if the amplitude contains one resonance we define $$\stackrel{~}{y}(z)=\frac{Y(z)B_\lambda (z)}{g(z)}=\underset{k=\mathrm{}}{\overset{\mathrm{}}{}}\stackrel{~}{y}_kz^k\mathrm{where}B_\lambda (z)=\frac{(z\lambda )(z\lambda ^{})}{(1z\lambda )(1z\lambda ^{})}.$$ (5) $`B_\lambda `$ is the one resonance Blaschke pole killing factor. Now the quantity $`\chi _1^2=_{k=1}^{\mathrm{}}\stackrel{~}{y}_k^2`$ will have a $`\chi ^2`$distribution. ## 4 Inputs The experimental inputs into our calculation are the $`\pi K`$ scattering magnitudes, $`a(s)`$, and phases, $`\varphi (s)`$, as measured, for example, by the LASS experiment. The amplitude is normalised such that $$f^I(s)=\frac{a(s)e^{i\varphi (s)}}{\rho (s)},\rho (s)=\frac{2q(s)}{\sqrt{s}}$$ (6) where $`q(s)`$ is the centre of mass 3-momentum of the $`\pi K`$ system and the superscript indicates we are on the physical sheet. In order to detect a resonance we must move onto the relevant unphysical sheet, this is done by swapping the sign of the phase. After mapping we have discrete data points $`f_i(z)`$ and errors $`\mathrm{\Delta }_i`$ around a portion of the circle. Where the data is more densely packed the amplitude is more tightly controlled, so we weight the errors by the density of data points in that region, i.e. we define $`ϵ_i=\mathrm{\Delta }_i\sqrt{\|\delta z_i\|/2\pi }`$. Between threshold and the start of the data we use the LASS effective range fit to create a few guide points. In the unphysical region we simply guess a few widely spaced guide points and assign large errors. The large spacing and errors will hopefully deweight these points so that they do not unduly affect the results of the calculation. With these three different inputs and using the Schwarz Reflection Principle we can cover the circle with data and then we interpolate to give the continuous functions that we need. A suitable form for the weight function is $`g(z)=\mathrm{exp}_{n=1}^Mc_nz^n`$ where the $`c_n`$ can be found from a Fourier cosine fit to $`\mathrm{log}ϵ(z)`$. We can now calculate the singular coefficients of our Laurent expansions in the usual way, i.e. $$y_k=\frac{1}{2\pi }_c\frac{Y(z)z^k}{g(z)}\|\mathrm{d}z\|\mathrm{and}\stackrel{~}{y}_k=\frac{1}{2\pi }_c\frac{Y(z)B_\lambda (z)z^k}{g(z)}\|\mathrm{d}z\|$$ (7) We want to find the value of $`\lambda `$ that minimises $`\chi _1^2`$ and then compare the final value with $`\chi _0^2`$. We can carry out a similar procedure to test for two resonances and by comparing the various $`\chi ^2`$’s we can determine the number of resonances present in the channel. ## 5 Results In order to assess its capabilities, we began by testing the method explained previously with a trial scattering amplitude. This amplitude was created using Jost functions and had two resonances present, a $`\kappa _1`$ with mass 900 MeV and width 500 GeV and a $`\kappa _2`$ with mass 1400 MeV and width 300 MeV. The results are shown in Table 1. We can see clearly that the $`\chi ^2`$ falls significantly as the number of resonances is increased and that changing the size of the errors in the unphysical region has very little effect on either of the pole positions obtained. Thus we can confidently claim that the amplitude contains two resonances and furthermore the parameters of these resonances have been quite accurately determined. Turning now to experimental data, Table 2 shows the results obtained from the LASS total S-wave data. We see that the fall in $`\chi ^2`$ in going from one to two resonances is negligible compared to the fall between zero and one resonance. Also the parameters for the second resonance are quite sensitive to the size of the errors in the unphysical region. This suggests that the second resonance listed in Table 2 is merely an artifact of looking for two poles and is not really present in the data. ## 6 Conclusions In this talk I have tried to outline a calculation to determine the number of resonances present in strange scalar channel. This method is most sensitive to lighter resonances and is capable of identifying two resonances in a channel even when they are broad and overlapping. The experimental $`\pi K`$ scattering data is found to contain only one resonance below 1.8 GeV and this state is readily identifiable with the $`K_0^{}(1430)`$. There is no evidence for a $`\kappa (900)`$. I would like to thank the organisers for the opportunity to attend such an interesting and enjoyable meeting.
warning/0007/cond-mat0007233.html	ar5iv	text	# Excited stationary states of trapped Bose-Einstein condensates ## Abstract We investigate the excited stationary states of Bose-Einstein condensates trapped in harmonic potentials. We derive simple analytical approximations of the first few eigenstates of the associated time-independent one-dimensional Gross-Pitaevskii equation and their energies. Our results are excited state generalizations of the Thomas-Fermi approximation of the ground state. As opposed to traditional studies of liquid helium, the recent experimental work on Bose-Einstein condensation centers on inhomogeneous mesoscopic systems. However, such systems are highly nonlinear and no exact analytical expressions for the wavefunctions and associate energies are yet known. In order to understand the behavior of Bose-Einstein condensates (BEC) it would be beneficial to have at least simple approximate expressions for the condensate wave functions. These could be used to study the qualitative behavior of the condensate and, given the approximate expressions are precise enough, might even be useful for the determination of overlap factors between different wave functions, spectra and other quantities of interest. Yukalov et al. have investigated this problem using a perturbation method. Also, Kivshar et al. recently highlighted the connection between this problem and that of dark solitons and gave numerical solutions to the first few excited stationary states. In this communication we derive simple analytical approximations to the stationary states $`\mathrm{\Phi }_n`$ $`(n=0,1,2,\mathrm{})`$ of the one-dimensional Gross-Pitaevskii equation for a BEC confined in a harmonic trap : $$\left(\frac{^2}{z^2}+\lambda ^2z^2+\frac{Q}{2}\|\mathrm{\Phi }_n(z)\|^22\epsilon _n\right)\mathrm{\Phi }_n(z)=0.$$ (1) Here $`z`$ is the longitudinal trap coordinate, the aspect ratio, $`\lambda =\omega _z/\omega 1`$, is given by the ratio of the longitudinal trap frequency $`\omega _z`$ to the transversal trap frequency $`\omega `$, $`\lambda `$ formally is the effective spring constant of the 1-dimensional problem. $`\lambda `$ is assumed to be much smaller than unity in order to yield a cigar-shaped condensate which can be modeled by a 1-dimensional problem. $`Q/2=4\pi aN/a_0`$ is the effective non-linearity for the interaction of the $`N`$ trapped atoms with mass $`m`$ and scattering length $`a`$ in the trap with an effective diameter $`a_0=\sqrt{\mathrm{}/(m\omega )}`$ . The states $`\mathrm{\Phi }_n`$ are normalized to unity, but, owing to the non-linearity of the Gross-Pitaevskii-equation, they do not form an orthogonal set, generally $`\mathrm{\Phi }_n\|\mathrm{\Phi }_m0`$. It is convenient to rescale the Gross-Pitaevskii-equation to make the coefficients of the kinetic and the non-linear term equal to unity. Using the rescaling $`Y_n=\mathrm{\Phi }_n\frac{1}{2}\sqrt{Q/\epsilon _0}`$ and $`x=z\sqrt{2\epsilon _0}`$ Eq. (1) becomes $$\left(\frac{^2}{x^2}+\lambda ^2x^2+\|Y_n(x)\|^2ϵ_n\right)Y_n(x)=0.$$ (2) The advantage of this scaling is that the solutions are parameterized solely by $`\lambda `$. The energy eigenvalue $`ϵ_n=\epsilon _n/\epsilon _0`$ is now scaled in units of the Thomas-Fermi energy $`\epsilon _0=(3^{2/3}/8)(Q\lambda )^{2/3}`$ of the lowest stationary state and the spring constant is expressed in terms of $`\lambda ^2=(\omega _z/\omega )^2=(512/9)\epsilon _0^3/Q^2`$. Our choice of rescaling results in the normalization $`\mathrm{\Phi }_n\|\mathrm{\Phi }_n=1=\frac{4\sqrt{\epsilon _0}}{\sqrt{2}Q}Y_n\|Y_n`$ which can be written as $$Y_n\|Y_n=\frac{4}{3\lambda }.$$ (3) The Thomas-Fermi approximation of $`Y_0`$ entirely neglects the contribution from the kinetic energy term and yields the standard approximate solution of (2): $$Y_0(x)=\sqrt{1\lambda ^2x^2},\text{ for }\|x\|<X_0=1/\lambda $$ (4) and $`Y_0(x)=0`$ otherwise, as illustrated in Fig. 1. To generalize the Thomas-Fermi approach to excited stationary states, $`Y_n`$ for $`n>0`$, we need to include the kinetic energy term. Our approach to deriving a solution is based on a piecewise ansatz which is motivated as follows. Over small regions about some fixed point $`x=x^{}`$ the potential energy term $`\lambda ^2x^2`$ can be considered slowly varying compared to the other terms in Eq. (2). The local solution about $`x^{}`$ is then given approximately by the solution of the nonlinear Schrödinger equation $`\left({\displaystyle \frac{^2}{x^2}}+\|Y_n(x)\|^2ϵ_n^{}\right)Y_n(x)=0.`$ (5) where $`ϵ_n^{}=ϵ_n\lambda ^2x^2`$. This equation is well known in soliton theory; the solutions with finite amplitude are given by the Jacobi-sine function ‘sn’: $$\sqrt{ϵ_n^{}}A\text{ sn}(\sqrt{ϵ_n^{}(1A^2/2)}(xa),\frac{A}{\sqrt{2A^2}})$$ (6) where $`a`$ is an arbitrary displacement along the $`x`$ axis and $`0<A1`$. The wavelength of the function in (6) is given by $$\mathrm{\Lambda }=\frac{4}{\sqrt{ϵ_n^{}(1A^2/2)}}\times K(\frac{A}{\sqrt{2A^2}}),$$ (7) where $`K`$ is a complete elliptic integral of the first kind . Imagine joining the solutions for two different values of $`x^{}`$ smoothly at some intermediate point along the $`x`$ axis. As an example, let this intermediate point be a zero of the two solutions. The slope of (6) at a zero is given by $`ϵ_n^{}A\sqrt{1A^2/2}`$ and so a smooth join requires the value of $`A`$ to be larger for the solution with the larger value of $`\|x^{}\|`$ to compensate for the smaller value of $`ϵ_n^{}=ϵ_n\lambda ^2x^2`$. Repeating this smooth joining of solutions for successive pairs of $`x^{}`$ values along the $`x`$ axis leads to a set of piecewise solutions for which the value of $`A`$ increases with $`\|x^{}\|`$. At the critical value $`A=1`$ the wavelength becomes infinite and the solution of Eq. (5) is given by the corresponding limit of (6) namely $$\sqrt{ϵ_n^{}}\mathrm{tanh}\left(\sqrt{ϵ_n^{}/2}(xa)\right).$$ (8) This piecewise analysis implies that the general finite-amplitude solution of Eq. (2) will be oscillatory near the origin and change to a tanh-like function some distance away. Moreover, the $`\mathrm{tanh}`$ function in (8) is relatively flat for $`\|xa\|`$ significantly different from zero, and so (8) is approximated well by just the prefactor $`\sqrt{ϵ_n^{}}=\sqrt{ϵ_n\lambda ^2x^2}`$. We note that this prefactor is the Thomas-Fermi solution (4) in terms of the local variable $`x^{}`$. We can reduce the number of piecewise segments needed in our ansatz using energy considerations . The potential energy and its variation are smallest in the center of the trap and to a good approximation we can neglect the trap potential in this region. This amounts to replacing $`ϵ_n^{}`$ with $`ϵ_n`$ and holding $`A`$ constant in our oscillatory solution (6) to give an oscillating function $`Y_{\mathrm{osc}}`$ which has a fixed wavelength and amplitude. Further away from the center of the trap the potential energy increases at the expense of the kinetic energy to the extent that the later becomes negligible. The wave function is then well described by the conventional Thomas-Fermi solution , $`\sqrt{ϵ_n\lambda ^2x^2}`$. Thus, replacing $`x^{}`$ with $`x`$ in the prefactor $`\sqrt{ϵ_n\lambda ^2x^2}`$ of (8) yields a function $`Y_{\mathrm{tf}}`$ which provides a transition between $`Y_{\mathrm{osc}}`$ and the Thomas-Fermi solution. We will call $`Y_{\mathrm{tf}}`$ the “tail function” as it gives the shape of the condensate in its outer regions. Our ansatz is to match these two segments $`Y_{\mathrm{osc}}`$ and $`Y_{\mathrm{tf}}`$ smoothly at points $`x=\pm T_n`$ which lie symmetrically about the origin. In analogy to Eq. (4) we also set $`Y_n(x)=0`$ for $`\|x\|>X=\sqrt{ϵ_n}/\lambda `$ where $`X`$ is the half-width of the condensate in the Thomas-Fermi approximation. Thus our ansatz is given by $$Y_n(x)=\{\begin{array}{cc}Y_{\mathrm{osc}}(x),\hfill & \|x\|T_n\hfill \\ Y_{\mathrm{tf}}(x),\hfill & T_n<\|x\|X\hfill \\ 0,\hfill & X<\|x\|\hfill \end{array}$$ (9) where $`Y_{\mathrm{osc}}(x)`$ $`=`$ $`\sqrt{ϵ_n}A`$ $`\times \text{sn}(\sqrt{ϵ_n(1A^2/2)}(xT_n),{\displaystyle \frac{A}{\sqrt{2A^2}}}),`$ $`Y_{\mathrm{tf}}(x)`$ $`=`$ $`\left({\displaystyle \frac{x}{\|x\|}}\right)^n\sqrt{ϵ_n\lambda ^2x^2}`$ $`\times \mathrm{tanh}\left(\sqrt{(ϵ_n\lambda ^2T^2)/2}(\|x\|T_n)\right).`$ The points $`x=\pm T_n`$, where the solution switches between the oscillatory part and the tail function is given by $`T_n=(n1)\mathrm{\Lambda }/4`$. The forms of $`T_n`$ and $`Y_{\mathrm{osc}}(x)`$ ensure that $`Y_n(x)`$ is either symmetric or antisymmetric for $`n`$ even or odd. The parameters $`A`$, $`T_n`$ and $`ϵ_n`$ are determined by the requirements that the two segments match smoothly and that the solution is normalized according to Eq. (3). The first requirement is satisfied when the gradients, $`\frac{}{x}Y_{\mathrm{osc}}\|_{x=T_n}`$ and $`\frac{}{x}Y_{\mathrm{tf}}\|_{x=T_n}`$ are equal, i.e. when $$ϵ_nA\sqrt{1A^2/2}=(ϵ_n\lambda ^2T_n^2)/\sqrt{2}.$$ (10) For notational simplicity we concentrate on the second excited state $`Y_2`$, i.e. $`n=2`$. The expression for $`T_2`$ can be approximated by $`T_2\sqrt{2/ϵ_2}`$ arctanh$`(\sqrt{(1+A^2)/2})`$ $`\sqrt{2/ϵ_2}\text{ln}(2/\sqrt{1A})`$. This can be used to solve Eq. (10) for $`A1\lambda ϵ_2W(8ϵ_2/\lambda )/2`$ where $`W`$ is the Lambert function defined by $`W\mathrm{exp}(W)=x`$. This, in turn, can be approximated to give $`A1\lambda ϵ_2\mathrm{ln}(8ϵ_2/\lambda )/2`$ and therefore $`T_2\mathrm{ln}(\lambda ϵ_2\mathrm{ln}(8ϵ_2/\lambda )/8)/\sqrt{2ϵ_2}`$. Thus we now have a simple expression for $`T_2`$ in terms of $`ϵ_2`$. The second requirement, i.e. that $`Y_2(x)`$ be normalized according to Eq. (3), determines the value of $`ϵ_2`$. The contribution from the oscillatory part of the wave function is $`Y_{\mathrm{osc}}\|Y_{\mathrm{osc}}`$ $`=`$ $`2[\sqrt{2}\text{ arctanh}(\sqrt{(1+A^2)/2})`$ $`\sqrt{(1+A^2)/ϵ_2}]`$ and the contribution from the tails is approximately $`Y_{\mathrm{tf}}\|Y_{\mathrm{tf}}`$ $``$ $`2{\displaystyle _T^X}dx[\lambda ^2(T^2x^2)`$ (12) $`+(ϵ_n\lambda ^2T^2)\mathrm{tanh}^2(\sqrt{(ϵ_n\lambda ^2x^2)/2})].`$ We use tanh($`\frac{\sqrt{ϵ}X}{\sqrt{2}})1`$ to find a simplified expression for (12). Inserting the above expression for $`A`$ and $`T`$ into $`Y_2\|Y_24/(3\lambda )`$, treating $`\lambda `$ and $`\delta _2=ϵ_2ϵ_0`$ as small expansion parameters and performing a somewhat tedious calculation using an expansion in $`\delta _2`$ followed by an expansion in $`\lambda `$ finally yields $`ϵ_21+2\sqrt{2}\lambda `$. The derivation for the general case $`n=1,2,3,\mathrm{}`$ is very similar. We find that $`ϵ_n`$ $``$ $`1+n\sqrt{2}\lambda ,`$ (13) $`A`$ $``$ $`1\lambda ϵ_n\mathrm{ln}(8ϵ_n/\lambda )/2,`$ (14) $`\text{and }T_n`$ $``$ $`(1n)\mathrm{ln}(\lambda ϵ_n\mathrm{ln}(8ϵ_n/\lambda )/8)/\sqrt{2ϵ_n}.`$ (15) Eq. (9) together with Eqs. (13-15) are the main results of this communication and give our analytical approximations of the stationary states of a harmonically trapped condensate. We note, however, that the greater the value of $`n`$, the more stringent the requirements on the aspect ratio $`\lambda `$ to be small. We should expect, and do obtain, best results for low order $`n`$ because in this case the above approximations are least critical. In a more detailed paper we will present a self-consistency argument to justify our approach. Figs. 2-4 compare our analytical and the exact (numerically determined) wave functions for low-order excited states and illustrate the increasing deviation from the fixed wavelength approximation for more highly excited states. As in the Thomas-Fermi solution in Fig. 1, the omission of the kinetic energy for the outer regions of the condensate leads to the artificial truncation at the outer edge $`x=X`$ of the wave functions rather than a Gaussian tail. Also Fig. 4 shows that the fixed wavelength approximation of Eq. (9) for the oscillatory segment of the solution breaks down for high excitation numbers. Interestingly we have found that the overlap between the numerical and the analytical solutions to $`Y_4`$ is higher for larger values of $`\lambda `$ (e.g. $`\lambda =1/20`$). We believe this is a consequence of the fixed wavelength approximation, we will explore this point elsewhere. These excited stationary states may soon be realized experimentally. Various possibilities to generate excited states have been discussed in the literature . Moreover, Burger et al. recently produced a slowly-moving dark soliton using the ‘phase imprinting method’ . It should be possible with a slight modification of their experiment to generate the first excited state $`Y_1(x)`$ as a stationary dark soliton. The generation of higher excited states should also be possible in a similar manner. In conclusion, we have derived simple, approximate analytic expressions of the excited stationary solutions of the 1-dimensional Gross-Pitaevskii equation using an ansatz based on piecewise solutions. Our analytic solutions agree very well with numerical solutions for a sufficiently small aspect ratio. Generalizations of the approach presented here should be rather straightforward, for example, one can introduce a coordinate dependent wavelength of the oscillatory part of the solution rather than using the fixed wavelength approximation. Further details will be explored elsewhere . We are grateful to D. Richards and A. Chefles for helpful discussions. This work was supported by the Engineering and Physical Sciences Research Council (EPSRC).
warning/0007/hep-ph0007248.html	ar5iv	text	# Measurements of the 𝑄²-Dependence of the Proton and Neutron Spin Structure Functions 𝑔₁^𝑝 and 𝑔₁^𝑛 ## Abstract The structure functions $`g_1^p`$ and $`g_1^n`$ have been measured over the range $`0.014<x<0.9`$ and $`1<Q^2<40`$ GeV<sup>2</sup> using deep-inelastic scattering of 48 GeV longitudinally polarized electrons from polarized protons and deuterons. We find that the $`Q^2`$ dependence of $`g_1^p`$ ($`g_1^n)`$ at fixed $`x`$ is very similar to that of the spin-averaged structure function $`F_1^p`$ ($`F_1^n`$ ). From an NLO QCD fit to all available data we find $`\mathrm{\Gamma }_1^p\mathrm{\Gamma }_1^n=0.176\pm 0.003\pm 0.007`$ at $`Q^2=5`$ GeV<sup>2</sup>, in agreement with the Bjorken sum rule prediction of $`0.182\pm 0.005`$. preprint: SLAC-PUB-7994, July 2000 T/E The spin-dependent structure function $`g_1(x,Q^2)`$ for deep-inelastic lepton-nucleon scattering is of fundamental importance in understanding the quark and gluon spin structure of the proton and neutron. The $`g_1`$ structure function depends both on $`x`$, the fractional momentum carried by the struck parton, and on $`Q^2`$, the squared four-momentum of the exchanged virtual photon. The fixed-$`Q^2`$ integrals (or first moments) $`\mathrm{\Gamma }_1^p(Q^2)=_0^1g_1^p(x,Q^2)𝑑x`$ for the proton and $`\mathrm{\Gamma }_1^n(Q^2)=_0^1g_1^n(x,Q^2)𝑑x`$ for the neutron are related to the net quark helicity $`\mathrm{\Delta }\mathrm{\Sigma }`$ in the nucleon. Measurements of $`\mathrm{\Gamma }_1^p`$ , $`\mathrm{\Gamma }_1^d`$ for the deuteron (which essentially measures the average of the proton and neutron), and $`\mathrm{\Gamma }_1^n`$ have found $`\mathrm{\Delta }\mathrm{\Sigma }`$ between 0.2 and 0.3, significantly less than the prediction that $`\mathrm{\Delta }\mathrm{\Sigma }=0.58`$ assuming zero net strange quark helicity and SU(3) flavor symmetry in the baryon octet. A fundamental sum rule originally derived from current algebra by Bjorken predicts $`\mathrm{\Gamma }_1^p(Q^2)\mathrm{\Gamma }_1^n(Q^2)=g_A/6g_V`$. Recent measurements are in agreement with this sum rule prediction when perturbative QCD (pQCD) corrections are included. According to the DGLAP equations , $`g_1`$ is expected to evolve logarithmically with $`Q^2`$, and in the case of $`g_1^p`$ to increase with $`Q^2`$ at low $`x`$, and decrease with $`Q^2`$ at high $`x`$ . A similar $`Q^2`$-dependence has been observed in the spin-averaged structure functions $`F_1(x,Q^2)`$, while the ratio $`g_1/F_1`$ has been found to be approximately independent of $`Q^2`$ . The precise behavior is sensitive to the underlying spin-dependent quark and gluon distribution functions. Fits to data for $`g_1`$ using NLO pQCD allow determinations of the first moments (from which the Bjorken sum rule can be tested) as well as the valence quark, sea quark, and gluon spin contributions. The goal of the present experiment (SLAC E155) was to make precise measurements over a wide range of $`Q^2`$ in a single experiment to further constrain these quantities. The ratio of polarized to unpolarized structure functions can be determined from measured longitudinal asymmetries $`A_{}`$ using $$g_1/F_1=A_{}/d+(g_2/F_1)[(2Mx)/(2E\nu )],$$ (1) where $`d=[(1ϵ)(2y)]/\{y[1+ϵR(x,Q^2)]\}`$, $`y=\nu /E`$, and $`\nu =EE^{}`$, where $`E`$ is the incident and $`E^{}`$ is the scattered electron energy in the lab frame, $`ϵ^1=1+2[1+\gamma ^2]\mathrm{tan}^2(\theta /2)`$, $`\gamma ^2=Q^2/\nu ^2`$, $`\theta `$ is the electron scattering angle, $`M`$ is the nucleon mass, and $`R(x,Q^2)=[F_2(x,Q^2)(1+\gamma ^2)]/[2xF_1(x,Q^2)]1`$ is typically 0.2 for the kinematics of this experiment . For the contribution of the transverse spin structure function $`g_2`$ we used the twist-two result of Wandzura and Wilczeck ($`g_2^{WW}`$) $$g_2^{WW}(x,Q^2)=g_1(x,Q^2)+_x^1g_1(\xi ,Q^2)𝑑\xi /\xi ,$$ (2) evaluated using the empirical fit to $`g_1/F_1`$ given below (Eqs. 5 and 6). The $`g_2^{WW}`$ model is in good agreement with existing data . Using other reasonable models for $`g_2`$ that agree with existing data makes negligible changes to the extracted $`g_1/F_1`$ values due to suppression of the $`g_2`$ contribution by the factor $`2Mx/(2E\nu )`$. The $`g_1`$ and $`g_2`$ structure functions are related to the virtual photon asymmetry $`A_1=(g_1/F_1)\gamma ^2(g_2/F_1)`$ (which is bounded by $`\|A_1\|1`$). In this Letter we report new measurements of $`g_1^p`$ made using a 48.35 GeV polarized electron beam at SLAC. The new data extend to higher $`Q^2`$ (40 GeV<sup>2</sup>) and lower $`x`$ (0.014) than previous high statistics SLAC measurements . Combined with measurements of $`g_1^d`$ made in this same experiment using a <sup>6</sup>LiD target , we can extract $`g_1^n`$ and compare with E154 which measured $`g_1^n`$ at similar kinematics using a polarized <sup>3</sup>He target as a source of polarized neutrons. Longitudinally polarized electrons were produced by photoemission from a strained-lattice GaAs crystal. Beam pulses were typically 0.3 $`\mu `$s long, contained 2–4$`\times 10^9`$ electrons, and were delivered at a rate of 120 Hz. The helicity was selected randomly on a pulse-to-pulse basis to minimize instrumental asymmetries. The longitudinal beam polarization $`P_b`$ was measured using Møller scattering from thin, magnetized ferromagnetic foils, periodically inserted about 25 m before the polarized target used to measure $`g_1`$. Results from two detectors (one detecting a single final-state electron, the other detecting two electrons in coincidence) agreed within errors, yielding $`P_b=0.813\pm 0.020`$. As in E143 , the 3-cm-long polarized target cell contained pre-irradiated granules of <sup>15</sup>NH<sub>3</sub> immersed in liquid He at 1 K in a uniform magnetic field of 5 T. Microwaves near 140 GHz were used to drive the hyperfine transition which aligns (or anti-aligns) the nucleon spins with the magnetic field, producing proton polarizations of typically 90% in 10 to 20 minutes. The polarization slowly decreased due to radiation damage, and was periodically restored by annealing the target at about 80 K. The 2-3 mm diameter electron beam spot was rastered over the 3 cm<sup>2</sup> front surface of the target to uniformly distribute beam heating and radiation damage. To study possible experimental biases, the target polarization direction was periodically reversed using slight adjustments to the microwave frequency. Also, the direction of the magnetic field was reversed several times during the experiment. Final asymmetry results were consistent for the four polarization/field direction combinations. The target polarization $`P_t`$ was monitored with the same NMR Q-meter system as was used in experiment E143 . The E143 design of target cell was modified for this experiment to improve the target polarization (average value of $`P_t`$ was about 0.8 for E155) and this change had unforeseen effects on the performance of the NMR system when it was used to measure the proton polarization. Consequently, the NMR system was operated outside its design envelope, resulting in a significant degree of non-linear behavior. This problem is now qualitatively understood but insufficient information about the NMR RF circuit parameters is available to allow adequate corrections for these non-linear effects to be calculated. Therefore the polarization data, for the proton target only, was extracted using the observed dependence of the polarization on the integrated beam dose deposited in the target material obtained primarily from experiment E155x . This method leads to a larger systematic error in the proton polarization measurements (typically 7%) than would have been obtained using the standard NMR technique. It should be emphasized that this problem is unique to this particular set of proton experimental data and was eliminated in experiment E155x by a further target cell design change. Scattered electrons were detected in three independent magnetic spectrometers centered at angles of 2.75, 5.5, and 10.5 degrees. The two small angle spectrometers were the same as in E154 , while the large angle spectrometer was new for this experiment. It was composed of a single dipole magnet and two quadrupoles, and covered $`7<E^{}<20`$ GeV, $`9.6^{}<\theta <12.5^{}`$, and $`18<\varphi <18`$ mr, for a maximum solid angle of 1.5 msr at 8 GeV. Electrons were separated from a much larger flux of pions by using a gas Cherenkov counter and a segmented lead glass electromagnetic calorimeter. The experimental asymmetries $`A_{}`$ were determined from $$A_{}=\left(\frac{N_{}N_+}{N_{}+N_+}\right)\frac{C_N}{fP_bP_tf_{RC}}+A_{RC},$$ (3) where the target polarization is parallel to the beam direction, $`N_{}`$ ($`N_+`$) is the number of scattered electrons per incident charge for negative (positive) beam helicity, $`C_N0.985`$ is a correction factor for the polarized nitrogen nuclei, $`f`$ is the dilution factor representing the fraction of measured events originating from polarizeable hydrogen within the target, and $`f_{RC}`$ and $`A_{RC}`$ take into account radiative corrections. The dilution factor $`f`$ varied with $`x`$ between 0.13 and 0.17; it was determined from a detailed model of the number of measured counts expected from each component of the target, including <sup>15</sup>NH<sub>3</sub>, various windows, NMR coils, liquid helium, etc. A typical target contained about 13% free protons, 66% <sup>15</sup>N, 10% <sup>4</sup>He, 6% Al, and 5% Cu-Ni by weight. The relative systematic error in $`f`$ ranges from 2.2% to 2.6%. A correction to the asymmetries was made for hadrons misidentified as electrons (typically 2% of electron candidates, but up to 15% in the lowest $`x`$ bin of the 10.5 spectrometer). The correction used the asymmetry measured for a large sample of inclusive hadrons, which was found to be close to zero at all kinematics. An additional correction was made for electrons from pair-symmetric processes (such as $`e^+/e^{}`$ pair production from photons) measured by reversing the spectrometer polarity. The measured pair-symmetric $`A_{}`$ was consistent with zero at all kinematics, so the correction is equivalent to a dilution factor correction of typically 10% at the lowest $`E^{}`$ of each spectrometer, decreasing rapidly to a negligible correction at higher $`E^{}`$. Corrections were applied for the rate-dependence of the detector response, which changed the measured asymmetries by less than 1%. Corrections for kinematic resolution were generally a few percent or less, except for $`x>0.6`$ where corrections to the measured asymmetries were as large as 15%. The internal radiative corrections for $`A_{}`$ were evaluated using the formulae of Kuchto and Shumeiko . The cross sections entering the asymmetry were ‘externally radiated’ according to Tsai . Comparison of Born and fully radiated asymmetries allowed us to extract the asymmetry corrections $`f_{RC}`$ and $`A_{RC}`$. By splitting the radiative correction into these two parts, we can propagate consistently the experimental error to the extracted Born asymmetries for the corresponding kinematic bins, in the presence of ‘dilution’ from elastic and inelastic radiative tails. Previous analyses (including E143) have used values of $`f_{RC}`$ closer to unity by taking only quasi-elastic radiative tails into account, leading to smaller error bars at low $`x`$. Our new treatment is based on a definition of $`f_{RC}`$ that insures that the additive correction $`A_{RC}`$ is statistically independent from the data point to which it is applied. Values for $`f_{RC}`$ range from 0.45 at the lowest $`x`$-bin to greater than 0.9 for $`x>0.15`$, similar to the results in E154. However, the resulting net correction of the measured asymmetries is relatively small (0.01 to 0.02). The E155 radiative corrections are based on an iterative global fit to all available data, in which all previous SLAC data were re-corrected in a self-consistent way. The E155 results for $`g_1^p/F_1^p`$ and $`g_1^n/F_1^n`$ are shown in Figs. 1 and 2 as a function of $`Q^2`$ at eleven values of $`x`$, and are listed in Table I. The neutron results were obtained from the proton results and E155 deuteron results using $$g_1^n=\frac{g_1^d}{11.5\omega _D}\frac{F_1^n+F_1^p}{F_1^d}g_1^p$$ (4) For the deuteron D-state probability we use $`\omega _D=0.05\pm 0.01`$, and $`F_1^p/F_1^n`$ was obtained from the NMC fit . Slight changes to the data of Refs. were made to use the $`g_2^{WW}`$ model for $`g_2`$ instead of assuming $`A_2=0`$. Data from all experiments have been matched to the $`x`$ bins in Figs. 1 and 2 using the simple fit below for small bin centering corrections. For the present experiment, most systematic errors (beam polarization, target polarization, fraction of polarizeable nucleons in the target) for a given target are common to all data and correspond to an overall normalization error of about 7.6% for the proton data. The remaining systematic errors (model dependence of radiative corrections, model uncertainties for $`R(x,Q^2)`$, resolution corrections) vary smoothly with $`x`$ in a locally correlated fashion, ranging from a few percent for mid-range $`x`$ bins, up to 15% for the highest and lowest bins. Given the relatively large overall normalization uncertainty, the E155 data are in good agreement with the average of world data . If we were to allow an overall normalization factor for our proton data, we would find a value of $`1.08\pm 0.03`$(stat)$`\pm 0.07`$(syst). In any given $`x`$ bin, there is no evidence of strong $`Q^2`$ dependence for the ratio $`g_1/F_1`$. A simple parametric fit to world data with $`Q^2>1`$ GeV<sup>2</sup> (to mitigate possible higher twist contributions ) and missing mass $`W>2`$ GeV (to avoid complications from the resonance region), shown as the dashed curves in all three figures, is given by $$\frac{g_1^p}{F_1^p}=x^{0.700}(0.817+1.014x1.489x^2)(1\frac{0.04}{Q^2})$$ (5) $$\frac{g_1^n}{F_1^n}=x^{0.335}(0.0130.330x+0.761x^2)(1+\frac{0.13}{Q^2}).$$ (6) This fit has an acceptable $`\chi ^2`$ of 478 for 483 degrees of freedom. The coefficients of $`0.04\pm 0.06`$ ($`0.13\pm 0.45`$) for the overall proton (neutron) $`Q^2`$ dependence are small and consistent with zero. To examine the $`x`$ dependence of $`g_1`$ at fixed $`Q^2`$, we averaged the E155 results over $`Q^2`$ assuming the $`Q^2`$ dependence of the fit above and use $`F_1`$ from to obtain results for $`g_1^p`$ and $`g_1^n`$ at a fixed $`Q^2=5`$ GeV<sup>2</sup>, shown in Fig. 3. The proton data suggest $`g_1`$ is approximately constant or slightly rising as $`x0`$, but the neutron data are consistent with the trend of the E154 data to become increasingly negative at low $`x`$. The difference $`g_1^pg_1^n`$ (which enters into the Bjorken sum rule) is theoretically expected to be well-behaved as $`x0`$ compared to either $`g_1^p`$ or $`g_1^n`$. This is because if isospin is a good symmetry, the sea quark and gluon contributions cancel, leaving only the difference of $`u`$ and $`d`$ quark valence distributions. The errors on the present data are too large to clearly support or contradict this expectation (see Fig. 3c). The choice of low$`x`$ extrapolation has a large impact on the the evaluation of the first moment of $`g_1`$. To be consistent with other analyses of $`g_1`$, we have made a NLO pQCD fit in the $`\overline{MS}`$ scheme to all recent data, using assumptions similar to those in . The polarized parton distributions were parameterized as $$\mathrm{\Delta }f(x,Q_0^2)=A_fx^{\alpha _f}f(x,Q_0^2),$$ (7) where $`\mathrm{\Delta }f=\mathrm{\Delta }u_v`$, $`\mathrm{\Delta }d_v`$, $`\mathrm{\Delta }\overline{Q}`$, and $`\mathrm{\Delta }G`$ are the polarized valence, sea, and gluon distributions, and the $`f(x,Q_0^2)`$ are the unpolarized parton distributions at $`Q_0^2=0.40`$ GeV<sup>2</sup> from Ref. . The positivity constraint $`\left\|\mathrm{\Delta }f\right\|<f`$ was imposed in this fit, as well as the requirement $`\alpha _f>0`$. The sea quark distributions were parameterized as $`\mathrm{\Delta }\overline{Q}=\frac{1}{2}(\mathrm{\Delta }\overline{u}+\mathrm{\Delta }\overline{d})+\frac{1}{5}\mathrm{\Delta }\overline{s}`$. We assumed a symmetric quark sea for this analysis. We have not fixed the normalization of the non-singlet distributions, so that the fit results test the Bjorken sum rule. However, $`\alpha _s(M_Z^2)`$ has been fixed at 0.114 for consistency with the unpolarized distributions that were used . The fit results are: $`A_u=0.95`$, $`A_d=0.42`$, $`A_Q=0.01`$, $`A_g=0.50`$, $`\alpha _u=0.57`$, $`\alpha _d=0.0`$, $`\alpha _Q=1.00`$, and $`\alpha _g=0.02`$. The overall $`\chi ^2`$/d.f. is $`1.10`$ using statistical errors only. Evaluations of the fit are plotted as the solid curves in Figs. 1-3, and indicate only a slight dependence on $`Q^2`$ for $`g_1/F_1`$ in the $`x`$ region where there are high statistics data. For $`x<0.014`$, the proton and neutron fits become increasingly negative at fixed $`Q^2`$ (see Fig. 3), although the difference stays closer to zero and makes only a small contribution to the Bjorken sum rule. Using the NLO pQCD fit, we find the quark singlet contribution $`\mathrm{\Delta }\mathrm{\Sigma }=0.23\pm 0.04`$(stat)$`\pm 0.06`$(syst) at $`Q^2=5`$ GeV<sup>2</sup>, well below the Ellis-Jaffe prediction of 0.58. We find $`\mathrm{\Gamma }_1^p=0.118\pm 0.004\pm 0.007`$, $`\mathrm{\Gamma }_1^n=0.058\pm 0.005\pm 0.008`$, and $`\mathrm{\Gamma }_1^p\mathrm{\Gamma }_1^n=0.176\pm 0.003\pm 0.007`$, in good agreement with the Bjorken sum rule prediction of $`0.182\pm 0.005`$ evaluated with up to third order corrections in $`\alpha _s`$ . For the first moment of the gluon distribution we obtain $`\mathrm{\Delta }G=1.6\pm 0.8\pm 1.1`$. The error on this quantity is too large to significantly constrain the gluon contribution to the nucleon spin sum rule. In summary, the new data on $`g_1^p`$ and $`g_1^n`$ extend the range of high statistics electron scattering results to lower $`x`$ and higher $`Q^2`$ than previous data, improving the errors obtained from NLO pQCD fits to world data. The Bjorken sum rule prediction is validated within errors, while the extracted quark singlet contribution is small at approximately 0.2. This work was supported by the Department of Energy (TJNAF, FIU, Massachusetts, ODU, SLAC, Stanford, Virginia, Wisconsin, and William and Mary); by the National Science Foundation (American, Kent, Michigan, and ODU); by the Schweizersche Nationalfonds (Basel); by the Kent State University Research Council (GGP); by the Commonwealth of Virginia (Virginia); by the Centre National de la Recherche Scientifique and the Commissariat a l’Energie Atomique (French groups).
warning/0007/hep-th0007008.html	ar5iv	text	# Quantum Mechanical Symmetries and Topological Invariants ## 1 Introduction Supersymmetric quantum mechanics was originally introduced as a toy model used to study some of the features of supersymmetric field theories . This simple toy model has, however, proven to be a very useful tool in dealing with a variety of problems in quantum and statistical mechanics . Supersymmetric quantum mechanics has also been used to derive some of the very basic results of differential topology. Among these are the supersymmetric derivation of the Morse inequalities and supersymmetric proofs of the Atiyah-Singer index theorem . The relationship between supersymmetry and topological invariants such as the indices of the elliptic operators is our main motivation for seeking general quantum mechanical symmetries with topological properties similar to those of supersymmetry. Various generalizations of supersymmetry have been considered in the literature . Among these only extended and generalized supersymmetries and a certain class of $`p=2`$ parasupersymmetries are known to share the topological characteristics of supersymmetry. The strategy pursued in this article is as follows. First, we introduce the notion of a topological symmetry (TS) by formulating a set of basic principles that ensure the desired topological properties. Then, we investigate the underlying algebraic structure of these symmetries. This is necessary for seeking a mathematical interpretation of the corresponding topological invariants. We have recently reported our preliminary results on $`ZZ_2`$-graded TSs of type $`(1,1)`$ and $`(2,1)`$ in . The purpose of the present article is to generalize the results of to arbitrary $`ZZ_n`$-graded TSs. The organization of the article is as follows. In section 2, we give the definition of a general $`ZZ_n`$-graded TS and introduce the associated topological invariants. In section 3, we consider the case of $`n=2`$ and derive the algebra of a $`ZZ_2`$-graded TS of arbitrary type $`(m_+,m_{})`$. In particular, we show that for $`m_{}=1`$, the algebra can be reduced to the algebra of supersymmetry or $`p=2`$ parasupersymmetry. In section 4, we consider the $`ZZ_n`$-graded TSs. Here we discuss the properties of the grading operator and derive the algebra of $`ZZ_n`$-graded TSs of arbitrary type $`(m_1,m_2,\mathrm{},m_n)`$. In section 5, we comment on the mathematical interpretation of the topological invariants associated with TSs. In section 6, we study some concrete examples of quantum systems possessing TSs. In section 7, we summarize our results and present our concluding remarks. The appendix includes the proof of some of the mathematical results that we use in our analysis. ## 2 $`ZZ_n`$-Graded TSs In order to describe the concept of a topological symmetry, we first give some basic definitions. In the following we shall only consider the quantum systems with a self-adjoint Hamiltonian $`H`$. We shall further assume that all the energy levels are at most finitely degenerate. * Definition 1: Let $`n`$ be an integer greater than 1, $``$ denote the Hilbert space of a quantum system, and $`_1,_2,\mathrm{},_n`$ be (nontrivial) subspaces of $``$. Then a state vector is said to have definite color $`c_{\mathrm{}}`$ iff it belongs to $`_{\mathrm{}}`$. * Definition 2: A quantum system is said to be $`ZZ_n`$-graded iff its Hilbert space is the direct sum of $`n`$ of its (nontrivial) subspaces $`_{\mathrm{}}`$, and its Hamiltonian has a complete set of eigenvectors with definite color. * Definition 3: Let $`m_{\mathrm{}}`$ be positive integers for all $`\mathrm{}\{1,2,\mathrm{},n\}`$, and $`m:=_{\mathrm{}=1}^nm_{\mathrm{}}`$. Then a quantum system is said to possess a $`ZZ_n`$-graded topological symmetry of type $`(m_1,m_2,\mathrm{},m_n)`$ iff the following conditions are satisfied. + The quantum system is $`ZZ_n`$-graded; + The energy spectrum is nonnegative; + For every positive energy eigenvalue $`E`$, there is a positive integer $`\lambda _E`$ such that $`E`$ is $`\lambda _Em`$-fold degenerate, and the corresponding eigenspaces are spanned by $`\lambda _Em_1`$ vectors of color $`c_1`$, $`\lambda _Em_2`$ vectors of color $`c_2`$, $`\mathrm{}`$, and $`\lambda _Em_n`$ vectors of color $`c_n`$. * Definition 4: A topological symmetry is said to be uniform iff for all positive energy eigenvalues $`E`$, $`\lambda _E=1`$. * Theorem 1: Consider a quantum system possessing a $`ZZ_n`$-graded topological symmetry of type $`(m_1,m_2,\mathrm{},m_n)`$, and let $`n_{\mathrm{}}^{(0)}`$ denote the number of zero-energy states of color $`c_{\mathrm{}}`$. Then for all $`i,j\{1,2,\mathrm{},n\}`$, the integers $$\mathrm{\Delta }_{ij}:=m_in_j^{(0)}m_jn_i^{(0)}$$ (1) remain invariant under continuous symmetry-preserving changes of the quantum system.<sup>1</sup><sup>1</sup>1Here we consider quantum systems whose energy spectrum depends on a set of continuous parameters. These parameters may be identified with coupling constants or geometric quantities entering the definition of the Hamiltonian and/or the Hilbert space. A continuous change of the system corresponds to a continuous change of these parameters. * Proof: The proof of this theorem is essentially the same as the proof of the topological invariance of the Witten index of supersymmetry. A symmetry-preserving continuous change of the quantum system will preserve the particular degeneracy and grading structures of positive energy levels. Because under such a change an initial zero energy eigenstate can only become a positive energy eigenstate and vice versa, the only possible change in the number of zero energy eigenstates are the ones involving the changes of $`n_i^{(0)}`$ of the form $$n_i^{(0)}\stackrel{~}{n}_i^{(0)}:=n_i^{(0)}+km_i,$$ (2) where $`k`$ is an integer greater than or equal to $`n_i^{(0)}/m_i`$. Moreover, such a change must occur simultaneously for all $`n_i^{(0)}`$’s, i.e., the transformation (2) is valid for all $`i\{1,2,\mathrm{},n\}`$. Therefore under such a symmetry-preserving change of the system, $$\mathrm{\Delta }_{ij}\stackrel{~}{\mathrm{\Delta }}_{ij}:=m_i\stackrel{~}{n}_j^{(0)}m_j\stackrel{~}{n}_i^{(0)}=m_i(n_j^{(0)}+km_j)m_j(n_i^{(0)}+km_i)=m_in_j^{(0)}m_jn_i^{(0)}=\mathrm{\Delta }_{ij},$$ i.e., $`\mathrm{\Delta }_{ij}`$’s remain invariant. $`\mathrm{}`$ A direct consequence of Theorem 1 is that (the value of) any function of $`\mathrm{\Delta }_{ij}`$’s is a topological invariant of the system. In particular, $`\mathrm{\Delta }_{ij}`$’s are the basic topological invariants. A typical example of a derived invariant is $$\mathrm{\Delta }:=\frac{1}{2}\underset{i,j=1}{\overset{n}{}}(\mathrm{\Delta }_{ij})^2.$$ (3) Note that $`\mathrm{\Delta }`$ is a measure of the existence of the zero-energy states. Let us next observe that the topological symmetries are simple generalizations of supersymmetry. First we recall that the Hilbert space of a supersymmetric system is $`ZZ_2`$-graded. We can relate the $`ZZ_2`$-grading of the Hilbert space $``$ to the existence of a ‘parity’ or ‘grading’ operator $`\tau :`$ satisfying $`\tau ^2`$ $`=`$ $`1,`$ (4) $`\tau ^{}`$ $`=`$ $`\tau ,`$ (5) $`[H,\tau ]`$ $`=`$ $`0.`$ (6) Here we use a $``$ to denote the adjoint of the corresponding operator. We can identify the subspaces $`_1`$ and $`_2`$ with the eigenspaces of $`\tau `$, $$_1=_+,_2=_{},_\pm :=\{\psi \|\tau \psi =\pm \psi \}.$$ (7) Now, consider the superalgebra $`[H,𝒬]`$ $`=`$ $`0,`$ (8) $`{\displaystyle \frac{1}{2}}\{𝒬,𝒬^{}\}`$ $`=`$ $`H,`$ (9) $`𝒬^2`$ $`=`$ $`0,`$ (10) of supersymmetric quantum mechanics with one nonself-adjoint symmetry generator $`𝒬`$ satisfying $$\{𝒬,\tau \}=0.$$ (11) It is well-known that using the superalgebra (8) – (10) together with the properties of the grading operator (4) – (6) and the symmetry generator (11), one can show that the conditions a) – c) of Definition 3, with $`n=2`$ and $`m_1=m_2=1`$, are satisfied. Therefore, supersymmetry is a $`ZZ_2`$-graded TS of type $`(1,1)`$. For a $`ZZ_2`$-graded TS of type $`(1,1)`$, there is a single basic topological invariant namely $`\mathrm{\Delta }_{11}`$. This is precisely the Witten index. ## 3 Algebraic Structure of $`ZZ_2`$-Graded TSs In this section we shall explore the algebraic structure of the $`ZZ_2`$-graded TSs that fulfil the following conditions. * The $`ZZ_2`$-grading is achieved by a grading operator $`\tau `$ satisfying (4) – (6); * There is a single nonself-adjoint symmetry generator $`𝒬`$; * $`𝒬`$ is an odd operator, i.e., it satisfies Eq. (11). We shall only treat the case of the uniform TSs. The algebraic structure of nonuniform TSs is easily obtained from that of the uniform topological symmetries (UTSs). In fact, the algebraic relations defining uniform and nonuniform TSs of the same type are identical. In order to obtain the algebraic structures that support TSs, we shall use the information on the degeneracy structure of the corresponding systems and the properties of the grading operator and the symmetry generator to construct matrix representations of the relevant operators in the energy eigenspaces $`_E`$ with positive eigenvalue $`E`$. We shall use the notation $`O^E`$ for the restriction of an operator $`O`$ onto the eigenspace $`_E`$. Throughout this article $`E`$ stands for a positive energy eigenvalue. The zero-energy eigenspace (kernel of $`H`$) is denoted by $`_0`$. In view of Eqs. (6) and (8), $`\tau ^E`$ and $`𝒬^E`$ are $`m\times m`$ matrices acting in $`_E`$. We also have the trivial identity: $`H^E=EI_m`$, where $`I_m`$ denotes the $`m\times m`$ unit matrix. Next, we introduce the self-adjoint symmetry generators $$Q_1:=\frac{1}{\sqrt{2}}(𝒬+𝒬^{})\mathrm{and}Q_2:=\frac{i}{\sqrt{2}}(𝒬𝒬^{})$$ (12) where $`i:=\sqrt{1}`$. Note that because $`\tau `$ is self-adjoint, we have $$\{Q_j,\tau \}=0\mathrm{for}j\{1,2\}.$$ (13) Now in view of Eq. (6), we can choose a basis in $`_E`$ in which $`\tau `$ is diagonal. Then using Eqs. (4), (12), (13), and the self-adjointness of $`Q_j`$, we obtain the following matrix representations for $`\tau ^E`$, $`Q_j^E`$, and $`𝒬^E`$. $`\tau ^E`$ $`=`$ $`\mathrm{diag}(\underset{m_+\mathrm{times}}{\underset{}{1,1,\mathrm{},1}},\underset{m_{}\mathrm{times}}{\underset{}{1,1,\mathrm{},1}}),`$ (14) $`Q_j^E`$ $`=`$ $`\left(\begin{array}{cc}0& A_j\\ A_j^{}& 0\end{array}\right),`$ (17) $`𝒬^E`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}\left(\begin{array}{cc}0& A_1+iA_2\\ A_1^{}+iA_2^{}& 0\end{array}\right),`$ (20) where ‘diag$`(\mathrm{})`$’ stands for a diagonal matrix with diagonal entries ‘$`\mathrm{}`$’, $`0`$’s denote appropriate zero matrices, and $`A_j`$ are $`m_+\times m_{}`$ complex matrices. The next step is to find general identities satisfied by $`Q_j^E`$ and $`𝒬^E`$ for all $`E>0`$. In order to derive the simplest such identities we appeal to the Cayley-Hamilton theorem of linear algebra. This theorem states that an $`m\times m`$ matrix $`Q`$ satisfies its characteristic equations, $`𝒫_Q(Q)=0`$, where $`𝒫_Q(x):=det(xI_mQ)`$ is the characteristic polynomial for $`Q`$. Using this theorem we can prove the following lemma. The proof is given in the appendix. * Lemma 1: Let $`m_\pm `$ be positive integers, $`m:=m_++m_{}`$, and $`Q`$ be an $`m\times m`$ matrix of the form: $$Q=\left(\begin{array}{cc}0& X\\ Y& 0\end{array}\right),$$ (21) where $`X`$ and $`Y`$ are $`m_+\times m_{}`$ and $`m_{}\times m_+`$ complex matrices. Let $`𝒫_{XY}(x)`$ and $`𝒫_{YX}(x)`$ denote the characteristic polynomials for $`XY`$ and $`YX`$, respectively. Then $`𝒫_{YX}(Q^2)Q=𝒫_{XY}(Q^2)Q=0`$. Furthermore, if $`m_+=m_{}`$, then $`𝒫_{YX}(Q^2)=𝒫_{XY}(Q^2)=0`$. Applying this lemma to $`Q_j^E`$ and $`𝒬^E`$, we find for $`m_+=m_{}`$ $`𝒫_j[(Q_j^E)^2]`$ $`=`$ $`0,`$ (22) $`𝒫[(𝒬^E)^2]`$ $`=`$ $`0,`$ (23) and for $`m_+>m_{}`$ $`𝒫_j[(Q_j^E)^2]Q_j^E`$ $`=`$ $`0,`$ (24) $`𝒫[(𝒬^E)^2]𝒬^E`$ $`=`$ $`0,`$ (25) where $`𝒫_j(x)`$ and $`𝒫(x)`$ denote the characteristic polynomials of $`A_j^{}A_j`$ and $`(A_1^{}+iA_2^{})(A_1+iA_2)/2`$, respectively, For $`m_{}>m_+`$, the roles of $`A_j`$ and $`A_j^{}`$ are interchanged. We shall, therefore, restrict our attention to the case where $`m_+m_{}`$. We can write Eqs. (22) – (25) in terms of the roots of the characteristic polynomials $`𝒫_j(x)`$ and $`𝒫(x)`$. This yields $`\left[(Q_j^E)^2\mu _{j1}^EI_m\right]\left[(Q_j^E)^2\mu _{j2}^EI_m\right]\mathrm{}\left[(Q_j^E)^2\mu _{jm_{}}^EI_m\right](Q_j^E)^{1\delta (m_+,m_{})}`$ $`=`$ $`0,`$ (26) $`\left[(𝒬^E)^2\kappa _1^EI_m\right]\left[(𝒬^E)^2\kappa _2^EI_m\right]\mathrm{}\left[(𝒬^E)^2\kappa _m_{}^EI_m\right](𝒬^E)^{1\delta (m_+,m_{})}`$ $`=`$ $`0,`$ (27) where $`\mu _j\mathrm{}^E`$ and $`\kappa _{\mathrm{}}^E`$ are the roots<sup>2</sup><sup>2</sup>2Note that the roots are not necessarily distinct. of $`𝒫_j(x)`$ and $`𝒫(x)`$, respectively, and $$\delta (m_+,m_{})=\delta _{m_+,m_{}}:=\{\begin{array}{ccc}1& \mathrm{for}& m_+=m_{}\\ 0& \mathrm{for}& m_+m_{}.\end{array}$$ In order to promote Eqs. (26) and (27) to operator relations, we introduce the operators $`M_j\mathrm{}`$ and $`𝒦_{\mathrm{}}`$ (for each $`j\{1,2\}`$ and $`\mathrm{}\{1,2,\mathrm{},m_{}\}`$) which commute with $`H`$ and have the representations: $$M_j\mathrm{}^E=\mu _j\mathrm{}^EI_m\mathrm{and}𝒦_{\mathrm{}}^E=\kappa _{\mathrm{}}^EI_m$$ (28) in $`_E`$. We then deduce from Eqs. (26), (27) and (28) that $`M_j\mathrm{}`$ and $`𝒦_{\mathrm{}}`$ must commute with $`\tau `$ and $`Q_j`$. Furthermore, they should satisfy the algebra $`(Q_j^2M_{j1})(Q_j^2M_{j2})\mathrm{}(Q_j^2M_{jm_{}})Q_j^{1\delta (m_+,m_{})}`$ $`=`$ $`0,`$ (29) $`(𝒬^2𝒦_1)(𝒬^2𝒦_2)\mathrm{}(𝒬^2𝒦_m_{})𝒬^{1\delta (m_+,m_{})}`$ $`=`$ $`0.`$ (30) Note that the roots $`\kappa _{\mathrm{}}^E`$ and $`\mu _j\mathrm{}^E`$ are defined using the matrices $`A_j`$. This suggests that the operators $`M_j\mathrm{}`$ and $`𝒦_{\mathrm{}}`$ are not generally independent. Furthermore, because $`A_j^{}A_j`$ are Hermitian matrices, the roots $`\mu _j\mathrm{}^E`$, which are in fact the eigenvalues of $`A_j^{}A_j`$, are real. This in turn suggests that the operators $`M_j\mathrm{}`$ are self-adjoint. In summary, the algebra of general $`ZZ_2`$-graded topological symmetry of type $`(m_+,m_{})`$ which is generated by one odd nonself-adjoint generator $`𝒬`$ is given by Eqs. (29) and (30) where $`m_+`$ is assumed (without loss of generality) not to be smaller than $`m_{}`$, the operators $`M_j\mathrm{}`$ and $`𝒦_{\mathrm{}}`$ commute with $`H`$, $`\tau `$ and $`𝒬`$, and $`M_j\mathrm{}`$ are self-adjoint. Moreover, $`M_j\mathrm{}`$ and $`𝒦_{\mathrm{}}`$ have similar degeneracy structure as the Hamiltonian<sup>3</sup><sup>3</sup>3The degeneracy structure of these operators will be the same as that of the Hamiltonian, if their eigenvalues $`\mu _j\mathrm{}^E`$ and $`\kappa _{\mathrm{}}^E`$ are distinct for different $`E`$. (at least for positive energy eigenvalues). In particular, it might be possible to express $`H`$ as a function of $`M_j\mathrm{}`$ and $`𝒦_{\mathrm{}}`$. In order to elucidate the role of the operators $`M_j\mathrm{}`$ and $`𝒦_{\mathrm{}}`$ and their relation to the Hamiltonian, we shall next consider the $`ZZ_2`$-graded UTSs of type $`(m_+,1)`$. If $`m_{}=1`$, then $`A_j^{}A_j`$ and $`(A_1^{}+iA_2^{})(A_1+iA_2)/2`$ are respectively real and complex scalars. In this case, $`𝒫_j(x)=x\mu _j^E`$ and $`𝒫(x)=x\kappa ^E`$, where $$\mu _j^E=A_j^{}A_j,\kappa ^E=\frac{1}{2}(A_1^{}+iA_2^{})(A_1+iA_2)=\frac{1}{2}[(A_1^{}A_1A_2^{}A_2)+i(A_1^{}A_2+A_2^{}A_1)],$$ (31) and the algebra (29) and (30) takes the form $`(Q_j^2M_j)Q_j^{1\delta (m_+,1)}`$ $`=`$ $`0,`$ (32) $`(𝒬^2𝒦)𝒬^{1\delta (m_+,1)}`$ $`=`$ $`0.`$ (33) Here we have used the abbreviated notation: $`M_j=M_{j1}`$ and $`𝒦=𝒦_1`$. Next, we define the self-adjoint operators $$K_1=𝒦+𝒦^{}\mathrm{and}K_2=i(𝒦𝒦^{}).$$ (34) In view of Eqs. (28) and (31) we have $$M_2=M_1K_1.$$ (35) In the following we shall consider the cases $`m_+=1`$ and $`m_+>1`$ separately. ### 3.1 $`ZZ_2`$-Graded UTS of Type $`(1,1)`$ Setting $`m_+=1`$ in Eqs. (32) and (33), we find $`Q_j^2`$ $`=`$ $`M_j,`$ (36) $`𝒬^2`$ $`=`$ $`𝒦.`$ (37) If we express $`𝒬`$ in terms of $`Q_j`$ and use Eqs. (34) and (35), we can write Eqs. (36) and (37) in the form $`Q_1^2`$ $`=`$ $`M_1,`$ (38) $`Q_2^2`$ $`=`$ $`M_1K_1,`$ (39) $`\{Q_1,Q_2\}`$ $`=`$ $`K_2.`$ (40) Now, we observe that Eqs. (38) – (40) remain form-invariant under the linear transformations of the form $$\begin{array}{c}Q_1\stackrel{~}{Q}_1=aQ_1+bQ_2,\\ Q_2\stackrel{~}{Q}_2=cQ_1+dQ_2,\end{array}$$ (41) where $`a,b,c`$ and $`d`$ are self-adjoint operators commuting with all other operators. More specifically, $`\stackrel{~}{Q}_j`$ satisfy Eqs. (38) – (40) provided that $`M_1`$ and $`K_j`$ are transformed according to $`M_1`$ $``$ $`\stackrel{~}{M}_1:=(a^2+b^2)M_1b^2K_1+abK_2,`$ (42) $`K_1`$ $``$ $`\stackrel{~}{K}_1:=(a^2+b^2c^2d^2)M_1++(d^2b^2)K_1+(abcd)K_2,`$ (43) $`K_2`$ $``$ $`\stackrel{~}{K}_2:=2(ac+bd)M_12bdK_1+(ad+bc)K_2.`$ (44) In particular, there are transformations of the form (41) for which $`\stackrel{~}{K}_j=0`$. These correspond to the choices for $`a,b,c`$ and $`d`$ that satisfy (either of) $$\frac{a+ic}{b+id}=\frac{K_2}{2M_1}\pm i\sqrt{1\frac{K_1}{M_1}\frac{K_2^2}{4M_1^2}}.$$ (45) One can use the representations of $`K_j`$ and $`M_1`$ in the eigenspaces $`_E`$ to show that the terms in the square root in (45) yield a positive self-adjoint operator, provided that the kernel of $`M_1`$ is a subspace of the zero-energy eigenspace $`_0`$. The above analysis shows that we can reduce the general algebra (38) – (40) to the special case where $`K_j=0`$. Writing this algebra in terms of $`𝒬`$, we obtain the superalgebra (8) – (10) with $`M_1`$ replacing $`H`$. In other words, if we identify the Hamiltonian with $`M_1`$, which we can always do, the algebra of $`ZZ_2`$-graded topological symmetry of type $`(1,1)`$ reduces to that of supersymmetry. ### 3.2 $`ZZ_2`$-Graded UTSs of Type $`(m_+,1)`$ with $`m+>1`$ If $`M_+>1`$, then Eqs. (32) and (33) take the form $`Q_j^3`$ $`=`$ $`M_jQ_j,`$ (46) $`𝒬^3`$ $`=`$ $`𝒦𝒬.`$ (47) Again we express $`𝒬`$ in terms of $`Q_j`$ and use Eqs. (34) and (35) to write (46) and (47) in the form $`Q_1^3=M_1Q_1,`$ (48) $`Q_2^3=(M_1K_1)Q_2,`$ (49) $`Q_2Q_1Q_2+\{Q_1,Q_2^2\}=(M_1K_1)Q_1+K_2Q_2,`$ (50) $`Q_1Q_2Q_1+\{Q_2,Q_1^2\}=M_1Q_2+K_2Q_1.`$ (51) It is remarkable that these relations are also invariant under the transformations (41) and (42) – (44). Therefore, again we can reduce our analysis to the special case where $`K_j=0`$. Substituting zero for $`K_j`$ in Eqs. (48) – (51), and writing them in terms of $`𝒬`$, we obtain $`[M_1,𝒬]=0,`$ (52) $`\{𝒬^2Q^{}\}+𝒬𝒬^{}𝒬=2M_1𝒬,`$ (53) $`𝒬^3=0.`$ (54) This is precisely the algebra of $`p=2`$ parasupersymmetry of Rubakov and Spiridonov with $`H`$ replaced by $`M_1/2`$. Hence, if we identify $`H`$ with $`M_1/2`$, which we can always do, the algebra of $`ZZ_2`$-graded topological symmetry of type $`(m_+,1)`$ with $`m_+>1`$ reduces to that of the $`p=2`$ parasupersymmetry. As shown in Ref. , one can use the algebra (52) – (54) of $`p=2`$ parasupersymmetry and properties of the grading operator (4) – (6) and (13) to obtain the general degeneracy structure of a $`p=2`$ parasupersymmetric system. In general the algebra of $`p=2`$ parasupersymmetry does not imply the particular degeneracy structure of the $`ZZ_2`$-graded UTS of type $`(m_+,1)`$, even for $`m_+=2`$. Therefore, the $`ZZ_2`$-graded UTS of type $`(2,1)`$ is a subclass of the general $`p=2`$ parasupersymmetries. As argued in Refs. and , these are parasupersymmetries for which an analog of the Witten index can be defined. In Ref. it is also shown that the positive energy eigenvalues of a $`p=2`$ parasupersymmetric system can at most be triply degenerate, provided that the eigenvalues of $`Q_1^E`$ for all $`E>0`$ are nondegenerate. This means that the $`ZZ_2`$-graded TSs of type $`(m_+,1)`$ with $`m_+>2`$ occur only if $`Q_1^E`$ have degenerate eigenvalues for all $`E>0`$. This suggests the presence of further (even) symmetry generators $`L_a`$ which would commute with $`Q_1`$ and label the basis eigenvectors within the degeneracy subspaces of $`Q_1`$. The existence of these generators is an indication that the $`ZZ_2`$-graded TSs of type $`(m_+,1)`$ with $`m_+>2`$ are not uniform. ### 3.3 Special $`ZZ_2`$-Graded TSs The analysis of the $`ZZ_2`$-graded TSs of type $`(m_+,1)`$ shows that the corresponding algebras can be reduced to a simplified special case by a redefinition of the symmetry generators $`Q_j`$. This raises the question whether this is also possible for the general $`ZZ_2`$-graded TSs of type $`(m_+,m_{})`$. The reduction made in the case of $`ZZ_2`$-graded TSs of type $`(m_+,1)`$ has its roots in the form of the matrix representation of the corresponding symmetry generators in $`_E`$. For a $`ZZ_2`$-graded UTS of arbitrary type $`(m_+,m_{})`$, a similar reduction, which eliminates the operators $`𝒦_{\mathrm{}}`$ in Eq. (30), is possible, if we can find a transformation $`Q_j\stackrel{~}{Q}_j`$ which satisfies the following conditions. * The transformed generators $`\stackrel{~}{Q}_j`$ have the representation $$\stackrel{~}{Q}_j^E=\left(\begin{array}{cc}0& \stackrel{~}{A}_j\\ \stackrel{~}{A}_j^{}& 0\end{array}\right)$$ in $`_E`$. * For all energy eigenvalues $`E>0`$, the corresponding matrices $`\stackrel{~}{A}_j`$ which define $`\stackrel{~}{Q}_j^E`$ satisfy $`\stackrel{~}{A}_2=U\stackrel{~}{A}_1`$, where $`U`$ is an $`m_+\times m_+`$ unitary and anti-Hermitian matrix. The first condition is necessary for the invariance of the algebra (29) – (30). It is satisfied by the linear transformations: $`Q_j^E\stackrel{~}{Q}_j^E=T_jQ_j^ET_j^{}`$ where $`T_j`$ are $`m\times m`$ matrices of the form $$T_j=\left(\begin{array}{cc}T_{j+}& 0\\ 0& T_j\end{array}\right),$$ and $`T_{j\pm }`$ are $`m_\pm \times m_\pm `$ matrices. Under such a transformation $`A_j`$ transform according to $`A_j\stackrel{~}{A}_j:=T_{j+}A_jT_j^{}`$. The second condition implies that the transformed matrices $`\stackrel{~}{A}_j`$ satisfy $$(\stackrel{~}{A}_1^{}+i\stackrel{~}{A}_2^{})(\stackrel{~}{A}_1+i\stackrel{~}{A}_2)=A_1^{}(I_{m_+}+iU^{})(I_{m_+}+iU)A_1=A_1^{}[(I_{m_+}U^{}U)+i(U+U^{})]A_1=0,$$ (55) where we have used the unitarity and anti-Hermiticity of $`U`$. The latter relation indicates that for the transformed system $`\kappa _{\mathrm{}}^E=0`$ and $`\mu _2\mathrm{}^E=\mu _1\mathrm{}^E`$, for all $`E`$ and $`\mathrm{}`$. Hence $`𝒦_{\mathrm{}}`$ can be identified with the zero operator and $`M_2\mathrm{}=M_1\mathrm{}=:M_{\mathrm{}}`$. Furthermore, one can check that (55) implies $`(\stackrel{~}{𝒬}^E)^{3\delta (m_+,m_{})}=0`$. Therefore, the transformed generators satisfy $`(\stackrel{~}{Q}_j^2M_1)(\stackrel{~}{Q}_j^2M_2)\mathrm{}(\stackrel{~}{Q}_j^2M_m_{})\stackrel{~}{Q}_j^{1\delta (m_+,m_{})}=0,`$ (56) $`\stackrel{~}{𝒬}^{3\delta (m_+,m_{})}=0.`$ (57) We shall term such $`ZZ_2`$-graded TSs, the special $`ZZ_2`$-graded TSs. ### 3.4 $`ZZ_2`$-graded TSs with one Self-Adjoint Generator The above analysis of $`ZZ_2`$-graded TSs can be easily applied to the case where there is a single self-adjoint generator $`Q`$. In fact, one can read off the corresponding algebra from Eqs. (29) and (30). The result is $$(Q^2M_1)(Q^2M_2)\mathrm{}(Q^2M_m_{})Q^{1\delta (m_+,m_{})}=0,$$ (58) where $`M_{\mathrm{}}`$ are self-adjoint operators commuting with $`H,Q`$ and $`\tau `$. For $`m_\pm =1`$, this equation reduces to that of the $`N=\frac{1}{2}`$ supersymmetry, provided that we identify $`M_1`$ with the Hamiltonian $`H`$. ## 4 Algebraic Structure of $`ZZ_n`$-Graded TSs In the preceding section, we have used Definition 3 and the properties (4) – (6) and (11) of the $`ZZ_2`$-grading operator $`\tau `$ to obtain the algebra of the $`ZZ_2`$-graded TSs with one generator. The main guideline for postulating these properties were Definition 2 and the known grading structure of supersymmetric systems. Similarly, in order to obtain the algebraic structure of the $`ZZ_n`$-graded TSs, with $`n>2`$, we must first postulate the existence of an appropriate $`ZZ_n`$-grading operator. In view of Definitions 1, 2 and 3, such a grading operator $`\tau `$ must commute with the Hamiltonian — Eq. (6) holds — and have $`n`$-distinct eigenvalues $`c_{\mathrm{}}`$ with eigenspaces $`_{\mathrm{}}`$ satisfying $`=_1_2\mathrm{}_n`$. The simplest choice for $`c_{\mathrm{}}`$ are the $`n`$-th roots of unity, e.g., $`c_{\mathrm{}}=q^{\mathrm{}}`$ where $`q=e^{2\pi i/n}`$. This choice suggests the following generalization of Eqs. (4) and (5). $`\tau ^n`$ $`=`$ $`1`$ (59) $`\tau ^{}`$ $`=`$ $`\tau ^1`$ (60) Note that for $`n=2`$, according to Eq. (4), $`\tau ^1=\tau `$ and Eq. (60) coincides with (5). In the following we shall only consider $`ZZ_n`$-graded TSs with one symmetry generator $`𝒬`$. In order to proceed along the same lines as in the case of $`n=2`$, we need to impose a grading condition on $`𝒬`$ similar to (11). We first consider the simplest case, namely $`ZZ_n`$-graded UTS of type $`(1,1,\mathrm{},1)`$. ### 4.1 $`ZZ_n`$-Graded UTS of Type $`(1,1,\mathrm{},1)`$ For $`n=2`$, condition (11) implies that the action of $`𝒬`$ on a definite color (parity) state vector changes its color (parity) – a bosonic state changes to a fermionic state and vice versa. The simplest generalization of this statement to the case $`n>2`$ is that the action of $`𝒬`$ must change the color of a definite color state by one unit, i.e., $$\tau \psi =c_{\mathrm{}}\psi \mathrm{implies}\tau (𝒬\psi )=c_{\mathrm{}+1}Q\psi .$$ (61) This condition is consistent with Eq. (59). If we impose this condition on the representations $`𝒬^E`$ and $`\tau ^E`$ in $`_E`$, then in a basis in which $`\tau ^E`$ is diagonal we have $$\tau ^E=\mathrm{diag}(q,q^2,\mathrm{},q^{n1},q^n=1),$$ (62) and $$𝒬^E=\left(\begin{array}{ccccccc}0& 0& 0& \mathrm{}& 0& 0& a_n\\ a_1& 0& 0& \mathrm{}& 0& 0& 0\\ 0& a_2& 0& \mathrm{}& 0& 0& 0\\ \mathrm{}& \mathrm{}& \mathrm{}& & \mathrm{}& \mathrm{}& \mathrm{}\\ 0& 0& 0& \mathrm{}& a_{n2}& 0& 0\\ 0& 0& 0& \mathrm{}& 0& a_{n1}& 0\end{array}\right),$$ (63) where $`a_{\mathrm{}}`$ are complex numbers depending on $`E`$. A simple calculation shows that $`\tau ^E`$ and $`𝒬^E`$ $`q`$-commute, i.e. $`[\tau ^E,𝒬^E]_q=0`$, where the $`q`$-commutator is defined by $`[O_1,O_2]_q:=O_1O_2qO_2O_1`$. Generalizing this property of $`\tau ^E`$ and $`𝒬^E`$ to $`\tau `$ and $`𝒬`$, we find $$[\tau ,𝒬]_q=0.$$ (64) This relation is the algebraic expression of the condition (61). It reduces to Eq. (11) for $`n=2`$. Another consequence of Eq. (63) is that a symmetry generator of a $`ZZ_n`$-graded UTS of type $`(1,1,\mathrm{},1)`$ with $`n>2`$ which satisfies (61) cannot be self-adjoint. Furthermore, one can easily check that $$(𝒬^E)^n=a_1a_2\mathrm{}a_nI_{n\times n}.$$ (65) We can generalize this equation to the whole Hilbert space and write it in the operator form: $$𝒬^n=𝒦.$$ (66) Here $`𝒦`$ is an operator that commutes with all other operators in the algebra. Next, we seek for the algebraic relations satisfied by the self-adjoint generators: $$Q_1:=\frac{1}{\sqrt{2}}(𝒬+𝒬^{})\mathrm{and}Q_2:=\frac{i}{\sqrt{2}}(𝒬𝒬^{}).$$ (67) Using the results reported in the appendix, namely Eq. (185), one can show that in an energy eigenspace $`_E`$, with $`E>0`$, $`Q_1`$ and $`Q_2`$ satisfy $`(Q_1^E)^n+\alpha _{n2}(Q_1^E)^{n2}+\mathrm{}`$ $`=`$ $`({\displaystyle \frac{1}{\sqrt{2}}})^nR_1,`$ (68) $`(Q_2^E)^n+\alpha _{n2}(Q_2^E)^{n2}+\mathrm{}`$ $`=`$ $`({\displaystyle \frac{1}{\sqrt{2}}})^nR_2,`$ (69) where $`\alpha _{\mathrm{}}`$’s are functions of $`\|a_i\|^2`$ and $`R_1`$ and $`R_2`$ are defined by $$R_1:=\underset{k=1}{\overset{n}{}}a_k+\underset{k=1}{\overset{n}{}}a_k^{},R_2:=\underset{k=1}{\overset{n}{}}(ia_k)+\underset{k=1}{\overset{n}{}}(ia_k^{}).$$ Eqs. (68) and (69) can be written in the operator form according to $`Q_1^n+M_{n2}Q_1^{n2}+\mathrm{}`$ $`=`$ $`({\displaystyle \frac{1}{\sqrt{2}}})^n(𝒦+𝒦^{}),`$ (70) $`Q_2^n+M_{n2}Q_2^{n2}+\mathrm{}`$ $`=`$ $`({\displaystyle \frac{1}{\sqrt{2}}})^n(i^n𝒦^{}+(i)^n𝒦),`$ (71) where $`M_i`$s are self-adjoint operators commuting with all other operators. We can rewrite Eqs. (70) and (71) in the following more symmetric way. * For $`n=2p`$, $$\begin{array}{c}(Q_1^2_1)(Q_1^2_2)\mathrm{}(Q_1^2_p)=(\frac{1}{2})^p(𝒦+𝒦^{}),\\ (Q_2^2_1)(Q_2^2_2)\mathrm{}(Q_2^2_p)=(\frac{1}{2})^p(𝒦+𝒦^{}).\end{array}$$ (72) * For $`n=2p+1`$, $$\begin{array}{c}(Q_1^2_1)(Q_1^2_2)\mathrm{}(Q_1^2_p)Q_1=(\frac{1}{\sqrt{2}})^{2p+1}(𝒦+𝒦^{}),\\ (Q_2^2_1)(Q_2^2_2)\mathrm{}(Q_2^2_p)Q_2=(\frac{i}{\sqrt{2}})^{2p+1}(𝒦^{}𝒦).\end{array}$$ (73) Here $`_i`$ are also operators that commute with all other operators. Note also that for even $`p`$’s, the algebraic relations for $`Q_1`$ and $`Q_2`$ coincide. Eqs. (66), (72), and (73) are the defining equations of the algebra of $`ZZ_n`$-graded TS of type $`(1,1,\mathrm{},1)`$. For $`n=2`$, they reduce to the familiar algebra of $`ZZ_2`$-graded TS of type $`(1,1)`$. In this case, as we discussed in section 3.1, we can perform a linear transformation on the self-adjoint generators that eliminates the operator $`𝒦`$. Next, consider the case where for all $`E>0`$, $`𝒦^E0`$ and introduce the transformed generator $`\stackrel{~}{𝒬}`$ by $$\stackrel{~}{𝒬}^E=\left(\frac{E}{\|\kappa ^E\|}\right)^{1/n}e^{i\varphi ^E/n}𝒬^E,$$ (74) where $$\kappa ^E:=a_1a_2\mathrm{}a_n\mathrm{and}e^{i\varphi ^E}:=\frac{\kappa ^E}{\|\kappa ^E\|}.$$ (75) In view of Eqs. (65) and (74), it is not difficult to see that $$(\stackrel{~}{𝒬}^E)^n=EI_n,$$ (76) Writing this equation in the operator form, we are led to $$H=\stackrel{~}{𝒬}^n.$$ (77) This is precisely the algebra of fractional supersymmetry of order $`n`$ . Next, we examine whether Eq. (66) guarantees the desired degeneracy structure of the $`ZZ_n`$-graded TS of type $`(1,1,\mathrm{},1)`$. In order to address this question, we suppose that the kernel of $`𝒦`$ is a subset of $`_0`$. Then for all $`E>0`$, $`\kappa ^E0`$. In view of Eq. (66), the eigenvalues of $`𝒬^E`$ are of the form $`q^{\mathrm{}}(\kappa ^E)^{1/n}`$, with $`\mathrm{}\{1,2,\mathrm{},n\}`$. Now, let $`\|q^{\mathrm{}},\nu _{\mathrm{}}`$ denote the corresponding eigenvectors, where $`\nu _{\mathrm{}}`$’s are degeneracy labels. We can easily show using Eq. (64) that for all positive integers $`s`$, $`\tau ^s\|q^{\mathrm{}},\nu _{\mathrm{}}`$ are eigenvectors of $`𝒬^E`$ with eigenvalue $`q^\mathrm{}s(\kappa ^E)^{1/n}`$. This is sufficient to conclude that all the eigenvalues of $`𝒬^E`$ are either nondegenerate or have the same multiplicity $`N_E`$. If for all $`E>0`$, $`N_E=1`$, then we will have the desired degeneracy structure of the $`ZZ_n`$-graded UTS of type $`(1,1,\mathrm{},1)`$. If there are $`E>0`$ for which $`N_E>1`$, then we have a nonuniform $`ZZ_n`$-graded TS of type $`(1,1,\mathrm{},1)`$. We conclude this section by noting that if we put $`𝒦=0`$ in the algebraic relations for $`ZZ_n`$-graded UTS of type $`(1,1,\mathrm{},1)`$, we will obtain the algebraic relations for $`ZZ_2`$-graded TS. This is not so strange, because putting $`𝒦=0`$ means that one of the $`a_i`$’s (say $`a_n`$) in $`𝒬^E`$ is zero. In this case, as we explain in the appendix, there is a unitary transformation that transforms $`Q_1^E`$ and $`Q_2^E`$ to off block diagonal matrices. In view of the analysis of section 3, this implies that the generators of this kind of $`ZZ_n`$-graded UTSs satisfy the algebra of $`ZZ_2`$ graded UT. ### 4.2 $`ZZ_n`$-Graded TS of Arbitrary Type $`(m_1,m_2,\mathrm{},m_n)`$ In order to obtain the algebraic structure of the $`ZZ_n`$-graded TS of arbitrary type $`(m_1,m_2,\mathrm{},m_n)`$, we need an appropriate grading operator. We shall adopt Eqs. (6), (59), (60) and (64) as the defining conditions for our $`ZZ_n`$-grading operator. We shall further assume that $`m_1m_2\mathrm{}m_n`$. This ordering can always be achieved by a reassignment of the colors. Again we shall confine our attention to the uniform $`ZZ_n`$-graded TS. The algebras of uniform and nonuniform TS of the same type are identical. Working in an eigenbasis in which $`\tau ^E`$ is diagonal and using Eq. (64), we have $$\tau ^E=\mathrm{diag}(\underset{m_1\mathrm{times}}{\underset{}{q,q,\mathrm{},q}},\underset{m_2\mathrm{times}}{\underset{}{q^2,q^2,\mathrm{},q^2}},\mathrm{},\underset{m_n\mathrm{times}}{\underset{}{q^n,q^n,\mathrm{},q^n}}).$$ (78) In view of this equation and Eq. (64), we obtain the following matrix representation for $`𝒬^E`$. $$𝒬^E=\left(\begin{array}{ccccccc}0& 0& 0& \mathrm{}& 0& 0& A_n\\ A_1& 0& 0& \mathrm{}& 0& 0& 0\\ 0& A_2& 0& \mathrm{}& 0& 0& 0\\ \mathrm{}& \mathrm{}& \mathrm{}& & \mathrm{}& \mathrm{}& \mathrm{}\\ 0& 0& 0& \mathrm{}& A_{n2}& 0& 0\\ 0& 0& 0& \mathrm{}& 0& A_{n1}& 0\end{array}\right),$$ (79) where $`A_{\mathrm{}}`$ (with $`\mathrm{}\{1,2,\mathrm{},n1\}`$) are complex $`m_{\mathrm{}+1}\times m_{\mathrm{}}`$ matrices, $`A_n`$ is a complex $`m_1\times m_n`$ matrix, and $`0`$’s are appropriate zero matrices. Next, we compute the $`n`$-th power of $`𝒬`$. The result is $$(𝒬^E)^n=\left(\begin{array}{ccccc}A_nA_{n1}\mathrm{}A_2A_1& 0& 0& \mathrm{}& 0\\ 0& A_1A_nA_{n1}\mathrm{}A_2& 0& \mathrm{}& 0\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ 0& 0& 0& \mathrm{}& A_{n1}A_{n2}\mathrm{}A_1A_n\end{array}\right).$$ (80) In order to find the most general algebraic identity satisfied by $`𝒬`$ we appeal to the following generalization of Lemma 1. The proof is given in the appendix. * Lemma 2: Let $`(m_1,m_2,\mathrm{},m_n)`$ be an $`n`$-tuple of positive integers satisfying $`m_1m_2\mathrm{}m_n`$, $`m:=_{\mathrm{}=1}^nm_{\mathrm{}}`$, $`\delta `$ is the number of times $`m_1`$ appears in $`(m_1,m_2,\mathrm{},m_n)`$, $`Q`$ is an $`m\times m`$ matrix of the form (79), and $`𝒫(x)`$ is the characteristic polynomial of the $`m_1\times m_1`$ matrix $`A_nA_{n1}\mathrm{}A_2A_1`$. Then $`Q`$ satisfies $$𝒫(Q^n)Q^{n\delta }=0.$$ (81) Substituting $`𝒬^E`$ for $`Q`$ in Eq. (81) and writing $`𝒫(x)`$ in terms of its roots $`\kappa _k^E`$, we find $$\left[(𝒬^E)^n\kappa _1^E\right]\left[(𝒬^E)^n\kappa _2^E\right]\mathrm{}\left[(𝒬^E)^n\kappa _{m_1}^E\right](𝒬^E)^{n\delta }=0.$$ (82) Next, we introduce the operators $`𝒦_k`$ for $`k\{1,2,\mathrm{},m_1\}`$ which commute with the Hamiltonian and have the representation: $$𝒦_k^E=\kappa _k^EI_m$$ (83) in $`_E`$. In view of Eqs. (82) and (83), we obtain $$(𝒬^n𝒦_1)(𝒬^n𝒦_2)\mathrm{}(𝒬^n𝒦_{m_1})𝒬^{n\delta }=0.$$ (84) The operators $`𝒦_k`$ have a similar degeneracy structure as the Hamiltonian (at least for the positive energy eigenvalues). Therefore, the Hamiltonian might be expressed as a function of $`𝒦_k`$. For a $`ZZ_n`$-graded UTS of type $`(1,1,\mathrm{},1)`$, $`\delta =n`$ and Eq. (84) reduces to (66). However, one can show that in general Eq. (84) does not ensure the desired degeneracy structure of the general $`ZZ_n`$-graded TSs. This is true even for the case $`n=2,m_+=2,m_{}=1`$ considered in section 3. In general, the $`ZZ_n`$-graded TSs correspond to a special class of symmetries satisfying (84). Finally, we wish to note that for a general $`ZZ_n`$-graded TS the algebraic relations satisfied by the self-adjoint generators are extremely complicated. We have not been able to express them in a closed form. ## 5 Mathematical Interpretation of $`\mathrm{\Delta }_{i,j}`$ In order to obtain the mathematical interpretation of the invariants $`\mathrm{\Delta }_{i,j}`$ of TSs, one must express the Hamiltonian in terms of the symmetry generators. This can be easily done using the defining algebra for the $`ZZ_n`$-graded TSs of type $`(1,1,\mathrm{},1)`$. In the following, we discuss the mathematical meaning of the topological invariants associated with these symmetries. We know that for a $`ZZ_n`$-graded TS of type $`(1,1,\mathrm{},1)`$ the operator $`𝒦`$ has a similar degeneracy structure as the Hamiltonian. Therefore, we may set $`H=f(𝒦)`$ where $`f`$ is a function mapping the eigenvalues $`\kappa ^E`$ to $`E`$. Next, suppose that the kernels of $`𝒦`$ and $`H`$ also coincide. Then as far as the general properties of the symmetry is concerned, we can confine our attention to the special case where $`𝒦=H`$, i.e., the fractional supersymmetry. Note that in this case, $`𝒬^n`$ is necessarily self-adjoint. Alternatively, we can use the rescaled symmetry generator $`\stackrel{~}{𝒬}`$ that does satisfy $`\stackrel{~}{𝒬}^n=H`$. In order to obtain the mathematical interpretation of the topological invariants $`\mathrm{\Delta }_{i,j}`$, we use an $`n`$-component representation of the Hilbert space in which a state vector $`\psi `$ is represented by a column of $`n`$ colored vectors $`\psi _{\mathrm{}}_{\mathrm{}}`$. In this representation the grading operator is diagonal and the generator of the symmetry and the Hamiltonian are respectively expressed by $`n\times n`$ matrices of operators according to $`𝒬`$ $`=`$ $`\left(\begin{array}{ccccccc}0& 0& 0& \mathrm{}& 0& 0& D_n\\ D_1& 0& 0& \mathrm{}& 0& 0& 0\\ 0& D_2& 0& \mathrm{}& 0& 0& 0\\ \mathrm{}& \mathrm{}& \mathrm{}& & \mathrm{}& \mathrm{}& \mathrm{}\\ 0& 0& 0& \mathrm{}& D_{n2}& 0& 0\\ 0& 0& 0& \mathrm{}& 0& D_{n1}& 0\end{array}\right),`$ (91) $`H`$ $`=`$ $`𝒬^n=\left(\begin{array}{cccccc}H_1& 0& 0& \mathrm{}& 0& 0\\ 0& H_2& 0& \mathrm{}& 0& 0\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ 0& 0& 0& \mathrm{}& 0& H_n\end{array}\right).`$ (96) Here $`D_n:_n_1`$ and $`D_{\mathrm{}}:_{\mathrm{}}_{\mathrm{}+1}`$, with $`\mathrm{}\{1,2,\mathrm{},n1\}`$, are operators and $`H_1`$ $`:=`$ $`D_nD_{n1}\mathrm{}D_2D_1,`$ $`H_2`$ $`:=`$ $`D_1D_nD_{n1}\mathrm{}D_2,`$ $`\mathrm{}`$ $`\mathrm{}`$ $`H_n`$ $`:=`$ $`D_{n1}D_{n2}\mathrm{}D_1D_n.`$ The condition that $`H`$ is self-adjoint takes the form $`H_{\mathrm{}}^{}=H_{\mathrm{}}`$, or alternatively $$D_{\sigma (n)}D_{\sigma (n1)}\mathrm{}D_{\sigma (2)}D_{\sigma (1)}=D_{\sigma (1)}^{}D_{\sigma (2)}^{}\mathrm{}D_{\sigma (n1)}^{}D_{\sigma (n)}^{},$$ (97) for all cyclic permutations $`\sigma `$ of $`(n,n1,n2,\mathrm{}1)`$. In addition, the assumption that $`H`$ has a nonnegative spectrum further restricts $`D_{\mathrm{}}`$. Note also that the algebraic relations satisfied by the self-adjoint generators also put restrictions on the choice of the operators $`D_{\mathrm{}}`$. This is because the operators $`M_i`$ appearing in Eqs. (70) and (71) involve $`D_{\mathrm{}}`$. The condition that $`M_i`$ commute with $`𝒬`$ leads to a set of compatibility relations among $`D_{\mathrm{}}`$. In view of Eq. (96) and the fact that $`n_{\mathrm{}}^{(0)}`$ is the dimension of the kernel of $`H_{\mathrm{}}`$, we can easily express the invariant $`\mathrm{\Delta }_{ij}`$ in the form: $$\mathrm{\Delta }_{i,j}=\mathrm{dim}(\mathrm{ker}H_j)\mathrm{dim}(\mathrm{ker}H_i).$$ (98) Suppose that the subspaces $`_{\mathrm{}}`$ are all identified with a fixed Hilbert Space. Consider the special case where $$D_3=D_4=\mathrm{}=D_n=1,D_2=D_1^{},$$ (99) and $`D_1`$ is a Fredholm operator, then $`H_1=D_1^{}D_1`$, $`H_2=D_1D_1^{}`$, and there is one independent invariant, namely $$\mathrm{\Delta }_{1,2}=\mathrm{dim}(\mathrm{ker}D_1^{}D_1)\mathrm{dim}(\mathrm{ker}D_1D_1^{})=\mathrm{dim}(\mathrm{ker}D_1)\mathrm{dim}(\mathrm{ker}D_1^{}).$$ This is just the analytic index of $`D_1`$. This example shows that the above construction has nontrivial solutions. In general, the operator $`D_{\mathrm{}}`$ need not satisfy (99). They are however subject to the above-mentioned compatibility relations. In order to demonstrate the nature of these relations, we consider $`ZZ_3`$-graded UTS of type $`(1,1,1)`$ with the symmetry generator $$𝒬=\left(\begin{array}{ccc}0& 0& D_3\\ D_1& 0& 0\\ 0& D_2& 0\end{array}\right).$$ (100) In this case, the algebraic relations (70) and (71) satisfied by the self-adjoint generators involve a single commuting operator which we denote by $`M`$. Assuming that this operator has the form $$M=\frac{1}{2}\left(\begin{array}{ccc}M_1& 0& 0\\ 0& M_2& 0\\ 0& 0& M_3\end{array}\right)$$ (101) and enforcing $$[𝒬,M]=0,Q_1^3+MQ_1=\frac{1}{\sqrt{2}}H,$$ (102) we find $`M_1D_3`$ $`=`$ $`D_3M_3`$ (103) $`M_2D_1`$ $`=`$ $`D_1M_1`$ (104) $`M_3D_2`$ $`=`$ $`D_2M_2`$ (105) One can manipulate these relations to obtain the following compatibility relations for $`D_{\mathrm{}}`$. $`D_2^{}D_2D_1D_3`$ $`=`$ $`D_1D_3D_2D_2^{},`$ (106) $`D_3^{}D_3D_2D_1`$ $`=`$ $`D_2D_1D_3D_3^{},`$ (107) $`D_1^{}D_1D_3D_2`$ $`=`$ $`D_3D_2D_1D_1^{}.`$ (108) Furthermore, under these conditions on $`M_{\mathrm{}}`$ and $`D_{\mathrm{}}`$, one can check that the relation for $`Q_2`$, i.e., $`Q_2^3+MQ_2=0`$, is identically satisfied. Next, consider the case where one of the $`D_i`$’s, say $`D_3`$, is $`1`$. Then we can use Eqs. (103) – (105) to express $`M_i`$ in terms of $`D_1`$ and $`D_2`$. This yields $$M_1=M_3=(D_1^{}D_1+D_2D_2^{}+1),M_2=(D_2^{}D_2+D_1D_1^{}+1).$$ (109) Note that we can easily satisfy Eqs. (106) – (108), if $`D_2=D_1^{}`$, i.e., when (99) is satisfied. In this case, $$M_1=M_2=(2D_1^{}D_1+1)=(2H_1+1),M_3=(2D_1D_1^{}+1)=(2H_2+1),$$ (110) and $$M=(H+\frac{1}{2}),$$ (111) where we have used Eqs. (96) and (101). ## 6 Examples In this section, we examine some examples of quantum systems possessing UTSs. ### 6.1 A System with UTS of Type $`(2,1)`$ Consider the Hamiltonian $$H=\frac{1}{2}(p^2+x^2)+\frac{1}{2}\left(\begin{array}{ccc}1& 0& 0\\ 0& 1& 0\\ 0& 0& 1\end{array}\right),$$ (112) where $`x`$ and $`p`$ are respectively the position and momentum operators. This Hamiltonian was originally considered in . It is not difficult to show that it commutes with $$𝒬=\left(\begin{array}{ccc}0& 0& \frac{1}{\sqrt{2}}(pix)\\ 0& 0& 0\\ 0& \frac{1}{\sqrt{2}}(p+ix)& 0\end{array}\right).$$ (113) Furthermore, $`𝒬`$, $`H`$, and the the self-adjoint generators $`Q_j`$ satisfy the algebra of the UTS of type $`(2,1)`$, i.e., $$Q_1^3=HQ_1,Q_2^3=HQ_2,𝒬^3=0.$$ For this system the grading operator is given by $`\tau =\mathrm{diag}(1,1,1)`$; there is a nondegenerate zero-energy ground state; and the positive energy eigenvalues are triply degenerate. Therefore, this system has a UTS of type $`(2,1)`$. If we denote by $`\|n`$ the normalized energy eigenvectors of the harmonic oscillator with unit mass and frequency, then a set of eigenvectors of the Hamiltonian (112) are given by $$\|\varphi _0=\left(\begin{array}{c}0\\ 0\\ \|0\end{array}\right),$$ (114) for $`E=0`$, and $$\|\varphi _n,1=\left(\begin{array}{c}\|n1\\ 0\\ 0\end{array}\right),\|\varphi _n,2=\left(\begin{array}{c}0\\ \|n1\\ 0\end{array}\right),\|\varphi _n,3=\left(\begin{array}{c}0\\ 0\\ \|n\end{array}\right),$$ (115) for $`E=n>0`$. Next, we check the action of the symmetry generator and the grading operator on $`\|\varphi _0`$ and $`\|\varphi _n,a`$. This yields $`𝒬\|\varphi _0=0,\tau \|\varphi _0=\|\varphi _0,`$ (116) $`𝒬\|\varphi _n,1=0,𝒬\|\varphi _n,2\|\varphi _n,3,𝒬\|\varphi _n,3\|\varphi _n,1,`$ (117) $`\tau \|\varphi _n,1=\|\varphi _n,1,\tau \|\varphi _n,2=\|\varphi _n,2,\tau \|\varphi _n,3=\|\varphi _n,3.`$ (118) Here we have used $``$ to denote equality up to a nonzero multiplicative constant. As one can see from Eqs. (116) and (118), the ground state has negative parity and the positive energy levels consist of two positive and one negative parity states. The topological invariants for this system are $`\mathrm{\Delta }_{2,1}=\mathrm{\Delta }_{1,2}=2`$. ### 6.2 A System with UTS of Type $`(1,1,1)`$ Consider the Hamiltonian $$H=𝒬^3=\frac{1}{2}(p^2+x^2)+\frac{1}{2}\left(\begin{array}{ccc}1& 0& 0\\ 0& 1& 0\\ 0& 0& 1\end{array}\right),$$ (119) which is precisely the Hamiltonian (112) written in another basis.<sup>4</sup><sup>4</sup>4We have used some of the results of to obtain this Hamiltonian. It possesses a symmetry generated by $$𝒬=\left(\begin{array}{ccc}0& 0& 1\\ \frac{1}{\sqrt{2}}(p+ix)& 0& 0\\ 0& \frac{1}{\sqrt{2}}(pix)& 0\end{array}\right).$$ (120) In fact, it is not difficult to show that $$𝒬^3=H.$$ (121) Furthermore, one can check that the self-adjoint generators $`Q_j`$ satisfy $$Q_1^3+MQ_1=\frac{1}{\sqrt{2}}H,Q_2^3+MQ_2=0,$$ (122) where the operator $`M`$ is given by $$M=\frac{1}{2}(p^2+x^2)+\left(\begin{array}{ccc}1& 0& 0\\ 0& 0& 0\\ 0& 0& 1\end{array}\right)=(H+\frac{1}{2}).$$ (123) This is in agreement with the more general treatment of section 5. In particular, Eq. (123) is a special case of Eq. (111). Obviously, $`M`$ commutes with $`H`$ and $`𝒬`$. In view of this observation and Eqs. (121), (122), (66), (70), and (71) we can conclude that this system has a $`ZZ_3`$-graded UTS of type $`(1,1,1)`$. A complete set of eigenvectors of the Hamiltonian (112) are given by $$\|\varphi _0=\left(\begin{array}{c}0\\ \|0\\ 0\end{array}\right),$$ (124) for $`E=0`$,and $$\|\varphi _n,1=\left(\begin{array}{c}\|n1\\ 0\\ 0\end{array}\right),\|\varphi _n,2=\left(\begin{array}{c}0\\ \|n\\ 0\end{array}\right),\|\varphi _n,3=\left(\begin{array}{c}0\\ 0\\ \|n1\end{array}\right),$$ (125) for $`E=n>0`$. The grading operator is $`\tau =\mathrm{diag}(q,q^2,1)`$, where $`q:=e^{2\pi i/3}`$. The symmetry generator $`𝒬`$ and the grading operator $`\tau `$ transform the energy eigenvectors according to $`𝒬\|\varphi _0=0,`$ $`\tau \|\varphi _0=q^2\|\varphi _0,`$ (126) $`𝒬\|\varphi _n,1\|\varphi _n,2,`$ $`𝒬\|\varphi _n,2\|\varphi _n,3,𝒬\|\varphi _n,3\|\varphi _n,1.`$ (127) $`\tau \|\varphi _n,1=q\|\varphi _n,1,`$ $`\tau \|\varphi _n,2=q^2\|\varphi _n,2,\tau \|\varphi _n,3=\|\varphi _n,3.`$ (128) In particular, $`\|\varphi _0`$ and $`\|\varphi _n,a`$ have colors $`q^2`$ and $`q^a`$, respectively, and the topological invariants of the system are given by $$\mathrm{\Delta }_{1,2}=\mathrm{\Delta }_{2,1}=1,\mathrm{\Delta }_{2,3}=\mathrm{\Delta }_{3,2}=1,\mathrm{\Delta }_{1,3}=\mathrm{\Delta }_{3,1}=0.$$ (129) ### 6.3 A system with UTS of type$`(1,1,\mathrm{},1)`$ Consider the Hamiltonian $$H=\frac{1}{2}(p^2+x^2)+\frac{1}{2}\mathrm{diag}(\underset{\mathrm{n}1\mathrm{t}\mathrm{i}\mathrm{m}\mathrm{e}\mathrm{s}}{\underset{}{1,1,\mathrm{},1}},1).$$ (130) This Hamiltonian has a symmetry generated by $$𝒬=\left(\begin{array}{ccccc}0& 0& \mathrm{}& 0& \frac{1}{\sqrt{2}}(pix)\\ 1& 0& \mathrm{}& 0& 0\\ 0& 1& \mathrm{}& 0& 0\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ 0& 0& \mathrm{}& \frac{1}{\sqrt{2}}(p+ix)& 0\end{array}\right),$$ (131) because $`𝒬^n=H`$. One can also check that the self-adjoint generators $`Q_j`$ satisfy $`Q_1^n+M_{n2}Q_1^{n2}+\mathrm{}`$ $`=`$ $`({\displaystyle \frac{1}{\sqrt{2}}})^n(2H),`$ (132) $`Q_2^n+M_{n2}Q_2^{n2}+\mathrm{}`$ $`=`$ $`({\displaystyle \frac{i}{\sqrt{2}}})^n(1+(1)^n)H,`$ (133) where $`M_{n2k}`$ are given by $$M_{n2k}=(1)^k\left[\frac{1}{2^k}\left(\genfrac{}{}{0pt}{}{nk1}{k}\right)+\frac{1}{2^{k1}}\left(\genfrac{}{}{0pt}{}{nk1}{k1}\right)H\right],$$ (134) and $`\left(\begin{array}{c}a\\ b\end{array}\right):=\frac{a!}{b!(ab)!}`$. It is not difficult to see that this system has a zero-energy ground state and that the positive energy levels are $`n`$-fold degenerate. A complete set of energy eigenvectors are given by $$\|0=\left(\begin{array}{c}0\\ 0\\ \mathrm{}\\ 0\\ \|0\end{array}\right),$$ (135) for $`E=0`$, and $$\|m,1=\left(\begin{array}{c}\|m1\\ 0\\ \mathrm{}\\ 0\\ 0\end{array}\right),\mathrm{},\|m,n1=\left(\begin{array}{c}0\\ 0\\ \mathrm{}\\ \|m1\\ 0\end{array}\right),\|m,n=\left(\begin{array}{c}0\\ 0\\ \mathrm{}\\ 0\\ \|m\end{array}\right),$$ (136) for $`E=m>0`$: These observations indicate that the quantum system defined by the Hamiltonian (130) has a $`ZZ_n`$-graded UTS of type $`(1,1,\mathrm{},1)`$. Clearly, the grading operator is $`\tau =\mathrm{diag}(q,q^2,\mathrm{},q^{n1},q^n=1)`$, the ground state has color $`c_n=q^n=1`$, and the nonvanishing topological invariants are $`\mathrm{\Delta }_{\mathrm{},n}=\mathrm{\Delta }_{n,\mathrm{}}=1`$, where $`\mathrm{}\{1,2,\mathrm{},n1\}`$. Next, we wish to comment that if we change the sign of the term involving the matrix $`\mathrm{diag}(\underset{\mathrm{n}1\mathrm{t}\mathrm{i}\mathrm{m}\mathrm{e}\mathrm{s}}{\underset{}{1,1,\mathrm{},1}},1)`$ in the Hamiltonian (130), then we obtain another quantum system with a $`ZZ_n`$-graded UTS of type $`(1,1,\mathrm{},1)`$ that is generated by $$𝒬=\left(\begin{array}{ccccc}0& 0& \mathrm{}& 0& \frac{1}{\sqrt{2}}(p+ix)\\ 1& 0& \mathrm{}& 0& 0\\ 0& 1& \mathrm{}& 0& 0\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ 0& 0& \mathrm{}& \frac{1}{\sqrt{2}}(pix)& 0\end{array}\right),$$ (137) This system has an $`(n1)`$-fold degenerate zero-energy ground state. ## 7 Conclusions We have introduced a generalization of supersymmetry that shares its topological properties. We gave a complete description of the underlying algebraic structure and commented on the meaning of the corresponding topological invariants. We showed that the algebras of the $`ZZ_2`$-graded TSs of type $`(m_+,1)`$ coincide with the algebras of supersymmetry or parasupersymmetry of order $`p=2`$. The algebraic relations obtained for the $`ZZ_2`$-graded TSs of type $`(m_+,m_{})`$ with $`m_{}>1`$ include as special cases the algebras of higher order parasupersymmetry advocated by Durand and Vinet . We also pointed out that the algebra of $`ZZ_n`$-graded TS of type $`(1,1,\mathrm{},1)`$ is related to the algebra of fractional supersymmetry of order $`n`$. Our approach in developing the concept of a TS differs from those of the other generalizations of supersymmetry in the sense that we introduce TSs in terms of certain requirements on the spectral degeneracy properties of the corresponding quantum systems, whereas in the other generalizations of supersymmetry such as parasupersymmetry and fractional supersymmetry one starts with certain defining algebraic relations. These relations are usually obtained by generalizing the relations satisfied by the generators of symmetries that relate degrees of freedom with different statistical properties in certain simple models. For example the Robakov-Spiridonov algebra of parasupersymmetry of order $`p=2`$ was originally obtained by generalizing the algebra of symmetry generators of an oscillator involving a bosonic and a ($`p=2`$) parafermionic degree of freedom . In order to investigate the topological content of these (statistical) generalizations of supersymmetry, one is forced to study the spectral degeneracy structure of the corresponding systems. The derivation of the degeneracy structure using the defining algebraic relations is usually a difficult task. In fact, for parasupersymmetries of order $`p>2`$ this problem has not yet been solved. Even for the parasupersymmetries of order $`p=2`$ the solution requires a quite lengthy analysis , and the defining algebra does not guarantee the existence of any topological invariants. This in turn raises the question of the classification of the $`p=2`$ parasupersymmetries that do have topological properties similar to supersymmetry . The analysis of the topological aspects of supersymmetry and $`p=2`$ parasupersymmetry shows that the information about the topological properties is contained in the spectral degeneracy structure of the corresponding systems. This is the main justification for our definition of a TS. Like any other quantum mechanical symmetry, a TS also possesses an underlying operator algebra. This algebra contains more practical information about the systems possessing the symmetry. As we showed in the preceding sections, the operator algebras associated with TSs can be obtained using the defining conditions on the spectral degeneracy structure of these systems. This observation may be viewed as another indication that, as far as the topological aspects are concerned, the spectral degeneracy structure is more basic than the algebraic structure. ## Acknowledgments A. M. wishes to thank Bryce DeWitt for introducing him to supersymmetric quantum mechanics and encouraging him to work on the topological aspects of supersymmetry more than ten years ago. We would also like to thank M. Khorrami and F. Loran for interesting discussions and fruitful comments. ## Appendix In this appendix we give the proofs of some of the mathematical results we use in sections 3 and 4. In the following we shall denote the characteristic polynomial of a matrix $`M`$ by $`𝒫_M(x)`$, i.e., $`𝒫_M(x)=det(xIM)`$. * Lemma 0: Let $`X`$ and $`Y`$ be $`m\times n`$ and $`n\times m`$ matrices respectively. Then $$𝒫_{YX}(\lambda )=\lambda ^{nm}𝒫_{XY}(\lambda ).$$ (138) * Proof: Let $$M:=\left(\begin{array}{cc}I_m& X\\ Y& \lambda I_n\end{array}\right).$$ (139) Then using the well-known properties of the determinant, we have $`det(M)`$ $`=`$ $`det\left(\begin{array}{cc}I_m& X\\ Y& \lambda I_n\end{array}\right)`$ (142) $`=`$ $`det\left(\begin{array}{cc}I_m& X\\ 0& \lambda I_nYX\end{array}\right)`$ $`=`$ $`det(\lambda I_nYX)`$ $`=`$ $`𝒫_{YX}(\lambda ).`$ (146) Similarly, we can show that $`det(M)`$ $`=`$ $`det\left(\begin{array}{cc}I_m& X\\ Y& \lambda I_n\end{array}\right)`$ (149) $`=`$ $`\lambda ^ndet\left(\begin{array}{cc}I_m& X\\ \frac{1}{\lambda }Y& I_n\end{array}\right)`$ (152) $`=`$ $`\lambda ^{nm}det\left(\begin{array}{cc}\lambda I_m& \lambda X\\ \frac{1}{\lambda }Y& I_n\end{array}\right)`$ (155) $`=`$ $`\lambda ^{nm}det\left(\begin{array}{cc}\lambda I_mXY& 0\\ \frac{1}{\lambda }Y& I_n\end{array}\right)`$ $`=`$ $`\lambda ^{nm}det(\lambda I_mXY)`$ $`=`$ $`\lambda ^{nm}𝒫_{XY}(\lambda ).`$ (159) Eqs. (146) and (159) yield the identity (138). $`\mathrm{}`$ * Lemma 1: Let $`m_\pm `$ be positive integers, $`m=m_++m_{}`$, and $`Q`$ be an $`m\times m`$ matrix of the form $$Q=\left(\begin{array}{cc}0& X\\ Y& 0\end{array}\right),$$ (160) where $`X`$ and $`Y`$ are $`m_+\times m_{}`$ and $`m_{}\times m_+`$ complex matrices. Then $`𝒫_{XY}(Q^2)Q=𝒫_{YX}(Q^2)Q=0`$. Furthermore, if $`m_+=m_{}`$, then $`𝒫_{XY}(Q^2)=𝒫_{YX}(Q^2)=0`$. * Proof: Let $`a_k`$ denote the coefficients of $`𝒫_{YX}(x)`$, i.e., $`𝒫_{YX}(x)=_{k=0}^na_kx^k`$. According to the Cayley-Hamilton theorem, $$𝒫_{YX}(YX)=\underset{k=0}{\overset{n}{}}a_k(YX)^k=0.$$ (161) In view of this identity, we can easily show that for any positive integer $`k`$, $`Q^{2k}`$ $`=`$ $`\left(\begin{array}{cc}(XY)^k& 0\\ 0& (YX)^k\end{array}\right),`$ (164) $`𝒫_{YX}(Q^2)`$ $`=`$ $`{\displaystyle \underset{k=0}{\overset{n}{}}}a_kQ^{2k}=\left(\begin{array}{cc}_{k=0}^na_k(XY)^k& 0\\ 0& _{k=0}^na_k(YX)^k\end{array}\right)`$ (167) $`=`$ $`\left(\begin{array}{cc}_{k=0}^na_k(XY)^k& 0\\ 0& 0\end{array}\right)`$ (170) $`𝒫_{YX}(Q^2)Q`$ $`=`$ $`\left(\begin{array}{cc}0& _{k=0}^na_k(XY)^kX\\ 0& 0\end{array}\right)`$ (173) $`=`$ $`\left(\begin{array}{cc}0& X_{k=0}^na_k(YX)^k\\ 0& 0\end{array}\right)=0.`$ (176) A similar calculation yields $`𝒫_{XY}(Q^2)Q=0`$. Furthermore, using Lemma 0 one can see that if $`m_+=m_{}`$, then $`𝒫_{XY}(x)=𝒫_{YX}(x)`$. In view of this identity and Eq. (170), we have (for the case $`m_+=m_{}`$) $`𝒫_{XY}(Q^2)=𝒫_{YX}(Q^2)=0`$$`\mathrm{}`$ * Corollary 1: Consider the $`n\times n`$ matrix $`M`$ whose elements are given by $$M_{ij}:=a_j\delta _{i,j+1}+a_i^{}\delta _{i+1,j},i,j\{1,2,\mathrm{},n\}$$ (177) i.e., $$M:=\left(\begin{array}{cccccc}0& a_1^{}& 0& \mathrm{}& 0& 0\\ a_1& 0& a_2^{}& \mathrm{}& 0& 0\\ 0& a_2& 0& \mathrm{}& 0& 0\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ 0& 0& 0& \mathrm{}& 0& a_{n1}^{}\\ 0& 0& 0& \mathrm{}& a_{n1}& 0\end{array}\right)$$ (178) Then the characteristic polynomial of $`M`$ is given by $$𝒫_M(x)=\{\begin{array}{cc}𝒫_{AA^{}}(x^2)\hfill & \text{for }n=2p\hfill \\ 𝒫_{AA^{}}(x^2)x\hfill & \text{for }n=2p+1\text{,}\hfill \end{array}$$ (179) where $`A`$ is the matrix with entries $$A_{ij}=a_{2i1}\delta _{ij}+a_{2i}^{}\delta _{i+1,j}.$$ (180) It is a $`p\times p`$ matrix for $`n=2p`$ and a $`p\times (p+1)`$ matrix for $`n=2p+1`$. * Proof: Consider the following unitary transformation $$M\stackrel{~}{M}=U^{}MU,$$ (181) where $`U`$ is defined by $$U_{ij}=\{\begin{array}{cc}\delta _{i,2j}\hfill & \text{for }jp\hfill \\ \delta _{i,2j2p1}\hfill & \text{for }j>p.\hfill \end{array}$$ (182) Substituting this equation in (181), we find $$\stackrel{~}{M}=\left(\begin{array}{cc}0& A\\ A^{}& 0\end{array}\right).$$ (183) Now applying Lemma 1, we obtain Eq. (179). Furthermore, the coefficeints of the characteristic polynomial of $`M`$ are functions of $`\|a_i\|^2`$s only. $`\mathrm{}`$ * Corollary 2: Consider the $`n\times n`$ complex matrix $$Q:=\left(\begin{array}{cccccc}0& a_1^{}& 0& \mathrm{}& 0& a_n\\ a_1& 0& a_2^{}& \mathrm{}& 0& 0\\ 0& a_2& 0& \mathrm{}& 0& 0\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ 0& 0& 0& \mathrm{}& 0& a_{n1}^{}\\ a_n^{}& 0& 0& \mathrm{}& a_{n1}& 0\end{array}\right).$$ (184) Then the characteristic polynomial of $`Q`$ is given by $$𝒫_Q(x)=(\underset{k=1}{\overset{n}{}}a_k+\underset{k=1}{\overset{n}{}}a_k^{})+x^n+\beta _{n2}x^{n2}+\mathrm{}+\beta _{n2k}x^{n2k}+\mathrm{}$$ (185) where $`\beta _{n2k}`$s are functions of $`\|a_i\|^2`$s. * Proof: A straightforward application of the properties of the determinant, one can show that $$𝒫_Q(x)=(\underset{k=1}{\overset{n}{}}a_k+\underset{k=1}{\overset{n}{}}a_k^{})+𝒫_M(x)\|a_n\|^2𝒫_M^{}(x),$$ (186) where $`M^{}`$ is obtained from $`M`$ by removing the first and last rows and columns. Now, since $`𝒫_M(x)`$ is an odd (even) polynomial for odd (even) $`n`$, $`𝒫_Q(x)`$ will have the form given by Eq. (185). $`\mathrm{}`$ * Lemma 2: Let $`(m_1,m_2,\mathrm{},m_n)`$ be an $`n`$-tuple of positive integers satisfying $`m_1m_2\mathrm{}m_n`$, $`m:=_{\mathrm{}=1}^nm_{\mathrm{}}`$, $`\delta `$ is the number of times $`m_1`$ appears in $`(m_1,m_2,\mathrm{},m_n)`$, $`Q`$ is an $`m\times m`$ matrix of the form (79), and $`𝒫(x)`$ is the characteristic polynomial of the $`m_1\times m_1`$ matrix $`A_nA_{n1}\mathrm{}A_2A_1`$. Then $`Q`$ satisfies $$𝒫(Q^n)Q^{n\delta }=0.$$ (187) * Sketch of Proof: The proof of this lemma is very similar to the proof of Lemma 1. The idea is to multiply the block diagonal matrix $`𝒫(Q^n)`$ with a power of $`Q`$ so that one obtains a matrix whose entries have a factor $`𝒫(A_nA_{n1}\mathrm{}A_1)`$. Then one uses the Cayley-Hamilton theorem to conclude that the resulting matrix must vanish. One can show by inspection that the smallest nonnegative integer $`r`$ for which a factor of $`𝒫(A_nA_{n1}\mathrm{}A_1)`$ occurs in all the entries of $`𝒫(Q^n)Q^r`$ is $`n\delta `$.
warning/0007/astro-ph0007026.html	ar5iv	text	# Hubble Space Telescope imaging survey of sub-mJy star-forming galaxies I: morphologies at 𝑧∼0.2 ## 1 Introduction The study of the observational constraints on the cosmic star formation history is currently among the most active fields in observational cosmology. The most widely used tracer of the comoving volume-averaged star formation rate (SFR) is the UV luminosity density (e.g. Madau et al. 1996, Steidel et al. 1996, 1999), found to peak at $`z12`$. However, little is known about the history of star formation in the Universe beyond its global average (e.g. Abraham et al. 1999). In this paper we present the first results from an on-going program with Hubble Space Telescope (HST) to study the morphology of starforming galaxies, using a sample of galaxies which dominate the comoving SFR and which were selected in a manner free from obscuration biases. An important caveat to the SFR constraints is that the UV luminosity of star-forming galaxies is dominated by the lowest-extinction regions. This leads to extreme sensitivity to obscuration corrections (e.g. Pettini et al. 1998, Meurer et al. 1997), large enough to eliminate the evidence for a redshift cut-off in the SFR. Indeed, SCUBA observations of the Hubble Deep Field (HDF) (Hughes et al. 1998) detected $`5`$ ultraluminous galaxies at $`z>2`$ (the redshift constraints come mainly from the radio-FIR correlation and not from the uncertain HDF IDs), which together implied an obscured star formation rate at least as large as the integrated de-reddened UV-derived rate. Moreover, the H$`\alpha `$-derived star formation rate from the Canada-France Redshift Survey (CFRS, Tresse & Maddox 1997) was a factor $`\times 3`$ larger than the UV estimate in the same sample, but comparable to those derived from the mid-IR for the CFRS (Flores et al. 1999) or the HDF (Rowan-Robinson et al. 1997). However, a serious problem is that optical-UV light is strongly skewed to low extinction regions, while the reverse is true for mid and far-IR selected samples, leading to problematic obscuration-related effects in either case. An unbiased technique for selecting star-forming galaxies is to sample at decimetric radio wavelengths, where both obscured and unobscured star forming galaxies contribute (Condon 1992). Several groups (including ourselves) have begun to exploit this technique to trace the comic star formation history, and to study the obscuration effects in star forming galaxies (e.g. Cram et al. 1998, Serjeant et al. 1998, Oliver et al. 1998, Haarsma & Partridge 1998, Cram 1998, etc). We have embarked on a programme to image almost all star-forming galaxies selected at decimetric radio wavelengths, using the WFPC2 in snapshot mode on the HST. In this paper we carry out a morphological study of a small subset of these galaxies at $`z0.2`$, from our first cycle $`8`$ observations in this programme. (Although these are at low redshifts by many standards, they are not local, with the volume-averaged SFR being a factor of $`\times 2`$ higher than in the local Universe for $`(1+z)^3`$ luminosity evolution models (Table 1).) The galaxies in this sample are selected so that they dominate the (radio-derived) star formation history of the Universe at this epoch, i.e. close to the $`L_{}`$ characteristic luminosity in the radio luminosity function. The sample selection is discussed in Section 2, which also presents the HST observations and data reduction. These data are then analysed in Section 3, followed by a discussion of the results in Section 4. ## 2 Sample selection and observations The following criteria are considered in selecting the sample for HST observations: 1. The objects are sub-mJy radio sources, selected at 1.4 GHz (Benn et al (1993); Georgakakis et al (2000)). This avoids obscuration-related selection biases. Decimetric radio fluxes are much less affected by AGN contribution than shorter wavelength observations (compare e.g. the ongoing contoversy over the AGN fraction in higher frequency sample with Hammer et al. (1995) finding $`50\%`$ AGN fraction at $`58`$ GHz while Windhorst et al. (1995) find only $`5\%`$). 2. Spectroscopic data are used to select emission-line star-forming galaxies. 3. The galaxies must not lie close to bright stars to avoid HST roll angle constraints. 4. Galaxies selected to lie close to the $`1.4`$GHz $`L_{}`$, which dominate the volume averaged SFR (Table 1). 5. Galaxies in a narrow redshift range around $`z0.2`$ to avoid differential evolution. Our primary sample satisfies (i), (ii) and (iii), and is the subject of our $`150`$ snapshots in cycles $`8`$ and $`9`$. Ten galaxies also satisfy conditions (iv) and (v); these are the cycle $`8`$ targets. The HST observations for four of these have now been completed. We selected the F814W (wide I band) filter. At this wavelength, the light is mostly dominated by old stellar population, so gives a more accurate measure of distortion in the underlying gravitational potential than e.g. UV observations. Details of the current sub-sample are presented in Table 2. The data were reduced using the Interactive Data Language (IDL), starting from the automatic pipeline products. Cosmic rays were identified as $`>3\sigma `$ differences between frames. For sky subtraction, we estimated the modal value of the underlying counts distribution using iterative fits to the readout histograms. The final HST images for the four sub-mJy radio sources are presented in Figure 1. ## 3 Results \[CM84\] 144 is extremely disturbed, with a tidal tail extending $`1.5^{\prime \prime }`$ North-East of the nucleus. Several secondary nuclei are also apparent, reminiscent of HST WFPC2 F814W imaging of ultra-luminous infrared galaxies (e.g. Borne et al. 2000) such as Arp 220 (Borne & Lucas 1997). \[CM84\] 074 is only mildly asymmetric. There are hints of face-on spiral structure, and an excess flux $`0.5^{\prime \prime }`$ West of the nucleus. There are two companion galaxies a few arcseconds to the West, also visible in the digitised sky survey. (We reject these as identifications of the radio source based on the ID magnitude quoted in Benn et al. 1993.) Galaxy Phoenix-Deep-29 is again not as clearly disturbed. However, the inner isophotes are clearly offset in position and (tentatively) in orientation from the outer isophotes. This galaxy also shows signs of interaction with a companion $`5^{\prime \prime }`$ to the East. There is a clear tidal tail from the companion extending a few arcseconds North. Galaxy Phoenix-Deep-96 shows clear signs of morphological disturbance, with the structure dominated by a bright central bar-like feature. There are also hints of multiple nuclei inside the bar, and there is a secondary peak $`0.5^{\prime \prime }`$ South-West of the bar. There are however no other clearly associated companion galaxies visible in the HST image. To quantify the morphological disturbance, we follow the procedure adopted by Brinchmann et al. (1998) and others, by rotating each image through $`\pi `$ and subtracting it from the original image. Normalising this to the galaxy flux yields the fractional asymmetric flux: $$A=\frac{\mathrm{\Sigma }\|G_{ij}G_{ij}^{}\|}{\mathrm{\Sigma }G_{ij}}k$$ (1) where $`G`$ ($`G^{}`$) is the rotated (unrotated) image, $`k`$ is a correction for systematics such as sky gradients, and the sum is performed over pixels $`ij`$. The values of $`A`$ are fairly insensitive to the choice of aperture used for the self-subtraction, provided all pixels at $`\stackrel{>}{_{}}1.5\sigma `$ are included in the subimage. The magnitude of $`k`$ can be estimated by applying the technique to blank sky regions of the same size. The errors in $`A`$ are in practice dominated by the uncertainty in $`k`$: we find typically $`\mathrm{\Delta }A0.05`$. A further advantage of this statistic is that the $`I`$-selected Canada-France Redshift Survey (CFRS) and $`b_\mathrm{J}`$-selected LDSS surveys have extensive published HST morphologies quantified using the same statistic, yielding valuable control samples. The asymmetry measures for the sub-mJy star-forming galaxies, quantified as discussed above, are listed in Table 2. ## 4 Discussion and conclusions In Figure 2 we plot the asymmetry statistics of our galaxies, and compare them to a control sample of coeval galaxies from the CFRS sample (Brinchmann et al. 1998). The control extends to fainter $`I_{\mathrm{F814W}}`$ fluxes than our star-forming galaxy sample, though there is some overlap. Remarkably, our galaxies are far more asymmetric than the control sample of an (ostensibly) identical population of optically-selected star-forming galaxies of the same $`I_{\mathrm{F814W}}`$. The optical control must be well-matched with the radio sample in as many optical properties as possible. The radio sample is more asymmetric at a fixed optical luminosity, but is this also true at a fixed optically-derived SFR? I.e., could the morphological differences between the radio sample and the optically-selected control be instead because they are not well-matched in optically-derived star formation rates? To address this we must compare the H$`\alpha `$ luminosities of our radio sample with those of the control. The data for the radio samples are presented in Benn et al (1993) and Georgakakis et al (2000). The individual H$`\alpha `$ luminosities of the CFRS galaxies are not published, but Tresse & Maddox (1998) report a tight correlation between H$`\alpha `$ luminosities and $`M_\mathrm{B}(AB)`$ absolute magnitudes. In Figure 3 we use this relation to estimate the H$`\alpha `$-derived star formation rates, and compare the optical control sample with our sub-mJy star-forming galaxies. We adopt the following conversion between H$`\alpha `$ luminosities and star-formation rates (SFRs): $$\mathrm{SFR}(\mathrm{H}\alpha )=L(\mathrm{H}\alpha )/(1.41\times 10^{34}\mathrm{W})$$ (2) which is derived assuming Salpeter IMF between $`0.1`$ $`M_{}`$ and $`125`$ $`M_{}`$ (see Serjeant et al. 1998, Oliver et al. 1998). We de-redden our H$`\alpha `$ fluxes by $`A_\mathrm{V}=1`$ for consistency with optically-selected samples. (Tresse & Maddox 1998 note that the H$`\alpha `$ luminosity function is essentially unchanged if assuming this $`A_\mathrm{V}`$ throughout their sample, instead of using individual Balmer decrements.) There is a more substantial overlap of the H$`\alpha `$-derived SFR with the control sample, and it is again clear that at a fixed H$`\alpha `$ SFR the radio-selected sample is significantly more disturbed than the optically-selected sample of star forming galaxies. In summary, the radio-selected sources are typically more morphologically disturbed than the optically-selected control, well-matched in redshift, optical luminosity and/or H$`\alpha `$ SFR. Why this might be the case is hinted in Table 2, where we compare the star formation rates of our galaxies estimated from H$`\alpha `$ and the $`1.4`$GHz radio luminosity. To convert decimetric radio luminosities to SFRs we use $$\mathrm{SFR}(1.4\mathrm{GHz})=L_{1.4}/(7.63\times 10^{20}\mathrm{W}\mathrm{Hz}^1)$$ (3) appropriate for the same IMF as above (see Oliver et al. 1998, Serjeant et al. 1998). Our H$`\alpha `$ SFRs are comparable in most cases to $`L_{}`$ in the H$`\alpha `$ luminosity function (Tresse & Maddox 1997). However the radio-derived star formation rates are much larger (though still comparable with the radio $`L_{}`$, Table 1). This indicates that a large fraction of the star formation in these objects occurs in regions with $`A_\mathrm{V}1`$, since the observed H$`\alpha `$ luminosities will be dominated by regions with $`A_\mathrm{V}\stackrel{<}{_{}}1`$. Indeed if the dust is well-mixed with the H$`\alpha `$ emitting gas, the different optical depths for H$`\alpha `$ and H$`\beta `$ ensure that $`A_\mathrm{V}=1.1`$ would be derived for a simple screen Balmer decrement model regardless of the true extinction to the rear of the cloud. We therefore suggest that optically-selected samples may under-represent disturbed galaxies with large amounts of obscured star formation, in any SFR-weighted quantity. The galaxy morphologies of the sub-mJy sources are also markedly different from coeval AGN. McLure el al. (1999) and McLeod & Reike (1995) have found giant ellipticals hosting both radio-loud and radio-quiet quasars at these redshifts, and radiogalaxies are also early type galaxies at these epochs. The EMSS-selected sample of Schade et al. (1999) at $`z0.15`$ also finds no evidence for strong interaction/merger activity for any AGN in the sample. This may pose problems for models which propose causal links between the cosmic evolution of AGN and the cosmic star formation history. For example, a model in which both are driven by similar merger events must also provide a mechanism for delaying the onset of central engine fuelling, to account for the (relatively) more relaxed AGN host galaxy potentials. This is related to the more fundamental and unsolved problem of AGN evolution: how to drive the gas fuel down $`5`$ orders of magnitude in radius. We can also compare our sample with the ultraluminous galaxy samples of e.g. Borne et al. (2000), who also used HST WFPC2 F814W in snapshot mode. These authors found highly disturbed systems often in moderately rich environments, which they used to argue for an evolutionary sequence involving compact groups and ultraluminous galaxies. Although asymmetry statistics have not been calculated for the ultraluminous galaxies, our targets appear qualitatively less disturbed and the environments possibly less rich. (Ideally both ultraluminous galaxies and radio-selected star-forming galaxies should also have their environments quantified with for example the $`B_{\mathrm{gg}}`$ statistic; e.g. Yee & Green 1984, Hill & Lilly 1991, Wold et al. 2000) This may suggest a link between the level of star formation activity after interactions and the richness of the environment, perhaps via the number or the nature of the interaction events. Although our current sample is small, the asymmetric morphologies in our $`L_{}`$ targets suggest that galaxy interactions play a major role in the evolution of the star formation and metal production rates in the low-$`z`$ Universe. Further papers in this series will present results for larger samples of sub-mJy star-forming galaxies. However, it is worth keeping in mind that the formation and evolution of galaxies cannot be characterised by a single parameter, such as the volume-averaged star formation rate. In expanding on this simple first-order description, the obscuration-independent selection of galaxies will be essential, as will co-ordinated multi-wavelength follow-up (e.g. ELAIS, Oliver et al. 2000). ## Acknowledgements It is a pleasure to thank Patricia Royle for her help in the preparation of these observations. Based on observations with the NASA/ESA Hubble Space Telescope, obtained at the Space Telescope Science Institute, which is operated by the Association of Universities for Research in Astronomy, Inc. under NASA contract No. NAS5-26555. This work was supported by PPARC (grant number GR/K98728) and by the EC TMR Network programme (FMRX-CT96-0068).
warning/0007/math0007165.html	ar5iv	text	# 𝐺-actions on graphs ## 1. Introduction Let $`\mathrm{\Gamma }`$ be a finite $`d`$-valent graph and let $`G`$ be a $`n`$-dimensional torus. In this paper we will be concerned with objects (rings, modules, $`G`$-representations $`\mathrm{}`$) associated to an “action” of $`G`$ on $`\mathrm{\Gamma }`$. To define what we mean by this term, let $`V_\mathrm{\Gamma }=V`$ be the vertices of $`\mathrm{\Gamma }`$ and $`E_\mathrm{\Gamma }`$ the oriented edges. For each $`eE_\mathrm{\Gamma }`$ let $`i(e)`$ and $`t(e)`$ be the initial and terminal vertices of $`e`$, and let $`\overline{e}`$ be the edge, $`e`$, with its orientation reversed. (Thus $`i(e)=t(\overline{e})`$ and $`t(e)=i(\overline{e}`$.) ###### Definition 1.1. Let $`\varrho `$ be a map which assigns to each oriented edge, $`e`$, of $`\mathrm{\Gamma }`$ a one dimensional representation, $`\varrho _e`$, with character $$\chi _e:GS^1,$$ (1.1) let $`\tau `$ be a map which assigns to each vertex, $`p`$, of $`\mathrm{\Gamma }`$ a $`d`$-dimensional representation, $`\tau _p`$, and let $`G_e`$ be the kernel of (1.1). $`\varrho `$ and $`\tau `$ define an action of $`G`$ on $`\mathrm{\Gamma }`$ if they satisfy the axioms (1.2)–(1.4) below : $`\tau _p`$ $``$ $`{\displaystyle \underset{i(e)=p}{}}\varrho _e`$ (1.2) $`\varrho _{\overline{e}}`$ $``$ $`\varrho _e^{}`$ (1.3) $`\tau _{i(e)}\|_{G_e}`$ $``$ $`\tau _{t(e)}\|_{G_e}.`$ (1.4) ###### Remark. For the connection between this graph-theoretic notion of “$`G`$-action” and the usual notion of $`G`$-action, see Section 7 (or, for more details, \[GZ2, § 3.1\]). Let $`_G^{}`$ be the weight lattice of $`G`$, and let $`\alpha _e_G^{}`$ be the weight of the representation, $`\varrho _e`$, i.e. $$\chi _e=e^{2\pi i\alpha _e}.$$ (1.5) By (1.2) and (1.5) both $`\varrho `$ and $`\tau `$ are determined by the $`\alpha _e`$’s; so an action of $`G`$ on a graph, $`\mathrm{\Gamma }`$, can be thought of as a labeling of each edge, $`e`$, of the graph by a weight, $`\alpha _e`$. This labeling, however, will be forced by (1.2)–(1.4) to satisfy certain axioms. For instance, by (1.3) $$\alpha _{\overline{e}}=\alpha _e.$$ (1.6) We will say that an action is a GKM action if, for every pair of edges with the same initial vertex, $`p=i(e_1)=i(e_2)`$, either $`e_1=e_2`$ or $`\alpha _{e_1}`$ and $`\alpha _{e_2}`$ are linearly independent. (For the geometric interpretation of this property, see Section 7). All the actions we consider below will be assumed to be GKM actions. This paper is the fourth in a series of papers on the equivariant cohomology of graphs. The first three papers in this series were concerned with the equivariant cohomology ring, $`H_G(\mathrm{\Gamma })`$. In this paper we will be concerned with a slightly more complicated object: the equivariant “$`K`$-cohomology” ring of $`\mathrm{\Gamma }`$. However, to motivate its definition, we will first recall how $`H_G(\mathrm{\Gamma })`$ is defined: Denote by $`𝔤`$ and $`𝔤_e`$ the Lie algebras of $`G`$ and $`G_e`$, and let $`\mathrm{SS}(𝔤^{})`$ and $`\mathrm{SS}(𝔤_e^{})`$ be the symmetric algebras over the duals of $`𝔤`$ and $`𝔤_e`$. From the inclusion of $`𝔤_e`$ into $`𝔤`$ one gets a restriction map $$r_e:\mathrm{SS}(𝔤^{})\mathrm{SS}(𝔤_e^{}).$$ (1.7) ###### Definition 1.2. $`H_G(\mathrm{\Gamma })`$ is the set of all functions, $`f:V_\mathrm{\Gamma }\mathrm{SS}(𝔤^{})`$, which satisfy the compatibility conditions $$r_ef_{i(e)}=r_ef_{t(e)}$$ (1.8) for all edges, $`e`$ of $`\mathrm{\Gamma }`$. Following \[KR\] we will define the $`K`$-theory analog of $`H_G(\mathrm{\Gamma })`$ simply by replacing $`\mathrm{SS}(𝔤^{})`$ in the definition by the representation ring, $`R(G)`$, of $`G`$. ###### Definition 1.3. $`K_G(\mathrm{\Gamma })`$ is the set of all functions, $`f:V_\mathrm{\Gamma }R(G)`$, which satisfy the compatibility condition (1.8), $`r_e`$ being the restriction map, $`R(G)R(G_e)`$. ###### Remarks. 1. Since $`G`$ is an $`n`$-torus, the representation ring $`R(G)`$ can be identified with the character ring of $`G`$, i.e. the ring of all finite sums $$m_ke^{2\pi i\alpha _k}$$ (1.9) with $`m_k`$ and $`\alpha _k_G^{}`$. We will frequently use this identification, referring to a representation by indicating the element of the character ring it corresponds to and vice-versa. 2. Point-wise multiplication makes $`K_G(\mathrm{\Gamma })`$ into a ring. Moreover, since the constant functions satisfy (1.8), this ring contains the ring, $`R(G)`$, as a subring. Given $`fK_G(\mathrm{\Gamma })`$ let $$\chi (f)=\underset{pV}{}f_p\underset{i(e)=p}{}(1e^{2\pi i\alpha _e})^1.$$ (1.10) We will call $`\chi (f)`$ the Atiyah-Bott character of the class, $`f`$. The individual summands on the right hand side are elements of a quotient ring of $`R(G)`$; however, we will prove ###### Theorem 1.1. The sum (1.10) is an element of $`R(G)`$. Thus (1.10) defines a morphism of $`R(G)`$-modules $$\chi :K_G(\mathrm{\Gamma })R(G)$$ which we will call the character map. A helpful way of looking at this map is in terms of virtual representations. Namely, to each $`pV`$, one can attach an infinite-dimensional virtual representation, $`Q(\tau _p)`$, the “$`\text{spin}^{}`$-quantization” of the action, $`\tau _p`$, of $`G`$ on $`^d`$, and Theorem 1.1 asserts that the sum $$Q(f)=\underset{pV}{}Q(\tau _p)f_p$$ (1.11) is a finite dimensional virtual representation and that its character is given by (1.10). Suppose, in particular, that $`f`$ has the form $$f_p=e^{2\pi i\alpha _p},\alpha _p_G^{}.$$ (1.12) Then by (1.8) $$\alpha _q\alpha _p=m_e\alpha _e$$ (1.13) for every pair of vertices, $`p`$ and $`q`$, and edge, $`e`$, joining $`p`$ to $`q`$. ###### Definition 1.4. $`f`$ is symplectic if $`m_e>0`$ for all $`e`$. If $`f`$ is symplectic, the representation (1.11) has the following convexity property. (Compare with \[GS, Theorem 6.3\].) ###### Theorem 1.2. If $`\alpha `$ is a weight of $`Q(f)`$ then $`\alpha `$ is in the convex hull of $`\{\alpha _p;pV\}`$. Let’s denote this convex hull by $`\mathrm{\Delta }`$. We will call a weight, $`\alpha `$, an extremal weight if it is a vertex of $`\mathrm{\Delta }`$. For these weights we will prove ###### Theorem 1.3. If $`\alpha `$ is extremal, it occurs in $`Q(f)`$ with multiplicity $`1`$. For non-extremal weights we will prove a more refined result. Fix a vector, $`\xi `$, in $`𝔤`$ with the property $`\alpha _e(\xi )0`$ for all edges, $`e`$, of $`\mathrm{\Gamma }`$; given a vertex, $`p`$, let $$_p=\{eE_\mathrm{\Gamma };p\text{ is a vertex of }e\text{ and }\alpha _e(\xi )>0\},$$ and let $`\sigma _p`$ be the number of edges $`e_p`$ for which $`p=t(e)`$. For $`e_p`$ define $$(1)^e=\{\begin{array}{cc}1\hfill & \text{ if }p=i(e)\hfill \\ 1\hfill & \text{ if }p=t(e)\hfill \end{array},$$ and let $`(1)^p`$ $`=`$ $`{\displaystyle \underset{e_p}{}}(1)^e=(1)^{\sigma _p}`$ $`\delta _p`$ $`=`$ $`\frac{1}{2}{\displaystyle \alpha _e},e_p`$ $`\delta _p^\mathrm{\#}`$ $`=`$ $`\frac{1}{2}{\displaystyle (1)^e\alpha _e},e_p.`$ ###### Definition 1.5. The Kostant partition function $$N_p:_G^{}$$ is the function which assigns to every weight, $`\alpha `$, the number of distinct ways in which $`\alpha `$ can be written as a sum $$\alpha =n_e\alpha _e,e_p$$ with non-negative integer coefficients. ###### Theorem 1.4. The multiplicity with which a weight, $`\alpha `$, occurs in $`Q(f)`$ is equal to $$\underset{p}{}(1)^pN_p(\alpha \alpha _p+\delta _p^\mathrm{\#}\delta _p).$$ (1.14) (Compare with \[GLS, (1.13)\].) The next results which we will describe involve a graph-theoretical analog of the notion of “reduction by a circle action” in symplectic geometry. Let $`T`$ be a circle subgroup of $`G`$ which is not contained in any of the groups, $`G_e`$. Then if $`\xi `$ is the infinitesimal generator of $`T`$ $$\alpha _e(\xi )0$$ (1.15) for all $`e`$. A function, $`\varphi :V`$ is called a $`T`$-moment map if for all edges $`eE_\mathrm{\Gamma }`$ $$\frac{\varphi (t(e))\varphi (i(e))}{\alpha _e(\xi )}>0.$$ (1.16) We recall (\[GZ1, § 2.2\]) that there is a simple necessary and sufficient condition for the existence of such a map. By (1.6) one can orient $`\mathrm{\Gamma }`$ by assigning to each unoriented edge the orientation for which $`\alpha _e(\xi )>0`$. Then, for the existence of a $`T`$-moment map, it is necessary and sufficient that this graph have no oriented cycles. We will call the numbers, $`\varphi (p)`$, the critical values of $`\varphi `$. By perturbing $`\varphi `$ slightly one can arrange that these $`\varphi (p)`$’s are all distinct. Let $`c`$ be a regular (non-critical) value of $`\varphi `$, and let $`V_c`$ be the set of all oriented edges, $`e`$, of $`\mathrm{\Gamma }`$ with $`\varphi (t(e))>c>\varphi (i(e))`$. One can make $`V_c`$ into the set of vertices of a new object, $`\mathrm{\Gamma }_c`$, and this object is our graph-theoretical “reduction of $`\mathrm{\Gamma }`$ at $`c`$”. (Unfortunately, $`\mathrm{\Gamma }_c`$ is not a graph. It is a slightly more complicated object: a “hypergraph”. For details see \[GZ3, § 3\].) Now fix an element, $`f`$, of $`K_G(\mathrm{\Gamma })`$, and for every edge, $`e`$, in $`V_c`$ let $`p=i(e)`$ and let $$\widehat{f}_e=f_p\underset{e^{}}{}(1e^{2\pi i\alpha _e^{}})^1$$ (1.17) the product being over all $`e^{}`$ with $`i(e^{})=p`$ and $`e^{}e`$. By composing the inclusion map of $`G_c`$ into $`G`$ with the projection of $`G`$ onto $`G/T`$, one gets a surjective finite-to-one map $`\pi _e:G_eG/T`$ and hence a “push-forward” in $`K`$-theory (see Section 3) $$(\pi _e)_{}:R(G_e)R(G/T).$$ This can be formally extended to elements of the quotient ring of $`R(G_e)`$ of the form (1.17), and by applying it to (1.17) one gets, for every vertex of $`\mathrm{\Gamma }_c`$, an element $$f_c^\mathrm{\#}(e)=(\pi _e)_{}r_e\widehat{f}_e$$ (1.18) of a quotient ring of $`R(G/T)`$. ###### Theorem 1.5. The sum $$\chi _c(f)=\underset{eV_c}{}f_c^\mathrm{\#}(e)$$ is in $`R(G/T)`$. We will prove this by proving a stronger result. Let $$f_p^\mathrm{\#}=f_p\underset{i(e)=p}{}(1e^{2\pi i\alpha _e})^1$$ be the $`p`$$`^{\text{th}}`$ summand on the right hand side of (1.10); and, for $`gG`$, consider the integral over $`T`$ $$f_p^\mathrm{\#}(gt)𝑑t.$$ (1.19) We will see in Section 5 that the integrand has poles at a finite number of points, $`t_iT`$ so this integral as it stands isn’t well defined. However, one can “regularize” it by moving the contour of integration to a curve in $`T^{}`$ which surrounds the $`t_i`$’s; and, denoting this regularized integral by $`\text{Res}_Tf_p^\mathrm{\#}`$ we will prove: ###### Theorem 1.6. $`\text{Res}_Tf_p^\mathrm{\#}`$ is an element of $`R(G/T)`$. Our strengthened version of Theorem 1.5 asserts: ###### Theorem 1.7. $`\chi _c(f)`$ is equal to the sum $$\underset{\varphi (p)>c}{}Res_Tf_p^\mathrm{\#}.$$ (1.20) Next we will explain what “quantization commutes with reduction” translates into the context of graphs. Recall that an element, $`f`$, of $`K_G(\mathrm{\Gamma })`$ of the form $$f_p=e^{2\pi i\alpha _p},\alpha _p_G^{}$$ is symplectic if $$\alpha _q\alpha _p=m_e\alpha _e,m_e>0$$ for every pair of vertices, $`p`$ and $`q`$, and edge, $`e`$ joining $`p`$ to $`q`$. For $`f`$ symplectic, the map $$\varphi :V,p\alpha _p(\xi )$$ is a $`T`$-moment map. Assume zero is a regular value of this map, i.e. $`\alpha _p(\xi )0`$ for all $`p`$; and let $`\mathrm{\Gamma }_{red}=\mathrm{\Gamma }_0`$ and $`\chi _{red}=\chi _0`$ . ###### Theorem 1.8. Let $`Q(\mathrm{\Gamma })`$ be the virtual representation of $`G`$ with character, $`\chi (f)`$ and $`Q(\mathrm{\Gamma }_{red})`$ the virtual representation of $`G/T`$ with character, $`\chi _{red}(f)`$. Then, as virtual representations of $`G/T`$ $$Q(\mathrm{\Gamma }_{red})=Q(\mathrm{\Gamma })^T.$$ (1.21) Finally in the last section of this paper we will show that if $`M`$ is a GKM manifold and $`\mathrm{\Gamma }`$ is its “one-skeleton”, these theorems about graphs have $`K`$-theoretic implications for $`M`$ (thanks to a beautiful recent result of Allen Knutson and Ioanid Rosu which asserts that $`K_G(M)K_G(\mathrm{\Gamma })`$). ## 2. Some algebraic preliminaries We will collect in this section some elementary facts about lattices and tori which will be needed in the proofs. Let $`V`$ be an $`n`$-dimensional real vector space and let $`L`$ be a rank $`n`$ lattice sitting inside $`V`$. Let $$L^{}=\{\alpha V^{};\alpha (v)\text{ for all }vL\}$$ be the dual lattice in $`V^{}`$. An element of $`L`$ is primitive if it is not of the form, $`kv`$, with $`vL`$ and $`\|k\|>1`$. ###### Lemma 2.1. $`vL`$ is primitive if and only if there is an $`\alpha L^{}`$ with $`\alpha (v)=1`$. ###### Lemma 2.2. $`v`$ is primitive if and only if there exists a basis $`v_1,\mathrm{},v_n`$ of $`L`$ with $`v=v_1`$. Now let $`G`$ be an $`n`$-dimensional torus and let $`𝔤`$ be its Lie algebra. ###### Definition 2.1. The group lattice of $`G`$, $`_G`$, is the kernel of the exponential map, $`\mathrm{exp}:𝔤G`$ and its dual, $`_G^{}`$, is the weight lattice of $`G`$. In particular $$G=𝔤/_G$$ and the exponential map is just the projection of $`𝔤`$ onto $`𝔤/_G`$. Given a weight, $`\alpha _G^{}`$, let $`\chi _\alpha `$ be the character of $`G`$ defined by $$\chi _\alpha (g)=e^{2\pi i\alpha (x)},g=\mathrm{exp}(x).$$ ###### Proposition 2.1. If $`\alpha `$ is primitive, the subgroup $$G_\alpha =\{gG;\chi _\alpha (g)=1\}$$ (2.1) is connected, i.e. is an $`(n1)`$-dimensional subtorus of $`G`$. More generally, if $`\beta `$ is primitive and $`\alpha =k\beta `$, $`k>1`$, the identity component of $`G_\alpha `$ is $`G_\beta `$ and $`G_\alpha /G_\beta `$ is a finite cyclic group of order $`k`$. Let $`\xi `$ be a primitive element of $`_G`$ and let $$G_\xi =\{\mathrm{exp}(t\xi );0t<1\}.$$ (2.2) then $`G_\xi `$ is a closed, connected one-dimensional subgroup of $`G`$. ###### Proposition 2.2. If $`\alpha (\xi )=0`$, then $`G_\xi G_\alpha `$ and if $`\alpha (\xi )0`$, then $`G_\xi G_\alpha `$ is a finite cyclic subgroup of order $`\|\alpha (\xi )\|`$. Let $`G_1=G/G_\xi `$ and let $`\gamma :G_\alpha G_1`$ be the composition of the inclusion, $`G_\alpha G`$, and the projection, $`GG_1`$. ###### Corollary 2.1. The map $`\gamma `$ is surjective and its kernel is a cyclic subgroup of $`G_\alpha `$ of order $`\|\alpha (\xi )\|`$. ## 3. The representation ring $`R(G)`$ The groups, $`G`$, in this section will be compact commutative Lie groups. For such a group every irreducible representation is one-dimensional, i.e. is defined by a homomorphism of $`G`$ into $`S^1`$. Thus the elements of $`R(G)`$ can be identified with the character ring of $`G`$: all finite sums of the form $$m_i\chi _i$$ (3.1) $`m_i`$ being an integer and $`\chi _i`$ a homomorphism of $`G`$ into $`S^1`$ (or “character”.) Hence, if $`G`$ is an $`n`$-torus, (3.1) is a sum of the form (1.9). In this section we will discuss some functorial properties of this ring. First we note that $`R(G)`$ is naturally a contravariant functor, i.e. if $`\gamma :GH`$ is a homomorphism of Lie groups, then a representation of $`H`$ can be converted, by composition with $`\gamma `$, into a representation of $`G`$; so there is a natural map $$\gamma ^{}:R(H)R(G)$$ (3.2) and it is easy to see that this is an algebra homomorphism. A much more interesting object for us will be a map in the opposite direction $$\gamma _{}:R(G)R(H)$$ (3.3) which we will define here modulo the assumption $$()\text{the kernel and cokernel of}\gamma \text{are finite.}$$ First let’s assume that $`\gamma `$ is surjective, i.e. that $`H=G/W`$ and that $`W`$ is a finite subgroup of $`G`$. Let $`\rho `$ be a representation of $`G`$ on a vector space, $`V`$, and let $`V^W`$ be the vectors in $`V`$ which transform trivially under $`W`$. Then the restriction of $`\rho `$ to $`V^W`$ is a representation, $`\rho ^W`$, of $`G/W`$ and $`\gamma _{}`$ is the map defined by $`\rho \rho ^W`$. Next assume that $`\gamma `$ is injective, i.e. that $`G`$ is a closed subgroup of $`H`$ and $`G\backslash H`$ is finite. Given a representation, $`\rho `$, of $`G`$ on a vector space, $`V`$, let $`\rho _{ind}`$ be the induced representation of $`H`$ (i.e. let $`V_{ind}`$ be the vector space consisting of maps $`f:HV`$ which satisfy $`f(gh)=\rho (g)f(h)`$ and let $$(\rho _{ind}f)(k)=f(kh^1)$$ for all $`kH`$). In this case $`\gamma _{}`$ is the map defined by $`\rho \rho _{ind}`$. Finally if $`\gamma `$ is neither injective nor surjective, let $`G_1`$ be the image of $`\gamma `$. Then $`\gamma `$ factors into the submersion $`\gamma _1:GG_1`$, composed with the inclusion, $`\gamma _2:G_1H`$, and one defines $$\gamma _{}=(\gamma _2)_{}(\gamma _1)_{}.$$ (3.4) This map is unfortunately not a ring homomorphism, but it is a morphism of $`R(H)`$-modules: for $`\chi R(G)`$ and $`\tau R(H)`$ $$\gamma _{}(\chi \gamma ^{}\tau )=(\gamma _{}\chi )\tau .$$ (3.5) We will mostly be interested in the case when $`\gamma `$ is a submersion, i.e. when $`H=G/W`$. In this case one has an alternative way of looking at $`\gamma `$: ###### Lemma 3.1. Let $`\rho `$ be a unitary representation of $`G`$ on a complex vector space, $`V`$. Then the orthogonal projection of $`V`$ onto $`V^W`$ is given by the operator $$P=\frac{1}{\|W\|}\underset{wW}{}\rho (w).$$ (3.6) ###### Proof. If $`vV^W`$, $`\rho (w)v=v`$, so $`Pv=v`$. Moreover, for all $`vV`$ and $`aW`$, $$\rho (a)Pv=\frac{1}{\|W\|}\underset{wW}{}\rho (aw)v=Pv,$$ so $`PvV^W`$. Finally, $$P^{}=\frac{1}{\|W\|}\underset{wW}{}\rho (w^1)=P,$$ so $`P`$ is the orthogonal projection of $`V`$ onto $`V^W`$. ∎ ###### Corollary 3.1. Let $`g`$ be an element of $`G`$ and let $`\overline{g}`$ be its image in $`G/W`$. Then $$(\gamma _{}\rho )(\overline{g})=\frac{1}{\|W\|}\underset{wW}{}\rho (gw).$$ (3.7) In particular, let $`f:G`$ be the function (3.1), i.e. $$f(g)=m_i\chi _i(g).$$ (3.8) Then $$\gamma _{}f(\overline{g})=\frac{1}{\|W\|}\underset{wW}{}f(gw).$$ (3.9) ## 4. Convexity and multiplicities ###### The proof of Theorem 1.1. : Let $`\alpha _1,\mathrm{},\alpha _N`$ be primitive vectors such that for every $`eE_\mathrm{\Gamma }`$ there exists a unique $`k\{1,\mathrm{},N\}`$ such that $`\alpha _e`$ is a multiple of $`\alpha _k`$. If $`m_1\alpha _1,\mathrm{},m_s\alpha _1`$ are all the occurrences of multiples of $`\alpha _1`$ among all the weights, let $`M_1=\text{l.c.m.}(m_1,\mathrm{},m_s).`$ Similarly we define $`M_2,\mathrm{},M_N`$. Then $$\chi (f)=\frac{g}{_{j=1}^N(1e^{2\pi iM_j\alpha _j})}$$ (4.1) with $`gR(G)`$. We will show that $`1e^{2\pi iM_1\alpha _1}`$ divides $`g`$ in $`R(G)`$. The vertices of $`\mathrm{\Gamma }`$ can be divided into two categories: 1. The first subset, $`V_1`$, contains the vertices $`pV_\mathrm{\Gamma }`$ for which none of the $`\alpha _e`$’s with $`i(e)=p`$, is a multiple of $`\alpha _1`$ 2. The second subset, $`V_2`$, contains the vertices $`pV_\mathrm{\Gamma }`$ for which there exists an edge $`e`$ such that $`i(e)=p`$ and $`\alpha _e`$ is a multiple of $`\alpha _1`$. (Notice that there will be exactly one such edge.) The part of (1.10) corresponding to vertices in the first category will then be of the form $$\underset{pV_1}{}f_p\underset{i(e)=p}{}(1e^{2\pi i\alpha _e})^1=g_1\underset{j=2}{\overset{N}{}}(1e^{2\pi iM_j\alpha _j})^1$$ (4.2) with $`g_1R(G).`$ If $`pV_2`$ then there exists an edge $`e`$ issuing from $`p`$ such that $`\alpha _e=m\alpha _1`$ with $`m\{0\}`$; let $`q=t(e)`$. Since $`\alpha _{\overline{e}}=\alpha _e`$ it follows that $`qV_2`$ as well and thus the vertices in $`V_2`$ can be paired as above. Let $`e_k,k=1,\mathrm{},d`$ and $`e_k^{},k=1,\mathrm{},d`$ be the edges issuing from $`p`$ and $`q`$ respectively, with $`e_d=e`$, $`e^{}=\overline{e}`$. Then, by (1.4), the $`e_k`$’s can be ordered so that $$r_e(e^{2\pi i\alpha _{e_k}})=r_e(e^{2\pi i\alpha _{e_k^{}}}),\text{ for all }k=1,\mathrm{},d1$$ which implies that $$1e^{2\pi i\alpha _{e_k}}1e^{2\pi i\alpha _{e_k^{}}}(mod1e^{2\pi i\alpha _e}).$$ (4.3) Similarly, from $$r_e(f_p)=r_e(f_q)$$ we deduce that $$f_qf_p(mod1e^{2\pi i\alpha _e}).$$ (4.4) The part of (1.10) corresponding to $`p`$ and $`q`$, $$f_p\underset{j=1}{\overset{d}{}}(1e^{2\pi i\alpha _{e_j}})^1+f_q\underset{j=1}{\overset{d}{}}(1e^{2\pi i\alpha _{e_j^{}}})^1,$$ can be expressed as $$\frac{f_p_{j=2}^N(1e^{2\pi i\alpha _{e_j^{}}})e^{2\pi im\alpha _1}f_q_{j=2}^N(1e^{2\pi i\alpha _{e_j}})}{(1e^{2\pi im\alpha _1})_{j=2}^N(1e^{2\pi i\alpha _{e_j}})_{j=2}^N(1e^{2\pi i\alpha _{e_j^{}}})}.$$ (4.5) From the congruences (4.3) and (4.4) we conclude that $`1e^{2\pi im\alpha _1}`$ divides the numerator of (4.5), so we deduce that $$f_p\underset{j=1}{\overset{d}{}}(1e^{2\pi i\alpha _{e_j}})^1+f_q\underset{j=1}{\overset{d}{}}(1e^{2\pi i\alpha _{e_j^{}}})^1=\frac{g_{p,q}}{_{j=2}^N(1e^{2\pi iM_j\alpha _j})}$$ with $`g_{p,q}R(G)`$. Therefore $$\underset{pV_2}{}f(p)\underset{i(e)=p}{}(1e^{2\pi i\alpha _e})^1=\frac{g_2}{_{j=2}^N(1e^{2\pi iM_j\alpha _j})}$$ (4.6) with $`g_2R(G)`$. Adding (4.2) and (4.6) we obtain $$\frac{g}{_{j=1}^N(1e^{2\pi iM_j\alpha _j})}=\frac{g_1+g_2}{_{j=2}^N(1e^{2\pi iM_j\alpha _j})}$$ with $`g_1+g_2R(G)`$, hence $`1e^{2\pi iM_1\alpha _1}`$ divides $`g`$. The same argument can be used to show that each $`1e^{2\pi iM_j\alpha _j}`$ divides $`g`$. The proof of the theorem now follows from: ###### Lemma 4.1. If $`PR(G)`$ and $`\alpha `$, $`\beta `$ are linearly independent weights such that $`1e^{2\pi i\alpha }`$ divides $`(1e^{2\pi i\beta })P`$, then $`1e^{2\pi i\alpha }`$ divides $`P`$. ###### Lemma 4.2. If $`PR(G)`$ and $`\beta _1,\mathrm{},\beta _k`$ are pairwise linearly independent weights such that $`1e^{2\pi i\beta _j}\text{ divides }P`$ for all $`j=1,\mathrm{},k`$ then $$(1e^{2\pi i\beta _1})\mathrm{}(1e^{2\pi i\beta _k})\text{ divides }P.\mathit{}$$ ###### The proof of Theorem 1.2. : Let $`\alpha `$ be a weight that is not in the convex hull of $`\{\alpha _p;pV_\mathrm{\Gamma }\}`$. Then there exists $`\xi 𝔤`$ and $`p_0V_\mathrm{\Gamma }`$ such that $`(\alpha \alpha _{p_0})(\xi )<0`$ and $`(\alpha _p\alpha _{p_0})(\xi )>0`$ for all $`pp_0`$. If $`eE_\mathrm{\Gamma }`$ and $`\alpha _e(\xi )<0`$ then $$(1e^{2\pi i\alpha _e})^1=e^{2\pi i\alpha _e}(1e^{2\pi i\alpha _{\overline{e}}})^1,$$ and using this we deduce that $$\chi (f)=\underset{pV}{}(1)^{\sigma _p}e^{2\pi i(\alpha _p^{}\alpha _e)}\underset{e_p}{}(1e^{2\pi i\alpha _e})^1,$$ (4.7) where $`^{}\alpha _e`$ in the exponent is the sum $$\underset{\begin{array}{c}i(e)=p\\ \alpha _e(\xi )<0\end{array}}{}\alpha _e=\underset{\begin{array}{c}t(e)=p\\ \alpha _e(\xi )>0\end{array}}{}\alpha _e=\delta _p^\mathrm{\#}\delta _p.$$ (4.8) From (4.7) and (4.8) we deduce that $$\chi (f)=\underset{pV}{}(1)^pe^{2\pi i(\alpha _p^{}\alpha _e)}\underset{e_p}{}(\underset{k_e0}{}e^{2\pi ik_e\alpha _e}).$$ (4.9) Suppose $`\alpha `$ is a weight of $`Q(f)`$; then there exists $`pV_\mathrm{\Gamma }`$ and non-negative integers $`\{k_e\}_{e_p}`$ such that $$\alpha =\alpha _p+\underset{\begin{array}{c}t(e)=p\\ \alpha _e(\xi )>0\end{array}}{}\alpha _e+\underset{e_p}{}k_e\alpha _e,$$ (4.10) which implies that $$\alpha \alpha _{p_0}=\alpha _p\alpha _{p_0}+\underset{\begin{array}{c}t(e)=p\\ \alpha _e(\xi )>0\end{array}}{}\alpha _e+\underset{e_p}{}k_e\alpha _e.$$ (4.11) But when we evaluate (4.11) at $`\xi `$, the right hand side is non-negative, while the left hand side is strictly negative ! This contradiction proves that $`\alpha `$ is not a weight of $`Q(f)`$. ∎ ###### The proof of Theorem 1.3. : Let $`\alpha =\alpha _{p_0}`$ be an extremal weight, i.e. a vertex of $`\mathrm{\Delta }`$. Then there exists $`\xi 𝔤`$ such that $`(\alpha _p\alpha _{p_0})(\xi )>0`$ for all $`pp_0`$. In this case (4.11) implies $$0=(\alpha _p\alpha _{p_0})(\xi )+\underset{\begin{array}{c}t(e)=p\\ \alpha _e(\xi )>0\end{array}}{}\alpha _e(\xi )+\underset{e_p}{}k_e\alpha _e(\xi ).$$ (4.12) Since each term on the right hand side is non-negative, (4.12) is only true if 1. $`p=p_0`$ (which also implies that $`\alpha _e(\xi )<0`$ for all $`e`$ with $`t(e)=p`$, i.e. that there are no terms in the first sum), and 2. $`k_e=0`$ for all $`e_p`$. This proves that the multiplicity with which $`\alpha `$ occurs in $`Q(f)`$ is 1. ∎ ###### The proof of Theorem 1.4. : From (4.8) and (4.10), the multiplicity with which a weight $`\alpha `$ appears in the term corresponding to the vertex $`p`$ is equal to $`(1)^p`$ times the number of distinct ways in which $`\alpha \alpha _p+\delta _p^\mathrm{\#}\delta _p`$ can be written as a sum $$\underset{e_p}{}k_e\alpha _e,$$ with $`k_e`$’s non-negative integers; and this number is $`N_p(\alpha \alpha _p+\delta _p^\mathrm{\#}\delta _p)`$. Counting the contributions given by all the vertices we obtain (1.14). ∎ ## 5. The residue operation Let $`G`$ be an $`n`$-dimensional torus, let $`T`$ be a circle subgroup of $`G`$, and let $$\chi _k=e^{2\pi i\alpha _k},k=1,\mathrm{},d$$ be characters of $`G`$ and $`f`$ an element of the character ring, $`R(G)`$. The goal of this section is to make sense of the integral $$_T\frac{f(gt)}{(1\chi _k(gt))}𝑑t$$ (5.1) as a function of $`gG`$. If the restriction of $`\chi _k`$ to $`T`$ is identically one, the denominator in the integrand is identically zero when $`\chi _k(g)=1`$. Hence, for (5.1) to make sense, we are forced to assume that the restriction of $`\chi _k`$ to $`T`$ is not identically one. Even with this assumption, however, the integrand has poles at the points where $`\chi _k(gt)=1`$; so to make sense of (5.1) we must “regularize” this integral and this we will do as follows. Fix a basis vector, $`\xi `$ of $`_T`$, and identify $`T`$ with $`S^1`$ via the map $$\mathrm{exp}(s\xi )e^{2\pi is}.$$ Then, with $`z=e^{2\pi ix}`$, the integrand of (5.1) becomes a meromorphic function $$f^\mathrm{\#}(gz)=f(gz)\underset{k=1}{\overset{d}{}}(1\chi _k(gz))^1$$ (5.2) on the complex plane with poles on the unit circle. Now move the contour of integration from the unit circle to a contour surrounding these poles, e.g. a circle of radius greater than one oriented in a counter-clock-wise sense plus a circle of radius less than one oriented in a clock-wise sense. In other words, replace (5.1) by the integral $$\frac{1}{2\pi i}_{C_+}f^\mathrm{\#}(gz)\frac{dz}{z}$$ (5.3) minus the integral $$\frac{1}{2\pi i}_C_{}f^\mathrm{\#}(gz)\frac{dz}{z},$$ (5.4) $`C_+`$ being a circle of radius greater than one and $`C_{}`$ a circle of radius less than one, both these circle being oriented in a counter-clock-wise sense. Let us denote this regularized integral, i.e. the difference of (5.3) and (5.4), by $`(Res_Tf^\mathrm{\#})(g)`$. It is easy to see that this function is $`T`$-invariant, $$(Res_Tf^\mathrm{\#})(gt)=(Res_Tf^\mathrm{\#})(g)$$ and hence defines a function on $`G/T`$. We will prove: ###### Theorem 5.1. $`Res_Tf^\mathrm{\#}`$ is an element of $`R(G/T)`$. ###### Remark. The definition of $`Res_Tf^\mathrm{\#}`$ depends on the identification of $`S^1`$ with $`T`$ given by $`\mathrm{exp}(x\xi )e^{2\pi ix}`$. If we replace $`\xi `$ by $`\xi `$, the orientations of the circles, $`C_+`$ and $`C_{}`$, will get reversed, and hence this will change the signs of (5.3) and (5.4). In proving Theorem 5.1, we can assume without loss of generality that $`f=e^{2\pi i\alpha }`$, $`\alpha _G^{}`$. Let $`e_1,..,e_n`$ be a basis of $`_G`$ with $`\xi =e_n`$ and let $`y_1,\mathrm{},y_{n1}`$ and $`x`$ be the coordinates on $`𝔤`$ associated with this basis. We can then write $$\alpha _i=k_ix+\beta _i(y)$$ (5.5) and $$\alpha =kx+\beta (y)$$ (5.6) with $`k_i=\alpha _i(\xi )`$ and $`k=\alpha (\xi )`$. Thus letting $`z`$ $`=`$ $`e^{2\pi ix}`$ (5.7) $`a_i`$ $`=`$ $`e^{2\pi i\beta _i(y)}`$ (5.8) $`b`$ $`=`$ $`e^{2\pi i\beta (y)}`$ (5.9) the integrals (5.3) and (5.4) become $$\frac{1}{2\pi i}_{C_+}\frac{bz^k}{(1a_iz^{k_i})}\frac{dz}{z}$$ (5.10) and $$\frac{1}{2\pi i}_C_{}\frac{bz^k}{(1a_iz^{k_i})}\frac{dz}{z}.$$ (5.11) Therefore, to prove the theorem, it suffices to show that each of these integrals individually is in $`R(G/T)`$. To verify this for (5.11), let’s order the factors in the denominator of the integrand so that $`k_i=k_i^{}<0`$ for $`1ir`$ and $`k_i>0`$ for $`r+1id`$. This integrand is then equal to $$\frac{bz^{k^{}1}}{_{i=1}^r(z^{k_i^{}}a_i)_{i=r+1}^d(1a_iz^{k_i})}$$ (5.12) with $`k^{}=kk_1\mathrm{}k_r`$. Hence, if $`k^{}>0`$, (5.12) is holomorphic at zero; so, in particular: ###### Lemma 5.1. The integral (5.4) is zero if $`k>k_1+\mathrm{}+k_r`$. For $`1ir`$ and $`z0`$, let $`a_i^{}=a_i^1`$ and let $$S_i(z)=\frac{1}{z^{k_i^{}}a_i}=\frac{a_i}{1a_i^{}z^{k_i^{}}}=a_i\underset{l=0}{\overset{\mathrm{}}{}}(a_i^{}z^{k_i^{}})^l,$$ (5.13) and, for $`r+1id`$, let $$S_i(z)=\frac{1}{1a_iz^{k_i}}=\underset{l=0}{\overset{\mathrm{}}{}}(a_iz^{k_i})^l.$$ (5.14) Then the integral (5.11) is just the degree -1 term in the Laurent series $$bz^{k^{}1}\underset{i=1}{\overset{d}{}}S_i(z)$$ (5.15) and this term is clearly a polynomial in $`b,a_1,..,a_r,a_1^1,\mathrm{},a_r^1`$, and $`a_{r+1},\mathrm{},a_d`$, with integer coefficients. Hence, by (5.8) and (5.9), it is clearly a trigonometric polynomial in $`y_1,..,y_{n1}`$. From this, together with Lemma 5.1, we conclude: ###### Theorem 5.2. The integral (5.4) is an element of $`R(G/T)`$. Moreover, if $`f=e^{2\pi i\alpha }`$ and $`k=\alpha (\xi )`$, this integral is zero if $`k>k_1+\mathrm{}+k_r`$. To evaluate the integral (5.10) we make the substitution, $`zz^1`$ and reduce this integral to an integral of the type we’ve just evaluated. We conclude ###### Theorem 5.3. The integral (5.3) is an element of $`R(G/T)`$. Moreover, if $`f=e^{2\pi i\alpha }`$ and $`k=\alpha (\xi )`$, this integral is zero if $`k<k_{r+1}+\mathrm{}+k_d`$. Since $`k_i`$ is negative for $`1ir`$ and positive for $`r+1id`$ we have, in particular: ###### Proposition 5.1. If $`f=e^{2\pi i\alpha }`$ with $`k=\alpha (\xi )`$, the integral (5.4) is zero when $`k`$ is positive and the integral (5.3) is zero when $`k`$ is negative. Moreover, if $`k=0`$, (5.4) is zero when $`r>0`$ and (5.3) is zero when $`dr>0`$. Suppose now that the weights, $`\alpha _i`$, $`i=1,..,d`$ are pairwise linearly independent, i.e. suppose that $`\alpha _i`$ and $`\alpha _j`$ are linearly independent for $`ij`$. Then the integrand in (5.10) - (5.11): $$bz^{k1}\underset{i=1}{\overset{d}{}}(1a_iz^{k_i})^1$$ (5.16) has simple poles on the unit circle for generic values of $`y`$. (Recall that since $`a_k=e^{2\pi i\beta _k(y)}`$, the location of these poles depends on $`y`$.) Thus, one can compute the difference between (5.10) and (5.11) by computing the residues of (5.16) at these poles. We will show that the sum of these residues, which is, by definition, the regularized integral (5.1), is given by an expression involving the $`K`$-theoretic push-forward which we described in Section 3. More explicitly, let $`G_i`$ be the kernel of the homomorphism $`\chi _i:GS^1`$ and let $`r_i`$ be the restriction map $`R(G)R(G_i)`$ and $`\pi _i`$ the projection of $`G_i`$ onto $`G/T`$. We will prove: ###### Theorem 5.4. Let $`f^\mathrm{\#}`$ be the function (5.2) and let $$\widehat{f_i}=f\underset{ji}{}(1\chi _j)^1.$$ Then $$Res_Tf^\mathrm{\#}=\underset{i=1}{\overset{r}{}}(\pi _i)_{}r_i\widehat{f_i}\underset{i=r+1}{\overset{d}{}}(\pi _i)_{}r_i\widehat{f_i}.$$ (5.17) ###### Proof. Let $`\theta _1,\mathrm{},\theta _d`$ be real numbers, let $`k_1,\mathrm{},k_d`$ be integers and let $`a_i=e^{2\pi i\theta _i}`$. As above we will order the $`k_i`$’s so that $`k_i<0`$ for $`1ir`$ and $`k_i>0`$ for $`r+1id`$. Let $`g(z)`$ be the function (5.16), i.e. $$g(z)=bz^{k1}\underset{i=1}{\overset{d}{}}(1a_iz^{k_i})^1.$$ ###### Lemma 5.2. Suppose that, for $`ij`$, $`\theta _i,\theta _j`$ and 1 are linearly independent over the rationals. Then $`g(z)`$ has simple poles on the unit circle. ###### Proof. Let $$\omega _i=e^{2\pi i/k_i}\text{ and }a_i^{1/k_i}=e^{2\pi i\theta _i/k_i}.$$ Then these poles are at the points $$\omega _i^la_i^{1/k_i},1lk_i,\mathrm{\hspace{0.33em}1}id;$$ (5.18) so if $`\theta _i,\theta _j`$ and 1 are linearly independent over the rationals these poles are distinct. ∎ Let us compute the residue of $`g(z)`$ at the pole (5.18). The quotient $$\frac{z\omega _i^la_i^{1/k_i}}{1a_iz^{k_i}}$$ evaluated at $`z=\omega _i^la_i^{1/k_i}`$ is equal, by l’Hopital’s rule, to: $$\frac{1}{a_ik_iz^{k_i1}}\text{ or, alternatively }\frac{z}{a_ik_iz^{k_i}}$$ evaluated at $`z=\omega _i^la_i^{1/k_i}`$, and since $`(\omega _i^la_i^{1/k_i})^{k_i}=a_i^1`$, this quotient is just $$\frac{1}{k_i}\omega _i^la_i^{1/k_i}.$$ (5.19) Thus the residue at $`z=\omega _i^la_i^{1/k_i}`$ of the function $$g(z)=\frac{1}{1a_iz^{k_i}}bz^{k1}\underset{ji}{}(1a_jz^{k_j})^1$$ is just $$\frac{b}{k_i}(\omega _i^la_i^{1/k_i})^k\underset{ji}{}(1a_j(a_i^{1/k_i})^{k_j})^1$$ which, if we set $$b_i=ba_i^{k/k_i}$$ (5.20) and $$a_{j,i}=a_ja_i^{k_j/k_i},$$ (5.21) can be written $$\frac{1}{k_i}(\omega _i^l)^kb_i\underset{ji}{}(1(\omega _i^l)^{k_j}a_{j,i})^1.$$ (5.22) We will now show that if we give $`b`$ and $`a_i`$ the values (5.8) - (5.9) the sum of these residues is identical with the right hand side of (5.17). If $`b`$ is equal to (5.9) and $`a_i`$ is equal to (5.8), then by (5.5) and (5.6) $$b_i=e^{2\pi i\sigma _i}$$ (5.23) and $$a_{j,i}=e^{2\pi i\alpha _{j,i}},$$ (5.24) where $$\sigma _i=\alpha \frac{k}{k_i}\alpha _i$$ (5.25) and $$\alpha _{j,i}=\alpha _j\frac{k_j}{k_i}\alpha _i.$$ (5.26) Let’s now give a more “intrinsic” definition of $`\sigma _i`$ and $`\alpha _{j,i}`$: Let $`𝔤_i`$ be the Lie algebra of the group, $`G_i`$, and $`𝔱`$ the Lie algebra of $`T`$. Since $`G_i`$ is by definition the kernel of the homomorphism, $`e^{2\pi i\alpha _i}:GS^1`$, $`𝔤_i`$ is the annihilator of $`\alpha _i`$; so, by (5.25), $`\sigma _i`$ is the unique element of $`𝔤^{}`$ which is annihilated by $`𝔱`$ and has the same restriction to $`𝔤_i`$ as $`\alpha `$. Similarly, $`\alpha _{j,i}`$ is the unique element of $`𝔤^{}`$ which is annihilated by $`𝔱`$ and has the same restriction to $`𝔤_i`$ as $`\alpha _j`$. Note, by the way, that since $`\sigma _i`$ and $`\alpha _{j,i}`$ are annihilated by $`𝔱`$, they are in the dual vector space to $`𝔤/𝔱`$; or, in other words, in the dual of the Lie algebra of $`G/T`$. Consider the kernel of the map $`G_iG/T`$. This consists of the elements $$\mathrm{exp}(\frac{l}{k_i}\xi ),l=1,..,k$$ and by (5.5) and (5.6) $$e^{2\pi i\alpha }(\mathrm{exp}(\frac{l}{k_i}\xi ))=(\omega _i^l)^k$$ (5.27) and $$e^{2\pi i\alpha _j}(\mathrm{exp}(\frac{l}{k_i}\xi ))=(\omega _i^l)^{k_j}.$$ (5.28) Thus the sum $$\frac{1}{k_i}\underset{l=1}{\overset{k_i}{}}(\omega _i^l)^ke^{2\pi i\sigma _i}\underset{ji}{}(1(\omega _i^l)^{k_j}e^{2\pi i\alpha _{j,i}})^1$$ of the residues of $`g(z)`$ over the poles $`(\omega _i^l)a_i^{1/k_i}`$, $`1lk_i`$ is by formula (3.9) identical to the expression $$(\pi _i)_{}r_i\frac{e^{2\pi i\alpha }}{_{ji}(1e^{2\pi i\alpha _j})}$$ if $`r+1id`$ (in which case $`k_i=\|k_i\|`$) and is equal to $$(\pi _i)_{}r_i\frac{e^{2\pi i\alpha }}{_{ji}(1e^{2\pi i\alpha _j})}$$ when $`1ir`$ (in which case $`k_i=\|k_i\|`$). ∎ ## 6. Quantization commutes with reduction We will prove below Theorems 1.7 and 1.8 of Section 1. As in Theorem 1.7, let $`f`$ be an element of $`K_G(\mathrm{\Gamma })`$, let $`\varphi :V`$ be a $`T`$-moment map, let $`c`$ be a regular value of $`\varphi `$ and let $`e`$ be an oriented edge of $`\mathrm{\Gamma }`$ with $`\varphi (q)>c>\varphi (p),`$ where $`p=i(e)`$ and $`q=t(e)`$ (i.e. $`e`$ corresponds to a vertex of the hypergraph, $`\mathrm{\Gamma }_c`$; we will denote this vertex by $`e`$, as well.) Consider the expressions $`\widehat{f_e}`$ $`=`$ $`f_p{\displaystyle \underset{e^{}}{}}(1e^{2\pi i\alpha _e^{}})^1`$ (6.1) $`\widehat{f_{\overline{e}}}`$ $`=`$ $`f_q{\displaystyle \underset{e^{\prime \prime }}{}}(1e^{2\pi i\alpha _{e^{\prime \prime }}})^1`$ (6.2) the product in (6.1) being over all edges, $`e^{}e`$, with $`i(e^{})=p`$, and the product in (6.2) being over all edges, $`e^{\prime \prime }\overline{e}`$, with $`i(e^{\prime \prime })=q`$. ###### Lemma 6.1. Let $`r_e=r_{\overline{e}}`$ be the restriction map $`R(G)R(G_e)`$. Then $$r_e\widehat{f_e}=r_{\overline{e}}\widehat{f_{\overline{e}}}.$$ (6.3) Let $`\pi _e=\pi _{\overline{e}}`$ be the projection of $`G_e`$ onto $`G/T`$. As a corollary of Lemma 6.1 we get two alternative ways of defining (1.18): $$(\pi _e)_{}r_e\widehat{f_e}=(\pi _{\overline{e}})_{}r_{\overline{e}}\widehat{f_{\overline{e}}}=f_c^\mathrm{\#}(e),$$ (6.4) and, as a consequence of (6.4), the following theorem: ###### Theorem 6.1. Let $`c`$ and $`c^{}`$ be regular values of $`\varphi `$. Suppose there exists just one vertex, $`p`$, with $`c<\varphi (p)<c^{}`$. Then $$\chi _c(f)\chi _c^{}(f)=Res_T\left(f_p\underset{i(e)=p}{}(1e^{2\pi i\alpha _e})^1\right)$$ (6.5) ###### Proof. If $`eV_c`$ and $`t(e)p`$, then $`eV_c^{}`$, and if $`eV_c^{}`$ and $`i(e)p`$, then $`eV_c`$. Moreover, in both cases, $$f_c^\mathrm{\#}(e)=f_c^{}^\mathrm{\#}(e),$$ (6.6) by (6.4). Thus, if $`e_i`$, $`i=1,..,r`$ are the elements of $`V_c`$ with $`t(e_i)=p`$, and $`e_i`$, $`i=r+1,\mathrm{},d`$, are the elements of $`V_c^{}`$ with $`i(e_i)=p`$, the difference between $`\chi _c(f)`$ and $`\chi _c^{}(f)`$ is, by (6.4), equal to $$\underset{i=1}{\overset{r}{}}f_c^\mathrm{\#}(e_i)\underset{i=r+1}{\overset{d}{}}f_c^{}^\mathrm{\#}(e_i),$$ or, also by (6.4), to $$\underset{i=1}{\overset{r}{}}(\pi _{e_i})_{}r_{e_i}\widehat{f}_{e_i}\underset{i=r+1}{\overset{d}{}}(\pi _{e_i})_{}r_{e_i}\widehat{f}_{e_i},$$ which, by (5.17), is identical with $$Res_T\left(f_p\underset{i(e)=p}{}(1e^{2\pi i\alpha _e})^1\right).\mathit{}$$ To prove Theorem 1.7, let $`c_0<c_1<\mathrm{}<c_N`$ be regular values of $`\varphi `$ with $`c_0=c`$, $`c_N`$ greater than $`\varphi _{max}`$, and with only one critical point, $`p_i`$, between $`c_i`$ and $`c_{i+1}`$. Then $$\chi _c(f)=\underset{i=0}{\overset{N}{}}(\chi _{c_i}(f)\chi _{c_{i+1}}(f))=\underset{\varphi (p)>c}{}Res_T\left(f_p\underset{i(e)=p}{}(1e^{2\pi i\alpha _e})^1\right),$$ proving Theorem 1.7. To prove Theorem 1.8, let $`f`$ be an element of $`K_G(\mathrm{\Gamma })`$ of the form (1.12) - (1.13) and let $`\varphi :V_\mathrm{\Gamma }`$, $`\varphi (p)=\alpha _p(\xi )`$. Then $$\chi (f)=\underset{p}{}e^{2\pi i\alpha _p}\underset{i(e)=p}{}(1e^{2\pi i\alpha _e})^1.$$ Identify $`T^{}`$ with $`0`$ and let $`C`$ be a circle in the complex plane with radius greater than one oriented in a counter-clock-wise sense. Then $$\frac{1}{2\pi i}_C\chi (f)(gz)\frac{dz}{z}=\underset{p}{}\frac{1}{2\pi i}_C\left(\frac{e^{2\pi i\alpha _p}}{_{i(e)=p}(1e^{2\pi i\alpha _e})}\right)(gz)\frac{dz}{z}.$$ (6.7) The right hand side of this identity is easy to evaluate: By Proposition 5.1, the summands with $`\alpha _p(\xi )<0`$ are zero, and the summands with $`\alpha _p(\xi )>0`$ are equal to $$Res_T\left(f_p\underset{i(e)=p}{}(1e^{2\pi i\alpha _e})^1\right),$$ so by Theorem 1.7 the right hand side is equal to $`\chi _{red}(f)`$. As for the the left hand side, by Theorem 1.1, $`\chi (f)`$ is in $`R(G)`$; so it is a finite sum of the form $$m_ke^{2\pi i\alpha _k}$$ with $`m_k`$ and $`\alpha _k_G^{}`$; and the integral over $`C`$ of the $`k`$-th term is zero except when $`e^{2\pi i\alpha _k}`$ doesn’t depend on $`z`$, in which case the integral is just $`2\pi ie^{2\pi i\alpha _k}`$. Hence the left hand side is equal to $$\underset{\alpha _k(\xi )=0}{}m_ke^{2\pi i\alpha _k},$$ which is the character of the representation, $`Q(\mathrm{\Gamma })^T`$. ## 7. GKM manifolds Let $`(M,\omega )`$ be a compact $`2d`$-dimensional symplectic manifold and $`\tau :G\times MM`$ a Hamiltonian action of $`G`$ on $`M`$. We will say that $`M`$ is a symplectic GKM manifold if $`M^G`$ is finite and if, for every $`pM^G`$, the weights $`\alpha _{i,p}_G^{}`$, $`i=1,..,d`$ of the isotropy representation of $`G`$ on $`T_pM`$ are pair-wise linearly independent. Let $$M^{(1)}=\{pM;dimG_pn1\}.$$ This set is called the one-skeleton of $`M`$; and $`M`$ is a GKM manifold if and only if $`M^{(1)}`$ consists of $`G`$-invariant imbedded 2-spheres, each of which contains exactly two fixed points. These 2-spheres can intersect at the fixed points; so the combinatorial structure of $`M^{(1)}`$ is that of a graph, $`\mathrm{\Gamma }`$, having the fixed points of $`\tau `$ as vertices and these 2-spheres as edges. For each oriented edge, $`e`$, of $`\mathrm{\Gamma }`$, let $`\varrho _e`$ be the isotropy representation of $`G`$ on the tangent space to this 2-sphere at the fixed point, $`t(e)`$; and for each vertex, $`p`$, of $`\mathrm{\Gamma }`$ let $`\tau _p`$ be the isotropy representation of $`G`$ on $`T_pM`$. It is easily checked that $`\varrho `$ and $`\tau `$ have properties (1.2) - (1.4) and hence define an action of $`G`$ on $`\mathrm{\Gamma }`$. For GKM manifolds the cohomology groups, $`H_G(\mathrm{\Gamma })`$ and $`K_G(\mathrm{\Gamma })`$, turn out to be equal to cohomology groups of $`M`$. More explicitly, let $`H_G(M)`$ be the equivariant cohomology ring of $`M`$ with complex coefficients and let $`K_G(M)`$ be the $`K`$-cohomology ring of $`M`$. Then there are ring homomorphisms $`H_G(M)`$ $``$ $`H_G(\mathrm{\Gamma }),(\text{see }\text{[GKM]})`$ (7.1) $`K_G(M)`$ $``$ $`K_G(\mathrm{\Gamma }),(\text{see }\text{[KR]}).`$ (7.2) With (7.1) and (7.2) as our point of departure, we will briefly describe some geometric implications of the theorems proved in this paper. The first of our results, Theorem 1.1, is a “combinatorial” explanation of why the right hand side of the Atiyah-Bott fixed point formula makes sense, i.e. why (1.10) does define a character of a virtual representation of $`G`$. Theorems 1.2 \- 1.4 are, in the manifold setting, well-known results about the “quantum” action of $`G`$ on $`M`$: Suppose $`[\omega ]H^2(M,)`$. Then there exists a line bundle, $`𝕃M`$, and a connection, $``$, on this bundle with $`curv()=\omega `$; and one says that the action, $`\tau `$, of $`G`$ on $`M`$ is pre-quantizable if it lifts to an action of $`G`$ on $`𝕃`$ preserving $``$. Now equip $`M`$ with a $`G`$-invariant Riemannian metric and let $$/_{}:S_{}^+S_{}^{}$$ be the $`spin^{}`$ Dirac operator. Given the connection, $``$, one can twist this operator with operator with $`𝕃`$ to get a Dirac operator $$/_{}^𝕃:S_{}^+𝕃S_{}^{}𝕃,$$ and the virtual vector space $$Q(M)=\text{kernel}(/_{}^𝕃)\text{cokernel}(/_{}^𝕃)$$ (7.3) is called the $`spin^{}`$-quantization of $`M`$. From the action of $`G`$ on $`𝕃`$, one gets a representation, $`\tau _Q`$, of $`G`$ on this space, and its character, $`\text{trace}\tau _Q`$, is equal, by the Atiyah-Bott formula, to the formal character, $`\chi (f)`$, defined by (1.10), $`f`$ being the element of $`K_G(\mathrm{\Gamma })`$ corresponding to $`[𝕃]`$ under the isomorphism (7.2). For $`\tau _Q`$, the convexity theorem (Theorem 1.2) is due to Guillemin and Sternberg, who pointed out in \[GS\] that it can be deduced from “quantization commutes with reduction” and the Atiyah-Guillemin-Sternberg convexity theorem for moment maps. (However, the simple proof of this theorem described in Section 4 seems to have eluded them.) As for Theorem 1.4, for co-adjoint orbits this is the celebrated Kostant Multiplicity Theorem. Our proof of it in Section 4 is modeled on Cartier’s proof of Kostant’s theorem in \[Ca\] and the symplectic version of the proof described in \[GLS\]. Let $`T`$ be a circle subgroup of $`G`$ and let $`M_c`$ be the reduction of $`M`$ with respect to $`T`$. For $`f=[𝕃]`$, the “reduced” character, $`\chi _c(f)`$, in Theorem 1.5 can be shown, by the orbifold version of Atiyah-Bott, to be equal to the character of the representation of $`G/T`$ on $`Q(M_c)`$. Our residue formula for it, (formula (1.20)) appears to be a new result even in the manifold case; however, the formula (1.21), which is a special case of this formula, is just the “quantization commutes with reduction” theorem for circle actions. A good reference for the long and entangled history of “$`[Q,R]=0`$” is the survey article \[Sj\]. For circle actions there are several relatively simple proofs, among them that of Duistermaat-Guillemin-Meinrenken-Wu (\[DGMW\]), Ginzburg-Guillemin-Karshon (\[GGK\]) and Metzler (\[Me\]). Of these, Metzler’s proof is probably the closest in spirit to our combinatorial proof of Theorem 1.8 in Section 6.
warning/0007/cond-mat0007256.html	ar5iv	text	# Soap froths and crystal structures \[ ## Abstract We propose a physical mechanism to explain the crystal symmetries found in macromolecular and supramolecular micellar materials. We argue that the packing entropy of the hard micellar cores is frustrated by the entropic interaction of their brush-like coronas. The latter interaction is treated as a surface effect between neighboring Voronoi cells. The observed crystal structures correspond to the Kelvin and Weaire-Phelan minimal foams. We show that these structures are stable for reasonable areal entropy densities. PACS numbers: 83.70.Hq 82.70.-y 61.50.Ah \] Dendritic polymers , hyper-branched star polymers and diblock copolymers represent a new class of molecular assemblies all of which form a variety of crystalline lattices, many of which are not close-packed. These assemblies are all characterized by compact cores and brush-like, soft coronas. These systems might be modeled by treating the micelles as sterically interacting hard spheres and it would follow that their crystalline phases should be stackings of hexagonal-close-packed (HCP) layers. Recently it has been shown that the face-centered cubic (FCC) lattice maximizes the total entropy and so hard-sphere crystals should form FCC structures. Note that the entropic difference between the various HCP lattices is a global issue: the local arrangement of spheres is the same for all close-packed variants and thus the lattice cannot be predicted from nearest-neighbor interactions. In order to understand the richness of crystal symmetries in the micellar systems, we propose an additional global consideration: we add an interaction proportional to the interfacial surface area between the cages which contain each micelle (Voronoi cells). Though approaches based on self-consistent field theory and two-body interactions can yield non-close-packed lattices , we propose a universal explanation for a host of new structures and present a new paradigm for the rational design and control of macromolecular assemblies . The interfacial interaction arises through the entropy of the brush-like coronas of the micelles. Because of constraints on their conformations, the brushes suffer an entropic penalty proportional to the interfacial area between the Voronoi cells surrounding each sphere. Thus they favor area-minimizing structures, precisely the type of structures that dry foams might make. Over a century ago, Lord Kelvin proposed that a body-centered-cubic (BCC) foam structure had the smallest surface-to-volume ratio but in 1994 Weaire and Phelan found that a structure based on the A15 lattice was more efficient. We note that neither the BCC nor A15 structures are close-packed and thus there is a fundamental frustration between the hard-core volume interaction and the surface interaction due to overlapping soft coronas. For concreteness, in this paper we focus on structures observed in a family of dendrimer compounds consisting of a compact poly(benzyl ether) core segment and a diffuse dodecyl corona . These conical dendrimers self-assemble in spherical micelles which subsequently arrange into the A15 lattice (Fig. 1). The interaction between the micelles is primarily steric, i.e., repulsive and short-range. The micellar architecture suggests that the potential is characterized by three regimes. At large distances, the micelles do not overlap and the interaction vanishes. As the coronas begin to overlap, the entropy of the brush-like coronas decreases, which gives rise to a soft repulsion between the micelles. Finally, at small separations the coronas begin to penetrate the compact cores: this is very unfavorable and gives rise to hard-core repulsion. This energy landscape is in qualitative agreement with recent, detailed molecular dynamics simulations . Although both originate in steric interaction, the two repulsive regimes are characterized by very different functional behaviors. The hard part of the potential results in a restricted positional entropy of the micelles which depends on the free volume, the difference between the actual and the hard-core volumes. The soft part comes from the decreased orientational entropy of the chains within the overlapping coronas. The matrix of overlapping coronas can be thought of as a compressed bilayer and thus the free volume may be written as a product of the interfacial area $`A`$ and the average spacing between the hard cores $`d`$ so that at any given density $$Ad=\mathrm{constant}.$$ (1) Though this approximation ignores the curvature of brush-like coronas, the dendrimers are relatively close and we expect this constraint to hold in this system. Since the repulsion decreases monotonically with distance, the system will favor a maximum thickness $`d`$ and will thus tend to minimize the interfacial area. Hence our proposed interfacial interaction, which is incompatible with the bulk free energy minimized by a close-packed arrangement of micelles. In the following, we compare the free energies of FCC, BCC, and A15 lattices and estimate the strength of the interfacial interaction such that the structure of the micellar crystal is dictated by the minimal-area principle. The calculation of the bulk free energies of condensed systems is fairly complicated even for hard-sphere systems and the best theoretical results are obtained numerically. It is interesting to note that elaborate analytic models, such as the high-density analog of the virial expansion and the weighted-density-functional approximation , are only slightly better than the simple cellular free-volume theory . The free-volume theory is a high-density approximation where each micelle is contained in a cell formed by its neighbors, and the communal entropy associated with the correlated motion of micelles is neglected. Within this theory, the positional entropy of a micelle is determined by the configurational space of its Voronoi or Wigner-Seitz cell. In the FCC lattice, the centers of mass of the micelles are within rhombic dodecahedra , while in the BCC lattice they are contained in regular octahedra although the BCC Voronoi cell is an orthic tetrakaidecahedron . For these lattices, the bulk free energy of a micelle is given by $$F_{\mathrm{bulk}}^X=k_BT\mathrm{ln}\mathbf{\left(}\alpha ^X\left(\frac{\beta ^X}{n^{1/3}}1\right)^3\mathbf{\right)},$$ (2) where $`X`$ is either FCC or BCC, $`n=\rho R^3`$ is the reduced number density, and $`R`$ is the hard-core radius of micelles. The coefficients $`\alpha ^{\mathrm{FCC}}=2^{5/2}`$ and $`\alpha ^{\mathrm{BCC}}=2^23^{1/2}`$ depend on the shape of the cells, whereas $`\beta ^{\mathrm{FCC}}=2^{5/6}`$ and $`\beta ^{\mathrm{BCC}}=2^{5/3}3^{1/2}`$ are determined by their size. The A15 lattice is somewhat more complicated: as shown in Fig. 1d the A15 unit cell includes 6 columnar sites, which make up 3 perpendicular interlocking columns, and 2 interstitial sites. A pseudo-Voronoi construction (subject to the constraint that all cells have equal volume) for this lattice leads to a partition consisting of irregular pentagonal dodecahedra and tetrakaidecahedra with two hexagonal and twelve pentagonal faces . Because of the irregularity of the cells, we calculate the bulk entropic free energy numerically, and the result is shown in Fig. 2. For our purposes we require an analytic form: by substituting the dodecahedra and tetrakaiecahedra by spheres and cylinders, respectively, and allowing for two adjustable parameters $`C`$ and $`S`$, which measure the deviation of the Voronoi cells from these spheres and cylinders, we find $`F_{\mathrm{bulk}}^{\mathrm{A15}}`$ $`=k_BT[{\displaystyle \frac{1}{4}}\mathrm{ln}\mathbf{\left(}{\displaystyle \frac{4\pi S}{3}}({\displaystyle \frac{\sqrt{5}}{4n^{1/3}}}1)^3\mathbf{\right)}`$ (4) $`+{\displaystyle \frac{3}{4}}\mathrm{ln}\mathbf{\left(}2\pi C({\displaystyle \frac{\sqrt{5}}{4n^{1/3}}}1)^2({\displaystyle \frac{1}{2n^{1/3}}}1)\mathbf{\right)}].`$ This form is within $`0.1\%`$ of the numerical result with $`S1.64`$ and $`C1.38`$. The interfacial free energy is minimized by the division of space with smallest area. The problem of finding the partition of space into equal-volume cells with the minimum interfacial area was first studied by Kelvin : he proposed a BCC lattice of orthic tetrakaidecahedra with slightly curved hexagonal faces to satisfy the Plateau rules . However, the Weaire-Phelan partition, which differs from the equal-volume Voronoi construction for the A15 lattice only in a delicate curvature of the pentagonal faces, is 0.3% more efficient . We note that the BCC and A15 structures are among the simplest tetragonal close-packed lattices , suggesting that other, more complex close-packed lattices might be more efficient still. However, the A15 structure appears to be the most efficient, although no proof of its supremacy exists. To argue that the A15 and BCC lattices are the equilibrium structures formed by micelles, we must estimate the entropy penalty per unit area and translate this into an entropy per dodecyl chain. The dodecyl bilayer is modeled as a polymer brush consisting of chain molecules attached to hard cores, and in the limit of high interdigitation its free energy consists solely of the excluded-volume repulsion of the chains: $$F_{\mathrm{surf}}=\frac{2\mathrm{}Nk_BT}{d},$$ (5) where $`d`$ is the layer thickness, $`\mathrm{}`$ is a parameter with the dimension of length and $`N`$ is the number of chains per micelle . Since the bilayer must fill the free volume, $`A_Md=2(n^14\pi /3)R^3`$, where $`A_M`$ is the interfacial area per micelle. Thus the interfacial free energy of a micelle is $$F_{\mathrm{surf}}^X=\frac{\mathrm{}Nk_BT}{R}\frac{\gamma ^Xn^{2/3}}{n^14\pi /3},$$ (6) where $`\gamma ^{\mathrm{FCC}}=2^{5/6}3=5.345`$, $`\gamma ^{\mathrm{BCC}}=5.306`$, and $`\gamma ^{\mathrm{A15}}=5.288`$ . We now calculate the range of $`\mathrm{}`$ such that the total free energy $$F^X=F_{\mathrm{bulk}}^X+F_{\mathrm{surf}}^X$$ (7) is minimized by the BCC and A15 lattices rather than by the naïve, close-packed, FCC lattice. In order to estimate the strength of the soft repulsion, we must first determine the actual reduced density $`n`$. Since the hard-core radius of the micelles is unknown, we limit $`n`$ by recognizing that it must be larger than the melting density, $`n0.120`$ for hard spheres , and that it must be smaller than the close-packing density of the A15 lattice, $`n=0.125`$. The most conservative lower bound of $`\mathrm{}`$ corresponds to the lowest possible density, i.e., the melting density. With $`N=162`$ chains per micelle , we find that at $`n=0.120`$ the FCC to BCC transition occurs for $`\mathrm{}0.1R`$ and the BCC to A15 transition occurs for $`\mathrm{}0.3R`$. This corresponds to an entropy per chain of $`0.5k_B`$ and $`1.5k_B`$, respectively. Both values are of the correct order of magnitude and the higher value of the latter is consistent with the relative rarity of the A15 phase. Since we expect that each chain has at least $`k_B`$ of entropy, we conclude that the energetics of the dendrimer micelles is dominated by interfacial effects. This is hardly surprising. The number of degrees of freedom of each micelle is quite large and the bulk free energy only depends on the position of the micelle as a whole. Since the micelles are soft, the internal degrees of freedom such as the chain conformations play an important role. This paradigm – which shows that the minimal surface problem can be fruitfully transplanted to the microscopic level – explains the morphology of a number of dense micellar systems. The same ideas can be applied to polymeric micelles made of, e.g., polystyrene-polyisoprene diblock copolymers dispersed in decane . In this case, the micelles are characterized by highly concentrated polystyrene core and diffuse polyisoprene corona, and they form BCC or FCC lattices, depending on the relative length of the polystyrene and polyisoprene chains. The BCC lattice is observed in diblock copolymers with similar lengths of core and coronal segments, whereas the FCC lattice occurs whenever the corona is thin compared to core. This is consistent with our model of the impenetrable core which is responsible for the hard part of the repulsion and which favors arrangements with large free volume. In addition, our model suggests that the A15 lattice is the ground state of an asymmetric diblock with an exceptionally large corona or, equivalently, a corona made of very floppy, “entropy-rich” chains. We note that distinguishing between A15 and BCC in powder-averaged diffraction is delicate: the first three BCC reflections are at $`\sqrt{2}`$, $`\sqrt{4}`$, and $`\sqrt{6}`$, while the first four A15 reflections are at $`\sqrt{2}`$, $`\sqrt{4}`$, $`\sqrt{5}`$, and $`\sqrt{6}`$, and thus a careful study would be necessary. The existence of the A15 lattice in the dendrimer aggregate also may be regarded as an experimental verification of the recent theoretical developments in minimal surfaces and, in particular, Weaire and Phelan’s conjecture that this structure solves the Kelvin problem. At this juncture, the presumably ideal A15 structure has not been observed unambiguously on a macroscopic scale in a soap froth . Last but not least, let us note that similar structures have been found in lyotropic materials, e.g., in lipid bilayers in water . In such systems the intermicellar potential also results in an effective interfacial free energy although it is not steric but substance-specific, and thus transcends the scope of this discussion. Our model may be further refined by including the effects of the curvature of the brush-like coats, the strain of the coronas into the interstitial regions, and solvent effects. In addition, the dual problem of determining the structure of foams might be amenable to our analysis through the introduction of excluded volume interactions between the bubbles . Recent work focusing on two-body interactions has shown that We hope that this study elucidates the relation between interaction and structure in supramolecular assemblies. By including an additional global contribution to the free energy we provide a rough yet universal quantitative guideline for the design of self-organized soft materials, which can be used for a number of applications such as photonic band-gap materials , Bragg switches , and porous microreactors . By tuning the ranges of hard and soft repulsion, one should be able to choose among the spectrum of symmetries from the lattice with minimal interfacial area to the lattice with maximal packing fraction and engineer the crystal structure most fitted for a particular application. It is a pleasure to acknowledge stimulating conversations with T. C. Lubensky and V. Percec. This work was supported in part by NSF Career Grant DMR97-32963. RDK was also supported by the Alfred P. Sloan Foundation.
warning/0007/gr-qc0007040.html	ar5iv	text	# Distinguishing marks of simply-connected universes ## 1 Introduction Current observational data favor the locally homogeneous and isotropic Friedmann-Lemaître-Robertson-Walker (FLRW) cosmological models as approximate descriptions of our universe at least since the recombination time. Thus in the framework of the general relativity theory it can be described through a Robertson-Walker (RW) metric $$ds^2=dt^2R^2(t)d\sigma ^2,$$ (1.1) where $`t`$ is a cosmic time, and $`d\sigma ^2=d\chi ^2+f^2(\chi )[d\theta ^2+\mathrm{sin}^2\theta d\varphi ^2]`$ with $`f(\chi )=\chi ,\mathrm{sin}\chi ,\mathrm{sinh}\chi ,`$ depending on the sign of the constant spatial curvature ($`k=0,\pm 1`$). These descriptions, however, are only local and do not fix the global shape (topology) of our universe. Despite the infinitely many possibilities for its global topology, it is often assumed that spacetime is simply-connected leaving aside the hypothesis that the universe may be multiply-connected, and compact (finite) even in the cases $`k=0`$ and $`k=1`$. In other words, it is often assumed that the $`t=const`$ spatial sections $`M`$ of a RW spacetime manifold are one of the following simply-connected spaces: Euclidean $`E^3`$ ($`k=0`$), elliptic $`S^3`$ ($`k=1`$), or the hyperbolic $`H^3`$ ($`k=1`$). However, the connectedness (either simply or multiply) for our three-space has not been settled by cosmological observations. Thus, the space $`M`$ where we live may also be any one of the possible multiply-connected quotient three-spaces $`M=\stackrel{~}{M}/\mathrm{\Gamma }`$, where $`\stackrel{~}{M}`$ stands for $`E^3`$, $`S^3`$ or $`H^3`$, and $`\mathrm{\Gamma }`$ is a discrete group of isometries of the covering space $`\stackrel{~}{M}`$ acting freely on $`\stackrel{~}{M}`$ . Whether we live in a simply or multiply-connected, finite (compact) or infinite (non-compact) space, and what is the size and the shape of the universe are open problems modern cosmology seeks to solve . The most immediate consequence of multiply-connectedness of the universe is that the sky may show multiple images of cosmic objects periodically distributed in the space. This periodic distribution of images arises from the correlations in their positions dictated by the discrete isometries of the covering group $`\mathrm{\Gamma }`$ of the three-manifold used to model its space section. One way to tackle the problems regarding the topology of the universe is through a suitable statistical analysis applied to catalogs of discrete cosmic sources to find out whether or not there are multiple correlated images of cosmic objects, and eventually determine the topological features of the universe from the pattern of images the sky shows. The correlations among the images of cosmic objects in multiply-connected universes can be couched in terms of distance correlations between the images. Indeed, one way of looking for distance correlations between cosmic images in multiply-connected universes is by using pair separations histograms (PSH), which are functions $`\mathrm{\Phi }(s_i)`$ that count the number of pair of images separated by a distance that lies in intervals (bins) $`J_i`$. The embryonic expectation was that the distance correlations would manifest as topological spikes in PSH’s, and that the spike spectrum would be a definite signature of the non-trivial topology . However, this initial expectation turned out to be false – . Nevertheless, the most striking evidence of multiply-connectedness in PSH’s is indeed the presence of such topological spikes, which arise from translational isometries $`g_t\mathrm{\Gamma }`$. The non-translational isometries $`g_{nt}\mathrm{\Gamma }`$, however, manifest as rather tiny deformations of the expected pair separation histogram $`\mathrm{\Phi }_{exp}^{sc}(s_i)`$ corresponding to the underlying simply-connected universe. However, from computer simulations it becomes clear that the expected pair separation histogram (EPSH) corresponding to a multiply-connected universe $`\mathrm{\Phi }_{exp}(s_i)`$, which is nothing but an PSH from which the statistical noise has been withdrawn, is not a suitable quantity for revealing the topology of multiply-connected universes . In a recent article, Gomero et al. (see also ) have proposed a way of extracting the topological signature of any multiply-connected universe of constant curvature by using a new quantity $`\phi ^{mc}(s_i)(n1)[\mathrm{\Phi }_{exp}(s_i)\mathrm{\Phi }_{exp}^{sc}(s_i)]`$, where $`n`$ is the number of images. Note, however, that this quantity cannot be used as distinguishing marks of simply-connected universes since it vanishes identically for such universes. This amounts to saying that the graphs of $`\phi ^{mc}(s_i)`$ for all three classes of RW simply-connected universes which arise from (simulated or real) catalogs exhibit nothing but statistical noise, and thus $`\phi ^{mc}(s_i)`$ should not be used as identifying markings in the simply-connected cases. As a matter of fact, the scheme discussed in as well as the approaches that make use of the cosmic microwave background radiation – were fundamentally devised to reveal the possible non-trivial topology of small universes. However, neither the multiply nor the simply-connectedness for our universe has been discarded or confirmed by the current astrophysical observations. In computer-aided simulations the histograms such as the PSH’s $`\mathrm{\Phi }(s_i)`$ contain statistical fluctuations, which can give rise to sharp peaks of statistical (non-topological) origin, or can hide (or mask) the tiny deformations due to non-translational isometries. The most immediate approach to cope with fluctuation problems in PSH’s is by using the mean pair separation histogram (MPSH) scheme to obtain the mean PSH $`<\mathrm{\Phi }(s_i)>`$ rather than a single PSH $`\mathrm{\Phi }(s_i)`$. In ref. they have used the MPSH technique to extract the topological signature of RW multiply-connected universes. This technique consists in the use of $`K`$ (say) computer-generated comparable catalogs to obtain the mean pair separation histograms $`<\mathrm{\Phi }(s_i)>`$ and $`<\mathrm{\Phi }^{sc}(s_i)>`$; and use them as approximations for $`\mathrm{\Phi }_{exp}(s_i)`$ and $`\mathrm{\Phi }_{exp}^{sc}(s_i)`$, to construct the topological signature $`\phi ^{mc}(s_i)(n1)[<\mathrm{\Phi }(s_i)><\mathrm{\Phi }^{sc}(s_i)>]`$. Obviously the greater is the number $`K`$ of catalogs the better are the approximations $`<\mathrm{\Phi }(s_i)>\mathrm{\Phi }_{exp}(s_i)`$ and $`<\mathrm{\Phi }^{sc}(s_i)>\mathrm{\Phi }_{exp}^{sc}(s_i)`$. In this article we point out that the statistical quantity $`\varphi ^{sc}(s_i)=\mathrm{\Phi }_{exp}^{sc}(s_i)`$ is indeed a suitable distinguishing mark of the simply-connected RW universes, and rederive its explicit expressions for Euclidean, hyperbolic and elliptic simply-connected universes (spherical balls $`_a`$ with radius $`a`$) fulfilled with an uniform distribution of cosmic objects. In doing so, on the one hand we obtain the exact (free from statistical fluctuation) expressions for the distinguishing mark $`\varphi ^{sc}(s_i)`$ of the three possible classes of simply-connected RW universes; on the other hand one attains a refined statistical meaning of the signature $`\phi ^{mc}(s_i)`$ and also obtains an improvement on the procedure to extract the topological signature of multiply-connected RW universes devised in ref. . In the next section we set our framework, define the basic notation, and derive the expressions for the distinguishing marks for the three possible classes of simply-connected RW universes (Euclidean, hyperbolic, elliptic). There we also present and analyze graphs of the distinguishing mark $`\varphi ^{sc}(s_i)`$ of simply-connected RW universes, and discuss the improvement we have obtained in the procedure to extract the topological signature of multiply-connected universes studied in ref. . In the last section we summarize and discuss our main results and present the concluding remarks. To close this section a word of clarification: although throughout this article we loosely use the terminology topological signature of a universe and/or of a manifold, it should be noted that the topological signature actually corresponds to an observed universe, which in this paper is a spherical ball $`_a\stackrel{~}{M}`$ of radius $`a`$, which contains the set of the observed images. ## 2 Distinguishing marks In this section we will first set the notation and then recast in a unified and compact way the explicit expressions for the probability densities obtained in so as to show that they can be used as distinguishing mark for the three possible classes ($`k=0,\pm 1`$) of simply-connected RW universes fulfilled with an uniform distribution of cosmic objects. Let us start by recalling that a catalog $`𝒞`$ is a set of observed images, subset of the set $`𝒪`$ of observable images ($`𝒞𝒪`$), which are clearly contained in the observable universe, which in turn is the part of the universal covering manifold $`\stackrel{~}{M}`$ causally connected to an image of a given observer. The observed universe is the part of the observable universe which contains all the sources registered in the catalog. Our observational limitations are formulated through selection rules which dictate how the subset $`𝒞`$ arises from $`𝒪`$. Catalogs whose images obey the same (well-behaved) distribution law and that follow the same selection rules are said to be comparable catalogs . It should be noted that in the process of construction of catalogs it is assumed a RW geometry (needed to convert redshift into distance) and that a particular type of sources (clusters of galaxies, quasars, etc) is chosen from the outset. So, for our purpose in the present work, in addition to the angular positions on the celestial sphere, the relevant information registered in a given catalog is the redshift corresponding to each image in the catalog. Consider a catalog $`𝒞`$ with $`n`$ cosmic images and denote by $`\eta (s)`$ the number of pairs of images whose separation is $`s`$. Consider also that our observed universe is a ball of radius $`a`$ and divide the interval $`(0,2a]`$ in $`m`$ equal subintervals $`J_i`$ of length $`\delta s=2a/m`$. Each of such subintervals has the form $$J_i=(s_i\frac{\delta s}{2},s_i+\frac{\delta s}{2}];i=1,2,\mathrm{},m,$$ (2.1) and is centered at $$s_i=(i\frac{1}{2})\delta s.$$ The PSH is a normalized function which counts the number of pair of images separated by a distance that lies in the subinterval $`J_i`$. Thus the function PSH is given by $$\mathrm{\Phi }(s_i)=\frac{2}{n(n1)}\frac{1}{\delta s}\underset{sJ_i}{}\eta (s),$$ (2.2) and is clearly subjected to the normalizing condition $$\underset{i=1}{\overset{m}{}}\mathrm{\Phi }(s_i)\delta s=1.$$ (2.3) If one considers an ensemble of comparable catalogs<sup>1</sup><sup>1</sup>1Note that a typical catalog of the ensemble reflects (corresponds to) a distribution of images in the observed universe $`_a`$. with the same number $`n`$ of images, and corresponding to the same three-manifold $`M`$ of constant curvature, one can compute probabilities and expected values of quantities which depend on the images in the catalogs of the ensemble. In particular, we can compute the expected number $`\eta _{exp}(s_i)`$ of pairs of cosmic images in a catalog $`𝒞`$ of the ensemble with separations in $`J_i`$. This quantity is quite relevant because from it one has the normalized expected pair separation histogram (EPSH) which clearly is given by $$\mathrm{\Phi }_{exp}(s_i)=\frac{1}{N}\frac{1}{\delta s}\eta _{exp}(s_i)=\frac{1}{\delta s}F(s_i),$$ (2.4) where obviously $`N=n(n1)/2`$ is the total number of pairs of cosmic images in $`𝒞`$, and $`F(s_i)=\eta _{exp}(s_i)/N`$ is the probability that a pair of images be separated by a distance that lies in the interval $`J_i`$. In what follows we shall consider that we have an ensemble of comparable catalogs whose underlying observed universe (spherical ball $`_a`$ with radius $`a`$) are simply-connected and fulfilled with an uniform distribution of pointlike objects. We will take $`\varphi ^{sc}(s_i)\mathrm{\Phi }_{exp}^{sc}(s_i)`$ as distinguishing mark for these three classes of simply-connected universes. Clearly for this uniform distribution of objects all separations $`\mathrm{\hspace{0.17em}0}<s_i2a`$ are allowed, so the identifying markings $`\varphi ^{sc}(s_i)`$ are continuous functions of $`s`$ given by $$\varphi ^{sc}(s)=\mathrm{\Phi }_{exp}^{sc}(s)=\frac{1}{\delta s}F_{sc}(s),$$ (2.5) where $`F_{sc}(s)`$ is the probability that a pair of images in a catalog $`𝒞`$, corresponding to a simply-connected universe, be separated by a distance $`s`$. For the sake of simplicity hereafter we will drop the subscript of $`F_{sc}(s)`$. To make explicit that the distinguishing mark depends upon the radius of the observed universe we rewrite (2.5) in the form $$\varphi ^{sc}(a,s)=\mathrm{\Phi }_{exp}^{sc}(a,s)=\frac{1}{\delta s}F(a,s)=(a,s),$$ (2.6) where $`(a,s)`$ clearly is the probability density, i.e. it is such that $`F(a,s)=(a,s)ds`$ gives the probability that two arbitrary points in the ball $`_a`$ be separated by a distance between $`s`$ and $`s+ds`$. Equation (2.6) makes apparent that the $`\varphi ^{sc}(a,s)`$ gives essentially the distribution of probability for all $`s`$ in the ball $`_a`$. Moreover, since the way one measures the distances varies for each constant curvature universe, it is clearly expected that the expression for distinguishing mark $`\varphi ^{sc}(a,s)`$ changes with the three-geometry of these simply-connected universes. In what follows we shall recast in a compact way the expressions of $`\varphi ^{sc}(a,s)`$ for Euclidean, hyperbolic and elliptic simply-connected universes . Consider in either of the simply-connected three-spaces a ball $`_a`$ centered at the origin $`O`$, and let $`P`$ and $`Q`$ be two arbitrary points in the ball. Denote by $`r[0,a]`$ the radial position of $`P`$, and by $`s2a`$ the distance from $`P`$ to $`Q`$ (see figure 1).<sup>2</sup><sup>2</sup>2In order to encompass the elliptic class in our compact approach we shall initially treat only the elliptic cases in which the radius $`a`$ of the universe $`_aS^3`$ is such that $`2a<\pi R`$, where $`R`$ is the radius of the curvature of the geometry, i.e. the scale factor of RW metric (1.1) for a given time $`t=t_0`$. Further, for the sake of simplicity and without loss of generality we shall also set $`R=1`$ for both the hyperbolic and elliptic cases. Consider now the quantity $`(a,r,s)drds`$, which is the probability that $`P`$ lies in a position between $`r`$ and $`r+dr`$, times the probability that the separation between $`P`$ and $`Q`$ lies between $`s`$ and $`s+ds`$. Clearly for the simply-connected cases we are concerned the probability density $`(a,r,s)`$ is proportional to the following two areas: (i) $`𝒜_S(r)`$ which is the area of the locus of the points $`P`$ located at a distance $`r`$ from the origin $`0`$; and (ii) the area of the locus of the points $`Q`$ that are separated from $`P`$ by $`s`$. Note, however, that when $`r+s<a`$ the latter locus is a two-sphere $`S^2`$ with area $`𝒜_S(s)`$, whereas when $`r+s>a`$ it changes into a spherical calotte (cap) with area $`𝒜_C(a,r,s)`$. To sum up we have that the expression for the probability density $`(a,r,s)`$ for the configuration in which $`P`$ is between $`r`$ and $`r+dr`$, and is separated from $`Q`$ by a distance between $`s`$ and $`s+ds`$, can be written in the general form $$(a,r,s)=\zeta 𝒜_S(r)\left[𝒜_S(s)\mathrm{\Theta }(asr)+𝒜_C(a,r,s)\mathrm{\Theta }(s+ra)\right],$$ (2.7) where $`\zeta `$ is a normalization constant and $`\mathrm{\Theta }`$ is the Heaviside function. Obviously, due to the cut-off effects of the function $`\mathrm{\Theta }`$ the first term of the sum in the right-hand side of (2.7) is nonzero only for $`r+s<a`$, whereas the second term is non-null only when $`r+s>a`$. Now since the area $`𝒜_S(r)`$ of the two-sphere as well as the area of the spherical calotte $`𝒜_C(a,r,s)`$ depend on what is the simply-connected three-space where the ball $`_a`$ is considered, then to obtain $`(a,r,s)`$ for each class of simply-connected universes we are interested one ought to: (i) calculate the areas $`𝒜_S(r)`$ and $`𝒜_C(a,r,s)`$; (ii) insert in (2.7) and integrate $`(a,r,s)`$ from $`r=0`$ to $`r=a`$; and (iii) impose the normalization condition $$_0^{2a}(a,s)𝑑s=_0^{2a}\varphi ^{sc}(a,s)𝑑s=1,$$ (2.8) to obtain the value of the normalization constant $`\zeta `$. In what follows we shall use this systematic scheme to determine $`(a,r,s)`$ for the three classes of simply-connected universes we are concerned. For Euclidean universe ($`_aE^3`$) one obviously has $`𝒜_S(r)=4\pi r^2`$, and straightforward calculations furnish $`𝒜_C(a,r,s)=(\pi s/r)[a^2(sr)^2]`$. According to the above-outlined scheme inserting these expressions in (2.7), integrating $`(a,r,s)`$ from $`r=0`$ to $`r=a`$, and using the condition (2.8) together with (2.6), one obtains that the expression for the distinguishing mark $`\varphi _E^{sc}(a,s)`$ of a Euclidean universe is given by $$\varphi _E^{sc}(a,s)=\frac{3}{16a^6}s^2(2as)^2(s+4a),$$ (2.9) which holds for $`s(0,2a]`$, and where the value of the normalization constant was found to be $`\zeta =9/(16\pi ^2a^6)`$. Before proceeding to the next class it should be noticed that the shape of the signature $`\varphi _E^{sc}(a,s)`$ does not depend on the value of the radius $`a`$. Indeed, in terms of a new variable $`s^{}=s/a`$ the expression (2.9) can be rewritten in the form $$\varphi _E^{sc}(a,s^{})=\frac{3}{16a}s^2(2s^{})^2(s^{}+4),$$ (2.10) which makes clear that for distinct radii $`a`$ one has different constant multiplying factors $`3/(16a)`$, but without changing the functional dependence of $`\varphi _E^{sc}(a,s)`$ with $`s`$. So, the shape of the graph of the distinguishing mark function for Euclidean universes does not depend on the value of the radius $`a`$. For hyperbolic universes ($`_aH^3`$) one obtains $`𝒜_S(r)=4\pi \mathrm{sinh}^2r`$ and $`𝒜_C(a,r,s)=2\pi \mathrm{sinh}s[\mathrm{sinh}s\mathrm{cosh}s\mathrm{coth}r+\mathrm{cosh}a\text{csch}r]`$. Again, through the second and third steps of the above-mentioned systematic scheme one finds the following expression for the distinguishing mark of a simply-connected hyperbolic universe fulfilled with an uniform distribution of objects: $$\varphi _H^{sc}(a,s)=\frac{8\mathrm{sinh}^2s}{(\mathrm{sinh}2a2a)^2}\left[\mathrm{cosh}a\text{sech}(s/2)\mathrm{sinh}(as/2)(as/2)\right],$$ (2.11) which holds for $`s(0,2a]`$, and where the value of the normalization constant in this case was found to be $`\zeta =[\pi (\mathrm{sinh}2a2a)]^2`$. As for the elliptic universes ($`_aS^3`$) when the diameter $`2a`$ is less than the separation $`\pi R`$ between antipodal points of $`S^3`$ the above scheme can be similarly used. For this case one obtains $`𝒜_S(r)=4\pi \mathrm{sin}^2r`$ and $`𝒜_C(a,r,s)=2\pi \mathrm{sin}s[\mathrm{sin}s\mathrm{cos}s\mathrm{cot}r+\mathrm{cos}a\mathrm{csc}r]`$. Using these expressions in (2.7) and following the above-outlined general procedure one finds $$\varphi _S^{sc}(a,s)=\frac{8\mathrm{sin}^2s}{(2a\mathrm{sin}2a)^2}\left[(as/2)\mathrm{cos}a\mathrm{sec}(s/2)\mathrm{sin}(as/2)\right],$$ (2.12) which hold for $`2a<\pi `$, where we have taken $`R=1`$. The value of the normalization constant in this case was found to be $`\zeta =[\pi (2a\mathrm{sin}2a)]^2`$. The elliptic universes ($`_aS^3`$) for which $`2a>\pi R`$ cannot be included in the above general scheme. They are trickier to be handled due to the connectivity of the spherical space $`S^3`$ and the additional requirement that $`s`$ must not exceed $`\pi R`$, which is needed to ensure that one is taking the shortest geodesic part between two points of $`S^3`$. For the sake of brevity and completeness we shall present here only the final expression for $`\varphi ^{sc}(a,s)`$. It turns out that the general expression of the distinguishing mark which holds for all elliptic universes with $`a(0,\pi ]`$ and fulfilled with a uniform distribution of objects is given by $`\varphi _S^{sc}(a,s)`$ $`=`$ $`{\displaystyle \frac{8\mathrm{sin}^2s}{(2a\mathrm{sin}2a)^2}}\{2a\mathrm{sin}2a\pi +\mathrm{\Theta }(2\pi 2as)[\mathrm{sin}2a+\pi `$ (2.13) $`as/2\mathrm{cos}a\mathrm{sec}(s/2)\mathrm{sin}(as/2)]\},`$ where $`s(0,\text{min}(2a,\pi )]`$. In what follows we shall present and analyze a few graphs of the distinguishing mark for each class or RW simply-connected universes. Figure 2 shows the distinguishing mark $`\varphi _E^{sc}(a,s)`$ for a Euclidean universe $`_a`$ with radius $`a=0.5`$. This marking also gives the probability distribution of the pair separation distance $`s`$ for $`s(0,2a]`$. A close inspection of this figure reveals that the most likely separation between two arbitrary pointlike objects in the Euclidean observed universe $`_a`$ is slightly greater than the radius $`a`$ of the ball. Figure 3 shows the distinguishing mark $`\varphi _H^{sc}(a,s)`$ for three values of the radius $`a`$. For $`a1`$ this function behaves approximately as that we have derived for the Euclidean universe (figure 2), as one would expect from the beginning. For increasing values of the radius $`a`$ of the observed universe the maximum of the mark $`\varphi _H^{sc}(a,s)`$ moves towards the greater values of $`s`$. In other words, the most likely value of $`s`$ increases (the maximum of $`\varphi _H^{sc}(a,s)`$ takes place later) for increasing values of the radius $`a`$. In figure 4 four graphs of the distinguishing mark $`\varphi _S^{sc}(a,s)`$ for different values of the radius $`a`$ of the universe are shown. For increasing values of $`a`$ from $`0`$ to $`\pi `$ the maximum of the signature $`\varphi _S^{sc}(a,s)`$ moves continuously towards the smaller values of $`s`$. Contrarily to the hyperbolic case, here for increasing values of the radius $`a`$ the maximum of $`\varphi _S^{sc}(a,s)`$ moves toward the origin (smaller values of $`s`$). This is also illustrated in figure 4, where the most likely value of the separation $`s`$ decreases (the maximum of $`\varphi _S^{sc}(a,s)`$ takes place earlier) for increasing values of $`a`$. It also should be mentioned that for an arbitrary radius of curvature of the geometry $`R`$ the expressions for $`\varphi _H^{sc}(a,s)`$ and $`\varphi _S^{sc}(a,s)`$ can be obtained by multiplying the right-hand side of (2.11) and (2.13) by $`1/R`$, and simultaneously by changing $`aa/R`$ and $`ss/R`$. In the remainder of this section we will discuss the improvement of the method to extract the topological signature $$\phi ^{mc}(s_i)=(n1)[\mathrm{\Phi }_{exp}(s_i)\mathrm{\Phi }_{exp}^{sc}(s_i)],$$ (2.14) of multiply-connected universes studied in (see also ). To this end, the relevant point to be noted is that an EPSH $`\mathrm{\Phi }_{exp}(s_i)`$ is essentially a typical PSH from which the statistical noise has been withdrawn. Hence we have $`\mathrm{\Phi }_{exp}(s_i)`$ $`=`$ $`\mathrm{\Phi }(s_i)\rho ^{mc}(s_i),`$ (2.15) $`\mathrm{\Phi }_{exp}^{sc}(s_i)`$ $`=`$ $`\mathrm{\Phi }^{sc}(s_i)\rho ^{sc}(s_i),`$ (2.16) where $`\rho ^{mc}(s_i)`$ and $`\rho ^{sc}(s_i)`$ represent the statistical noises that arise in the corresponding PSH’s. Using now the decompositions (2.15) and (2.16) together with (2.14) one readily obtains $$\phi ^{mc}(s_i)=(n1)[\mathrm{\Phi }(s_i)\mathrm{\Phi }^{sc}(s_i)+\rho ^{sc}(s_i)\rho ^{mc}(s_i)],$$ (2.17) which clearly gives the topological signature intermixed with two statistical fluctuations. Now, from equations (2.14) – (2.17) it is clear that one can approach the topological signature of multiply-connected universes $`\phi ^{mc}(s_i)`$ by reducing the statistical fluctuations, i.e. by making $`\rho ^{mc}(s_i)0`$ as well as $`\rho ^{sc}(s_i)0`$ through any suitable statistical method to lower the noises. An improvement of the method devised in to extract the topological signature $`\phi ^{mc}(s_i)`$ of multiply-connected universes with a uniform distribution of matter comes out from the very fact that having the derived expressions (2.9), (2.11) and (2.13) one has from the beginning $`\rho ^{sc}(s_i)=0`$ for those universes. Thus, for example, if the MPSH is the technique one uses to reduce the statistical fluctuations, the topological signature (2.14) in these cases reduces to the form $`\phi ^{mc}(s_i)(n1)[<\mathrm{\Phi }(s_i)>\mathrm{\Phi }_{exp}^{sc}(s_i)]`$, with the exact expression for $`\mathrm{\Phi }_{exp}^{sc}(s_i)`$ rather than the approximate mean $`<\mathrm{\Phi }^{sc}(s_i)>`$. ## 3 Concluding remarks To a certain extent it is well-known that RW geometry (1.1) does not fix the global shape (topology) of the spacetime, and that there is an infinite number of topologically distinct $`t=const`$ spatial sections $`M`$ for the RW spacetime manifold. Nevertheless, it is often (implicitly or explicitly) assumed that the $`t=const`$ spatial sections $`M`$ of a RW spacetime manifold are one of the following simply-connected spaces: $`E^3`$ ($`k=0`$), $`S^3`$ ($`k=1`$), or $`H^3`$ ($`k=1`$). However, this assumption of simply-connectedness for our three-space has not been settled by cosmological observations. As a matter of fact, neither the simply nor the multiply-connectedness for the three-space where we live has been discarded or confirmed by the available astrophysical data. The two main approaches to constrain or determine the topology of the our three-space rely on the existence of multiple (topological) images of either cosmic objects or spots of microwave background radiation, and thus they aim at non-trivial topology of small universes — a possible simply-connectedness (trivial topology) of the universe has not been suitably considered in these approaches. A special method to determine possible non-trivial topologies of RW universes, and which relies on the existence of multiple images, was recently discussed in . There it is suggested that the quantity $`\phi ^{mc}(s_i)(n1)[\mathrm{\Phi }_{exp}(s_i)\mathrm{\Phi }_{exp}^{sc}(s_i)]`$ is a suitable measure of the topological signature of the multiply-connected RW universes. However, $`\phi ^{mc}(s_i)`$ cannot be used as the topological signature for simply-connected universes, since it vanishes identically. This means that if we live in RW simply-connected universe the graphs of $`\phi ^{mc}(s_i)`$ which arise from real (or simulated) catalogs will exhibit nothing but statistical noise. Thus $`\phi ^{mc}(s_i)`$ can be used not only to extract the topological signature of multiply-connected universes but also to decide between multiply or simply-connectedness of the universe, since in this latter case it gives rise simply to statistical noise. One might think at first sight that the vanishing of $`\phi ^{mc}(s_i)`$ (which means that it gives rise to nothing but statistical fluctuations) would lead only to the simply-connectedness without separating among the three possible classes of simply-connected RW universes. In practice, though, the vanishing of $`\phi ^{mc}(s_i)`$ takes place for one underlying RW metric used to convert redshift into distance to have the pair separations, and thus it also gives the underlying manifold of the corresponding simply-connected universe, as there is a clear correspondence between geometry and the covering manifold in these cases. Nevertheless, since $`\phi ^{mc}(s_i)`$ vanishes identically (gives rise to nothing but statistical noise) for all classes of RW universes it is not an univocal (or unequivocal) distinguishing mark of these universes. Thus we have been led to take the quantity $`\varphi ^{sc}(a,s)`$ as distinguishing mark of the simply-connected RW universes. Clearly $`\varphi ^{sc}(a,s)`$: (i) does not vanish; (ii) can be used to separate the three possible classes of simply-connected RW universes; and (iii) it is the quantity which really matters (in this statistical context) for the simply-connected cases. Actually $`\varphi ^{sc}(a,s)`$ can also be used to distinguish RW universes with different radius $`a`$. Further, note that since the way one measures the distances varies for each constant curvature universe, it was really expected from the outset that the expression for the distinguishing mark $`\varphi ^{sc}(a,s)`$ would be distinct for different simply-connected universes. We have also presented the explicit expressions of the distinguishing marks $`\varphi _E^{sc}(a,s)`$, $`\varphi _H^{sc}(a,s)`$ and $`\varphi _S^{sc}(a,s)`$ for, respectively, Euclidean, hyperbolic and elliptic simply-connected RW universes fulfilled with an uniform distribution of cosmic objects. Besides, we have presented and analyzed graphs of this signature for the simply-connected RW universes, and discussed the improvement that these exact expressions bring to the method to extract the topological signature of multiply-connected universes discussed in . The distinguishing marks for the simply-connected RW universes $`\varphi _E^{sc}(a,s)`$, $`\varphi _H^{sc}(a,s)`$ and $`\varphi _S^{sc}(a,s)`$, which we have studied in is this work give, in each case, the probability distributions of the pair separation $`s(0,2a]`$. If one takes these probability distributions as ground distributions, then the topological signature of multiply-connected RW universes $`\phi ^{mc}(s_i)`$ studied in can be understood as a measure of the deviation between the pair separation probability distribution in the multiply-connected cases \[given by $`\mathrm{\Phi }_{exp}(s_i)`$\] and the corresponding ground pair separation probability distribution. The isometries $`g`$ of the covering group $`\mathrm{\Gamma }`$ modify the ground pair separation probability distribution, and the quantity $`\phi ^{mc}(s_i)`$ measures that deviation of topological origin. It is worth mentioning that in the cases of multiply-connected universes for which the smaller length of the fundamental polyhedron $`𝒫`$ of $`M`$ ($`𝒫\stackrel{~}{M}`$) is greater than the diameter $`2R_H`$ ($`R_H`$ is the particle horizon) of the observed universe $`_{R_H}\stackrel{~}{M}`$, no multiple images can be observed. These multiply-connected universes are therefore indistinguishable from the simply-connected universes with the same covering space, equal radius, and identical distribution of cosmic sources. In these multiply-connected cases in which the scale of the multiply-connectedness is greater than the radius $`R_H`$ universe $`_{R_H}`$ no sign of the multiply-connectedness will arise, and the distinguishing marks we have discussed in this work can play a relevant role, when there is a homogeneous distribution of matter, of course. To close this article it is worth mentioning that the ultimate goal in the statistical approaches to extract the topological signature (mark) is the comparison of the graphs (signature or mark) obtained either theoretically or from simulated catalogs against similar graphs obtained from real catalogs. To do so, one clearly has to have either the exact explicit expression or the simulated patterns of the topological signatures of the possible universes. The expressions we have found for the distinguishing mark of simply-connected RW universes with a uniform distribution of matter can certainly be used in such comparisons. Note, however, that even if the universe turns out to be simply-connected one still has to face the remaining problem of reducing the noise from just one or even a few real catalogs of cosmic sources. ## Acknowledgments I thank A.A.F. Teixeira and G.I. Gomero for stimulating and fertile discussions, for the careful reading of the manuscript and indication of relevant omissions. I also thank the scientific agency CNPq for financial support. ## Captions for the figures Two-dimensional schematic figure of observed universe: spherical ball $`_a`$ of radius $`a`$, which contains the set of the observed images. The circular arc with center in $`P`$ represents a spherical calotte (cap) that changes into a sphere $`S^2`$ when $`r+sa`$. The distinguishing mark $`\varphi _E^{sc}(a,s)`$ for a Euclidean simply-connected universe $`_a`$ for a radius $`a=0.5`$. The horizontal axis gives the pair separation $`s`$ while the vertical axis gives the normalized number of pairs. This curve also gives the probability distribution of the pair separation distance $`s`$ in this universe $`_a`$. A close inspection reveals that the most likely separation between two arbitrary images in a Euclidean universe $`_a`$ is slightly greater than the radius $`a`$ of the universe. Graphs of the distinguishing mark $`\varphi _H^{sc}(a,s)`$ of hyperbolic simply-connected universes $`_a`$ for three different values of the radius $`a`$. The normalized number of pairs in the vertical axis is given in unit of $`1/(2a)`$, while in the horizontal axis the unit of length is equal to $`2a`$. Note that for $`a1`$ the signature $`\varphi _H^{sc}(a,s)`$ behaves approximately as its Euclidean counterpart $`\varphi _E^{sc}(a,s)`$, whose graph is shown in figure 2. For increasing values of the radius $`a`$ of the universe $`_a`$ the maximum of the signature $`\varphi _H^{sc}(a,s)`$ shifts to the right. The most likely value of $`s`$ increases for increasing values of the radius $`a`$, and for $`a1`$ there is noticeable concentration of large values of $`s`$ near the extreme value $`s=2a`$. Graphs of the distinguishing mark $`\varphi _S^{sc}(a,s)`$ of elliptic simply-connected universes $`_a`$ for four different values of the radius $`a`$. The normalized number of pairs in the vertical axis is given in unit of $`1/(2a)`$, while in the horizontal axis the unit of length is equal to $`2a`$. For increasing values of $`a`$ from $`0`$ to $`\pi `$ the maximum of the signature $`\varphi _S^{sc}(a,s)`$ moves towards the smaller values of $`s`$. This behavior is the opposite of that corresponding to the hyperbolic case shown in figure 3. The most likely value of the separation $`s`$ decreases for increasing values of $`a`$.
warning/0007/math0007091.html	ar5iv	text	# Projective Dimension is a Lattice Invariant ## 1. Introduction. Back in the late 60’s, Roger Wiegand asked the following question in : > Let $`R`$ be a commutative \[von Neumann\] regular ring and $`J`$ an ideal of $`R`$ generated by a set $``$ of idempotents. Let $``$ be the Boolean algebra of all idempotents of $`R`$. Then is the projective dimension of $`J=R`$ as an $`R`$-module the same as the projective dimension of $``$ as a $``$-module? In this paper we show that the answer to this question is ‘yes’. Richard Pierce popularized this problem, and did some of the early work on it. It is not difficult to see that the answer is ‘yes’ if $`J`$ is projective. In Pierce showed that projective dimension of an ideal generated by an independent set of idempotents in a boolean algebra was $`\kappa `$ where the independent set had cardinality $`\mathrm{}_\kappa `$ (here $`\kappa \omega `$ is replaced by $`\mathrm{}`$ for projective dimension). Osofsky proved the same result for arbitrary commuting idempotents in any ring, so in the case of ideals generated by independent idempotents the answer to the Wiegand question is ‘yes’. Then Richard Pierce showed that it is ‘yes’ in case either the projective dimension of $`R`$ or the projective dimension of $``$ is one. Since then, the problem has been solved in some special cases with extra hypotheses on the idempotents forcing projective dimension to be the subscript of the minimal aleph of a generating set, although the general problem remained open. The essence of the problem is that the additive order of some of the idempotents in $``$ might be one prime (for example the prime 2 in case $`R=`$) and a different prime in another ring $`R^{}`$, or perhaps even infinite in a third ring $`R^{\prime \prime }`$. Here we conquer the problem of different primes by working in a (not regular) ring $`R`$ of characteristic 0. We show that the answer to Wiegand’s question is ‘yes’ in all cases. In Section 2, we prove a subtle but elementary result about free abelian groups, namely for any free abelian group $`G`$ and any direct sum decomposition of $`G/p^\nu G`$, this decomposition lifts to a direct sum decomposition of $`G`$. In Section 3 we apply this result to any commutative von Neumann regular ring $`R`$ containing a lattice of idempotents isomorphic to $``$. Unlike Pierce’s paper concerning the case of projective dimension 1 (), we do not give an internal characterization of projective dimension of ideals in a commutative von Neumann regular ring. However, there is a candidate for such a characterization in a series of papers by the author: , , and . ## 2. A theorem on abelian groups The aim in this section is to prove the following: ###### Theorem A. Let $`G`$ be a free abelian group and $`\{\overline{b_\alpha }:\alpha \}`$ a (free) basis for $`G/p^\nu G`$ with $`p`$ a prime. Then there exists a family of integers $`\{u_\alpha :\alpha \}`$, relatively prime to $`p`$, and a free basis of $`G`$, $`\{y_\alpha :\alpha \}`$, such that $`\overline{y_\alpha }=u_\alpha \overline{b_\alpha }`$ in $`G/p^\nu G`$ for all $`\alpha `$. A way of restating this theorem is that the direct sum decomposition $`G/p^\nu G=_\alpha \overline{b_\alpha }/p^\nu `$ lifts to a direct sum decomposition $`G=_\alpha y_\alpha `$. In fact, any direct sum decomposition of $`G/p^\nu G`$ will lift to a direct sum decomposition of $`G`$ by taking bases of each of the summands and lifting them. We use the fact that the ring $`/p^\nu `$ is local, that is, has a unique maximal ideal. If $`p^\nu `$ is replaced by an arbitrary integer which has at least two distinct prime factors, the result is false since $``$ is indecomposable whereas $`/n`$ decomposes if $`n`$ is a product of two relatively prime factors $`>1`$. ### Basic notation. We fix a prime power $`p^\nu `$. For any abelian group $`G`$, we denote the natural map from $`G`$ to $`G/p^\nu G`$ by an overline. If $`\overline{x}`$ is an element of $`\overline{G}=G/p^\nu G`$ we will assume from the notation that $`xG`$ is some preimage of $`\overline{x}`$. If $`G`$ is some free abelian group, we will denote some free basis for $`G`$ by $$𝔛=\{x_\sigma :\sigma 𝔎\}.$$ and we will denote a basis of $`\overline{G}`$ as a (free) $`\overline{}`$-module by $$𝔅=\{\overline{b_\alpha }:\alpha \}.$$ ### Reduction to the countable case. Much of this paper relies heavily on a beautiful paper by Kaplansky () for both technique and results. Here we adapt the basic technique of Kaplansky’s paper to get a specialized result on free abelian groups. We have the same objective as Kaplansky did, namely to reduce the question under study to the countable case. ###### Lemma 2.1. Let $`G`$ be a nonzero free abelian group with free basis $`𝔛`$, and let $`𝔅`$ be a basis of $`\overline{G}`$ as a (free) $`\overline{}`$-module. Let $`𝔠`$ be any countable subset of $`𝔅`$. Then there exists a nonzero countably generated direct summand $`H`$ of $`G`$ such that $$\overline{H}=\underset{i=0}{\overset{\mathrm{}}{}}\overline{b_{\alpha _i}}\overline{}$$ for $`\{\overline{b_{\alpha _i}}:i\omega \}`$ some countable subset of $`𝔅`$ containing $`𝔠`$. Moreover, $`H`$ itself is generated by a countable subset of $`𝔛`$. ###### Proof. We are given that $`𝔛=\{x_\sigma :\sigma 𝔎\}`$ is a free basis for $`G`$. Fix a lifting $`\left\{b_\alpha \right\}`$ of $``$. For any countable subset $`𝔠`$, let $`X_𝔠𝔎`$ be the smallest (necessarily countable) subset of $`𝔎`$ such that $`_{\alpha 𝔠}b_\alpha _{\sigma X_𝔠}x_\sigma `$. Similarly, for any countable subset $`𝔠^{}𝔎`$, let $`B_𝔠^{}`$ be the smallest (necessarily countable) subset of $``$ such that $`_{\sigma 𝔠^{}}x_\sigma _{\alpha B_𝔠^{}}b_\alpha `$. Now start with any nonempty countable set $`𝔠_0`$ such that $`𝔠𝔠_\mathfrak{0}`$. We use finite induction to define two sequences $`\{𝔠_i,𝔠_i^{}:i<\omega \}`$ of countable sets by $`𝔠_n^{}`$ $`=`$ $`X_{𝔠_𝔫}`$ $`𝔠_{n+1}`$ $`=`$ $`B_{𝔠_n^{}}.`$ In words, think of $`𝔅`$ as images of $`\{b_\alpha :\alpha \}`$. Starting with a countable subset $`𝔠_0`$ of the basis $`𝔅`$ of $`\overline{G}`$, use our lifting of $``$ to get an inverse image $`c_0G`$ and take the smallest countable subset $`𝔠_0^{}`$ of the basis $`𝔛`$ of $`G`$ whose span contains $`c_0`$. Now take images of $`𝔠_0^{}`$ modulo $`p^\nu `$ and find the smallest countable subset $`𝔠_\mathfrak{1}𝔠_0`$ of the basis $`𝔅`$ which span a group containing all of the elements of $`\overline{𝔠_i}`$. Iterate a countable number of times. We then have for all $`i`$, $`𝔠_i^{}𝔠_{i+1}^{}`$ and ($``$) $$\overline{F_n}=\underset{\alpha 𝔠_n}{}\overline{b_\alpha }\overline{}\overline{G_n}=\underset{\sigma 𝔠_n^{}}{}\overline{x_\sigma }\overline{}\overline{F_{n+1}}=\underset{\alpha 𝔠_{n+1}}{}\overline{b_\alpha }\overline{}.$$ Set $`H=_{\sigma _{n=0}^{\mathrm{}}𝔠_n^{}}x_\sigma `$. Clearly $`H`$ is a direct summand of $`G`$. Moreover, $`H`$ is countably generated since the indexing set is a countable union of countable sets. Equation ($``$) forces $$\overline{H}=\underset{\alpha _{i=0}^{\mathrm{}}F_i}{}\overline{b_\alpha }\overline{}.$$ ###### Lemma 2.2. Let $`G`$ be a nonzero free abelian group, and let $`𝔅`$ be a basis of $`\overline{G}`$ as a (free) $`\overline{}`$-module. Then $`G`$ is the union of a well-ordered (by inclusion) family $`\{H_\mu :\mu <\mathrm{\Omega }\}`$ of subgroups such that: $`H_\mu `$ and $`_{\kappa <\mu }H_\kappa `$ are direct summands of $`G`$ for every $`\mu `$ in the ordinal $`\mathrm{\Omega }`$; for each $`\mu `$, $`H_\mu /_{\kappa <\mu }H_\kappa `$ is countable; and each $`\overline{H_\mu }`$ is generated by some subset of the $`\{\overline{b_\alpha }:\alpha \}`$. ###### Proof. Fix a basis $`𝔛`$ of $`G`$. Well order $`𝔅`$. Assume we have $`H_\kappa `$ for all $`\kappa <\mu `$ such that: 1. Each $`H_\kappa `$ is generated by a subset of $`𝔛`$; 2. $`\overline{H_\kappa }`$ is generated by some subset of $`𝔅`$; and 3. $`H_\kappa H_\kappa ^{}`$ if $`\kappa >\kappa ^{}`$. 4. $`H_\kappa /_{\kappa ^{}<\kappa }H_\kappa ^{}`$ is countably generated. $`H_\kappa `$ and $`_{\mu <\kappa }H_\mu `$ are direct summands of $`G`$ since they are generated by subsets of our fixed basis. If $`_{\kappa <\mu }H_\kappa G`$, that union cannot map onto $`\overline{G}`$. Let $`\overline{b_\beta }`$ be the smallest element of $`𝔅`$ (under the well ordering of $`𝔅`$) not in $`\overline{_{\kappa <\mu }H_\kappa }`$. Apply Lemma 2.1 to get a countably generated subgroup $`K_\mu `$ generated by elements of $`𝔛`$ with $`\overline{b_\beta }\overline{K_\mu }`$ and $`\overline{K_\mu }`$ generated by a subset of $`𝔅`$. Set $`H_\mu =K_\mu +_{\nu <\mu }H_\nu `$. Since $`H_\mu `$ clearly has the required properties and this process must eventually give all of $`G`$ (at least by the order type of $`𝔅`$), by transfinite induction we are done. ∎ ###### Corollary 2.3. Assume that, for any countably generated free abelian group $`G`$ with $`𝔅`$ a basis for $`\overline{G}`$, there is a direct decomposition lifting of $$\overline{G}=\underset{\alpha }{}\overline{b_\alpha }\overline{}$$ to the direct decomposition $$G=\underset{\alpha }{}y_\alpha .$$ Then Theorem A is true for any free abelian group $`G`$. ###### Proof. Using the notation of Lemma 2.2, we let $`G=_{\mu <\mathrm{\Omega }}H_\mu `$ where for all $`\mu <\mathrm{\Omega }`$, $`H_\mu =K_\mu +_{\kappa <\mu }H_\kappa `$ with $`K_\mu `$ countably generated. For each $`b_{\alpha _i}`$ in $`K_\mu _{\kappa <\mu }H_\kappa `$, set $`b_{\alpha _i}=c_i+b_i^{}`$, where $`c_i`$ is the projection of $`b_{\alpha _i}`$ to $`_{\kappa <\mu }H_\kappa `$. If $`b_i^{}=0`$, ignore it and renumber. By assumption, we can lift the direct sum decomposition of the quotient $$\overline{H_\mu /\underset{\kappa <\mu }{}H_\kappa }\overline{K_\mu /K_\mu \underset{\kappa <\mu }{}H_\kappa }=\underset{i=0}{\overset{\mathrm{}}{}}\overline{b_i^{}}\overline{}$$ to a direct sum decomposition $$K_\mu /K_\mu \underset{\kappa <\mu }{}H_\kappa =\underset{b_{\alpha _i}K_\mu _{\kappa <\mu }H_\kappa }{}y_i^{}$$ with units $`\left\{u_i\right\}`$ such that $`y_i^{}u_ib_i^{}p^\nu G`$. Now set $`y_{\alpha _i}=y_i^{}+u_ic_i`$ so that $`y_{\alpha _i}`$ lifts $`b_{\alpha _i}`$. Assume for all $`\mu <\lambda `$, $`H_\mu =_{\kappa \mu }L_\kappa `$, where $`L_\kappa `$ is the free group generated by a lifting of the decomposition of $`H_\kappa /_{\kappa ^{}<\kappa }H_\kappa ^{}`$ generated by the appropriate subset of $`𝔅`$. Then we have $`_{\mu <\lambda }H_\mu =_{\mu <\lambda }L_\mu `$ and by the above, $`H_\lambda =L_\lambda _{\mu <\lambda }L_\mu `$. By transfinite induction we get $`G=_{\lambda <\mathrm{\Omega }}L_\lambda `$. ∎ ### Infinite Gaussian elimination modulo $`p^\nu `$. The reader is assumed thoroughly familiar with the details of Gaussian elimination as developed in an introductory linear algebra course. Infinite Gaussian elimination on a row finite $`\omega \times \omega `$ matrix can proceed very much like the algorithm on a finite matrix. As in , one looks for a pivot in a row rather than a column as in many texts and standard implementations of finite Gaussian elimination. That insures that only a finite number of entries need to be examined to either obtain a unit pivot or to know that no such pivot exists. Subtracting multiples of a pivot row from all other rows to make entries in the pivot column equal to 0 will, in general, involve an infinite number of operations before the algorithm is complete. To avoid this, in the infinite case, rows are included with previously obtained pivot rows one at a time, and one clears the previously obtained pivot columns in a row at the time that the row is included, and then finds a pivot if possible and clears above the pivot in the new pivot column. In the infinite case there is no LU decomposition or forward pass and back substitution because these might lead to rows changing infinitely often, and there are no row permutations because some row might conceivably be permuted to a higher numbered position an infinite number of times and thus never examined for a pivot. However, it is still the case that a row finite $`\omega \times \omega `$ matrix is invertible if and only if with these modifications of standard Gaussian elimination, infinite Gaussian elimination will row reduce the matrix to a matrix whose columns are a permutation of the columns of the identity matrix. We now modify infinite Gaussian elimination to produce an algorithm which we call infinite Gaussian elimination modulo $`p^\nu `$.<sup>1</sup><sup>1</sup>1The author has a working Maple V implementation of this algorithm. See the appendix in the copy of this paper archived on http://arXiv.org or URL http://www.math.rutgers.edu/pub/osofsky/getbasis.html This algorithm clearly also works if we have a finite matrix $`A`$. We indicate the variables needed in the algorithm with a little information about them, then give the steps of the algorithm, and then add a step by step explanation of what unusual steps do. We start with a row finite $`\omega \times \omega `$ matrix $`𝐀`$ with entries in $``$. In our proof of Theorem A, the rows of $`𝐀`$ will be some lifting of a given basis for $`\overline{}^{(\omega )}`$ to elements of $`^{(\omega )}`$. By the expression ‘principal submatrix’ of an infinite matrix, we will mean the submatrix obtained by taking the first $`n`$ rows and first $`k`$ columns of the matrix, where $`n`$ and $`k`$ are both finite. A ‘principal minor’ will be the determinant of a square principal submatrix. Additional variables are needed to perform the algorithm. We use a diagonal matrix $`𝐔`$ (or a countable row vector) to hold units modulo $`p^\nu `$. Multiplying row $`i`$ of $`𝐀`$ by an appropriate unit $`𝐔_{i,i}`$ enables us to make a crucial determinant 1. The actual row reduction is done in arbitrarily large but finite principal submatrices of an $`\omega \times \omega `$ matrix $`𝐑`$. Another $`\omega \times \omega `$ matrix $`𝐂`$ (for candidates) holds, in a finite principal submatrix, the current candidates for lifting basis elements times units. These candidates change during the elimination but each row only changes a finite number of times. As the algorithm progresses, we multiply (an initial segment of) row $`i`$ of $`𝐀`$ by the appropriate unit $`𝐔_{i,i}`$ (integer relatively prime to $`p^\nu `$) and then insert it into both $`𝐑`$ and $`𝐂`$. All changes to $`𝐂`$ other than the concatenation of rows from $`\mathrm{𝐔𝐀}`$ consist of adding multiples of $`p^\nu `$ to entries so nothing changes modulo $`p^\nu `$. In addition, we use a finite square matrix $`𝐌`$ which is generated from a submatrix of $`𝐂`$ and has determinant 1. At the end of each loop of this algorithm, the matrix $`𝐑`$ will be a row reduction of $`𝐂`$ with row operations captured by $`𝐌`$. Also, any entry of $`𝐑`$ which is a multiple of $`p^\nu `$ is 0; it is set to 0 before any arithmetic is done using it. At any given stage of the algorithm we work with finite matrices large enough to hold all nonzero entries in a finite number of rows. Moreover, the results of each loop of the algorithm applied to $`\overline{𝐀}`$ are identical with the results of applying normal infinite Gaussian elimination to $`𝐀`$. ###### Algorithm 1 (Infinite Gaussian elimination modulo $`p^\nu `$). We start with an $`\omega \times \omega `$ integer valued row finite matrix $`𝐀`$. * Initialize. Let your row index I be set to 0. Set up the matrix variables $`𝐌`$, $`𝐂`$, $`𝐑`$ and $`𝐔`$. Set up a row vector J to hold pivot columns. Read the $`0^{th}`$ row of $`𝐀`$ into $`𝐑`$, replacing any element divisible by $`p^\nu `$ with 0. * For K going from 0 to I – 1, subtract $`𝐑_{I,J(K)}`$ times row K of $`𝐑`$ from row I of $`𝐑`$. * Search row I of $`𝐑`$ for the first entry which is relatively prime to $`p`$. If no such element is found then STOP. The rows of $`𝐀`$ do not form a basis modulo $`p^\nu `$. Otherwise, let the first entry relatively prime to $`p`$ be in column J(I), and call column J(I) the I<sup>th</sup> pivot column. * Set $`𝐔_{I,I}`$ equal to an integer $`u`$ such that $`u𝐑_{I,J(I)}1modp^\nu `$. Multiply row I of $`𝐀`$ by $`u`$. If some entry in the resulting row is a multiple of $`p^\nu `$, set that entry to 0. Insert the result as row I in both $`𝐂`$ and $`𝐑`$. * For K going from 0 to I – 1, subtract $`𝐑_{I,J(K)}`$ times row K of $`𝐑`$ from row I of $`𝐑`$. * The pivot in row I of $`𝐑`$ is now congruent to 1 modulo $`p^\nu `$. Subtract a multiple of $`p^\nu `$ from it to make the pivot 1. Subtract the same multiple of $`p^\nu `$ from the $`(I,J(I))`$ entry of $`𝐂`$. * If any entry in row I of $`𝐑`$ is a multiple of $`p^\nu `$, subtract that multiple of $`p^\nu `$ from the corresponding entry in $`𝐂`$ and set the entry in $`𝐑`$ equal 0. * For K going from 0 to I – 1, subtract $`𝐑_{K,J(I)}`$ times row I of $`𝐑`$ from row K of $`𝐑`$ to clear every entry in column J(I) above the I<sup>th</sup> row. * If any entry in $`𝐑`$ is a multiple of $`p^\nu `$, then set that entry equal to 0. * Set $`𝐌`$ equal to the matrix $`\left[𝐂_{\mathrm{K},\mathrm{J}(\mathrm{K})}\right]_{0\mathrm{K}\mathrm{I}}`$. Set $`𝐂=\mathrm{𝐌𝐑}`$. * For each nonpivot column $`\mathrm{}`$ of $`𝐂`$, check to see if the first nonzero entry $`𝐂_{k,\mathrm{}}`$ is divisible by $`p^\nu `$. If so, form the set $`𝔖_{\mathrm{}}`$ consisting of all $`l_i`$ such that column $`l_i`$ is a pivot column, $`𝐂_{k,l_i}`$ is the first nonzero entry in column $`l_i`$, and $`𝐑_{l_i,\mathrm{}}0`$. If $`𝔖_{\mathrm{}}\mathrm{}`$, check if $`p^\nu `$ times the gcd of $`𝔖_{\mathrm{}}`$ divides $`𝐂_{k,\mathrm{}}`$. If so, express this gcd as a sum $`_𝔖_{\mathrm{}}𝐂_{k,l_i}b_{l_i}`$. Form a column vector with zeros everywhere except for $`b_{l_i}𝐂_{k,\mathrm{}}/d`$ in row $`l_i`$, and add this to column $`\mathrm{}`$ of $`𝐑`$. Premultiply by $`𝐌`$, and use the result as the new column $`\mathrm{}`$ of $`𝐂`$. The new $`𝐂_{k,\mathrm{}}`$ will be 0. * Read row I + 1 of $`𝐀`$ into $`𝐑`$, replacing multiples of $`p^\nu `$ by 0. * Increment I by 1 and GOTO Step 2. END That is the end of the algorithm. To get a picture of what is happening, at the end of the $`(n1)^{th}`$ loop at Step 13 the column permuted matrix $`𝐑`$ (picturing $`j(i)`$ as though it were $`i`$) looks like $$\begin{array}{c}\text{R}\end{array}=\left[\begin{array}{cccccccccccccccccccc}\begin{array}{cccccccccccccccccccc}1& 0& \mathrm{}& 0& & r_{0,n}& & & & & & & & & & & & & & \\ 0& 1& \mathrm{}& 0& & r_{1,n}& & & & & & & & & & & & & & \\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& & \mathrm{}& & & & & & & & & & & & & & \\ 0& 0& \mathrm{}& 1& & r_{n1,n}& & & & & & & & & & & & & & \\ & & & & & & & & & & & & & & & & & & & \\ r_{n,0}& r_{n,1}& \mathrm{}& r_{n,n1}& & r_{n,n}& & & & & & & & & & & & & & \end{array}& & \begin{array}{cccccccccccccccccccc}& & & & & & & & & & & & & & & & & & & \\ \text{B}\mathrm{}& & & & & & & & & & & & & & & & & & & \\ & & & & & & & & & & & & & & & & & & & \\ & & \multicolumn{-1}{c}{}& & & & & & & & & & & & & & & & & \\ \text{D}\mathrm{}& & & & & & & & & & & & & & & & & & & \end{array}& & & & & & & & & & & & & & & & & \end{array}\right]$$ for an appropriate finite matrix $`𝐁`$ and finite row $`𝐃`$, and all entries in $`𝐑`$ which are divisible by $`p^\nu `$ are 0. Now for a more detailed explanation of how this algorithm works. In the permuted matrix used in the discussion, $`j(i)`$ will be treated as though it were $`i`$ to aid in visualization of the progress of the algorithm. That is, we will pretend that we have permuted the columns of the matrix. Step 2 is the first pass at clearing already obtained pivot columns (which have pivot 1) in row $`i`$. It is used to get the unit mod $`p^\nu `$ we must multiply the $`i^{th}`$ row of $`𝐀`$ by to make sure that we can make the pivot in row $`i`$ equal to 1. It is not performed when $`i=0`$. Step 6 relies on the claim that the pivot is congruent to 1 modulo $`p^\nu `$. Why is that claim true? Adding one row of a matrix to another corresponds to premultiplication by a matrix of determinant 1. After Step 3, if we look at the principal minor of the column permuted matrix $`𝐑`$, it has determinant the $`(i,i)`$ entry of the permuted $`𝐑`$ because it is upper triangular with all other diagonal entries 1. When we multiply what was the last row before Step 3 by $`u`$, we make that determinant congruent to 1 modulo $`p^\nu `$. Now we redo the elementary row operations of determinant 1 to get an upper triangular matrix with element in the $`(i,j(i))`$ slot equal to the determinant. In Step 6, subtracting multiples of $`p^\nu `$ from the same entries in both $`𝐑`$ and $`𝐂`$ does not change $`\overline{𝐂}`$ and does insure that the elementary row operations we have done so far will reduce the new $`𝐂`$ to the new $`𝐑`$. Since we want entries in $`𝐑`$ congruent to 0 mod $`p^\nu `$ to be 0, we set them to 0 in Step 9. This can only affect entries in nonpivot columns. Now we must make sure that our $`𝐂`$ row reduces to the new $`𝐑`$. This is done in Step 10. At this stage, the appropriate principal submatrix of the column permuted matrix $`𝐑`$ is the identity matrix. So the row operations we have done have reduced the corresponding principal submatrix of the column permuted matrix $`𝐂`$ to the identity. By standard linear algebra, the matrix $`𝐌`$ is the inverse of the product of the elementary matrices which produce this elimination by premultiplication. Thus from $`𝐑=𝐌^1\mathrm{𝐌𝐑}`$ we see that setting $`𝐂=\mathrm{𝐌𝐑}`$ gives us a matrix which row reduces to the new $`𝐑`$, and since $`𝐑`$ did not change modulo $`p^\nu `$, neither did $`\mathrm{𝐌𝐑}`$. In Step 11, the algorithm bounds the power of $`p^\nu `$ that can divide entries of $`𝐂`$ after the corresponding row of $`𝐀`$ becomes all zeros. This step may change $`𝐂`$ and nonzero entries in $`𝐑`$ modulo $`p^\nu `$. If the first entry $`𝐂_{k,l}`$ in a nonpivot column of $`𝐂`$. The several imposed conditions on $`𝔖_{\mathrm{}}`$ insure that no zero entry of $`𝐑`$ becomes nonzero, and the divisibility property makes the added vector a multiple of $`p^\nu `$. If a nonzero entry appears in $`𝐂`$ after all the nonzero mod $`p^\nu `$ entries in its row occur in pivot columns, it may propagate, but that leads to entries in the row divisible by higher powers of $`p^\nu `$, and eventually Step 11 will make all of these entries zero. Thus Step 11 makes sure that no row has an infinite number of entries congruent to 0 modulo $`p^\nu `$. New row operations are only done to the rows above the pivot row when their entries in the current pivot column is nonzero. Hence once the finite set of rows of $`𝐑`$ from $`0`$ to $`i`$ have zero entries except for a pivot of 1, and there are no more nonzero multiples of $`p^\nu `$ in these rows of $`𝐂`$, those rows will no longer be affected by the elimination process. The last steps of the algorithm just set up for the next loop. ### The proof of Theorem A. ###### Proof. By Corollary 2.3, it is enough to show that, for a countably generated free abelian group $`G`$ with $`𝔅`$ a basis for $`\overline{G}`$, there is a direct decomposition lifting of $$\overline{G}=\underset{\alpha 𝔅}{}\overline{b_\alpha }\overline{}$$ to the direct decomposition $$G=\underset{\alpha 𝔅}{}y_\alpha $$ Form a matrix $`𝐀`$ whose rows are some lifting of $`𝔅`$. Do infinite Gaussian elimination modulo $`p^\nu `$ on $`𝐀`$. Since the rows of $`\overline{𝐀}`$ form a basis for $`\overline{}^{(\omega )}`$ and modulo $`p^\nu `$ this algorithm agrees with infinite Gaussian elimination, after a finite number of steps, the top $`i+1`$ rows of $`𝐑`$ will be rows of the identity and all rows of the identity will eventually arise as rows of $`𝐑`$. Since every entry of $`R`$ which is zero modulo $`p^\nu `$ is actually $`0`$, $`𝐂`$ is row reduced to the identity provided every row at some point stops changing in taking the product $`\mathrm{𝐌𝐑}`$. Since all of the entries of row $`n`$ of $`𝐂`$ which are not congruent to $`0`$ mod $`p^\nu `$ are contained in a finite number of columns, any row of $`𝐂`$ ceases to change when all the rows of the identity with $`1`$ in those columns have been obtained in the matrix $`𝐑`$. Hence after an infinite number of steps each row of $`𝐂`$ will have stabilized and the stabilized rows of $`𝐂`$ will form a basis for $`𝐙^{(\omega )}`$ which lifts the direct sum decomposition. ∎ ## 3. Lattices of commuting idempotents ### Definitions and notation. The following notation will be used, usually without comment, in the rest of this paper. Let $``$ be a lattice of commuting idempotents in a ring $`R`$ with $`1`$, that is, $``$ is closed under multiplication and addition of orthogonal idempotents. The idempotents in $``$ together with the identity generate a boolean algebra $``$ under multiplication as in $`R`$ but addition the symmetric difference $`e+_{}f=e\left(1f\right)+f\left(1e\right)`$. Let $`\left[\right]`$ be the semigroup algebra of $`,`$, that is, the free abelian group with basis the elements of $``$ and multiplication the multiplication as in $``$. Let $$𝒮=\left[\right]/(e+f)ef:ef=0.$$ $`𝒮`$ is a free lattice ring in the sense that it can be formed for any modular, complemented lattice and has appropriate universal properties with respect to embedding such lattices in rings. For convenience, we will assume that $``$ is a Boolean ideal, that is, if $`f=f^2R`$, then $`f`$. This does not change $`R`$. ### Elementary properties of $`𝒮`$. Much of the known material assumed in this subsection can be found in graduate level text books such as . The next proposition is essentially a sequence of remarks, included with short proofs. ###### Proposition 3.1. The following hold for the free lattice ring $`𝒮`$. 1. The additive group of $`𝒮`$ is torsionfree. 2. The lattice of idempotent generated ideals of $`𝒮`$ is isomorphic to $``$. 3. Any finitely generated ideal of $`𝒮`$ is cyclic and isomorphic to a sum $`_{i=1}^nf_i𝒮`$ for some set of orthogonal idempotents $`\left\{f_i\right\}`$. 4. $`R`$ is an $`𝒮`$-module under the map induced by the inclusion of $``$ in $`R`$. 5. The projective dimension of an idempotent generated ideal $`I`$ of $`𝒮`$ is greater than or equal to the projective dimension over $`R`$ of the module $`I_𝒮R`$. ###### Proof. 1. The kernel of the ring map from $``$ to $`𝒮`$ is generated by idempotents and so pure. 2. Any element of $`𝒮`$ is of the form $`_{i=1}^ne_in_i`$ where $`\left\{e_i\right\}`$ are pairwise orthogonal and $`n_i`$. Assume such an element is idempotent. By the torsionfree property of $`𝒮,+`$, the $`n_i`$ must be all $`1`$, and $`e=_{i=1}^ne_i`$. But then the symmetric difference of $`e`$ and $`f`$ is the same as in $``$. 3. Given a finite set of idempotents $`\{e_i:1in\}`$, the minimal nonzero idempotents in the lattice they generate will be pairwise orthogonal and generate the same lattice. Since $`𝒮`$ is a quotient of the ring $`\left[\right]`$, any element of $`𝒮`$ is of the form $`_{j=1}^me_jn_j`$. Moreover, if the $`\left\{e_j\right\}`$ happen to be orthogonal, $`\left(_{j=1}^me_jn_j\right)𝒮=_{j=1}^m\left(e_jn_j𝒮\right)`$. Now let $`I`$ be the finitely generated ideal $$I=\underset{i=1}{\overset{k}{}}\left(\underset{j=1}{\overset{l_j}{}}e_{i,j}n_{i,j}𝒮\right)𝒮.$$ Split each $`e_{i,j}`$ into an orthogonal sum of the nonzero minimal elements in the lattice generated by $`\{e_{i,j}:1jl_j,\mathrm{\hspace{0.17em}1}ik\}`$. Collecting multiples of each of these minimal elements, we get a generator for $`I`$ of the form $`_{i=1}^k^{}f_im_i`$ where the $`\left\{f_i\right\}`$ are pairwise orthogonal idempotents in $``$. But $`\left(_{i=1}^{k^{\prime \prime }}f_im_i\right)𝒮_{i=1}^{k^{\prime \prime }}f_i𝒮`$ if we ignore terms with $`m_i=0`$. 4. The obvious map $`\left[\right]R`$ is a ring homomorphism whose kernel contains $$\left(e+f\right)ef:ef=0.$$ 5. $`I`$ is a direct limit of idempotent generated cyclics and so flat. A projective resolution $$\mathrm{}P_iP_{i1}\mathrm{}P_0I0$$ is therefore pure exact. Moreover, since $`P_i`$ is a projective $`𝒮`$-module, $`P_i_𝒮R`$ is a projective $`R`$-module. Thus $$\mathrm{}P_i_𝒮RP_{i1}_𝒮R\mathrm{}P_0_𝒮RI_𝒮R0$$ is a projective resolution of $`I_𝒮R`$. If the kernel of a map $`P_iP_{i1}`$ is $`𝒮`$-projective, by pure exactness and the fact that tensoring preserves projectivity we see that the kernel of $`P_i_𝒮RP_{i1}_𝒮R`$ is $`R`$-projective. Thus the $`𝒮`$-projective dimension of $`I`$ is at most $`i`$ implies that the $`R`$-projective dimension of $`I_𝒮R`$ is also at most $`i`$. ∎ ###### Proposition 3.2. The additive group of $`𝒮`$ is a free abelian group. ###### Proof. Let $`𝔛`$ be the family of all subsets $`X`$ of $`\left\{0\right\}`$ such that whenever $`\left\{e_i\right\}`$ is a set of orthogonal idempotents in $`X`$, if $`\left\{f_j\right\}`$ is any set of orthogonal idempotents such that $`\left\{e_i\right\}\left\{f_j\right\}`$ and $`_ie_i=_jf_j`$, then at least one $`f_jX`$. $`𝔛`$ is an inductive poset under $``$, so by Zorn’s lemma there is a maximal element $`B`$ in $`𝔛`$. $`B`$ is $``$-linearly independent in $`𝒮`$ because the only relations on the $``$-linearly independent idempotents in $`\left[\right]`$ set an idempotent equal to an orthogonal sum of other idempotents. $`B`$ will be a vector space basis for $``$ over the field of 2 elements. Let $`f\left\{0\right\}`$. If $`fB`$, then $`B\left\{f\right\}𝔛`$. Hence there must be a set $`\left\{e_i\right\}`$ of orthogonal idempotents in $`B\left\{f\right\}`$ and a different set $`\left\{f_i\right\}B\left\{f\right\}`$ of orthogonal idempotents with $`_{i=1}^ne_i=_{j=1}^mf_j`$. If $`f\left\{e_i\right\}\left\{f_j\right\}`$ then we get $`_{e_if}e_i=_{f_jf}f_j`$ with all summands in $`B`$, a contradiction. Similarly, if $`f\left\{e_i\right\}\left\{f_j\right\}`$ we get a contradiction. Hence $`f`$ is in precisely one of the two sets, say $`f=e_1`$. Then $`f=_jf_j_{i=2}^ne_i`$ is in the span of $`B`$. ∎ Proposition 3.2 strongly reinforces the observation that $`𝒮`$ is a free object. The basis found for its additive group will be a basis for $`𝒮_𝒮F`$ over $`F`$ for any field $`F`$. In his proof in of the affirmative answer to the Wiegand question in the case $`n=1`$, R. S. Pierce proved the next lemma with completely different terminology. See \[9, Lemma 2.7\]. ###### Proposition 3.3. Let $`\left\{\kappa _\alpha \right\}`$ be a set of elements in a submodule of a free $`𝒮`$-module $`K`$, where the $`\left\{\kappa _\alpha 1\right\}`$ are all nonzero. Then if $`\left\{\kappa _\alpha _𝒮R\right\}`$ is $`R`$-independent in $`K_𝒮R`$, then $`\left\{\kappa _\alpha \right\}`$ is $`𝒮`$-independent in $`K`$. ###### Proof. Assume not. Then there is a shortest sum $`_{i=1}^n\kappa _{\alpha _i}s_i=0`$ where the summands are all nonzero in $`(S)`$. Considering elements of the free module $`K`$ as consisting of sums of idempotents times basis elements, we see that the annihilator of each $`\kappa _{\alpha _i}s_i`$ is generated by an idempotent $`\left(1\epsilon _i\right)`$. Since $`n`$ is the smallest number of summands that can give you a zero and $`_{i=1}^n\kappa _{\alpha _i}s_i\epsilon _1=0`$, we have $`\kappa _{\alpha _i}s_i\epsilon _10`$ for all $`i`$. Similarly $`\kappa _{\alpha _i}s_i\epsilon _1\epsilon _20`$ for all $`i`$. Continuing in this manner we get $`\kappa _{\alpha _i}s_i_{j=1}^n\epsilon _j0`$ for all $`i`$. Then $`_i\kappa _is_i_{j=1}^n\epsilon _j`$ has all summand nonzero and there is an integer $`m`$ such that $`_{i=1}^n\kappa _{\alpha _i}s_im^1_{i=1}^n\epsilon _i`$ is an element not divisible by any integers other than $`\pm 1`$ in the free abelian additive group of $`K`$. But then $`_{i=1}^n\kappa _{\alpha _i}s_im^1_{i=1}^n\epsilon _i1`$ is nonzero in $`K_𝒮R`$ and each of the summands is nonzero. ∎ We quote a Proposition due to Kaplansky that is basic to almost all studies of infinitely generated projective modules, with two consequences giving rise to the same result for von Neumann regular rings. ###### Proposition 3.4 (Kaplansky). A projective module over any ring is a direct sum of countably generated submodules. From this we obtain: 1. Any projective right module over a von Neumann regular ring is isomorphic to a direct sum of cyclic (idempotent generated) right ideals. 2. Any projective module over a commutative semihereditary ring is isomorphic to a direct sum of finitely generated right ideals. See for a proof. The proof of this theorem is the template on which the preliminary proofs in Section 2 are based. ### The proof of an affirmative answer to the Wiegand question. We now complete our work on the Wiegand question. ###### Proposition 3.5. Let $`R`$ be a commutative von Neumann regular ring. Let $`F`$ be a projective $`𝒮`$-module and let $`K`$ be any pure submodule of $`F`$. Then if $`K_𝒮R`$ is projective as an $`R`$-module, then $`K`$ is projective as an $`𝒮`$-module. ###### Proof. Since $`K_𝒮R`$ is a projective $`R`$-module, it is a direct sum of the form $`K_𝒮R=_\alpha x_\alpha R`$ where for each $`\alpha `$ there is an $`e_\alpha `$ such that $`x_\alpha Re_\alpha R`$. If any $`e_\alpha `$ is of finite but composite order, express it as an orthogonal sum of idempotents of prime power order by the Chinese Remainder Theorem. In the von Neumann regular case where there are no nilpotent elements, the prime power must be the prime itself. We can then divide the indexing set into a family of subsets $$𝔉_p=\{\alpha :char\left(e_\alpha _𝒮R\right)=p\}$$ for $`p`$ a prime or $`0`$. Consider the map $`K\stackrel{I_k1}{}K_𝒮R_{\alpha 𝔉_0}x_\alpha R`$. Its image is a projective $`𝒮`$-module, so it splits. Hence without loss of generality we can work with the kernel of this map in place of $`K`$ and assume that $`K_𝒮R`$ is torsion. But then it is the orthogonal sum of its $`p`$-primary components so we need only look at sums of the form $`_{\alpha 𝔉_p}x_\alpha R`$ for a fixed prime $`p`$. That is, without loss of generality, $`K_𝒮R`$ is $`p`$-primary. Since the additive group of $`𝒮`$ is free, the additive group of $`F`$ is free and hence $`K`$ is a subgroup of a free abelian group and so free. By Theorem A, there is a basis $`\left\{b_\lambda \right\}`$ of $`K`$ which lifts the direct sum decomposition $`G_p/pG_p=_{\alpha 𝔉_p}x_\alpha R`$ to a direct sum decomposition of $`K`$. For every $`\alpha `$, let $`𝔅_\alpha =\{b_\lambda :b_\lambda 1x_\alpha R\}`$. Let $`H_\alpha `$ be the $`𝒮`$-submodule of $`K`$ generated by $`𝔅_\alpha `$. Since the generators of $`H_\alpha `$ all map to $`x_\alpha R`$ under $`Id_K1_R`$, so must $`H_\alpha `$. Since $`H_\alpha `$ contains $`𝔅_\alpha `$ and $`_\alpha 𝔅_\alpha `$ is a basis for $`K`$, $`K=_\alpha H_\alpha `$. By Proposition 3.3, that sum is direct. Select any element $`y`$ in $`H_\alpha `$ which maps to $`x_\alpha `$. This $`y`$ is an element lying in a finitely generated free submodule of $`F`$. Hence it is of the form $`y=_{i=1}^m_{j=1}^{k_i}c_{i,j}e_{i,j}n_{i,j}`$ where the $`c_{i,j}`$ are basis elements of $`F`$, and we can use our little trick of decomposing into the minimal idempotents in a finite lattice to get that $`e_{i,j}`$ and $`e_{k,l}`$ are either the same idempotent or orthogonal. Because of the $``$-purity of $`K`$, we may find a $`y_\alpha H_\alpha `$ such that each sum of the form $`_{e_{i,j}=e_{k,l}}c_{i,j}e_{i,j}n_{i,j}`$ is of content 1 and hence this $`y_\alpha `$ generates a direct summand of $`F`$. But then $`y_\alpha 𝒮`$ is a direct summand of $`H_\alpha `$ which maps to the same submodule of $`K_𝒮R`$. We conclude that $`H_\alpha =y_\alpha S`$ for all $`\alpha `$. Thus $`K=_\alpha y_\alpha 𝒮`$ so $`K`$ is projective. ∎ ###### Corollary 3.6. Let $`F`$ be a projective $`𝒮`$-module of the form $$F=\underset{\alpha 𝔄}{}e_\alpha 𝒮$$ where each $`e_\alpha 𝒮`$ is isomorphic to an ideal of $`𝒮`$ contained in $`𝒮`$. Then for any pure submodule $`K`$ of $`F`$, $`pd_R\left(K_𝒮R\right)=pd_𝒮\left(K\right)`$. ###### Proof. We can take a short projective resolution of $`K`$ over $`𝒮`$, say $$0LPK0$$ is exact with $`P`$ projective and, like $`F`$, a direct sum of cyclic projectives of the form $`e𝒮`$ for some $`e`$. Then if we let $`\mathrm{}1=\mathrm{}`$, $`pd_𝒮\left(L\right)=pd_𝒮\left(K\right)1`$. This short exact sequence is pure, so tensoring with $`R`$ over $`𝒮`$ gives a short projective resolution of $`K_𝒮R`$ $$0L_𝒮RP_𝒮RK_𝒮R0$$ with $`pd_R\left(L_𝒮R\right)=pd_R\left(K_𝒮R\right)1`$. Induction on $`pd_𝒮\left(K\right)`$ completes the proof. ∎ ###### Theorem B (The answer to the Wiegand question). For any commutative von Neumann regular ring $`R`$ with a commuting set of idempotents $``$, $`pd_R\left(R\right)=pd_𝒮\left(𝒮\right)=pd_{}\left(\right)`$. ###### Proof. $`𝒮`$ has a projective resolution of the form required in Corollary 3.6. Then Corollary 3.6 gives the desired conclusion. ∎ One way to summarize this answer to the Wiegand question is to say that, when working in a submodule of a free module over a commutative regular ring, the lattice of direct summands carries all of the information about the module, and the coefficients essentially none. For example, note that in Theorem B, the lattices of direct summands in the three ideals $`R`$, $`𝒮`$, and $``$ are isomorphic, as they correspond to the idempotents themselves. However, as soon as one gets to free modules on more than one generator, that property fails. Since the number of one dimensional subspaces of a 2-dimensional vector space depends on the cardinality of the field, if $`R=𝒮/3𝒮`$ then the number of direct summands of $`eReR`$ isomorphic to $`eR`$ and the number of direct summands of $`ee`$ isomorphic to $`e`$ will always be different for any idempotent $`e`$. ## 4. Appendix Here is a Maple program which implements an algorithm similar to but not identical with the infinite gaussian elimination modulo $`p^\nu `$ of this paper. The $`𝐂`$ of this algorithm is the analogue of the $`M`$ in the algorithm here. Except for pivot columns, $`𝐑`$ is only determined modulo $`p^\nu `$. The # indicates a comment in the program. The program for the algorithm in this paper, as well as in this appendix, can be found via URL http://math.rutgers.edu/$``$osofsky in both .mws and .html formats. \# mgcdex IS A PROGRAM TO COMPUTE THE GREATEST COMMON DIVISOR OF A \# VECTOR OF INTEGERS, AND A LINEAR COMBINATION OF ENTRIES OF THE \# VECTOR WHICH GIVES THAT GCD. mgcdex:=proc(A,B) local i,j,a,b; with(linalg): B:=array(1..vectdim(A)+1); for i from 1 to vectdim(A) do B\[i\]:=1; B\[vectdim(A)+1\]:=0; od; for i from 1 to vectdim(A) do B\[vectdim(A)+1\]:=igcdex(B\[vectdim(A)+1\],A\[i\],’a’,’b’); B\[i\]:=b; for j from 1 to i-1 do B\[j\]:=a\B\[j\]; od; od; end: \# THE PROGRAM GetBasis IMPLEMENTING A VARIANT OF # GAUSSIAN ELIMINATION MOD $`p^\nu `$. \# The input consists of a finite matrix A and a prime power $`p`$. GetBasis:=proc(A,p) local ind, i, j, checkdet, k, temp, n, m, h, l, mat, Adj, u,check, mat1, getgcd, ell, hold, ii, V, B, fl0, sum: global R, C, U, mgcdex: with(linalg): \# INITIALIZE \# For checking purposes we will also hold the inverse of A in R. \# C is the matrix whose rows are the required basis. \# The second part of the augmented matrix B will be the \# inverse of C. It is not necessary to do this but it may help. R:=array(1..rowdim(A), 1..coldim(A)): copyinto(A,R,1,1): C:=array(1..rowdim(A),1..rowdim(A)); for i from 1 to rowdim(A) do for j from 1 to coldim(A) do R\[i,j\]:=mods(R\[i,j\],p): od:od: \# ind(ex) holds our column permutation. ind:=array(1..coldim(A)): for i from 1 to coldim(A) do ind\[i\]:=i; od; \# U is a diagonal matrix of units modulo p used to multiply \# rows and make pivots 1. U:=array(1..rowdim(A),1..rowdim(A)); for i from 1 to rowdim(A) do for j from 1 to rowdim(A) do if (i¡¿j) then C\[i,j\]:=0; U\[i,j\]:=0; else C\[i,j\]:=1; U\[i,j\]:=1; fi: od; od; R:=concat(R,C); \# checkdet holds candidates for the next pivot. checkdet:=array(1..coldim(A)): if (rowdim(A)¿coldim(A)) then RETURN(‘More rows than columns cannot form a basis.‘):fi: \# We need a temporary location to compute changes in C to avoid \# nonzero entries divisible by p. getgcd:=array(1..rowdim(A)); \# The index ’i’ will stand for the row currently being worked on. \# END INITIALIZE \# THE ACTUAL COMPUTATION \# The variable $`i`$ will denote the working row. \# Compute the determinant of the block to be used and the pivot \# column by looking for a unit mod $`p`$ to be the next pivot. for i from 1 to rowdim(A) do for m from 1 to coldim(A) do checkdet\[m\]:=R\[i,m\]: for k from 1 to i-1 do checkdet\[m\]:=checkdet\[m\]- R\[i,ind\[k\]\]\R\[k,m\]: od: od: for k from 1 to coldim(A) while (igcd(checkdet\[k\],p) $`<>`$ 1) do : od: if (k $`>`$ coldim(A)) then print(A,R,C):RETURN(‘No pivot. Not a basis mod p.‘): fi: \# If necessary, permute columns by permuting entries of ind. if (ind\[i\]¡¿k) then for n from 1 to coldim(A) while (ind\[n\] $`<>`$ k) do : od: temp:=ind\[i\]:ind\[i\]:=k:ind\[n\]:=temp: fi: \# Multiply the working row by the inverse of the pivot to \# make the pivot 1 mod p and clear below the diagonal. u:= (checkdet\[k\](̂-1) mod p) : U\[i,i\]:= mods(u,p): for n from 1 to coldim(A) do R\[i,n\]:= mods(u\R\[i,n\],p): od: \# Clear below the permuted diagonal. if (i¡¿1) then for h from 1 to i-1 do temp:=mods(R\[i,ind\[h\]\], p); R:=addrow(R,h, i, -temp): C:=addcol(C, i, h, temp) od: fi: \# Clear above the diagonal. if (i $`<>`$ 1) then for h from i-1 to 1 by -1 do temp:=mods(R\[h,ind\[i\]\],p); R:=addrow(R,i,h,-temp): C:=addcol(C,h, i, temp); od: fi: \# We now correct for some (enough) nonzero multiples of p which may \# occur in our candidate C for a lifting. m:=0; ell:=0; fl0:=0; if (mods(C\[k,ind\[i\]\],p)$`<>`$ 0) then fl0:=1; fi; if ((C\[k,ind\[i\]\] $`<>`$ 0)and(igcd(C\[k,ind\[i\]\],p)=p)) then for j from 1 to i-1 do sum:=0; if (i $`<>`$ 1) then for ii from 1 to k-1 do sum:=sum+C\[ii,j\]^2; od; fi; if (sum=0) then ell:=ell+1; getgcd\[ell\]:=ind\[j\];fi; od; if (ell $`<>`$ 0) then V:=array(1..ell); for j from 1 to ell do V\[j\]:=C\[k,getgcd\[j\]\]; od; mgcdex(V,B); if (mods(C\[k,ind\[i\]\],(B\[ell+1\]\p))=0) then m:=1; else fl0:=1; fi; fi; if (m=1) then k:=k-1; temp:=C\[k,ind\[i\]\]/B\[ell+1\]; for j from 1 to ell do C:=addcol(C,getgcd\[j\],i,(-B\[j\]\temp)); R:=addrow(R,i,getgcd\[j\],B\[j\]\temp); od; fi; \# Printouts added to observe progress. mat:=submatrix(C,1..i,1..i); mat1:=submatrix(R,1..i,1..coldim(A)+i); print(‘Row ‘,i,‘ C = ‘,mat,‘ Rowreduction = ‘,mat1): od: # This is the end of the working program. print(‘Orig A = ‘,A,‘ C = ‘,C,‘ U = ‘,U,‘R = ‘,R);end:
warning/0007/astro-ph0007052.html	ar5iv	text	# Astrometry with Hubble Space Telescope (version 1.2) ## 1 Introduction ### 1.1 Astrometry with Hubble Space Telescope In 1990 NASA launched Hubble Space Telescope. In addition to cameras and spectrographs usable from the far ultraviolet to the near-infrared, the observatory contains three white-light interferometers. As part of engineering and science support their primary task was telescope guiding; to position and hold science targets within the science instrument apertures with tolerances approaching 0$`.^{\prime \prime }`$1, equivalent to 100 milliseconds of arc (100 mas). Pointing and tracking requires two such Fine Guidance Sensors (FGS), leaving the third free for astrometry. The design goal for astrometry with an FGS was 3 mas precision over the entire field of regard. When designed and built in the early 1980’s, The FGS represented an order of magnitude improvement over existing ground-based techniques. HST launch delays provided sufficient time for ground-based techniques to equal and surpass this goal. Thus, our post-launch goal was redefined. This goal, 1 mas precision small-field astrometry, has been achieved, but not without significant challenges. An HST FGS will remain a competitive astrometric tool for faint targets in crowded fields and for faint small-separation binaries, until the advent of large-aperture ground-based and longer-baseline space-based interferometers. This article describes the science that can be done with the FGS and presents some recent science results. We outline how these data are modeled, acquired, and calibrated. We next show how the astrometer, FGS 3, works as an interferometer. Finally, we present a guide to the literature that provides additional detail for each item discussed in this article. Astrometry with 1 mas precision is obtainable with the HST Wide Field Planetary Camera. The techniques are similar to those used in ground-based CCD astrometry and should be discussed separately. ## 2 The choice of science targets The choice of HST astrometry science targets is rightly determined by what can be done from the ground. The unique capabilities of HST must remain reserved for projects demanding them. In this section we summarize the strengths of HST astrometry, pointing out a few examples of the kinds of objects for which HST and the FGS are ideally suited. Since the FGS interferometer offers two observational modes, fringe scanning and fringe tracking, we divide the remainder of this article by mode at each step. ### 2.1 Fringe Scanning: Using HST for Binary Star Astrometry Four characteristics of a high-resolution ($`<1\mathrm{"}`$) observing technique must be considered when attempting to resolve a target binary star: the source brightness limit, the resolving capability, the astrometric accuracy, and the magnitude difference between the components that can be observed. Although all four parameters are interdependent, we provide current limits for each characteristic separately in order to illustrate the efficacy of HST FGS3 as an astrometric instrument. For binary stars HST FGS3 observations provide results currently unachievable with ground-based techniques. HST FGS3 is accurate to 1-3 mas at a resolution limit of 15 mas for magnitude differences less than 2 mag. It can effectively observe targets to at least $`V15`$. At component separations greater than 200 mas, it can bridge brightness differences of at least $`\mathrm{\Delta }V=4`$ mag, as demonstrated by the detection of Gliese 623B. Future developments in ground-based interferometers should improve the accuracy of astrometric measurements and provide better resolution, but will still be limited to relatively bright targets and moderate magnitude differences. HST FGS3 is the only high-resolution instrument currently available that (1) can provide high-precision astrometry for very close binaries, (2) allows relatively faint targets to be observed, and (3) can bridge at least moderately large magnitude differences between the components in a binary system. ### 2.2 Fringe Tracking: Using HST to Obtain Parallaxes In the absence of distance, astrophysics lacks precision. Our demonstrated measurement precision for parallaxes is 0.5 mas or better, given six epochs of observation at maximum parallax factor. Often a parallax target is unmeasurable because reference stars are not located within the field of view of the measuring device, especially true for small-format (typically several square arcmin) CCD cameras. An FGS provides a large field of regard (Fig. 1), shaped somewhat like a pickle. This shape comes from the the pick-off mirror (Fig. 5) in the HST focal plane. Parallax observations are generally spaced by six months, during which time, due to HST solar array illumination constraints, the FGS field of view rotates by 180. The lens-shaped region (Fig. 6) in common to the two extreme orientations provides a relatively large field with short axis = 3$`.^{}`$5 and long axis = 14’. FGS 3 has a large dynamic range, able to obtain fringe tracking position measures for stars in the magnitude range $`4V17`$. This large dynamic range is provided by a neutral density filter that reduces the magnitude of bright stars by 5 mag. The unfiltered range is $`8.5V17`$. HST can obtain precise parallaxes for binaries whose components are separated $`0.^{\prime \prime }03<\rho <1.^{\prime \prime }0`$. Nearly simultaneous fringe scanning and tracking measurements provide the component positions relative to reference stars, essential in determining the center of mass of a system, and thus individual component masses. Lastly there is timeliness of result. One no longer need wait 3-6 years for the parallax of astrophysically interesting objects. The distance to a sufficiently interesting and important object can be obtained on the same time scale as other astrophysical information. If an object or class of objects is interesting now, a theory can be tested now. ## 3 Science Results Science targets have included nearby stars used as probes for an extrasolar planet search; stars whose astrophysics would be greatly aided by accurate distances; low-mass binary stars to define the mass-luminosity relation for the lower main sequence; members of the Hyades Cluster, key to the distance scale; extragalactic objects (QSO) to provide an inertial reference frame for HIPPARCOS; and the epitome of crowded fields, an extragalactic cluster, R136 in the Large Magellanic Cloud. Once again we divide by technique, discussing representative results from first fringe scanning, then fringe tracking. #### 3.0.1 Fringe scanning and binary stars Fringe scanning science depends on deconvolution and the assumption that all objects in the field are point sources. For a recent science result we turn to the low-mass binary Wolf 1062 = Gl 748, observed in support of a lower main sequence mass-luminosity project. Fig. 2 presents fringe scans of this binary along the two orthogonal axes of FGS 3. Since each observed fringe is a linear superposition of two fringes (one for each component star in the binary), it is modeled with two identical single star fringes. Their placement along the axis and relative amplitudes provide separation and brightness differences. Relative separation along the two axes provides position angle information. Fig. 3 shows all the measured separations and position angles on the best-fit orbit. The mean absolute difference between the observed and computed separations is only 1.1 mas, and only 0$`.^{}`$77 in the position angle. Residual vectors for all data points are drawn, but are smaller than the points. Perhaps it is even more impressive to realize that the box illustrated is only 0$`.^{\prime \prime }`$6 in size, so that good seeing from the ground would result in a stellar image the size of the entire figure. A combination of fringe scanning and fringe tracking observations for the low-mass binary, L722-22, has yielded a relative parallax $`\pi =0.^{\prime \prime }1656\pm 0.^{\prime \prime }0008`$ and component masses of 0.179 and 0.112 $`M_{}`$ with formal random errors for the mass as low as 1.5%. #### 3.0.2 Fringe tracking: positions in a reference frame Fringe tracking science is primarily that of relating the position of a target to positions of stars defining a reference frame. Here we discuss a Hyades parallax program and a fundamental astrometry project involving HIPPARCOS. Trigonometric parallax observations were obtained for seven Hyades members in six fields of view. These have been analyzed along with their proper motions to determine the distance to the cluster. Formal uncertainties on individual parallaxes average 1 mas. This relatively large error is due both to poor spatial distribution and small number of reference stars in each field. Knowledge of the convergent point and mean proper motion of the Hyades is critical to the derivation of the distance to the center of the cluster. Depending on the choice of the proper-motion system, the derived cluster center distance varies by 9%. Therefore, a full utilization of the HST FGS parallaxes awaits the establishment of an accurate and consistent proper-motion system. Observations of separations of HIPPARCOS stars from extragalactic objects have been made to determine the rotation of the HIPPARCOS instrumental system with respect to the ICRS-VLBI reference frame. A determination from 78 observations yields accuracies on the order of 2 mas rms in the coordinate rotational offsets near the mean epoch of the HST observations and 2 mas yr<sup>-1</sup> in the coordinate rotation rates. The main contributing sources of error are HST measurement errors and proper motion errors introduced by the three year time difference between the mean HIPPARCOS and HST observational epochs. ## 4 Data modeling Our two examples are a binary star for fringe scanning and a parallax field for fringe tracking. ### 4.1 Fringe scanning We consider a binary star as the simplest object requiring fringe scanning. We assume that the fringe produced by a binary star is a linear superposition of the fringes produced by two single stars. The interferometer response to an actual binary star is shown in Fig. 2. If $`F(x)`$ is the fringe produced by a single star on the x-axis, a binary star should be described by $$D(x)=A\times F(x+z_x)+B\times F(x+z_x+S_x)$$ (1) where A and B are the relative intensities of the two components (constraining $`A+B=1`$). The zero point offset, $`z_x`$, and a component separation, $`S_x`$, complete the model, there being a similar expression for the fringe produced along the y-axis. $`S_x`$ and $`S_y`$ yield the binary separation, and, once transformed to equatorial coordinates, the position angle. We obtain the magnitude difference from $$\mathrm{\Delta }m=2.5\times log(A/B)$$ (2) If a binary is observed in fringe tracking mode, the FGS will lock on an erroneous zero-crossing position which is generated by two closely overlapping s-curves. A fringe scanning mode observation gives relative positions of the two components. One can determine the relative positions with an accuracy and precision of 1 mas, once the fringes from the two stars are deconvolved. ### 4.2 Fringe tracking We consider a typical parallax target and associated reference stars as an example of fringe tracking. The primary science target is ideally surrounded by 5 - 10 other stars used as a reference frame relative to which we determine position and motion. Fringe tracking astrometry is generally a two-step process. We first determine the characteristics of the reference frame. One of the epochs of observation is chosen as the constraint plate. From these data we determine the scale and rotation relative to the constraint plate for each observation set within a single orbit. Since for many of our targets the observation sets span over two years, we also include the effects of reference star parallax ($`\pi `$) and proper motion ($`\mu `$). $$\xi =ax+by+c(P_x\pi +\mu _xt)$$ (3) $$\eta =dx+ey+f(P_y\pi +\mu _yt)$$ (4) The orientation relative to the celestial sphere is obtained from ground-based astrometry. Uncertainties in the field orientations are generally $`0.^{}03<ϵ_\theta <0.^{}09`$. We obtain the parallax factors, $`P_x`$ and $`P_y`$ from a JPL Earth orbit predictor. Finally for a rich-enough reference field ($`n>4`$ stars) we constrain $`\mathrm{\Sigma }\mu =0`$ and $`\mathrm{\Sigma }\pi =0`$ for the entire reference frame. The second step consists of applying the plate constants to the measurements of the science target. Plate coefficients $`a,b,\mathrm{},f`$ are applied as constants, while we solve for science target $`\pi `$ and $`\mu `$ in the above equations. ## 5 The astrometer ### 5.1 The design FGS 3 is an interferometer. Interference takes place in a prism that has been sliced in half, had a quarter-wave retarding coating applied, and then reassembled. Fig. 4 shows one of these Koester’s Prisms. Most of the FGS consists of supporting optics used to feed the Koester’s Prisms (Fig 5). In particular the star selectors walk the 5” instantaneous field of view throughout the interferometer field of regard shown in Fig 1. The output of each face (face A and face B) is measured by a PMT. These signals are combined $$S=\frac{AB}{A+B}$$ (5) to form a signal, S, that is zero for waves exactly vertically incident on the Koester’s Prism front face. Tilting the wavefront back and forth (equivalent to pointing the telescope slightly off, then on target, then slightly off to the other side) generates the fringe pattern seen in Fig. 2. A perfect instrument would generate a perfectly symmetric fringe pattern. The significant spherical aberration of the as-built HST primary mirror, in the presence of internal FGS misalignments, produces a signature in the fringe which mimics coma. Coma causes decreased modulation and multiple peaks and valleys in a fringe. FGS 3 produces the least complicated fringes over the widest area within the field of regard. A replacement FGS installed in 1997 contains an articulated fold flat that removes most of the internal misalignments. This FGS (FGS 1r) produces nearly perfect fringes and should at a minimum yield superior binary star results. FGS 3 will remain in servie until FGS 1r is fully calibrated in 2000. ### 5.2 FGS operating modes We discuss strategies for obtaining the highest quality data possible. Our goal, 1 mas precision small-field astrometry, has been achieved, but not without significant challenges. These included a mechanically noisy on-orbit environment, the self-calibration of FGS 3, and significant temporal changes in our instrument. Solutions include a denser set of drift check stars for each science observation, fine-tuning exposure times, overlapping field observations and analyses for calibration, and a continuing series of trend-monitoring observations. The single greatest contributor to data quality is to treat all data aquired in the same orbit as a unit observation, e.g., a ’plate’. ### 5.3 Fringe scanning The target is placed in the center of the pickle and the star selectors are commanded to move the instantaneous field of view across the target star image. This action produces a fringe. In practice 10 - 30 scans are obtained in a reciprocating pattern forwards and back across the star. The final fringe results from a reversal, shift and add process. We have evidence that over the span of an orbit the positions reported by FGS 1 and FGS 2 for the guide stars change. This results in a drift-like motion of FGS 3, the astrometer. Drift can exceed 30 mas over the span of 36 minutes. The x and y drift rates are generally dissimilar and the drift is not of constant rate. However, since a single scan across a science target requires about a minute, the drift per scan is reduced to less than 1 mas. The reciprocating data acquisition strategy is very nearly self-compensating for drift. ### 5.4 Fringe tracking For fringe tracking, onboard electronics locate the zero crossing between the highest positive and lowest negative fringe peak (see Fig. 2). The position of this zero crossing is determined at a 40 Hz rate during an observation time ranging $`10<t<300s`$. The median of $`>2400`$ zero crossings provides a robust position estimate. The star selectors are used to move the instantaneous field of view from one star to the next in the FGS field of regard shown in Fig. 1. Drift is correctable, but imposes additional overhead, reducing the time available within an orbit to observe the science target. An astrometric observation set must contain visits to one or more reference stars, multiple times during each observation sequence. Presuming no motion intrinsic to all reference and science target stars over a span of 40 minutes, one determines drift and corrects the reference frame and target star for this drift. As a result we reduce the error budget contribution from drift to less than 1 mas. In Fig. 1, the science target might be observed three times during an observing session (single orbit) and each reference star twice. ## 6 Astrometric Calibration The calibration of the fringe scanning mode has as its goal the determination of separation and magnitude differences between components of binary stars. As discussed earlier all fringe scanning targets are observed near the pickle center. In contrast fringe tracking takes place over the entire field of regard. Calibration requirements for these two modes of operation are quite different. ### 6.1 Fringe scanning calibration The most fundamental calibration for any astrometric instrument is to determine the scale in units of the celestial sphere. The FGS scale comes from measurements of several calibration binary stars frequently observed by speckle interferometry with 4-m class telescopes. On the basis of these speckle measures, extremely well-determined orbital elements were obtained. These binary star orbits yield, for any date of HST FGS observation, accurate angular separations of the binary, and hence, the scale. Secondary calibration consists of the development of a library of suitable single star fringe templates. This library is required since fringe morphology weakly depends on star color. From a library of templates the required number (one for each component of a binary or multiple star system) are chosen and fit to the binary star fringe using the least-squares algorithm, GaussFit, that allows for errors in both independent and dependent variables. ### 6.2 Fringe tracking calibration The Optical Telescope Assembly (OTA) of the HST is an aplanatic Cassegrain telescope of Ritchey-Chrètien design. Our initial problem included mapping and removing optical distortions whose effect on measured positions exceeded 0$`.^{\prime \prime }`$5. There was no existing star field with catalogued 1 mas precision astrometry, our desired performance goal, to use as a fiducial grid. Our solution was to use FGS3 to calibrate itself. As a result of this activity distortions are reduced to better than 2 mas over much of the FGS3 field of regard (Fig. 1). This model is called the Optical Field Angle Distortion (OFAD) calibration. To describe these distortions we have adopted a pre-launch functional form originally developed by the instrument builder, the Perkin-Elmer Corporation. These distortions can be described (and modeled to the level of one millisecond of arc) by the two dimensional fifth order polynomials: $`x^{}=a_{00}+a_{10}x+a_{01}y+a_{20}x^2+a_{02}y^2+a_{11}xy+a_{30}x(x^2+y^2)+a_{21}x(x^2y^2)`$ $`+a_{12}y(y^2x^2)+a_{03}y(y^2+x^2)+a_{50}x(x^2+y^2)^2+a_{41}y(y^2+x^2)^2`$ $`+a_{32}x(x^4y^4)+a_{23}y(y^4x^4)+a_{14}x(x^2y^2)^2+a_{05}y(y^2x^2)^2`$ $`y^{}=b_{00}+b_{10}x+b_{01}y+b_{20}x^2+b_{02}y^2+b_{11}xy+b_{30}x(x^2+y^2)+b_{21}x(x^2y^2)`$ $`+b_{12}y((y^2x^2)+b_{03}y(y^2+x^2)+b_{50}x(x^2+y^2)^2+b_{41}y(y^2+x^2)^2`$ $`+b_{32}x((x^4y^4)+b_{23}y(y^4x^4)+b_{14}x(x^2y^2)^2+b_{05}y(y^2x^2)^2`$ where x, y are the observed position within the FGS field of view, x’, y’ are the corrected position, and the numerical values of the coefficients $`a_{ij}`$ and $`b_{ij}`$ are determined by calibration. Ray-traces were used for the initial estimation of the OFAD. Gravity release, outgassing of graphite-epoxy structures within the FGS, and post-launch adjustment of the HST secondary mirror required that the final determination of the OFAD coefficients $`a_{ij}`$ and $`b_{ij}`$ be made by an on-orbit calibration. M35 was chosen as the calibration field. Since the ground-based positions of our target calibration stars were known only to 23 milliseconds of arc, the positions of the stars were estimated simultaneously with the distortion parameters. This was accomplished during a marathon calibration, executed on 10 January 1993 in FGS 3. The entire 19 orbit sequence is shown in Fig. 6. GaussFit was used to simultaneously estimate the relative star positions, the pointing and roll of the telescope during each orbit (by quaternions), the magnification of the telescope, the OFAD polynomial coefficients, and four parameters that describe the star selector optics inside the FGS. ### 6.3 Other fringe tracking calibrations Because each FGS contains refractive elements (Bradley et al. 1991), it is possible that the position measured for a star could depend on its intrinsic color. Changes in position would depend on star color, but the direction of shift is expected to be constant, relative to the FGS axes. This lateral color shift would be unimportant, as long as target and reference stars have similar color. However, this is certainly not the case for many of the science targets (e.g., very red stars such as Proxima Centauri and Barnard’s Star). Finally, to provide the large dynamic range, a neutral-density filter can be placed in front of the Koester’s Prism. As a consequence there will be a small but calibratable shift in position (due to filter wedge) when comparing the positions of the bright star to the faint reference frame. The shift is constant in direction (relative to FGS 3) and size, since the filter does not rotate within its holder. ### 6.4 Maintaining the fringe tracking calibration The FGS3 graphite-epoxy optical bench was predicted to outgas for a period of time after the launch of HST. The outgassing was predicted to change the relative positions of optical components on the optical bench. The result of whatever changes were taking place was a change in scale. The amount of scale change was far too large to be due to true magnification changes in the HST optical assembly. Two of the parameters that descibe the star selector optics cause a scale-like change, if they are allowed to vary with time. The solution was to revisit the M35 calibration field periodically to monitor these scale-like changes and other slowly varying non-linearities. Revisits will be required as long as it is desirable to do 1 mas precision astrometry with any FGS. The result of this activity judges the validity of the current OFAD coefficients and provides warning for the need for recalibration. With these data we remove the slowly varying component of the OFAD, so that uncorrected distortions remain below 2 mas for center of FGS 3. The character of these changes is generally monotonic with abrupt jumps in conjunction with HST servicing missions. ## 7 Bibliography Here are pointers to a number of key papers in the development of HST interferometric astrometry, and to some recent results. A detailed description of the FGS interferometer can be found in The flight hardware and ground system for Hubble Space Telescope astrometry. Bradley, A., Abramowicz-Reed, L., Story, D., Benedict, G. & Jefferys, W. 1991, PASP, 103, 317 and in Lupie, O., Nelan, E. P., 1998, FGS Instrument Handbook, version 7, (STScI) Our choice of astrometer (there being three from which to choose) is documented in Astrometric performance characteristics of the Hubble Space Telescope fine guidance sensors. Benedict, G. F., et al. 1992, PASP, 104, 958 Our primary calibration tool is GaussFit, a least-squares and robust estimation algorithm described in Jefferys, W., Fitzpatrick, J., and McArthur, B. 1988 Celest. Mech. 41, 39. Calibration and observing strategies for fringe scanning are discussed in Binary star observations with the Hubble Space Telescope Fine Guidance Sensors. I - ADS 11300. Franz, O. G., et al. 1991, ApJ, 377, L17 and A deconvolution technique for Hubble Space Telescope FGS fringe analysis. Hershey, J. L. 1992, PASP, 104, 592 Introductions to deriving and maintaining the fringe tracking calibration and observing strategies include Astrometry with Hubble Space Telescope Fine Guidance Sensor number 3: Position-mode stability and precision. Benedict, G. F., et al. 1994, PASP, 106, 327, Maintaining the FGS3 Optical Field Angle Distortion Calibration. McArthur, B., Benedict, G. F., Jefferys, W. H. & Nelan, E. 1997, The 1997 HST Calibration Workshop with a new generation of instruments. Editors Stefano Casertano, Robert Jedrzejewski, Charles D. Keyes, and Mark Stevens. Baltimore, MD : Space Telescope Science Institute 1997, p. 472. and Working with a space-based optical interferometer: HST Fine Guidance Sensor 3 small-field astrometry. Benedict, G. F., et al. 1998, Proc. SPIE, 3350, 229 Recent astrometric results naturally divide into fringe scanning and fringe tracking. For fringe scanning, Astrometric Companions Detected at Visible Wavelengths with the Hubble Space Telescope Fine Guidance Sensors. Franz, O. et al. 1994, American Astronomical Society Meeting, 185, 85.24 discusses the first results for large $`\mathrm{\Delta }m`$. The first definitive binary orbit determined with the Hubble Space Telescope Fine Guidance Sensors was for Wolf 1062 = Gliese 748 Franz, O. G., et al. 1998, AJ, 116, 1432 From the beginning, stellar interferometers have explored stellar diameters. The FGS has determined interferometric angular diameters of Mira variables Lattanzi, M. G., Munari, U., Whitelock, P. A. & Feast, M. W. 1997, ApJ, 485, 328 . FGS fringe scanning has contributed to the optical Mass-Luminosity Relation at the lower Main Sequence (.08 to .20 $`M_{}`$), discussed in Henry Henry, T. J. , Franz, O. G., Wasserman, L. H., Benedict, G. F., Shelus, P.J., Ianna, P.A., Kirkpatrick, J. D., & McCarthy, D. W. 1999, ApJ, 512, 864 Recent fringe tracking results include a report on the search for planets around Proxima Centauri Benedict, G. F., McArthur, B., Chappell, D. W., Nelan, E., Jefferys, W. H., Van Altena, W., Lee, J., Cornell, D., Shelus, P. J., Hemenway, P. D., Franz, O. G., Wasserman, L. H., Duncombe, R. L., Story, D., Whipple, A., & Fredrick, L. W. 1999, AJ, 118, 1086, parallaxes of two high-velocity stars Macconnell, D. J., Osborn, W. H. & Miller, R. J. 1997, AJ, 114, 1268, and the distance to the Hyades Cluster based on FGS parallaxes Van Altena, W. F. et al., 1997. ApJ, 486, L123. An FGS program to link the HIPPARCOS reference frame to an extragalactic reference system is detailed in Hemenway, et al 1997, AJ, 114, 2796 Finally a combination of fringe tracking and scanning yielded a parallax and component masses for the low-mass binary L722-22 Hershey, J. L. & Taff, L. G. 1998, AJ, 116, 1440 A snapshot, circa 1994, of the scientific capabilities of the FGS can be found in Astronomical and astrophysical objectives of sub-milliarcsecond optical astrometry: proceedings of the 166th Symposium of the International Astronomical Union held in the Hague, Netherlands, August 15-19, 1994 Edited by E. Hog, P. Kenneth Seidelmann, International Astronomical Union. Symposium no. 166, Kluwer Academic, Publishers, Dordrecht, 1995 Papers include Fringe Interferometry in Space: The Fine Structure of R136A with the Astrometer Guidance Sensor Aboard HST. Lattanzi M. G. et al., pg 95 and Hubble Space Telescope: A Generator of Submilliarcsecond Precision Parallaxes. Benedict, G. F., et al. , pg. 89 George F. (Fritz) Benedict This article relied heavily on inputs from and reviews by B. McArthur, E. Nelan, O. Franz, L. Wasserman, and W. Jefferys.
warning/0007/astro-ph0007016.html	ar5iv	text	# Integral field spectroscopy of the radio galaxy 3C 171 . Based on observations performed at the Canada France Hawaii Telescope ## 1 Introduction A number of radio galaxies and quasars at various redshifts exhibit line emission from ionized gas up to several tens or even hundreds of kiloparcsecs from the nucleus. The existence of such gaseous envelopes is most probably linked to the formation of galaxies and active nuclei. Conversely, the existence of an active nucleus influences the physical conditions of the gas, its kinematics, and star formation in the host galaxy. However, illumination of the gas by the ultraviolet radiation emitted by the active nucleus is probably not the only mechanism responsible for the ionization of the gas (Clark et al. 1998, hereafter C98, Villar-Martín et al. 1999, Tadhunter et al. 2000); the strong link between optical emission line and radio properties suggests interactions between the gas and the radio-emitting plasma, possibly because of shock heating and subsequent ionization of the gas. Radio-loud quasars at low redshift are well adapted to the study of the interactions between the gas and radiation from the active nucleus and/or the radio plasma. This indeed becomes much more difficult for more distant objects due to their smaller spatial extent and surface brightness dimming. Integral field spectroscopy is well suited for this purpose, since, contrary to long-slit spectroscopy, it allows to fully map the velocity field in one exposure, providing the size of the ionized nebulosity is compatible with the instrument field. Such a study was done for example for three quasars with redshifts between 0.268 and 0.370 by Durret et al. (1994), for a 0.734 redshift quasar by Crawford & Vanderriest (1997), for 3C 48 by Chatzichristou et al. (1999), for four ultraluminous IRAS galaxies by Wilman et al. (1999) and for six radio-loud quasars with redshifts between 0.26 and 0.60 by Crawford & Vanderriest (2000). 3C 171 is a radio-galaxy at a redshift z=0.238 (corresponding to a scale of 3.2 kpc/arcsec for H$`{}_{0}{}^{}=75`$ km s<sup>-1</sup> Mpc<sup>-1</sup> and q<sub>0</sub>=0.5). It is associated with an optical emission-line region extending over 6 arcsec on either side of the nucleus (Heckman et al. 1984). These authors found large velocity gradients and striking similarities between the ionized gas and radio distributions. The radio emission itself has two bright hot spots East and West of the nucleus, and fainter emission extending perpendicularly to the radio axis, North and South of these hot spots. The correspondence between the \[OIII\] and radio emissions was confirmed by Baum et al. (1988) and by Blundell (1996). HST spectroscopy by Hutchings et al. (1998) has revealed outward motions of the gas at several hundred km/s in the very inner region. A line-free image of the host galaxy of 3C 171 shows a moderate elongation in the North-South direction (Baum et al. 1988). Using long slit spectra along the radio axis (PA=102) C98 have shown evidence for shocks induced by jet-cloud interactions: close radio/optical association, ionization minima almost coincident with both radio hot spots, high velocity line splitting spatially associated with the two inner hot spots, large line widths in the external gas and anticorrelation between line width and ionization state in the external gas. In this paper we present 2D spectroscopic data on 3C 171, which allow us to spatially analyse line ratios and kinematics and, used together with radio images, to interpret them in terms of an important contribution of shocks in the regions close to radio hot spots, as in the scenario reported by C98. The data are presented in Section 2, the morphology and excitation of the ionized gas are studied in Section 3, the kinematics in Section 4 and the discussion and conclusions in Sections 5 and 6. ## 2 The data 3C 171 was observed on February 7-9, 1995 during a total exposure time of 4 hours and 13 minutes with the TIGER instrument at the 3.60m Canada France Hawaii telescope. Individual exposures were offset from one another in the East-West direction in order to cover completely the object. Airmasses were between 1.2 and 1.3. The grism used was R300, giving a spectral resolution of 6.9 Å FWHM in a spectral interval containing the HeII4686<sup>1</sup><sup>1</sup>1Note that this line is close to the filter edge and therefore the measured fluxes can be highly underestimated., H$`\beta `$ and \[OIII\]4959-5007 lines, corresponding to a velocity resolution of 350 km s<sup>-1</sup> at 6200 Å. The spatial sampling was 0.39 arcsec and each individual field about 7$`\times `$7 arcsec<sup>2</sup>; the seeing was 0.9 arcsec FWHM. The total spatial coverage is about 9$`\times `$4 arcsec<sup>2</sup>. References for a full description of the TIGER instrument and of the data reduction methods for this kind of object can be found in Durret et al. (1994). The first steps of the data reduction (from spectra extraction to wavelength calibration) strongly depend on the instrument optics and were thus performed with the TIGER data processing software. The following stages (cosmic ray removal, photometric calibration, exposure merging, image reconstruction) were achieved using the XOasis data reduction package, dedicated to the CFHT integral field spectrograph OASIS. The algorithms of this new package are indeed more powerful than the previous ones, especially for cosmic removal and exposure merging. The reader can find a complete description of this package in the XOasis user manual at http://www-obs.univ-lyon1.fr/$``$oasis/reduc/reduc\_tiger\_frames.html. We also make use of the HST image (two exposures of 300s taken by the WFPC2 with the red F702W filter retrieved from the archive, realigned and summed) published by Koff et al. (1996). We have retrieved the radio image published by Heckman et al. (1984) from the 3CRR atlas (http://www.jb.man.ac.uk). Optical and radio images have been aligned by assuming that the radio core is placed at the centroid of the brightest nuclear concentration in optical images. ## 3 Morphology and excitation of the ionized gas An image and contours of the \[OIII\] emission reconstructed from our spectra are displayed in Fig. 1. These contours clearly reveal the inhomogeneous structure of the nebulosity, with the ionized gas extending roughly along an East-West direction up to the radio hot spots, i.e. about 5 arcsec West and 4 arcsec East of the nucleus. Strong \[OIII\] emission is observed in the nuclear region, as well as in two blobs on either side of the nucleus at nuclear distances of $`\pm `$2.7 arcsec. The blob 2.7 arcsec West of the nucleus will hereafter be referred to as region W1. In Fig. 1 we also show the radio emission at 1441 MHz superimposed on our \[OIII\] image (total emission, i.e., computed by integrating the emission in the wavelength range covered by the two \[OIII\] lines). It can be seen that the peaks of the radio hot spots are placed at the edges of the elongated structure in \[OIII\], about 2 arcsec further out than the \[OIII\] maxima. A contour plot of the H$`\beta `$ emission reconstructed from our spectra is displayed in Fig. 2. A region of relatively strong H$`\beta `$ emission is observed almost 4 arcsec West of the nucleus, that is notably further out than region W1 observed in \[OIII\] West of the nucleus and therefore much closer to the radio hot spot. We will refer to it as region W2. We have also extracted the continuum image, reconstructed by integrating the whole wavelength range after a polynomial fit to the continuum. It only shows the central spot which is slightly resolved (FWHM $``$ 1.3 arcsec). We fit ellipses to the contours and obtained the resulting surface brightness profile shown in Fig. 3. A r<sup>1/4</sup> law fitting results in a scalelength (effective radius) of 4.8 arcsec, i. e., 15 kpc, a typical value for an elliptical galaxy. The HST image is displayed in Fig. 4 with \[OIII\] contours from our image superimposed. The HST image was also convolved by the PSF estimated from our direct image to create a “smoothed HST image” which is identical to our \[OIII\] image, with a bright central region and two blobs coinciding with those in our \[OIII\] image. The similarities are easily understood by taking into account the fact that the filter used for HST observations includes the emission line contributions from \[OIII\] and H$`\alpha `$. The emission observed in the HST image is more extended to the East, with an almost linear feature well aligned with the radio axis and a brightness enhancement placed at the same position as the eastern radio hot spot. A filamentary-like feature is also observed to the West of the nucleus, corresponding to our West blob but this time not exactly aligned with the radio emission. As shown in Section 4, the spectra obtained in some regions clearly indicate line splitting, with at least two components showing different excitations and kinematical behaviours: one about 3 arcsec to the West and another one 2.5 arcsec to the East. Two and three components are necessary to fit the spectra in the former and latter regions respectively; however, the signal to noise ratio (S/N) is not high enough to allow the individual fitting of each spectrum. These two regions are also marked in Fig. 1 with crosses. We have reconstructed two emission line images in order to separate each of the blue and red components of \[OIII\] and H$`\beta `$ in the region close to W1. They are shown in Figs. 5 and 6. We stress that the peak at 2.5 arcsec to the East in Fig. 5 is located in the region where three components are present. While the morphology of the blue component is quite similar for \[OIII\] and H$`\beta `$, the red components are notably different, with emission peaks shifted by about 2 arcsec, corresponding to regions W1 and W2, respectively. We have computed \[OIII\]/H$`\beta `$ ratios for the blue and red components, rejecting the spectra where the S/N for H$`\beta `$ is smaller than 1.5. The resulting images are displayed in Fig. 7. It can be noticed that, for the blue component, the regions of highest excitation are located almost symmetrically with respect to the nucleus at $``$ 1 arcsec in a direction almost perpendicular (PA=124) to that of the elongation of the central isophotes in Fig. 5 (PA=64). There are two local maxima, one with \[OIII\]/H$`\beta `$ = 5 corresponding to the 3$``$component region, and another one with \[OIII\]/H$`\beta `$ = 6 placed at ($`2,1`$) arcsec. For the red component, the highest excitation is reached in the region closest to the center of the radio galaxy, with decreasing values towards the outskirts. In order to have a high enough S/N ratio for H$`\beta `$ to be able to fit \[OIII\] and H$`\beta `$ without constraining their redshifts to be the same, we have filtered our data-cubes with a Gaussian of FWHM 1.2 arcsec. This allowed us to get more spatially extended information on \[OIII\]/H$`\beta `$. We then re-analysed the resulting spectra, now considering the spectra with S/N for H$`\beta `$ greater than 5. The results are shown in Fig. 8. As expected, the details have been smoothed out, and now the blue component shows a plateau in the central region, a local peak close to the three-component region and a general trend for decreasing ratios to the East. The maximum to the West (at (1.5,1) arcsec) is due to the contribution of the red component in this region (Fig. 8, top). For the red component (Fig. 8, bottom), the behaviour is essentially the same as for the non-convolved data, but extending somewhat further. ## 4 Kinematics of the ionized gas As already mentioned in section 3, some of the spectra show clear signs of line splitting. These are more evident in the region about 3 arcsec to the West, where we could fit the emission lines with two components. In Fig. 9 we show the resulting fit for the region 2.7 arcsec West of the nucleus covering 3 arcsec<sup>2</sup>. For the regions about 2.5 arcsec East of the nucleus, no fit is possible for each individual lens, but we have extracted a spectrum over about 1.4 arcsec<sup>2</sup> showing that at least 3 components are present (see Fig. 10). The spectrum of the nuclear region is shown in Fig. 11. At difference with Hutchings et al. (1998) we actually detect the MgI stellar absorption line feature, from which we determine the systemic velocity to be 63126$`\pm `$334 km s<sup>-1</sup> (the error bar is here the FWHM of the instrumental spectral line broadening, and is most probably overestimated). This value is quite close to the one derived from the \[OIII\] lines in the nuclear spectrum: 63069 km s<sup>-1</sup>. We will use this value hereafter when calculating velocities relative to the nucleus. We note that only one line of the MgI triplet is detected. This could be due to the presence of emission lines such as FeII which can be strong enough to fill the absorption lines (see e.g. Boroson & Green 1992). The low S/N ($``$ 6) precludes any accurate measurement of the velocity dispersion and equivalent width of the MgI line. Note that this absorption feature is also present in the nuclear low dispersion spectra of C98 as a faint dip just left of the \[NI\]5199 line. The velocity distributions have been obtained for the convolved data, for which the S/N is high enough to allow to fit \[OIII\] and H$`\beta `$ without constraining their redshifts to be the same. Velocities and FWHM distributions for the blue and red components for the two lines are shown in Figs. 12, 13 and 14. The kinematics of the blue component, for both H$`\beta `$ and \[OIII\] show a central region of about 1 arcsec with rather well organized motions and almost constant FHWM $`1315`$ Å (533 to 644 km/s once corrected from instrumental broadening). They resemble typical rotation, with a kinematical position angle (PA=60) in agreement with that of the elongations seen in the corresponding emission line maps (PA=64, but see below for other possible interpretations). Kinematics are more complicated for the region to the East, with velocities positive to the North and negative to the South. We note that the FWHM reaches two local maxima of about 25 Å (1160 km/s) corresponding to the 3-component region and to the position closest to the eastern radio hot spot. From 2 arcsec outwards to the East, velocities start reaching negative values with FWHM somewhat higher that in the central region. Note that velocities and FWHM maxima are higher for H$`\beta `$ than for \[OIII\]. The red component reaches positive velocities of about 600 km/s, but with different trends for the two lines. The peak velocity (575 km/s) and the minimum FWHM (13 Å) for H$`\beta `$ coincide with W1; velocities then decrease almost along PA=0 on either side of the nucleus and FWHM reach their maxima for region W2 and to the South of W2 (18 Å, i.e., 800 km/s). \[OIII\] velocities are the greatest (625-600 km/s) for the central 1 arcsec of the red component, with an almost flat distribution around W1 and a sharp decrease to the South-West over W2. The FWHM of \[OIII\] slightly increases from East to West, with maximum values (19 Å or 856 km/s) around W2. ## 5 Discussion As reported by Hutchings et al. (1998), the inner optical jet (inside $``$ 1 arcsec) of 3C 171 as traced by the HST image is not aligned with the radio axis. However, at larger scales the association between the optical emission lines and radio morphologies suggests that the processes producing both types of emissions are closely related. As confirmed by recent hydrodynamical simulations by Higgins et al. (1999), a collision between an extragalactic jet and a dense intergalactic cloud can lead to structures comparable to those observed in 3C171. We note that our data do not cover the tail towards the North reported by Tadhunter et al. (2000) in an H$`\alpha `$ emission line image. The close association between radio and emission line morphologies together with the ionization minima coincident with both radio hot spots, the high velocity line-splitting displaced by 2 arcsec behind the hot spots, the FWHM of about 1300 km/s and the anticorrelation between line width and ionization state, have been considered as evidence for shocks induced by jet-cloud interactions (C98). C98 derive these conclusions from high (about 2 Å) and low (about 8.5 Å) resolution long slit spectra along the radio axis (PA=122), with 1.3 arcsec spatial resolution. We use 2D spectroscopic information with better spatial resolution (0.9 arcsec) but with spectral resolution of 6.9 Å. We have extracted the kinematical and emission line values for a cut along PA=122 in order to compare with their low resolution data. A general agreement is obtained. The only noticeable difference is that in our data the \[OIII\] /H$`\beta `$ ratio reaches a local minimum at the center, followed by two maxima on either side at about 1 arcsec, probably due to a better spatial sampling (compare Figs. 7 and 8). With respect to the presence of various emission line components, we have only fit the regions where a single component fit was not satisfactory (see section 4). We have only used two components for each of the \[OIII\] 4959, \[OIII\] 5007 and H$`\beta `$ lines. C98 fit five components to their high resolution data, so a direct comparison is not possible. From the two component fitting we have applied to the region to the West, we may clearly separate two regions with different kinematics and ionization states, the red one clearly detached from the blue one, which is associated with the central region. Note that the red component corresponds to the linear feature on the HST image which is not well aligned with the radio axis. The central 2 arcsec show a velocity field that could be reminiscent of (asymmetric) rotation with an amplitude of 50 km/s. The kinematical position angle (PA=64) coincides with that of the central 1 arcsec elongated structure present in the HST image. \[OIII\] and H$`\beta `$ trace essentially the same kinematics. Similar rotations have been found in some radio-loud quasars (Chatzichristou et al. 1999; Crawford & Vanderriest 2000). To reproduce the gas rotation we have constructed a kinematical model representing a symmetric, inclined disk rotation as a solid body out to r=1.1 arcsec and then a flat rotation curve with V<sub>max</sub>=140 km/s (from r=1.1 to 1.4 arcsec). A non constant value of the major axis produces a better representation of the velocity field; we have fixed the kinematical major axis at PA=47° inside r=0.5 arcsec, then increasing linearly outwards to PA=67. The adopted inclination is $`i`$=82°, derived from the ellipticity of the central body. In Fig. 15 we show the resulting isovelocity contours. In spite of the simplicity of the kinematical model, the observed velocities are reasonably well reproduced in this region. \[OIII\] /H$`\beta `$ ratios in this central region show two maxima with \[OIII\] /H$`\beta `$ $``$ 10 at about 1 arcsec (3.2 kpc), at the ends of the kinematical minor axis. This result is due to the fact H$`\beta `$ contours are more elongated than \[OIII\] contours. Based on a long slit spectrum along PA=60 with the STIS on HST, Hutchings et al. (1998) have suggested that a central outflow in this direction could also explain the observed kinematics, consistently with the clumpy emission-line structures indicating outward motions of a few $``$ 100 km/s within a centrally illuminated and ionized biconical region. For the sake of comparison we have extracted a cut along this PA; we obtain a small amplitude (about 80 km/s), smooth velocity distribution, where the local peaks detected by Hutchings et al. at about $``$0.3, 0.1 and 0.5 arcsec, are absent. Their better spatial resolution ($``$ 0.1 arcsec) could explain the differences. However, such a model requires a mechanism for bending the jet from PA=60 to $``$100 over a distance of about 1 arcsec. Notice that the spatial resolution of the radio map is not sufficient to confirm the presence of such a central jet. The red component 3 arcsec to the West is detached by about 600 km/s from the blue one, and this cannot be reconciled with a gravitational origin. \[OIII\] and H$`\beta `$ kinematics are slightly different, with differences up to about 200 km/s for the two lines. We note that, with the exception of the nuclear region, we do not achieve high enough S/N to extract HeII 4686 fluxes that would allow to analyze HeII/H$`\beta `$ versus \[OIII\] /H$`\beta `$ in the context of the models of shock $`+`$ precursor of Dopita & Sutherland (1995), in the same way as Feinstein et al. (1999) for the radio galaxy 3C 299. The emission line FWHM reach values of about 860 km/s in the regions where \[OIII\] /H$`\beta `$ ratios are minimal, in agreement with the hypothesis of jet-cloud interaction processes, as decribed by C98 (see above). Such interactions are expected to be more important in rich environments, which seems to be the case at higher redshifts; in fact, 3C 171 has been proposed by C98 as an intermediate redshift prototype of high redshift radio galaxies. However, Baum et al. (1988) reported that 3C 171 is a very isolated object, with the closest possible companion at 200 kpc in projection. Moreover, from ROSAT X-ray images, McNamara et al. (1994) concluded that 3C 171 is not associated with a rich cluster of galaxies (whereas they did not exclude its association with a poor cluster or group). This also seems to be the case for the intermediate redshift FRII radio galaxy 3C 299; it also shows jet-cloud interaction producing shocks that ionize the gas and produce the radio optical alignement effect (Feinstein et al. 1999), and is also reported to reside in a non-cluster environment (Wan & Daly 1996; Zirbel 1997). ## 6 Conclusions We have mapped for the first time the extended ionized gas around 3C 171 in the \[OIII\] and H$`\beta `$ emission lines and derived the kinematical and physical properties. We have found that the properties of the central region can be interpreted in terms of those of a typical ENLR disk of radius 1 arcsec (3.2 kpc) following a low amplitude rotation. The continuum surface brightness profile follows an $`r^{1/4}`$ law, suggesting that the underlying galaxy is an elliptical with an effective radius of 15 kpc. The kinematics are much more complicated when approaching the radio hot spots, with clear line splitting. Two components can be fit in the West region, corresponding to an extension of the central region and to a detached blob at about 600 km/s. Line ratios and FWHM are compatible with the jet-cloud interaction scenario proposed by C98. 3C 171 is quite an isolated object, at most belonging to a poor cluster or group, with properties resembling those of high redshift radio galaxies. Such high redshift radio galaxies reside in much richer environments, which are invoked to explain the origin of such jet-cloud interaction. In intermediate redshift radio galaxies showing evidence for shocks produced by jet-cloud interactions as the origin of optical-radio alignements, the mechanisms are more likely to be related with the ambient gas, since a number of these objects neither have nearby companions nor rich cluster environments. ###### Acknowledgements. We acknowledge discussions with E. Emsellem, P. Ferruit, J. Masegosa and M. Villar-Martín. We are also very grateful to J. Perea for his help setting up his SIPL graphics package at the IAP. This work is financed by DGICyT grants PB93-0139 and PB96-0921. Financial support to develop the present investigation has been obtained through the Junta de Andalucía.
warning/0007/hep-th0007003.html	ar5iv	text	# 1 Introduction ## 1 Introduction The possible changes of space-time properties at small (Plank) scale now extensively discussed . Due to the absence of any experimental information it seems instructive to look for some directions exploring attentively the standard Quantum Physics space-time structure. Some years ago Aharonov and Kaufherr have shown that in nonrelativistic Quantum Mechanics (QM) the correct definition of physical reference frame (RF) must differ from commonly accepted one, which in fact was transferred copiously from Classical Physics . The main reason is that to perform exact quantum description one should account the quantum properties not only of studied object, but also RF, despite the possible practical smallness. The most simple of this RF properties is the existence of Schroedinger wave packet of free macroscopic object with which RF is usually associated . Then it introduces additional uncertainty into the measurement of object space coordinate in this RF. Furthermore this effect results in the states transformations between two such RFs which includes quantum corrections to the standard Galilean group transfromations . Algebraic and group theorettical structure of this transformation was studied in . In their work Aharonov and Kaufherr formulated Quantum Equivalence Principle (QEP) in nonrelativistic QM - all the laws of Physics are invariant under transformations between both classic and this finite mass RFs which called quantum RFs. The importance of RF quantum properties account was shown already in Quantum Gravity and Cosmology studies and will be considered here in connection with the time problem in quantum gravity. In this paper the consistent relativistic covariant theory of quantum RFs formulated, our preliminary results were published . In this theory no new ad hoc hypothesis introduced; all calculations are performed in the standard QM formalism. It will be shown that the transformation of the particle state between two quantum RFs obeys to relativistic invariance principles, but differs from standard Poincare Group transformations, due to quantum relativistic correction for RF motion. Solving the evolution equation for quantum clocks models the proper time in moving quantum RF calculated and the related effects of RF momentum quantum fluctuations revealed. This clocks model applied for the analysis of the space-time structure of canonical quantum gravity . In chap.2 canonical formalism for quantum RFs described. In chap.3 we study quantum clocks models and obtain relativistic proper time for quantum RF. In chap.4 the relativistic evolution equations and unitary transformations for quantum RFs described. We’ll consider also RF quantum motion in gravitation field where gravitational ’red shift’ results in additional clocks time fluctuations. ## 2 Quantum Coordinates Transformations For the beginning we’ll consider Quantum Measurements problems related to Quantum RFs model. In QM framework the system defined as RF should be able to measure the observables of studied quantum states and so include the measuring device - detector D. As the realistic example we can regard the photoemulsion plate or the diamond crystal which can measure the particle position and simultaneously record it. Despite the multiple proposals up to now the established theory of collapse doesn’t exist . Yet our problem premises doesn’t connected directly with any state vector collapse mechanism and and it’s enough to detailize standard QM collapse postulate of von Neumann measurement theory . We consider RF which consists of finite number of atoms (usually rigidly connected) and have the finite mass. It’s well known that the solution of Schroedinger equation for any free quantum system can be factorized as : $$\mathrm{\Psi }(t)=c_l\mathrm{\Phi }_l^c(\stackrel{}{R}_c,t)\varphi _l(u_k,t)$$ (1) where center of mass coordinate $`\stackrel{}{R}_c=m_i\stackrel{}{r}_i/M`$, $`c_l`$ are the partial amplitudes. $`u_k`$ describes the internal degrees of freedom, which for potential forces are reduced to $`\stackrel{}{r}_{i,j}=\stackrel{}{r}_i\stackrel{}{r}_j`$ . Here $`\mathrm{\Phi }_l^c`$ describes the c.m. motion of the system. It means that the evolution of the system is separated into the external evolution of pointlike particle M and the internal evolution defined by $`\varphi _l(u_k,t)`$. So the internal evolution is independent of whether the system is localized in some ’absolute’ reference frame (ARF) or not. Quantum Field Theory evidences that the factorization of c.m. and relative motion holds true even for nonpotential forces and variable $`N`$ in the secondarily quantized systems . Moreover this factorization expected to be correct for nonrelativistic systems where binding energy is much less then its mass $`M`$, which is characteristic for the real detectors and clocks. Consequently it’s reasonable to assume that this factorization fulfilled also for the detector states despite we don’t know their exact structure. For our problem it’s enough to assume that eq.(1) holds for RF state only in the time interval $`T`$ from RF preparation moment $`t_0`$ until the act of measurement starts , i.e. the measured particle $`n`$ wave packet $`\psi _n`$ impacts with D. If this factorization holds the space coordinate measured in this RF depends not only on $`\psi _n`$ but also on $`\mathrm{\Phi }_l^c`$ which permit in principle to study quantum RF effects. In this case the possible factorization violation at later time when the particle state collapse occured is unimportant for us. We regard in our model that all measurements are performed on the quantum pairs ensemble of particles $`G^2`$ and $`F^1`$. It means that each event is resulted from the interaction between the ’fresh’ RF and particle ,prepared both in the specified quantum states, alike the particle alone in the standard experiment. To illustrate the meaning of Quantum RF consider gedankenexperiment where in ARF the wave packet of RF $`F^1`$ described by $`\psi _1=\eta _1(x)\xi _1(y)\zeta _1(z)`$ at time moment $`t_0`$. The test particle $`n`$ with mass $`m_n`$ belongs to narrow beam which average velocity is orthogonal to $`x`$ axe and its wave function at $`t_0`$ is $`\psi _n=\eta _n(x)\xi _n(y)\zeta _n(z)`$. Before they start to interact this system wave function is the product of $`F^1`$ and $`n`$ packets. We want to find $`n`$ wave function for the observer in $`F^1`$ rest frame. In general it can be done by means of the canonical transformations described below, but if $`n`$ beam state is localized so that $`\psi _n`$ can be approximated by delta-function $`\delta (xx_b)\delta (yy_b)\delta (zz_b)`$ then $`n`$ wave function in $`F^1`$ can be easily calculated $`\psi _n^{}(\stackrel{}{r}_n^{})=\eta _1(x_n^{}x_b)\xi _1(y_n^{}y_b)\zeta _1(z_n^{}z_b)`$. It shows that if for example $`F^1`$ wave packet along $`x`$ axe have average width $`\sigma _x`$ then from the ’point of view’ of observer in $`F^1`$ each object localized in ARF acquires wave packet of the same width $`\sigma _x`$ in $`F^1`$ and any measurement in $`F^1`$ and ARF will confirm this conclusion. The generalized Jacoby canonical formalism will be applied in our model alternatively to Quantum Potentials used in . Consider the system $`S_N`$ of $`N`$ objects $`W^k`$ which include $`N_f`$ frames $`F^i`$ which have also some internal degrees of freedom and $`N_g=NN_f`$ pointlike ’particles’ $`G^i`$. At this stage we can regard both of them as equivalent objects in the relation to their c.m. motion. We’ll assume for the beginning that particles and RFs canonical operators $`\stackrel{}{r}_i,\stackrel{}{p}_i`$ are defined in absolute (classical) ARF - $`F^0`$ having very large mass $`m_0`$, but later this assumption can be abandoned. We’ll start with Jacoby canonical coordinates $`\stackrel{}{u}_j^l`$ associated with $`F^l`$ rest frame, which for $`l=1`$ equal : $`\stackrel{}{u}_i^1={\displaystyle \frac{\underset{j=i+1}{\overset{N}{}}m_j\stackrel{}{r}_j}{M_{i+1}}}\stackrel{}{r}_i^l;1i<N;\stackrel{}{u}_N=\stackrel{}{u}_s=\stackrel{}{R}_c`$ (2) where $`M_i=\underset{j=i}{\overset{N}{}}m_j`$. $`\stackrel{}{u}_i^l`$ can be obtained and is the linear combination of $`\stackrel{}{u}_i^1`$. Conjugated to $`\stackrel{}{u}_i^l`$ canonical momentums are : $$\stackrel{}{\pi }_i^1=\mu _i(\frac{\stackrel{}{p}_{i+1}^s}{M_{i+1}}\frac{\stackrel{}{p}_i}{m_i}),\stackrel{}{\pi }_N=\stackrel{}{p}_s=\stackrel{}{p}_1^s$$ (3) where $`\stackrel{}{p}_i^s=\underset{j=i}{\overset{N}{}}\stackrel{}{p}_j`$ ,and reduced mass $`\mu _i^1=M_{i+1}^1+m_i^1`$ . The relative coordinates $`\stackrel{}{r}_j\stackrel{}{r}_1`$ can be represented as the linear sum of $`\stackrel{}{u}_i^1`$. They don’t constitute canonical set due to the quantum motion of $`F^1`$ . The Hamiltonian of $`S_N`$ motion in ARF is expressed also via momentums $`\stackrel{}{\pi }_i^1`$ : $$\widehat{H}=\underset{i=1}{\overset{N}{}}\frac{\stackrel{}{p}_i^2}{2m_i}=\frac{\stackrel{}{p}_s^2}{2M}+\underset{j=1}{\overset{N1}{}}\frac{(\stackrel{}{\pi }_j^1)^2}{2\mu _j}=\widehat{H}_s+\widehat{H}_c$$ (4) In $`F^1`$ rest frame the true observables are $`\stackrel{}{\pi }_i^1,\stackrel{}{u}_i^1`$ and it’s impossible to measure $`S_N`$ observables $`\stackrel{}{p}_s`$ and $`\stackrel{}{R}_c`$. The true Hamiltonian of $`S_N`$ in $`F^1`$ should depend on the true observables only , so we can regard $`\widehat{H}_c`$ as the real candidate for its role. It results into modified Schroedinger equation which depends not only of particles masses ,but on observer mass $`m_1`$ also. Now we’ll regard here the alternatve form of this formalism which use Jacoby frame condition (JFC) and is more convenient for the relativistic problem. For the described system $`S_N`$ Langrangian in ARF $`L=\frac{m_i\dot{\stackrel{}{r}}_i^2}{2}`$ gives $`H`$ of (4) after Legandre transform. If one wish to include ARF motion in this formalism the simplest way is to define formally $`L^{}=L+\frac{m_0\dot{\stackrel{}{r}}_0}{2}`$. It gives $`N+1`$ canonical momentums : $`\stackrel{}{p}_j=\frac{L^{}}{\dot{\stackrel{}{r}}_j}`$. The new Langrangian $`L^{}`$ is formally symmetric relative to the frame choice and it gives the Hamiltonian $`H^{}=H+\frac{\stackrel{}{p}_0^2}{2m_0}`$ for $`H`$ of (4). Due to it to anchor this momentums and $`H^{}`$ to $`F^i`$ rest frame in which they acquire some values one must broke $`L^{}`$ symmetry introducing the frame condition (FC) or kinematical (holonomial) constraint . For ARF rest frame we choose FC $`\stackrel{}{p}_0^{\mathrm{\hspace{0.17em}2}n}0`$, where from the formal reasons $`n=2`$. It means that $`\dot{\stackrel{}{r}}_0=0`$ \- RF is at rest relative to itself which seems quite natural, yet it differs from FC used in . All Classical and QM results are reproduced in this scheme if ARF mass is taken infinite. $`S_N`$ quantization in $`F^1`$ performed with Hamitonian $`\widehat{H^{}}`$ and FC regarded as the operator which obeys to Dirack rules for the first order constraints . Galilean-like passive transformations from ARF to $`F^1`$ and back can be found introducing FC also for $`F^1`$ $`\stackrel{}{p}_{11}^{\mathrm{\hspace{0.17em}2}n}0`$, where $`\stackrel{}{p}_{1i}`$ are the canonical momentums in $`F^1`$. $`S_{N+1}`$ unitary transformation from ARF to $`F^1`$ is convenient to write via the $`F^{0,1}`$ total momentum $`\stackrel{}{p}_f=\stackrel{}{p}_0+\stackrel{}{p}_1`$ and $`F^0,F^1`$ relative momentum $`\stackrel{}{\pi }_f`$ conserving other momentums $`\stackrel{}{p}_i`$. Their conjugated coordinates $`\stackrel{}{r}_f,\stackrel{}{u}_f`$ have the standard form of (2). In this notations the transformation from $`F^0`$ to $`F^1`$ is equal to : $$U_{1,0}=P_fe^{ia_f\stackrel{}{p}_f\stackrel{}{r}_f}e^{i\stackrel{}{p}_f\stackrel{}{b}_s}\underset{i=2}{\overset{N}{}}e^{im_i\stackrel{}{r}_i\stackrel{}{\beta }}$$ (5) where $`a_f=\mathrm{ln}\frac{m_0}{m_1},\stackrel{}{b}_s=\frac{M}{m_0}\stackrel{}{R}_c`$. $`P_f`$ is $`\stackrel{}{r}_f`$ reflection (parity) operator. $`\stackrel{}{\beta }=\frac{\stackrel{}{p}_f}{m_1}`$ is the operator corresponding to the velocity parameter in Galilean transformation. Under this transformation $`\stackrel{}{p}_f`$ transformed to $`\stackrel{}{p}_{1f}=\stackrel{}{p}_{10}+\stackrel{}{p}_{11}`$ and $`\stackrel{}{\pi }_{1f}=\stackrel{}{\pi }_f`$. Alike the transformation from $`\stackrel{}{p}_j`$ to $`\stackrel{}{\pi }_i^1`$ obtained operator $`U_{1,0}`$ includes the dilatation transformation . For $`N=2`$ one obtains $`F^1`$ momentums and coordinates : $`\stackrel{}{p}_{10}=(1{\displaystyle \frac{m_0}{m_1}})\stackrel{}{p}_0{\displaystyle \frac{m_0}{m_1}}\stackrel{}{p}_1;\stackrel{}{r}_{10}={\displaystyle \frac{m_1\stackrel{}{r}_1+m_2\stackrel{}{r}_2}{m_0}}+\stackrel{}{b}_s`$ $`\stackrel{}{p}_{11}=\stackrel{}{p}_0;\stackrel{}{r}_{11}=\stackrel{}{r}_0+(1{\displaystyle \frac{m_1}{m_0}})\stackrel{}{r}_1{\displaystyle \frac{m_2}{m_0}}\stackrel{}{r}_2+\stackrel{}{b}_s`$ (6) $`\stackrel{}{p}_{12}={\displaystyle \frac{m_2}{m_1}}(\stackrel{}{p}_0+\stackrel{}{p}_1)+\stackrel{}{p}_2;\stackrel{}{r}_{12}=\stackrel{}{r}_2`$ Results for $`N>2`$ can be easily deduced from this formulaes. It’s easy to see that ARF FC transformed into $`F^1`$ FC. All $`\stackrel{}{\pi }_i^1`$ are conserved and space shift on $`\stackrel{}{b}_s`$ conserves all the distances $`\stackrel{}{r}_i\stackrel{}{r}_j`$. In the limit where heavy $`F^1`$ moves nearly classically $`U_{01}`$ becomes the Galilean momentum transformation with the velocity $`\stackrel{}{\beta }`$. $`S_{N+1}`$ Hamitonian in $`F^1`$ also can be rewritten via new relative momentums $`\stackrel{}{\pi }_{1j}`$ which can be easily derived following (3) : $$\widehat{H}^1=\underset{i=0}{\overset{N}{}}\frac{\stackrel{}{p}_{1i}^2}{2m_i}=H_s^1+H_c^1=\frac{\stackrel{}{p}_{1s}^2}{2M_{N+1}}+\underset{j=2}{\overset{N+1}{}}\frac{\stackrel{}{\pi }_{1j}^2}{2\mu _{1j}}$$ (7) The term $`\widehat{H}_s^1`$ describes $`S_N`$ c.m. motion relative to $`F^1`$ which doesn’t influence on the evolution of $`S_{N+1}`$ true observables $`\stackrel{}{\pi }_{1i},\stackrel{}{u}_{1i}`$ or $`\stackrel{}{r}_i\stackrel{}{r}_j`$. $`\stackrel{}{r}_{1i},\stackrel{}{p}_{1i}`$ aren’t $`S_{N+1}`$ observables for $`F^1`$ observer, yet $`\stackrel{}{p}_{1i}`$ expectation values can be found from $`\stackrel{}{\pi }_i^1`$ measurements. Now we have quantum system $`S_{N+1}`$ which include ARF and in ARF rest frame we can ascribe to it without any contradictions with QM the state vector which for $`N=2`$ is equal : $`\psi _s(\stackrel{}{p}_0,\stackrel{}{p}_1,\stackrel{}{p}_2)=\phi (\stackrel{}{p}_1,\stackrel{}{p}_2)\|\stackrel{}{p}_0=0\|\stackrel{}{p}_1\|\stackrel{}{p}_2`$. After $`U_{1,0}`$ transformation it acquires the similar form in $`F^1`$ rest frame with $`\|\stackrel{}{p}_{11}=0`$. As the result of this transform we obtain the new canonical coordinates referred to finite mass $`F^1`$ rest frame. They permit to factorize internal $`S_N`$ motion and ARF motion and dropping ARF term in $`H^1`$ of (7) we obtain $`S_N`$ Hamitonian. Remind that active transformation shifts $`G^2`$ state $`\psi _2`$ on the distance $`\stackrel{}{a}`$ and velocity $`\stackrel{}{\beta }`$ relative to RF. Passive $`G^2`$ transformation means the transition from one RF to another, but for quantum RF with state $`\psi _s`$ it can’t be described by any state shift on $`\stackrel{}{a},\stackrel{}{\beta }`$ and have more complicated form. $`U_{1,0}`$ is such passive transformation and active $`G^2`$ transformation is the standard Galilean one even in $`F^1`$ . In general the quantum transformations in 2 or 3 dimensions should also take into account the possible rotation of quantum RF axes relative to ARF, which introduce additional angular uncertainty into objects coordinates. Thus after performing coordinate transformation $`\widehat{U}_{A,1}`$ from ARF to $`F^1`$ c.m. we must rotate all the objects (including ARF) around it on the uncertain polar and azimuthal angles $`,\varphi _1,\theta _1`$ which are $`F^1`$ internal degrees of freedom. We can imagine $`F^1`$ axes as some solid rods which orientation this angles describe. As the result the complete transformation is: $`\widehat{U}_{A,1}^T=\widehat{U}_{A,1}^R\widehat{U}_{A,1}`$. Such rotation transformation operator commutes with $`\widehat{H}_c`$ and due to it can’t change the evolution of the transformed states . ## 3 Quantum Clocks Models To construct the relativistic covariant formalism of quantum RFs it’s necessary first to define the time in such RFs. In nonrelativistic mechanics time $`t`$ is universal and is independent of observer, while in relativistic case each observer in principle has its own proper time $`\tau `$. We don’t know yet the nature of time , but phenomenologically it can be associated with the clock hands motion or some other relative motion of the system parts . In Special Relativity the time in moving frame $`F^1`$ can be defined by external observer at rest measuring the state of $`F^1`$ comoving clocks. We’ll consider the same procedure in relativistic QM i.e. some clock observable being measured at some time from the rest frame gives the estimate of proper time of moving quantum RF $`F^1`$. For some clocks models $`F^1`$ internal evolution which define $`F^1`$ clocks motion and consequently its proper time $`\tau _1`$ can be factorized from $`F^1`$ c.m. motion. Its quantum c.m. motion described by the relativistic Schrodinger equation for massive boson. This is Klein-Gordon square root (KGR) equation in which only positive root will be regarded for initial positive energy state . Solving Dirack constraints it was shown recently that this first order equation is completely equivalent to free Field secondary quantization . For our relativistic model we should regard more strictly the features of reference frames and clocks, taking into account the internal motion. Consider the evolution of some system $`F^1`$ where the internal interactions described by the Hamiltonian $`\widehat{H}_c`$ are nonrelativistic , which as was discussed in chap.1 is a reasonable approximation for the measuring devices or clocks. We’ll use the parameter $`\alpha _I=\frac{\overline{H}_c}{m_1}`$ ,where $`m_1`$ is $`F^1`$ constituents total rest mass. In $`F^1`$ c.m. $`\alpha _I=m_1^1\phi _c\|\widehat{H}_c\|\phi _c`$ where $`\phi _c`$ is $`F^1`$ internal state of (1). It describes the relative strength of the internal $`F^1`$ interactions and for the realistic clocks is of the order $`10^{10}`$. In addition we’ll assume that all RF constituents spins and orbital momentums are compensated so that its total orbital momentum is zero, like in $`\alpha `$-particle ground state. In this case the system $`F^1`$ c.m. motion can be reduced to the motion of the spinless boson with the mass $`m_1`$ and in the next order the mass operator $`m_t=m_1+H_c`$ will be used. We’ll start the proper time study with the simple models of quantum RFs with clocks, yet we expect its main results to be true also for the more sophisticated models. To introduce our main idea let’s regard the dynamics of the moving clocks in Special relativity . We’ll suppose that the proper (clocks) time is defined by the coordinate $`\theta `$ describing some internal system motion independent of its c.m. motion. For the simplicity assume that Hamiltonian of clocks $`H_c`$ results in the trajectory $`\theta (t)=\omega t+\theta _0`$ of the clocks canonical observable $`\theta `$, which renormalized into the time observable $`\tau =\frac{\theta }{\omega }`$. This is the property which is expected from ideal clocks and the simplest example of such system is the motion of free particle relative to observer $`\tau =\frac{x}{v}`$ . For this and some other clocks models described below the Hamitonian of clocks with mass $`m_1`$ which c.m. moves with momentum $`\stackrel{}{p}_1`$ relative to ARF : $$H_T=(m_t^2+\stackrel{}{p}_1^2)^{\frac{1}{2}}$$ where $`m_t=m_1+H_c`$. If $`\theta ,\stackrel{}{p}_1`$ commutes, solving Hamilton equations in ARF time $`\tau _0`$ one obtains $`\theta (\tau _0)=B_1\omega \tau _0+\theta _0^{}`$, where $`B_1=\frac{m_t}{H_T}`$ coincides with Lorentz boost value. So as expected , if $`\theta `$ is measured by the observer at rest he finds the proper time $`\tau _1=B_1\tau _0`$ of moving frame. Yet we’ll show that the quantum fluctuations of RF motion results in the principally new additional effects. One of the most simple and illustrative quantum clocks models is the quantum rotator proposed by Peres . The rotator Hamiltonian $`\widehat{H}_c=2\pi \omega i\frac{}{\theta }`$ , where $`\theta `$ is the rotator’s polar angle. Preparing the special initial state $`\phi _c(\theta )=\|v_J^0`$ at $`t=0`$, where $`J`$ is its maximum orbital momentum one obtains the close resemblance of the classical clocks hand motion. The clocks state $`\phi _c(\theta 2\pi \omega t)`$ for large $`J`$ has the sharp peak at $`\overline{\theta }=2\pi \omega t`$ with the uncertainty $`\mathrm{\Delta }_\theta =\pm \frac{\pi }{N}`$ and can be visualized as the constant hand motion on the clocks circle. Our main clocks model - $`C_x`$ exploits the nonrelativistic particle motion relative to observer with Hamitonian $`H_c=\frac{\stackrel{}{p}^{\mathrm{\hspace{0.17em}2}}}{2m}`$ . Let’s consider the particle 3-dimensional motion, but choose as its initial state at $`t=0`$ the Gaussian packet factorized in $`x`$ direction which momentum state vector is : $$\varphi _c(\stackrel{}{p})=A\varphi (p_y,p_z)e^{\frac{\sigma _x^2}{2}(\overline{p}_xp_x)^2}$$ (8) for which $`\overline{p}_x0`$. $`\sigma _x`$ is the initial wave packet spatial spread. Then the simplest Hermitian observable which gives the time estimate is $`\widehat{\tau }=\frac{mx}{\overline{p_x}}`$ \- the particle’s position on the arbitrary $`x`$ axe. It describes the nonshifted measurement with $`\overline{\tau }=t`$ and the finite dispersion $`D_0(t)`$ for $`0<t<\mathrm{}\text{[17]}`$. In fact in $`C_x`$ model $`\widehat{\tau }`$ is the clocks hand position operator or the pseudotime operator, and not a time operator in a strict sense . So from all sides $`C_x`$ can be regarded as the realistic clocks model in which measuring $`\widehat{\tau }`$ one obtains the correct $`t`$ estimate with some statistical error having quantum origin. $`C_x`$ wave function $`\phi _c(x,t)`$ evolution can be factorized as the packet centre of gravity motion with the constant velocity $`\frac{p_x}{m}`$ and the packet smearing around it. For the given initial state there is unambiguous correspondence between the state vector $`\|\phi _c(x,t)`$ and time $`t`$, so the quantum clocks synchronization at $`t=0`$ means the preparation of the state $`\phi _c(x,0)`$. From the corresponding Heisenberg equation one can find Heisenberg position operator for the Hamiltonian $`H_c`$ : $$x(t)=(\frac{p_xt}{m}+x_0)$$ (9) where $`x_0=x(0)`$ is Schrodinger position operator If $`\overline{x}_0=0`$ the corresponding clock time operator, which will be extensively used in relativistic theory can be decomposed as : $$\widehat{\tau }=t+\frac{p_x\overline{p}_x+x_0m}{\overline{p}_x}$$ The first term gives the time expectation value and the rest gives the clocks dispersion $`D(t)`$. To simplify our discussion we’ll consider also the clocks model $`C_0`$ with the linear approximation of the position operator $`x(t)=\omega t+x_0`$ where parameter $`\omega =\frac{\overline{p}_x}{m}`$ which is the analog of Peres clocks for unbounded motion. $`C_0`$ Hamiltonian $`H_c^0=\omega p_x`$ is unbounded from below for the continuous spectra, but for the interpretation of the relativistic clock effects it’s unimportant. Any initial $`C_0`$ state (8) evolves as $`\phi _c^0(x\omega t)`$, so the initial form of wave function is conserved and only its centre of gravity moves. Now we’ll consider the relativistic $`C_x`$ model in which RF $`F^1`$ and the particle $`G^2`$ system $`S_2`$ motion is relativistic. We’ll suppose that ARF proper time $`\tau _0`$ is defined also by some quantum clocks ,which dispersion is so small that can be neglected and $`\tau _0`$ is the parameter. If $`F^1`$ internal interactions neglected $`F^1`$ c.m. motion described by the massive boson wave packet evolution and $`S_2`$ Hamiltonian $`H_T`$ in ARF is the sum of two KGR Hamiltonians for the positive energy states : $$H_T=(m_1^2+\stackrel{}{p}_1^{\mathrm{\hspace{0.17em}2}})^{\frac{1}{2}}+(m_2^2+\stackrel{}{p}_2^{\mathrm{\hspace{0.17em}2}})^{\frac{1}{2}}=(s+\stackrel{}{p}_s^{\mathrm{\hspace{0.17em}2}})^{\frac{1}{2}}$$ (10) , where $`\stackrel{}{p}_s=\stackrel{}{p}_1+\stackrel{}{p}_2`$ and $`s`$ is invariant mass square. $`\sqrt{s}`$ can be regarded as the Hamiltonian of two objects $`G^2,F^1`$ relative motion in their c.m.s. equal to system $`S_2`$ mass operator : $$m_t=\sqrt{s}=(m_1^2+\stackrel{}{q}^{\mathrm{\hspace{0.17em}2}})^{\frac{1}{2}}+(m_2^2+\stackrel{}{q}^{\mathrm{\hspace{0.17em}2}})^{\frac{1}{2}}$$ (11) where $`\stackrel{}{q}`$ is $`G^2`$ relative invariant momentum . If $`\|\overline{q}\|`$ is small we can choose as $`p_x`$ \- clock momentum $`\stackrel{}{q}`$ projection along any suitable direction for which $`\overline{q}_x0`$. In this case $`F^1,G^2`$ relative motion can be regarded as nonrelativistic and $`F^1`$ mass operator approximated : $$m_tm_s+\frac{q_x^2}{2\mu _{12}}+E_k(q_y,q_z)$$ ,where $`\mu _{12}`$ is $`G^1,F^2`$ reduced mass, $`m_s=m_1+m_2`$ is $`S_2`$ rest mass. In this case $`E_k`$ is small and can be omitted in the calculations. Like in nonrelativistic case $`F^1`$ proper time in this $`C_x`$ relativistic model can be estimated measuring in ARF the distance $`x=x_2x_1`$ between $`F^1`$ and the particle $`G^2`$ which operator is equal to : $`x=i\frac{}{q_x}`$. For the obtained $`m_t`$ $`S_2`$ Hamiltonian $`H_T`$ can be formally rewritten : $$H_T=[(m_s+H_c)^2+\stackrel{}{p}_s^{\mathrm{\hspace{0.17em}2}}]^{\frac{1}{2}}$$ (12) where $`H_c=\frac{q_x^2}{2\mu _{12}}`$. Moreover it is reasonable to assume that this square root Hamitonian can describe the evolution of any clocks model with nonrelativistic interactions $`H_c`$ i.e. for $`\alpha _I1`$ . Here and below the algebraic operations with the operators (if they don’t result into singularities) means Tailor raw decomposition. If $`F^1,G^2`$ relative motion is nonrelativistic we can assume for the beginning that $`F^1`$ and $`S_2`$ c.m.s. proper time practically coincide. For the classical motion $`F^1`$ Lorentz factor in $`S_2`$ c.m.s. $`(1+\frac{\stackrel{}{q}^{\mathrm{\hspace{0.17em}2}}}{m^2})^{\frac{1}{2}}`$ and below we’ll show that in quantum case their difference is also negligible. It’s impossible to resolve in analytical form the Schrodinger equation for $`H_T`$ of (12) , only some approximated solutions discussed below can be found. $`S_2`$ observables evolution can be found solving Heisenberg equation for the Hamitonian $`H_T`$ of (12) or for exact Hamitonian of (10) as will be done below . After the simple algebra one obtains $`x`$ evolution in ARF proper time $`\tau _0`$ : $`\dot{x}=i[x,H_T]={\displaystyle \frac{im_t}{(m_t^2+\stackrel{}{p}_s^2)^{\frac{1}{2}}}}[x,H_c]=iB_1[x,H_c]`$ (13) We’ll call the operator $`B_1(\stackrel{}{p}_s,m_t)`$ the time boost operator, which interpretation will be discussed after some calculations. The clock observables we obtain in this clock models are the functions of canonical momentums only and due to it their factor ordering is unimportant for our problem. After the commutators calculations we can approximate operator $`m_t`$ by the parameter $`m_tm_s+\frac{\overline{q}_x^2}{2m}`$. The operator $`x`$ easily restored from $`\dot{x}`$ : $$x(\tau _0)=B_1(\stackrel{}{p}_s,m_t)\frac{q_x\tau _0}{\mu _{12}}+x_0$$ where $`x_0`$ is Schroedinger position operator for $`\tau _0=0`$. If we take that $`\overline{x}_0=0`$ it results into $`F^1`$ proper time operator : $$\widehat{\tau }_1=B_1(\stackrel{}{p}_s,m_t)\frac{q_x}{\overline{q}_x}\tau _0+\frac{\mu _{12}x_0}{\overline{q}_x}$$ (14) Its meaning will be discussed after some calculations, but formally it’s $`F^1`$ moving clocks hand position measured in ARF at the moment $`\tau _0`$. $`\tau _1`$ operator in $`C_0`$ model have the simpler form which prompts its interpretation : $$\widehat{\tau }_1=B_1(\stackrel{}{p}_s,m_t)\tau _0+\frac{x_0^{}}{\omega }$$ (15) If $`\overline{x}_0^{}=0`$ $`C_0`$ $`\widehat{\tau }_1`$ expectation value $`\overline{\tau }_1=\overline{B}_1\tau _0`$ coincides with the classical Lorentz time boost value. Its dispersion have the form : $`D_\tau =D_L(\tau _0)+D_c=D_B\tau _0^2+\overline{D}_2\tau _0+D_0`$ (16) where $`D_B=\overline{B}_1^2(\overline{B}_1)^2`$ and $`D_0=\frac{x_0^{}_{}{}^{}2}{\omega ^2}`$ is the clocks mechanism dispersion, which for $`C_0`$ is time independent. Operator $`D_2`$ is equal to : $$D_2=\frac{B_1x_0+x_0B_1}{\omega }$$ (17) The numerical calculations show that for $`C_0`$ localized states $`D_2`$ expectation value is very small and can be neglected. If $`D_0`$ is small $`\tau _1`$ fluctuations are defined mainly by $`D_L(\tau _0)`$ Lorentz boost dispersion stipulated by $`\stackrel{}{p}_s`$ fluctuations in $`F^1`$ wave packet. It’s independent of the clocks mechanism and demonstrates that the proper time measurement have the principal quantum uncertainty growing unrestrictedly proportional to $`\tau _0^2`$. For $`C_x`$ model the factor $`\frac{q_x}{\overline{q_x}}`$ in (14) produces additional $`\widehat{\tau }_1`$ fluctuations. Due to it Lorentz boost expectation value differs only for the small factor of the order $`\alpha _I`$ : $$\overline{\tau _1}=\tau _0\overline{B}_1[1+\frac{\overline{B}_1}{\sigma _x^2\mu _{12}m_s}(1\overline{B}_1^2)]$$ It results from $`m_t`$ dependence on $`p_x`$ and reflects influence of clocks energy on total mass. We’ll neglect this effect in $`C_x`$ dispersion also described by ansatz (16), but with different parameters : $`D_2={\displaystyle \frac{\mu _{12}}{\overline{q}_x^2}}(q_xB_1x_0+x_0q_xB_1);`$ (18) $`D_B={\displaystyle \frac{\overline{q}_x^2}{(\overline{q}_x)^2}}\overline{B}_1^2(\overline{B}_1)^2;D_0={\displaystyle \frac{\mu _{12}^2\sigma _x^2}{\overline{q}_x^2}}`$ (19) Here $`\overline{D}_2=0`$ for the gaussian wave packets (8) and any other localizable states. Due to $`q_x`$ fluctuations absent in $`C_0`$ model the part of $`D(\tau _1)`$ : $$D_x=D_0+\frac{\overline{q}_x^2(\overline{q}_x)^2}{(\overline{q}_x)^2}(\overline{B}_1)^2\tau _0^2$$ can be related to the packet smearing along $`x`$ coordinate, regarded as the clocks mechanism uncertainty. To illustrate the physical meaning of this time operator let’s consider the corresponding approximate solutions of $`F^1`$ state evolution equation for Hamiltonian (12). For $`\alpha _I0`$ we can decompose $`H_T`$ of (12) in the first $`\alpha _I`$ order : $$i\frac{d\mathrm{\Psi }_s}{d\tau _0}[(m_s^2+\stackrel{}{p}_s^2)^{\frac{1}{2}}+\frac{m_1\widehat{H}_c}{(m_s^2+\stackrel{}{p}_s^2)^{\frac{1}{2}}}]\mathrm{\Psi }_s$$ (20) Here the first term is independent of $`H_c`$ which permit to represent $`\mathrm{\Psi }_s`$ as the sum of factorized states. The second term is in fact the product of clock Hamitonian and Lorentz boost $`B_1`$. Let’s choose the initial $`F^1`$ state $`\mathrm{\Psi }_s(0)=\mathrm{\Phi }_s(\stackrel{}{p}_s)\phi _c(x,0)`$ and $`\mathrm{\Phi }_s=c_l\|\stackrel{}{p}_{sl}`$, where the sum denotes the integral over $`\stackrel{}{p}_s`$. From our definition of quantum clocks synchronization it follows that $`\mathrm{\Psi }_s(0)`$ describes $`F^1`$ clocks synchronized with ARF clocks at $`\tau _0=0`$. Solving equation (20) one finds : $$\mathrm{\Psi }_s(\tau _0)=c_l\phi _c(x,B_l\tau _0)\|\stackrel{}{p}_{sl}e^{iE(\stackrel{}{p}_{sl})\tau _0}$$ (21) where $`E(\stackrel{}{p})=(m_s^2+\stackrel{}{p}^2)^{\frac{1}{2}}`$ , $`B_l=B_1(\stackrel{}{p}_{sl},m_s)`$. For linear clock $`C_0`$ Hamiltonian $`H_c=H_c^0`$ for small $`\alpha _I`$ this state can be rewritten : $$\mathrm{\Psi }_s(\tau _0)=\underset{l}{}c_l\phi _c^0(x\omega B_l\tau _0)\|\stackrel{}{p}_{sl}e^{iE(\stackrel{}{p}_{sl})\tau _0}$$ (22) To make the situation more clear suppose that $`\phi _c^0(0)=\delta (x)`$, which evolves at rest into $`\delta (x\omega \tau _0)`$ . Then $`x`$ measurement defines the time $`\tau `$ of quantum clocks at rest unambiguously and with zero dispersion, but $`\mathrm{\Psi }_s`$ of (22) in general isn’t $`x`$ eigenstate. It means that at any $`\tau _0>0`$ $`\mathrm{\Psi }_s`$ is the entangled superposition of the states $`\phi _c^0`$ which $`F^1`$ clocks acquires at the consequent $`\tau _1`$ moments. As was shown there is one-to one correspondence between clock state $`\phi _c(x,t)`$ and the time moment $`t`$ and in some sense it can be regarded as the ’superposition’ of $`F^1`$ proper time moments, or more precisely $`F^1`$ states existed at this moments. For example $`F^1`$ clocks hand can show 3,4 and 5 o’clocks simultaneously which can be tested by $`x`$ measurement at some $`\tau _0`$ in ARF. This spread corresponds to $`D_B`$ dispersion term resulting from the $`F^1`$ momentum $`\stackrel{}{p}_s`$ uncertainty. For the realistic clocks their $`x`$ dispersion given by $`D_0`$ isn’t zero even at rest and this two terms added as statistically independent effects. $`\mathrm{\Psi }_s`$ for $`C_x`$ Hamitonian is given by (21) and admits the same interpretation. It corresponds to the more complicated form of time dependent dispersion (19) which can be eventually factorized into the same two parts - relativistic and clock mechanism. So we conclude that the interpretation which follows from the approximate Schrodinger equation agrees well with Heisenberg operator calculus. In fact operator $`\tau _1`$ describes $`F^1`$ proper time in the limit when this clock dispersion is very small and the clock energy is much less then $`F^1`$ total mass energy i.e. $`\alpha _I0`$. Obtained results suppose that the proper time of any quantum RF being the parameter in it simultaneously will be the operator from the ’point of view’ of other RF. Qualitatively the appearance of RF proper time fluctuations can be understood considering the superposition of momentum eigenstates $`\|\stackrel{}{p}_{si}`$ in $`S_2`$ wave packet as the superposition of $`S_2`$ velocities $`\stackrel{}{\beta }_i`$ and corresponding Lorentz factors $`\gamma _1(\stackrel{}{\beta }_i)`$. In Special Relativity $`F^1`$ proper time $`\tau _1`$ measured at the same $`\tau _0`$ in ARF depends on $`\gamma _1`$. If we formally extends this dependence on $`F^1`$ wave packet motion we get that the proper time will fluctuate proportionally to $`\gamma _1`$ spread. So $`F^1`$ clocks measurement in ARF shows how much time passed in $`F^1`$ in this particular event and can give the different value for another event of the same ensemble. It means that the time moments in different RFs corresponds only statistically with the dispersion $`D_\tau `$ in ARF given by (16). It differs from Special Relativity where one to one correspondence between $`\tau _1,\tau _0`$ time moments always exists , but can be incorporated into relativistic QEP if we find the analogous time relations between two quantum RFs of finite mass. In fact $`\tau _1`$ is more correct to relate to $`S_2`$ c.m.s. rest frame, but regarding the difference between $`F^1`$ and $`S_2`$ c.m.s. proper time operators $`\tau _1^{},\tau _1`$ it’s easy to show that they coincide if $`\overline{q}_x0`$. From it we conclude that the principal part of the relativistic time operator, independent of any particular clocks mechanism features have the form in the limit $`\alpha _I0`$ : $$\widehat{\tau }_1^{}=B_1(\stackrel{}{p}_1,m_1)\tau _0$$ (23) Moreover this formulae permits to define formally the time operator for any object including the single massive particle. This operator form of $`\tau _1^{}`$ is closely connected with Fock-Shwinger proper time $`\tau _F`$ formalism interpretation and will be discussed in detail in the forcoming paper . Note only that $`\widehat{\tau }_1^{}(\tau _0)`$ measurement gives $`F^1`$ proper time $`\tau _F`$ estimate at $`\tau _0`$ moment of ARF time. On the opposite in Fock-Shwinger formalism $`\tau _F`$ is the parameter time to which particular values operators $`\widehat{\tau }_0(\tau _F),\stackrel{}{r}_1(\tau _F)`$ related. In distinction with our formalism it makes $`\tau _F`$ interpretation confusing, because $`\stackrel{}{r}_1`$ and other $`F^1`$ operators are measured in ARF, hence the time of measurement defined in $`F^1`$ to which as we have shown in quantum case they related only statistically. The practical realization of $`x`$ measurement in ARF can be the intricated procedure, which scheme we don’t intend to discuss here. Note only that to perform it one should measure simultaneously the distance between $`F^1`$ and $`G^2`$ and their total momentum giving total velocity and this two operators commute. Some examples of the analogous nonlocal observables measurements are described in . The most disputable question here is the relativistic particle coordinate measurements. Yet in the considered case, when the relative $`F^1,G^2`$ average velocity is small then $`x`$ is the nonrelativistic coordinate operator. Yet to prove the quantum equivalence principle it’s necessary to perform the full relativistic calculations. We’ll present such completely relativistic results for $`C_x`$ model using Newton-Wigner Hermitian operator of the space coordinate which is the direct analog of nonrelativistic operator $`x_1`$ : $$\widehat{x}_{NW}^1=i\frac{d}{dp_{x1}}i\frac{p_{x1}}{2(m_1^2+\stackrel{}{p}_1^2)}$$ (24) The operator of two objects relative coordinates conjugated to c.m. momentum $`q_x`$ can be derived from this objects c.m. Hamiltonian (11) : $$\widehat{x}_{NW}=x+F(\stackrel{}{q})=i\frac{d}{dq_x}i\frac{q_x}{\sqrt{s}}(\frac{1}{w_1}+\frac{1}{w_2})$$ (25) where $`w_i=(m_i^2+\stackrel{}{q}^2)^{\frac{1}{2}}`$. The clocks time observable in $`F^1`$ rest frame is proportional to $`x_{NW}`$ : $$\tau =\frac{x_{NW}\overline{x}_{NW}(0)}{\overline{\beta }_x}$$ where $`\beta _x=q_x(w_1^1+w_2^1)`$ is $`F^2,G^1`$ relative velocity, If we choose $`\overline{x}_{NW}(0)=0`$ , then solving Heisenberg equation in ARF for the Hamiltonian of (10) we find the resulting $`F^1`$ time operator : $$\widehat{\tau }_1=\frac{B_1(\stackrel{}{p}_s,m_t)\tau _0\beta _x+x_{NW}(0)}{\overline{\beta }_x}$$ (26) where in $`B_1`$ $`m_t=\sqrt{s}`$. This is the exact relativistic expression for $`\tau _1`$ without assumption of $`q_x`$ smallness. $`\overline{\tau }_1`$ corresponds to Lorentz boost value $`\overline{B}_1`$ which depends both on $`\stackrel{}{p}_s`$ and $`\stackrel{}{q}`$. It’s easy to note that the momentum dependent part of $`x_{NW}`$ is constant in time and consequently can only enlarge the clocks mechanism dispersion $`D_0`$. Due to it the dispersion structure is the same as for nonrelativistic relative motion of (19) but its members are described by the more complicated formulaes omitted here. In fact this calculations evidence that $`x_{NW}`$ meausurements introduces only additional time-independent clock dispersion of the order of $`G^2`$ Compton wavelength without changing our previous conclusions about time operator properties. In fact $`F^1`$ proper time measurement in ARF can be performed by two different methods which equivalence must be proved. In the first method described above the detector $`D_0`$ installed in ARF measures $`\tau _1`$ and induces $`C_x`$ state collapse. In the second one the detector $`D_1`$ installed in $`F^1`$ measures the clock state and after it $`D_1`$ signal transfered to ARF. In this case we should consider the collapse in the moving frame , which is difficult to describe. But we must note that independently of its mechanism such interaction happens after this clocks evolves to this state and so can’t influence directly on their evolution, so it seems correct to neglect it at this stage. Obtained time-fluctuation effect reminds the well-known life-time dilatation for the relativistic unstable particles . In this framework such particle can be regarded as the elementary binary clock having only two states. Obtained results evidence that the proper time in Quantum RF depend on the RF quantum state, but doesn’t prove QEP directly. To do it we must consider two finite mass RFs on equal ground and to find the time transformation between them. ## 4 Relativistic Quantum Frames To calculate the time operator between two RFs of finite mass it’s necessary first to find the particle evolution equation in quantum RF rest frame. In general the system Poincare group irreducible representations contain the information which permit to describe its evolution completely, but due to appearance of time operators to find this representations for quantum RF is quite a problem. Therefore we choose another route; first we’ll find the free particle evolution equation and corresponding proper time operator from Dirack constraints quantization. After it we’ll investigate Poincare transformations for quantum RFs with the clues prompted by this Hamiltonian ansatz. Dirack constraint formalism which permit to define free particle/antiparticle positive Hamiltonian was developed by Gitman and Tiutin . They’ve shown that starting from free scalar particle action $`S=m𝑑s`$ Dirack constraint quantization of $`p_\mu ^2m^2`$ initial superhamiltonian results into positive square root Hamiltonian $`H_p`$ as function of 3-momentum $`\stackrel{}{p}`$ plus additional charge $`\xi =\pm 1`$ discriminating antiparticles. In quantum case it was shown to result in Klein-Gordon square root (KGR) equation for both $`\xi `$ $$i\frac{d\psi }{d\tau _0}=\sqrt{\stackrel{}{p}^2+m^2}\psi $$ Following this approach we start from classical two particles action defined in ARF $`F^0`$ with time parameter $`\tau _0`$ : $$S(\tau _0)=L𝑑\tau _0=m_1[(\dot{x}_{10}^2\dot{x}_{1i}^2)^{\frac{1}{2}}+m_2(\dot{x}_{20}^2\dot{x}_{2i}^2)^{\frac{1}{2}}]𝑑\tau _0$$ From it one finds 4-momentums $`\pi _{1\mu },\pi _{2\mu }`$ which satisfy to superhamiltonian constraints $`\pi _{j\mu }^2=m_j^2`$. Due to $`m_1,m_2`$ dynamics independence $`S_2`$ system Dirack quantization results in double number of 3- momentums $`\stackrel{}{p}_j`$ and charges $`\xi _j`$. After simple calculations repeating Gitman-Tiutin anzats one obtains system Hamiltonian : $$H_T=\sqrt{\stackrel{}{p}_1^2+m_1^2}+\sqrt{\stackrel{}{p}_2^2+m_2^2}$$ For two particles system quantization we can use this Hamiltonian in Schrodinger equation. But if we consider $`m_1`$ as quantum RF then like in nonrelativistic case, before quantize it we must put additional constraints corresponding to $`F^1`$ choice as rest frame and defining $`m_2`$ operators transformations into it. Analogously to nonrelativistic case we put constraint on $`F^1`$ momentum in its rest frame $`\stackrel{}{p}_{11}^20`$, meaning that RF don’t move relative to itself. For charges we put constraint $`\xi _1^{}=\xi _1^2`$ \- .i.e. RF can’t be antiparticle for itself, and correspondingly $`\xi _2^{}=\xi _1\xi _2`$. Below for simplicity we’ll describe in detail only restricted Hilbert space sector without antiparticles . Corresponding to $`H_T`$ two particles state vector isn’t interpereted by us as their wave function in $`F^0`$. For two particles correlations it can result into contradictions connected with nonlocalities. But we don’t need it in such role and will study only reduced state vector in quantum RF $`F^1`$. To find the transformations from ARF to $`F^1`$ we consider the system $`S_2`$ of RF $`F^1`$ and particle $`G^2`$ which momentums $`\stackrel{}{p}_i`$, energies $`E_i`$ are defined in ARF. For the constraint described above we choose ARF FC $`\stackrel{}{p}_0^{\mathrm{\hspace{0.17em}2}}0`$ and $`S_2`$ Hamitonian : $`H_A^0=(m_0^2+\stackrel{}{p}_0^2)^{\frac{1}{2}}+(m_1^2+\stackrel{}{p}_1^2)^{\frac{1}{2}}+(m_2^2+\stackrel{}{p}_2^2)^{\frac{1}{2}}`$ (27) if $`G_2`$ is boson. Like in nonrelativistic case all $`\stackrel{}{p}_i`$ are the operators and state vector ascribed also to ARF $`\|\stackrel{}{p}_0=0`$. This ARF constraint formalism reproduces all relativistic QM results for $`m_0\mathrm{}`$. Performing transformations to $`F^1`$ rest frame we assume that the proper time parameter $`\tau _1`$ can be defined in it from $`F^1`$ clocks measurements extrapolation as was described in previous chapter. Then from $`\stackrel{}{p}_1`$ constraint and correspondence with Lorentz momentum transformations we phenomenologically find $`m_i`$ momentums: $`\stackrel{}{p}_{12}=\stackrel{}{p}_i+{\displaystyle \frac{(\stackrel{}{n}_1\stackrel{}{p}_i)(E_1m_1)\stackrel{}{n}_1E_i\stackrel{}{p}_1}{m_1}}+\stackrel{}{F}_i(\stackrel{}{p}_0)`$ (28) where $`\stackrel{}{n}_1=\stackrel{}{p}_1\|\stackrel{}{p}_1\|^1`$. $`\stackrel{}{F}_i`$ are undefined at this stage operators for which $`\stackrel{}{F}_i=0`$ and can be neglected in the following calculations. This transformation results in $`\stackrel{}{p}_{11}=0`$ if $`\stackrel{}{p}_0=0`$ and will advocate its form below where we’ll discuss Poincare group for quantum RFs. Until then this is phenomenological transformation which for definite RFs momentum and velocities reproduces Lorentz transformations. If $`G^2`$ have spin zero then the Hamiltonian $`H`$ transformed from ARF to $`F^1`$ is equal : $`H_{Tot}^1=H_0^1+H_1^1+H_2^1={\displaystyle \underset{i=0}{\overset{2}{}}}(m_i^2+\stackrel{}{p}_{1i}^{\mathrm{\hspace{0.17em}2}})^{\frac{1}{2}}`$ (29) For classical Special Relativity where normally RF supposed to have infinite mass $`\stackrel{}{p}_{1i},H_R^1`$ corresponds to the canonical momentums for finite mass RFs . We see that $`F^0`$ motion Hamiltonian is factorized and so we can drop it and regarding $`S_2`$ motion can use $`S_2`$ Hamiltonian $`H^1=H_1^1+H_2^1`$. In quantum case in $`H^1`$ we can’t simply omit $`\stackrel{}{p}_{11}`$ because now it’s operator. So in $`F^1`$ proper time $`\tau _1`$ $`S_2`$ evolution equation is : $$i\frac{d\psi ^1}{d\tau _1}=[(m_1^2+\stackrel{}{p}_{11}^{\mathrm{\hspace{0.17em}2}})^{\frac{1}{2}}+(m_2^2+\stackrel{}{p}_{12}^{\mathrm{\hspace{0.17em}2}})^{\frac{1}{2}}]\psi ^1=(s(\stackrel{}{q})+\stackrel{}{p}_s^{\mathrm{\hspace{0.17em}2}})^{\frac{1}{2}}\psi ^1$$ (30) where $`S_2`$ c.m. observables $`\stackrel{}{q},\stackrel{}{p}_s`$ defined in chap.3. Solutions of this equation describe $`G^2`$ normalized free wave packet localizable relative to $`F^1`$ rest frame : $$\mathrm{\Psi }^1(\tau _1)=\phi _2(\stackrel{}{p}_{12})e^{iE^1\tau _1}\|m_2,\stackrel{}{p}_{12}\|m_1,\stackrel{}{p}_{11}=0=\phi _2^{}(\stackrel{}{q})e^{iE^1\tau _1}\|\sqrt{s},\stackrel{}{p}_s=\stackrel{}{p}_{12}\|m_1,\stackrel{}{q}\|m_2,\stackrel{}{q}$$ (31) expressed also via $`S_2`$ c.m. observables. Here $`E^1=E_1^1+E_2^1`$ are $`H^1`$ eigenvalues. They differ from the standard KGR energy only on $`m_1`$ and so we can use in $`F^1`$ rest frame the standard KGR momentum spectral decomposition and the states scalar product . In $`F^1`$ rest frame together with its proper time $`\tau _1`$ the space coordinate can be defined. We choose arbitrarily as $`G^2`$ coordinate (nonhermitian) operator in $`F^1`$ : $`\widehat{x}_{12}=i\frac{}{q_x}`$ and corresponding Hermitian Newton-Wigner operator can be easily derived. Note that $`x_q`$ defined in $`F^1`$ differs from the same operator defined in c.m.s., yet our following results doesn’t depend on the particular form of this operator. $`x_{12}`$ also differs from the operator $`x_p=i\frac{}{p_{12x}}`$ which corresponds to the classical distance between $`F^1`$ and $`G^2`$. They coincide only in the limit $`m_1\mathrm{}`$ or in nonrelativistic case. Now we can calculate $`F^2`$ proper time operator as function of the proper time in $`F^1`$. To perform it we assume again that $`F^2`$ c.m. motion is equivalent to the spinless particle $`G^2`$ motion. In the described framework the Hamiltonian of $`F^2`$ with $`C_0`$ or $`C_x`$ clocks in $`F^1`$ rest frame can be obtained substituting in $`\widehat{H}^1`$ of (30) $`m_2=m_2^{}+\widehat{H}_c`$. $`\widehat{\tau }_2`$ can be found solving Heisenberg equation for $`F^2`$ clocks coordinate $`\dot{x}=i[x,H^1]`$ analogously to (13). If we omit analogously to (23) the members describing the clocks mechanism fluctuations the $`F^2`$ proper time operator $`\widehat{\tau }_2`$ is equal : $$\widehat{\tau }_2=\frac{m_2\tau _1}{(m_2^2+\stackrel{}{p}_{12}^2)^{\frac{1}{2}}}\frac{m_2^{}\tau _1}{(m_2^2+\stackrel{}{p}_{12}^2)^{\frac{1}{2}}}=\widehat{B}_1(\stackrel{}{p}_{12},m_2^{})\tau _1$$ (32) This formalism is completely symmetrical and the operator obtained from (32) exchanging indexes 1 and 2 relates the time $`\widehat{\tau _1}`$ in $`F^1`$ and $`F^2`$ proper time \- parameter $`\tau _2`$. The Special Relativity limit when $`\tau _2`$ becomes the parameter is obvious and analoguosly to it the average time boost depends on whether $`F^1`$ measures $`F^2`$ clocks observables, as we consider or vice versa, and this measurement makes $`F^1`$ and $`F^2`$ nonequivalent . The new effect will be found only when $`F^1`$ and $`F^2`$ will compare their initially synchronized clocks. In QM formalism this synchronization means that $`F^2`$ state prepared at the moment $`\tau _0`$ can be factorized as $`\mathrm{\Phi }_2(\stackrel{}{p}_{12})\phi _c(x,0)`$ analogous to (21). If this $`F^2`$ time measurements repeated several times (to perform quantum ensemble) it’ll reveal not only classical Lorentz time boost , but also the statistical spread having quantum origin with the dispersion given in (16). Obtained relation between two finite mass RFs proper times evidence that Quantum Equivalence principle can be correct also in relativistic case. If the number of particles $`N_g>1`$ then for the system state description the clasterization formalism can be used . According to it for $`N=3`$ Hamiltonian in $`F^1`$ of two free particles $`G^2,G^3`$ rewritten through the system canonical observables acquires the form : $$\widehat{H^1}=(m_1^2+\stackrel{}{p}_{11}^2)^{\frac{1}{2}}+(s_{23}+\stackrel{}{p}_{1,23}^2)^{\frac{1}{2}}=(s+\stackrel{}{p}_s^2)^{\frac{1}{2}}$$ (33) ,where $`\sqrt{s}_{23}`$ is $`G^2,G^3`$ invariant mass, $`\sqrt{s},\stackrel{}{p}_s`$ are the system total invariant mass and momentum. In clasterization formalism at the first level the relative motion of $`G^2,G^3`$ defined by $`\stackrel{}{q}_{23}`$ their relative momentum is considered. At the second level we regard them as the single quasiparticle \- cluster $`C_{23}`$ with mass $`\sqrt{s}_{23}`$ and momentum $`\stackrel{}{q}`$ in the system c.m.s. It transformed to $`\stackrel{}{p}_{1,23}`$ momentum in $`F^1`$ and so at any level we can regard the relative motion of two objects only. This procedure can be extended in the obvious inductive way to arbitrary $`N`$. If we have two reference frames $`F^1,F^2`$ and $`N_g0`$ then their relative momentums can be also described by the cluster formalism. Due to appearance of the time operator between two RFs to find Poincare group transformations for quantum RFs $`\widehat{U}_{2,1}^s(\tau _2,\tau _1)`$ is quite a problem and here we can present it only phenomenologically for some simple examples. We don’t include rotations into consideration, so this results are suitable completely only for 1-dimensional case and for more dimensions they give only partial description of the Lorentz transformations. Consider first the case $`N=2`$ when $`S_2`$ include $`F^1,F^2`$ only and its state in $`F^1`$ rest frame $`\mathrm{\Psi }^1(\tau _1)`$ is the solution (31) of eq. (30). We’ll take that it transformed by $`U_{2,1}^F`$ into state $`\mathrm{\Psi }^2(\tau _2)`$ in $`F^2`$ rest frame. If $`F^1,F^2`$ clocks are synchronized at $`\tau _1=\tau _2=0`$ then for this time moment $`\mathrm{\Psi }^2(0)=\widehat{U}_{2,1}^F(0,0)\mathrm{\Psi }^1(0)`$ and from $`F^1,F^2`$ symmetry it follows : $`\|\mathrm{\Psi }^2(0)=\phi _1^{}(\stackrel{}{p}_{21})\|m_2,\stackrel{}{p}_{22}=0\|m_1,\stackrel{}{p}_{21}`$. $`F^{1,2}`$ internal wave functions $`\phi _c^{1,2}(x,0)`$ at $`\tau _1=0`$ are obviously invariant and so omitted here. Like in nonrelativistic case we introduce $`\stackrel{}{p}_f=\stackrel{}{p}_{11}+\stackrel{}{p}_{12}`$, $`\stackrel{}{p}_f^{}=\stackrel{}{p}_{21}+\stackrel{}{p}_{22}`$ and conjugated $`\stackrel{}{r}_f,\stackrel{}{r}_f^{}`$. From the correspondence with Lorentz transformations it should give $`\stackrel{}{p}_{12}=\frac{m_2}{m_1}\stackrel{}{p}_{21}`$ and if to demand that fro relative momentum in $`F^2`$ $`\stackrel{}{q}_2=\stackrel{}{q}`$ must be fulfilled then the simplest transformation is : $$\widehat{U}_{2,1}^F(0,0)=P_fe^{ia_f\stackrel{}{p}_f\stackrel{}{r}_f}$$ (34) where $`a_f=\mathrm{ln}\frac{m_1}{m_2}`$, $`P_f`$ is $`\stackrel{}{r}_f`$ reflection (parity) operator. We see $`\widehat{U}_{21}^F(0,0)`$ ansatz practically coincides with nonrelativistic transform of (5) for $`N=1`$. The passive $`S_N`$ transformation for spinless $`G^i`$ also found from the correspondance principle as the minimal extension of standard Poincare transformations : $$\widehat{U}_{2,1}^s(0,0)=U_{21}^F(0,0)\underset{j=3}{\overset{N}{}}e^{i\stackrel{}{\beta }_f\stackrel{}{N}_j^{}}$$ (35) where velocity operator $`\stackrel{}{\beta }_f=\stackrel{}{p}_f(H_2^1)^1`$, $`\stackrel{}{N}_i^{}=H_i^1\frac{}{\stackrel{}{p}_{1i}}+\frac{}{\stackrel{}{p}_{1i}}H_i^1`$ are $`G^i`$ Poincare generators in $`F^1`$ which coincide with standard ansatz. Then the transformation operator for arbitrary $`\tau _1,\tau _2`$ is : $$\widehat{U}_{21}^s(\tau _1,\tau _2)=\widehat{W}_2(\tau _2)\widehat{U}_{21}^s(0,0)\widehat{W}_1^1(\tau _1)$$ (36) , where $`\widehat{W}_{1,2}(\tau _{1,2})=exp(i\tau _{1,2}\widehat{H}^{1,2})`$ are $`S_N`$ evolution operators and $`H^{1,2}`$ \- $`S_N`$ Hamiltonians in $`F^1,F^2`$ rest frames. It means that despite $`\tau _2`$ and $`\tau _1`$ are correlated only statistically through $`\widehat{\tau }_2`$ nevertheless $`S_N`$ state vectors for free motion in $`F^2,F^1`$ at this moments are related unambiguously. Transformed $`S_N`$ momentums are : $`\stackrel{}{p}_{21}={\displaystyle \frac{m_2}{m_1}}\stackrel{}{p}_{12}+d_1\stackrel{}{p}_{11},\stackrel{}{p}_{22}=\stackrel{}{p}_{11}+d_2\stackrel{}{p}_{11}`$ (37) $`\stackrel{}{p}_{2i}=\stackrel{}{p}_{1i}+{\displaystyle \frac{(\stackrel{}{n}_{12}\stackrel{}{p}_{1i})(E_2^1m_2)\stackrel{}{n}_{12}E_i^1\stackrel{}{p}_{12}}{m_2}}+d_i\stackrel{}{p}_{11}`$ (38) , where $`\stackrel{}{n}_{12}=\frac{\stackrel{}{p}_{12}}{\|p_{12}\|}`$, $`E_i^1`$ are $`G^i`$ energies in $`F^1`$. If to demand that all relative momentums $`\stackrel{}{q}_{ij}`$ conserved (or reflected), then $`d_i`$ can be calculated, but due to their unimportance we omit it here. It’s easy to see that $`\stackrel{}{p}_{1i}`$ of (28) for ARF to $`F^1`$ transform follows from $`U_{2,1}^s`$ after the simple substitutions, and so the semiqualitative Hamiltonian derivation of (28) was consistent. We see that the passive spinless $`G^3`$ transformation differs from the standard one only by the change of velocity parameter to the operator $`\stackrel{}{\beta }`$ which commutes with $`G^3`$ Hamiltonian. It was argued that RF quantum properties can become important in Quantum Gravity , where in principle one should quantize the field, matter and RF simultaneously . In principle our approach permits to calculate the time operator $`\widehat{\tau }_1`$ for RF $`F^1`$ moving in the external gravitational field $`g_{\mu \nu }(x)`$. We assume that ARF is located in the region where this field is weak and so we can take $`\tau _0=x_0`$ \- world time. Analogously to (12) $`F^1`$ clocks Hamiltonian in ARF (for $`g_{oa}=0`$ gauge) : $$\widehat{H}_T=[g_{00}(m_1+H_c)^2+g_{00}g_{ab}p_1^ap_1^b]^{\frac{1}{2}}$$ (39) ,where $`a,b=1,3`$ . Now $`H_T`$ depends on $`x_\mu `$ and due to it solving Heisenberg equation (13) for the clocks hand coordinate $`x_c`$ one obtains the differential relation for $`\tau _1`$ : $$d\widehat{\tau }_1=\frac{\sqrt{g_{00}}(m_1+H_c)d\tau _0}{[(m_1+H_c)^2+g_{ab}p_1^ap_1^b]^{\frac{1}{2}}}=\sqrt{g_{00}}B_g(x,\stackrel{}{p}_1)d\tau _0$$ (40) In this case $`\widehat{\tau }_1`$ becomes the integral operator , where integral is taken over $`\tau _0`$ interval. If $`g_{\mu \nu }`$ is the classical metrics then this relation contains no new physics , except the additional gravitational ’red shift’ time boost proportional to $`\sqrt{g_{00}}`$ . But in Quantum Gravity $`g_{\mu \nu }(x)`$ becomes the operator and its fluctuations can induce the additional quantum fluctuations of the measured $`F^1`$ clocks time. Despite that this fluctuation calculations are quite complicated we can expect from the general Quantum Statistics rules that they can be factorized from the considered Lorentz boost fluctuations induced by the $`F^1`$ momentum fluctuations : $$D_T=D_G(\tau _0)+D_L(\tau _0)+D_o(\tau _0)$$ From this rules we can expect also that for $`F^1`$ motion in the homogeneous gravitation field $`D_G`$ will grows proportionally to $`\tau _0`$ analogous to QED fluctuations (Brownian motion effects). Note that this fluctuations must be independent of RF mass. This approach can give some new insight into the famous time problem of Quantum Gravity which we discuss here briefly. In this aspect the situation in Classical and Quantum Gravity seems to differ principally. Strictly speaking if the metrics becomes the operator it stops to be space-time metrics which unambiguously defines the space-time geometry. Due to it the observer can correctly use only the operational definition of physical space-time by means of clocks and other measurements. In gravity this operational time can originate from some evolving observable of gravitation field or to be the operator describing the time measurement for some free matter object carrying some nongravitational ’foreign’ clocks. The idea that space-time events can be described by their relation with some distributed system or media was extensively explored for long time . The most close to our purposes is the incoherent dust system, each piece of it carrying clocks. Gravity ADM quantization for such system with selfgravitation account permits to extract positive Schrodinger hamiltonian as was shown by Brown and Kuchar . Hence the dust pieces motion in their model was described only semiclassically. Introduction of ’dust space’ $`\stackrel{}{z}`$ permit to quantize the gravitation field. Yet the free quantum motion of dust pieces transforms $`\stackrel{}{z}`$ into the operator on the initial space-time manifold $`x_\mu `$ which makes this quantization procedure contradictory. We describe here briefly the model of dust RFs quantum motion where in the first approximatio its selfgravitation neglected. Let’s consider first classical RF $`F^1`$ free falling in external gravitational field. In $`F^1`$ comoving ’Gaussian’ frame where frame conditions imposed before the field variation we have $`g_{00}^{}=1,g_{0a}^{}=0`$. In this RF for the classical field gravity constraints fulfilled $`H_a(x)=0,H_0(x)=0`$ which permit to calculate $`g_{ab}^{},p_{ab}^{}`$ evolution for $`F^1`$ clock time solving corresponding Hamilton equations for $`H_0`$ . In quantum case this vacuum field constraints results in Wheeler - deWitt equation $`\widehat{H}_0\mathrm{\Psi }=0`$ from which Schrodinger Hamiltonian can’t be derived easily. Now let’s account RF quantum motion and suppose that this constraints holds true also in quantum $`F^1`$ comoving frame. $`F^1`$ proper time for the external observer is given by the operator analogous to (40), but in comoving frame $`\tau _1`$ is just the parameter. In this case we can calculate field observables evolution in $`F^1`$ clocks time from Heisenberg equations for $`H_0`$ vacuum constraint : $$\dot{g}_{ab}^{}(x)=i[g_{ab}^{}(x),H_0(x)]$$ where the commutator in general is nonzero. Note that this equation is obviously local, so to calculate $`g_{ab}^{}(x,\tau _1)`$ we must define $`g_{ab}^{},p_{ab}^{}`$ only on a small spacelike surface region around $`x`$ at a preceding moment $`\tau _1d\tau _1`$. Space coordinates $`x_a`$ supposedly can be defined at least in the close vicinity of $`F^1`$ analoguosly to the definition given above for the flat space-time. Obviously this approach have many associated problems some of which are the construction of multifingered time for quantum RF dust and the field theoretical behavior of such commutators, despite it seems to deserve additional study. For the conclusion we can claim that the extrapolation of QM laws on free macroscopic objects regarded as RFs prompt to change the common approach to the space-time which was taken copiously from Classical Physics. In this paper the relativistic covariant theory of quantum RFs constructed and at least in flat space-time it agrees with the principle of equivalence for quantum RFs. The quantum RF momentum uncertainty results in the quantum statistical fluctuations of Lorentz boost which relates the proper times in two RFs. So in this model each observer has its proper time - parameter and euclidian coordinate space which can’t be related unambiguously with the another observers space-time and in this sense is local.
warning/0007/astro-ph0007476.html	ar5iv	text	# The Cosmological Mass Function in the Zel’dovich Approximation ## 1. Introduction The derivation of the distribution of masses in gravitationally bound objects is one of the principle goals of the theory of the structure formation (for a review see \[Monaco 1998\]). Comparing the theoretical mass function with observations provides important constraints on the cosmological models (see e.g. Bond & Myers 1996, \[Bahcall & Fan 1998\], \[Reichart et al. 1999\]). Rich clusters of galaxies represent a particular interesting class of objects for two reasons. Firstly, they are the largest gravitationally bound objects in the universe and therefore represent rare events. As one of the consequences of being rare events clusters are particularly sensitive to some parameters of the cosmological models ($`\mathrm{\Omega }_m`$ and $`\sigma _8`$). Secondly, the formation of clusters is relatively simple process since it is primarily determined by the gravitational dynamics while other processes (hydro, thermal, etc) are less important than e.g. in the process of galaxy formation. As a result the numerical simulations of clusters of galaxies are more realistic and reliable than simulations of galaxy formation. Measuring the mass function of galaxy clusters is not easy but recently a certain progress has been achieved for both optically (see e.g. \[Bahcall & Cen 1993, Girardi et al. 1998\]) and X-ray (\[Reiprich & Böhringer 1999\]) selected samples. Although there are systematic differences between mass functions obtained by different groups there is a general agreement in a broad sense. Most of theoretical derivations of the cosmological mass function are based on the ideas of \[Press & Schecter (1974)\] that can be summarized as follows: * The mass fraction ( $`F(>M)`$ ) in gravitationally bound objects with masses greater than $`M`$ can be estimated as the fraction of mass satisfied the collapse condition at this scale ( $`\mathrm{\Psi }(\delta _M>\delta _c)`$ ): $`F(>M)=2\mathrm{\Psi }(\delta _M>\delta _c)`$. * The collapse condition is local i.e. it can be expressed in terms of the quantities at one point. * The quantity that determines the collapse is the linearly extrapolated filtered density contrast $`\delta _M\delta _c`$ at a given point. * The value of the threshold $`\delta _c=3/20(12\pi )^{2/3}1.69`$ that corresponds to the collapse of a spherical top-hat model with the similar initial density contrast. It was assumed that the collapse of the spherical top-hat model approximately corresponded to the virialization of the gravitationally bound clump. The mathematical aspects of the Press-Schechter formalism is outlined in the following section. Here I would like to discuss briefly some of ideas suggested since the formalism was proposed in 1974. The excursion set approach (\[Peacock & Heavens 1990\], \[Bond et al. 1991\]) justified the assumption that $`F(>M)=2\mathrm{\Psi }(\delta _M>\delta _c)`$ in the case of a sharp $`k`$-space filter. Many realized that the threshold $`\delta _c=1.69`$ does not provide the best fit to N-body simulations. Although some authors used the canonical value (e.g. \[Bond et al. 1991\], \[Efstathiou et al. 1988\]) others preferred the lower values: $`\delta _c=1.58`$ (Bond & Myers 1996), $`\delta _c=1.44`$ (\[Carlberg & Couchman 1989\]), or even as low as $`\delta _c=1.33`$ (\[Efstathiou & Rees 1988\] and \[Klypin et al. 1995\]). Recently Shapiro et al. (1999) showed that the virialization of the top-hat model occurs when linear extrapolation of the density contrast reaches $`\delta _c1.52`$. One of the major efforts in reduction of the discrepancy of the theory with simulations has been related to incorporating the anisotropic character of gravitational collapse. \[Bond & Myers (1996)\] developed a model that incorporated both the Zel’dovich approximation on large scales and the collapse of a homogeneous ellipsoid on the nonlinear scale. \[Monaco (1995)\] suggested a different collapse condition that corresponded to the collapse along the first axis in the Zel’dovich approximation. \[Audit et al. (1997)\] incorporated some of the nonlinear effects into an anisotropic collapse model. \[Lee & Shandarin (1998a)\] suggested to use the collapse condition corresponding to the collapse along all three axes as described by the extrapolation of the Zel’dovich approximation. \[Sheth & Tormen (1999)\] obtained an analytic fit to the numerical mass function in the SCDM, OCDM and $`\mathrm{\Lambda }`$CDM models and then \[Sheth, Mo & Tormen (1999)\] provided a semianalytic derivation of the formula assuming an anisotropic collapse an incorporating some nonlocal effects. All but one models mentioned above assumed that the formation of a gravitationally bound object is related to the collapse along three axes. Only \[Monaco (1995)\] assumed the collapse condition corresponding to the collapse along only the first axis. In this talk I briefly review the Press-Schechter formalism. Then I describe the derivation of the mass function ($`\lambda _3`$-function) in the Zel’dovich approximation. I compare the result with the standard Press-Schechter model and the model suggested by \[Sheth, Mo & Tormen (1999)\]. I briefly discuss the results of comparison of the $`\lambda _3`$-function with N-body simulations. Finally, I discuss the effect of the initial gravitational potential on the cosmological mass function and show that the clusters have a significant tendency to form in the troughs of the initial gravitational potential. ## 2. The Press-Schechter Formalism The mass function $`n(M)`$ is the number density of gravitationally bound clumps with masses between $`M`$ and $`M+dM`$. Let $`F(>M)`$ be the fraction of the mass contained in the gravitationally bound objects with masses greater than $`M`$. Press and Schechter (1974) suggested the fraction $`F(>M)`$ and the mass function $`n(M)`$ can be related as $$n(M)=\frac{\overline{\rho }}{M}\frac{F}{M},$$ (1) where $`\overline{\rho }`$ is the mean mass density in the universe and the minus sign reflects the fact that $`F`$ is a decreasing function of $`M`$. Press and Schechter also made the assumption that the fraction of mass $`F(>M)`$ can be estimated as a fraction of mass $`\mathrm{\Psi }(\delta _M>\delta _c)`$ in the initial density field filtered with the window function $`W`$ (corresponding to $`\stackrel{~}{W}`$ in $`k`$-space) $$\delta _M(𝐱,t)=D(t)\delta _{in}(𝐱^{})W(\|𝐱^{}𝐱\|/R)d^3x^{}$$ (2) satisfying the collapse condition $`\delta _M>\delta _c`$. Here $`\delta =(\rho \overline{\rho })/\overline{\rho }`$ is the density contrast, $`D(t)`$ is linear growth factor, $`𝐱`$ is the comoving coordinate. The mass $`M`$ and the linear scale $`R`$ of the filter are related as $$M=f_W\frac{4\pi }{3}R^3\overline{\rho },$$ (3) where $`f_W`$ is a factor depending on the shape of the smoothing filter $`W`$. For a sharp k-space filter adopted here $`f_W=9\pi /214.1`$ and thus $`M=6\pi ^2R^3\overline{\rho }`$ (see e.g. \[Lacey & Cole (1994)\]). Assuming that the initial density contrast is Gaussian random field its pdf (probability distribution function) is $$p(\delta _M)=\frac{1}{\sqrt{2\pi }\sigma _M}\mathrm{exp}\left[\frac{\delta _M^2}{2\sigma _M^2}\right],$$ (4) where the variance $`\sigma _M^2`$ is a function of mass $`M`$ $$\sigma _M^2=\frac{d^3k}{(2\pi )^3}P(k)\stackrel{~}{W}^2(kR),$$ (5) where $`P(k)=\|\delta _k\|^2`$ is the initial spectrum of perturbations and $`\stackrel{~}{W}(kR)`$ is the window function in the k-space. Press and Schechter argued that a fluid element becomes a part of a gravitationally bound object of mass $`M`$ when its linearly extrapolated density contrast $`\delta _M`$ reaches the critical value $`\delta _c=3/20(12\pi )^{2/3}1.69`$. This corresponds to the collapse of the top-hat spherical perturbation having the initial density contrast similar to the fluid element in question. The collapse of a spherical top-hat model has been assumed approximately to correspond the virialization of the clump. Recently Shapiro et al. (1999) showed that the virialization corresponds to $`\delta _c1.52`$ rather than to $`\delta _c1.69`$. The fraction of mass satisfying the collapse condition on scale $`M`$ is $`\mathrm{\Psi }(\delta _M>\delta _c)`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2\pi }\sigma (M)}}{\displaystyle _{\delta _c}^{\mathrm{}}}\mathrm{exp}\left[{\displaystyle \frac{\delta _M^2}{2\sigma ^2(M)}}\right]𝑑\delta _M`$ (6) $`=`$ $`{\displaystyle \frac{1}{2}}\mathrm{erfc}\left[{\displaystyle \frac{\delta _c}{\sqrt{2}\sigma (M)}}\right],`$ where erfc$`(x)`$ is the complementary error function. Assuming that $`F(M)\mathrm{\Psi }(\delta _M>\delta _c)`$ one easily obtains the mass function $`n(M)`$ (eq. 1 ). One obvious problem with this result is that the normalization integral $$_0^{\mathrm{}}𝑑F\mathrm{\Psi }(\delta _{M=\mathrm{}}>\delta _c)=\frac{1}{2}$$ (7) meaning that only a half of the mass is contained in the gravitationally bound clumps. Press and Schechter renormalized $`n(M)`$ by introducing an additional factor of 2 ($`F(M)=2\mathrm{\Psi }(\delta _M>\delta _c`$ ) $`n_{ps}(M)`$ $`=`$ $`2{\displaystyle \frac{\overline{\rho }}{M}}{\displaystyle \frac{\mathrm{\Psi }}{M}}=2{\displaystyle \frac{\overline{\rho }}{M}}{\displaystyle \frac{d\sigma }{dM}}{\displaystyle \frac{\mathrm{\Psi }}{\sigma }}`$ (8) $`=`$ $`\sqrt{{\displaystyle \frac{2}{\pi }}}{\displaystyle \frac{\overline{\rho }}{M}}{\displaystyle \frac{d\sigma }{dM}}{\displaystyle \frac{\delta _c}{\sigma ^2(M)}}\mathrm{exp}\left[{\displaystyle \frac{\delta _c^2}{2\sigma _M^2}}\right].`$ Later, the normalization problem was correctly resolved in the frame of the excursion set model (\[Peacock & Heavens 1990\] and \[Bond et al. 1991\]). The derivation of Press and Schechter did not take into account the so called cloud-in-cloud problem. Function $`\mathrm{\Psi }(\delta _M>\delta _c)`$ in eq.6 gives the fraction of mass that satisfies the collapse condition at the filtering scale $`M`$. However, some of the fluid particles may satisfy the collapse condition at larger filtering scales. In the correct model the fluid elements must be assigned to the clumps of mass $`M_1`$ being equal to the largest filtering mass at which the collapse condition is fulfilled. In the excursion set formalism this corresponds to the first crossing of the collapse threshold $`\delta _c`$ while $`\delta `$ evolves with the growth of $`\sigma `$. An elegant method to normalize the mass function was suggested by \[Jedamzik (1995)\] (see also the discussion in \[Yano et al. (1996)\]) who derived the integral equation $$\mathrm{\Psi }(\delta _M>\delta _c)=_M^{\mathrm{}}𝑑M_1n(M_1)\frac{M_1}{\overline{\rho }}P(M,M_1).$$ (9) that relates the fraction of the fluid elements firstly crossed the collapse threshold at the filtering scale $`M_1`$ ( $`dM_1n(M_1)M_1/\overline{\rho }`$ ) and the fraction of mass satisfying the collapse condition at filtering scale $`M`$ ( $`\mathrm{\Psi }(\delta _M>\delta _c)`$ ). Function $`P(M,M_1)`$ is the probability that a fluid particle firstly crossed the collapse threshold at the scale $`M_1`$ satisfies the collapse condition at the scale $`M`$. In the case of the sharp k-space filter and Gaussian $`\delta _M`$ this probability is exactly equal to $`1/2`$ for all $`M_1>M`$. Thus, the integral equation (9) can be immediately solved for the mass function $`n(M)`$. The solution is the correctly normalized mass function of eq. 8. ## 3. Mass Function in the Zel’dovich Approximation The simplest theory describing the anisotropic character of the gravitational collapse in a generic case of random initial condition is the Zel’dovich approximation (\[Zel’dovich 1970\], see also \[Shandarin & Zel’dovich 1989\] for a discussion). In particular, the Zel’dovich approximation provides a formula for an anisotropic collapse of a fluid element $$\rho (𝐪,t)=\frac{\overline{\rho }}{[1D(t)\lambda _1(𝐪)][1D(t)\lambda _2(𝐪)][1D(t)\lambda _3(𝐪)]},$$ (10) where $`D(t)`$ is the linear growth function and $`\lambda _1(𝐪),\lambda _2(𝐪)`$ and $`\lambda _3(𝐪)`$ are the eigenvalues of the initial deformation tensor. Using the ordering convention $`\lambda _1(𝐪)>\lambda _2(𝐪)`$ and $`\lambda _2(𝐪)>\lambda _3(𝐪)`$ the condition $`1D(t)\lambda _1(𝐪)=0`$ has been interpreted as a collapse of a fluid particle along one principle axis (\[Zel’dovich 1970\]). Similarly the conditions $`1D(t)\lambda _i(𝐪)=0`$ ($`i=2,3`$) can be interpreted as collapses along the second and third principle axes. \[Shandarin & Klypin (1984)\] showed that the formation of gravitationally bound clumps was the best correlated with the maxima of the smallest eigenvalue ($`\lambda _3`$) of the initial deformation tensor. Although the formation of the clumps may be also related to other pointlike singularities (\[Arnol’d et al. 1982\]) here we assume that a fluid particle becomes a part of a gravitationally bound clump of mass $`M`$ when its smallest eigenvalue $`\lambda _3`$ reaches the critical value $`\lambda _c`$ at the largest filtering scale $`M`$ (Lee & Shandarin 1998a). The Zel’dovich approximation (eq.10) predicts that the collapse condition is $`\lambda _c=1`$ (it is assumed that $`D(t)`$ normalized to $`D(t_0)=1`$, where $`t_0`$ is the present time). However, because of multistreaming effect all fluid particles (except the set of measure zero) enter the multi-stream flow regions before they collapse. We approximately incorporate this complex effect by reducing the threshold $`\lambda _c`$ to a smaller value. The comparison with the Press-Schechter mass function as well as with the numerical mass function suggests that $`\lambda _c=0.37`$ is a reasonable choice. The derivation of the mass function in the Zel’dovich approximation is similar to the Press-Schechter derivation except that the collapse condition is $`\lambda _3(M)=\lambda _c`$ instead of $`\delta _M=\delta _c`$. \[Doroshkevich (1970)\] derived the joint pdf of three eigenvalues $$p(\lambda _1,\lambda _2,\lambda _3)=\frac{3375}{8\sqrt{5}\pi \sigma ^6}\mathrm{exp}\left(\frac{3I_1}{\sigma ^2}+\frac{15I_2}{2\sigma ^2}\right)(\lambda _1\lambda _2)(\lambda _2\lambda _3)(\lambda _1\lambda _3),$$ (11) where $`I_1=\lambda _1+\lambda _2+\lambda _3`$, $`I_2=\lambda _1\lambda _2+\lambda _2\lambda _3+\lambda _3\lambda _1`$ and $`\sigma ^2`$ is the density contrast variance as defined in eq. 5. Integrating $`p(\lambda _1,\lambda _2,\lambda _3)`$ over two eigenvalues one can obtain the pdf of one of the eigenvalues. We are interested in the collapse along the third axis and therefore $`p(\lambda _3)`$ is of primary interest $`p(\lambda _3)`$ $`=`$ $`{\displaystyle \frac{\sqrt{5}}{12\pi \sigma }}\{3\sqrt{3\pi }\mathrm{exp}({\displaystyle \frac{15\lambda _3^2}{4\sigma ^2}})\mathrm{erfc}\left({\displaystyle \frac{\sqrt{3}\lambda _3}{2\sigma }}\right)`$ (12) $`+`$ $`\sqrt{2\pi }\left(20{\displaystyle \frac{\lambda _3^2}{\sigma ^2}}1\right)\mathrm{exp}\left({\displaystyle \frac{5\lambda _3^2}{2\sigma ^2}}\right)\mathrm{erfc}\left(\sqrt{2}{\displaystyle \frac{\lambda _3}{\sigma }}\right)`$ $``$ $`20{\displaystyle \frac{\lambda _3}{\sigma }}\mathrm{exp}({\displaystyle \frac{9\lambda _3^2}{2\sigma ^2}})\}.`$ Repeating the derivation of the previous section using the pdf of eq.12 instead of eq.4 one arrives to an analog of the normalization integral equation (eq.9) $$\mathrm{\Psi }(\lambda _3(M)>\lambda _c)=_M^{\mathrm{}}𝑑M_1n(M_1)\frac{M_1}{\overline{\rho }}P(M,M_1),$$ (13) here $`\mathrm{\Psi }(\lambda _3(M)>\lambda _c)`$ is the fraction of mass where $`\lambda _3(M)>\lambda _c`$ on the filter scale $`M`$. Solving exactly eq.13 is much more difficult than eq.9 because now $`P(M,M_1)`$ is not a constant. In the limit $`M_1MM`$ the probability $`P=0.5`$ as in eq.9 but in the limit $`M_1M`$ it drops to $`P=0.08`$. \[Lee & Shandarin (1998a)\] used the limiting value $`P=0.08`$ and analytically derived the mass function in the Zel’dovich approximation $`n(M)`$ $`=`$ $`{\displaystyle \frac{25\sqrt{5}}{24\pi }}{\displaystyle \frac{\overline{\rho }}{M}}{\displaystyle \frac{d\sigma }{dM}}{\displaystyle \frac{\lambda _{3c}}{\sigma _M^2}}`$ (14) $`\{`$ $`3\sqrt{3\pi }\mathrm{exp}\left({\displaystyle \frac{15\lambda _3^2}{4\sigma ^2}}\right)\mathrm{erfc}\left({\displaystyle \frac{\sqrt{3}\lambda _3}{2\sigma }}\right)`$ $`+`$ $`\sqrt{2\pi }\left(20{\displaystyle \frac{\lambda _3^2}{\sigma ^2}}1\right)\mathrm{exp}\left({\displaystyle \frac{5\lambda _3^2}{2\sigma ^2}}\right)\mathrm{erfc}\left(\sqrt{2}{\displaystyle \frac{\lambda _3}{\sigma }}\right)`$ $``$ $`20{\displaystyle \frac{\lambda _3}{\sigma }}\mathrm{exp}({\displaystyle \frac{9\lambda _3^2}{2\sigma ^2}})\}.`$ In the following sections I compare the obtained result with the Press-Schechter and Sheth-Mo-Tormen mass functions as well as with numerical simulations. ## 4. Comparison of Three Analytic Mass Functions Both the Press-Schechter and $`\lambda _3`$-mass functions have a common factor $`\frac{\overline{\rho }}{M}\frac{d\sigma }{dM}`$ which depends on the initial spectrum and $`f(\sigma )F/\sigma `$ that completely characterizes a model. Thus, comparing different models is convenient by comparing $`f(\sigma )`$ as functions of $`\sigma `$. The Press-Schechter and $`\lambda _3`$-functions are $`f_{PS}(\sigma )`$ $`=`$ $`\sqrt{{\displaystyle \frac{2}{\pi }}}{\displaystyle \frac{\delta _c}{\sigma ^2}}\mathrm{exp}\left({\displaystyle \frac{\delta _c^2}{2\sigma ^2}}\right),`$ $`f_{\lambda _3}(\sigma )`$ $`=`$ $`{\displaystyle \frac{25\sqrt{5}}{24\pi }}{\displaystyle \frac{\lambda _{3c}}{\sigma ^2}}\{3\sqrt{3\pi }\mathrm{exp}({\displaystyle \frac{15\lambda _3^2}{4\sigma ^2}})\mathrm{erfc}\left({\displaystyle \frac{\sqrt{3}\lambda _3}{2\sigma }}\right)`$ (15) $`+`$ $`\sqrt{2\pi }\left(20{\displaystyle \frac{\lambda _3^2}{\sigma ^2}}1\right)\mathrm{exp}\left({\displaystyle \frac{5\lambda _3^2}{2\sigma ^2}}\right)\mathrm{erfc}\left(\sqrt{2}{\displaystyle \frac{\lambda _3}{\sigma }}\right)`$ $``$ $`20{\displaystyle \frac{\lambda _3}{\sigma }}\mathrm{exp}({\displaystyle \frac{9\lambda _3^2}{2\sigma ^2}})\}.`$ Sheth and Tormen (1999) and Sheth, Mo and Tormen (1999) derived a new mass function that fits better the results of the N-body simulations $$f_{SMT}(\sigma )=A\left[1+\left(\frac{a\delta _c^2}{\sigma ^2}\right)^q\right]\sqrt{\frac{2}{\pi }}\frac{\sqrt{a}\delta _c}{\sigma ^2}\mathrm{exp}\left(\frac{a\delta _c^2}{\sigma ^2}\right),$$ (16) here $`a=0.707`$, $`q=0.3`$ and the constant $`A=0.322`$ found from the normalization condition $$_0^{\mathrm{}}f(\sigma )𝑑\sigma =1.$$ (17) Choosing $`a=1`$, $`q=0`$ and the constant $`A=1`$ one obtains the Press-Schechter function. These three functions are shown in Fig. 1a. The small box shows the range of $`\sigma `$ where the theoretical mass functions were checked against N-body simulations by \[Sheth & Tormen (1999)\] and \[Lee & Shandarin (1999)\]. Fig. 1b shows the ratios $`f_{PS}/f_{SMT}`$ and $`f_{\lambda _3}/f_{SMT}`$. ## 5. Comparison with N-body Simulations Figure 2 shows the comparison of the $`\lambda _3`$-function with the numerical mass functions for the scale invariant initial spectra: $`P(k)k^n`$ with $`n=1`$ and $`n=0`$ (see for the details \[Tormen (1998)\]). The top panel ($`n=1`$) shows a quite good agreement of the $`\lambda _3`$-function with the numerical mass function, while in the $`n=0`$ case the agreement is significantly worse. Fig. 3 shows the comparison of the $`\lambda _3`$-function with the N-body simulation of the SCDM model (\[Governato et al. (1999)\]). At four epochs ($`z=1.86,1.14,0.43`$, and $`0`$) the $`\lambda _3`$-function is in a better agreement that the Press-Schechter mass function. A similar result has been reported by \[Sheth & Tormen (1999)\]. ## 6. Large-Scale Biasing It has been noticed for sometime that the initial gravitational potential may noticeably affect the large scale structure. Kofman and Shandarin (1988) showed that the adhesion approximation predicts that the formation of voids is associated with positive peaks of the primordial gravitational potential. Sahni et al. (1994) studied the effect and measured a significant correlation between the sizes of voids and the value of primordial gravitational potential in numerical simulations of the adhesion model. Recently, Madsen et al. (1998) have demonstrated by N-body simulations that the underdense and the overdense regions are closely linked to the regions with the positive and the negative gravitational potential respectively. \[Lee & Shandarin (1998b)\] showed that the initial potential also affects the masses of clusters. In order to incorporate the primordial gravitational potential fluctuations term into the derivation of the mass function, we first derive the conditional probability density distribution $`p(\delta \|\phi <\phi _c)`$ ($`\phi _c>0`$): $`p(\delta \|\phi <\phi _c)`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2\pi }\sigma _\delta }}\mathrm{exp}({\displaystyle \frac{\delta ^2}{2\sigma _\delta ^2}})[1\mathrm{erf}\left({\displaystyle \frac{\phi _c}{\sqrt{2}\sigma _\phi }}\right)]^1\times `$ (18) $`\left[1+\mathrm{erf}\left({\displaystyle \frac{\kappa \frac{\delta }{\sigma _\delta }\frac{\phi _c}{\sigma _\phi }}{\sqrt{2(1\kappa ^2)}}}\right)\right].`$ Here $`\sigma _\delta ^2`$, $`\sigma _v^2`$, and $`\sigma _\phi ^2`$ are the density, velocity and the potential variances respectively; $`\kappa =<\delta \phi >/\sigma _\delta \sigma _\phi =\sigma _v^2/\sigma _\delta \sigma _\phi `$ is the crosscorrelation coefficient of the the density contrast $`\delta `$ smoothed on the scale $`k_c`$ and the primordial (unsmoothed) potential fluctuations $`\phi `$. As a result, eq.1 for the conditional mass function $`n(M\|\phi <\phi _c)`$ becomes $$n(M\|\phi <\phi _c)=\frac{\overline{\rho }}{M}\left(\frac{F}{\sigma _\delta }\frac{d\sigma _\delta }{dM}+\frac{F}{\sigma _v}\frac{d\sigma _v}{dM}\right).$$ (19) The further calculation needs to be done numerically. Fig. 4 illustrates how the mass function depends on the initial potential in the CDM model with $`\mathrm{\Gamma }=\mathrm{\Omega }h=0.25`$ normalized to $`\sigma _8=1`$. The top panel shows the mass function for regions of positive and negative initial potential as well as unconditional mass function. The bottom panel show the ratio of conditional mass functions to unconditional one. We also calculate the probability that a clump with mass M is located in the potential regions satisfying the chosen condition, for instance, $`\phi <\phi _c`$ $$P(\phi <\phi _c\|M)=\frac{n(M\|\phi <\phi _c)}{n(M)}P(\phi <\phi _c),$$ (20) where $`P(\phi <\phi _c)`$ is the fraction of space satisfying the given condition (see Fig. 5). The scale of the initial potential $$R_\phi =\sqrt{3}\sigma _\phi /\sigma _\phi ^{}=\sqrt{3\frac{_{k_l}^{\mathrm{}}𝑑kk^2P(k)}{_0^{\mathrm{}}𝑑kP(k)}}120h^1\mathrm{Mpc}$$ (21) does not depend on any ad hoc scale; the dependence on $`k_l`$ is extremely weak ($`\sqrt{ln(1/k_l)}`$ for the Harrison-Zel’dovich spectra assumed here). The geometry of the gravitational potential does not evolve much on large scales (\[Kofman & Shandarin (1988)\], \[Madsen et al. (1998)\]). Therefore, the potential at present is very similar to the primordial one on scales much greater than the scale of nonlinearity. A simple explanation to this in the frame of the standard scenario of the structure formation is due to the fact that the mass has been displaced by the distance about $`10h^1\mathrm{Mpc}`$ (\[Shandarin 1993\]). Therefore, the potential on scales greater than, say, $`30h^1\mathrm{Mpc}`$ has been changed very little. For the model in question the scale of the primordial potential is found to be $`R_\phi 120h^1\mathrm{Mpc}`$. The scale of the density contrast field reaches this value $`R_\delta =\sqrt{3}\sigma _\delta /\sigma _\delta ^{}120h^1\mathrm{Mpc}`$ only after it is smoothed on $`k_c0.017h\mathrm{Mpc}^1`$. The corresponding density variance on this scale is $`\sigma _\delta (0.017h\mathrm{Mpc}^1)0.03`$. On the other hand, the number of clumps with masses $`10^{14}10^{15}h^1M_{}`$ can easily be 30% greater in the troughs of the potential than the mean density $`n(>M)=0.5[n(>M\|\phi <0)+n(>M\|\phi >0)]`$ (see Fig. 4). Thus, the bias factor $`b`$ (defined by the relation $`\mathrm{\Delta }n_{cl}/n_{cl}=b\mathrm{\Delta }\rho _m/\rho _m`$) reaches at least $`10`$ on the scale about $`120h^1\mathrm{Mpc}`$. Figure 5 demonstrates that the most massive clusters ($`M>10^{15}h^1M_{}`$) are almost certainly located in the the troughs in the initial potential. The bias defined as the density contrast of the clusters with respect to the mass density contrast $`b=\delta _{cl}/\delta \rho `$ reaches the value $`310`$ on the scale of the potential $`R_\varphi 120h^1Mpc`$ (Lee & Shandarin 1998b). ## 7. Summary In the talk I discussed new modifications of the Press-Schechter theory of the cosmological mass function. One assumes a different collapse condition that implies that a fluid particle becomes a part of a gravitationally bound object after it experiences collapses along three axes. The comparison with other models (Fig. 1) shows that it predicts about 25% more gravitationally bound clumps than the Sheth-Mo-Tormen model in the range $`0.45\sigma 3.1`$ where the comparison with the N-body simulations has been done. A direct comparison with the N-body simulations (Fig. 2 and 3) shows a quite good agreement although not as good as the Sheth-Mo-Tormen model. The $`\lambda _3`$-function based on the Zel’dovich approximation has been obtained analytically similarly to the Press-Schechter function. A drawback of the derivation is a quite crude approximation of the probability function $`P(M,M_1)`$ in the normalization integral eq.13. A more accurate normalization will be reported separately. The Sheth-Mo-Tormen model also suffer from a normalization problem: the shape of the mass function has been derived but the normalization has been enforced by demanding equality of eq.17 Another modification is the conditional mass function showing that the clusters of galaxies tend to form in the troughs of the initial gravitational potential and avoid the peaks of the potential. The gravitational potential field has a typical scale of about $`120h^1Mpc`$ and as a result has an advantage of being independent of the arbitrariness of the smoothing scale (if the filter scale is smaller than roughly $`50h^1Mpc`$) and at present it has almost same geometry as at the epoch of decoupling. Figures 4 and 5 quantify this large-scale biasing. ## 8. Acknowledgments I am grateful to Ravi Sheth for useful discussions during the workshop. This work has been partly supported by the University of Kansas GRF 2001 grant. Figure Captions Fig. 1. (a) The fraction of mass $`f=dF/d\sigma `$ in the gravitationally bound objects as a function of $`\sigma `$ as predicted by the Press-Schechter model (short dashed line), Sheth-Mo-Tormen model (solid line), and $`\lambda _3`$-model (long dashed line). The small box shows the range of $`\sigma `$ where the models have been checked against N-body simulations (see Fig. 2 in \[Sheth, Mo & Tormen (1999)\]) (b) The logarithm of the ratios $`f_{PS}/f_{SMT}`$ (short-dashed line) and $`f_{\lambda _3}/f_{SMT}`$ (long-dashed line). Fig. 2. The square dots represent the numerical mass function with poissonian error bars. The solid line is the $`\lambda _3`$-mass function with $`\lambda _{3c}=0.37`$ while the dashed, the dotted lines are the PS mass functions with $`\delta _c=1.69,1.5`$ respectively. The upper and the lower panels correspond to the $`n=1`$ and the $`n=0`$ power-law models respectively. See also the top left panel of Fig.2 in \[Tormen (1998)\]. Fig. 3. The square dots represent the numerical data for the case of SCDM model with $`\mathrm{\Omega }=1,h=0.5`$. The solid line is our mass function with $`\lambda _{3c}=0.37`$, and the dashed, the dotted lines are the PS mass functions with $`\delta _c=1.69,1.5`$ respectively. Fig. 4. In the upper panel the conditional cumulative mass function satisfying chosen potential condition is plotted. The solid, the long dashed, the dot-dashed, and the dashed lines correspond to the conditions $`\phi <\sigma _\phi `$, $`\phi <0`$, $`\phi >0`$, and $`\phi >\sigma _\phi `$ respectively, while the dotted line represents the unconditional cumulative PS mass function. The shaded area is $`1\sigma `$ fit to the observational cumulative mass function of rich clusters by Bahcall and Cen (1993). In the lower panel the ratio of the conditional cumulative mass functions to the unconditional one is plotted for each condition. The CDM spectrum with $`\mathrm{\Gamma }=0.25`$ normalized to $`\sigma _8=1`$ has been used. Fig. 5. The probability that a clump with mass M can be found in the regions satisfying chosen potential condition is plotted. The heavy solid, the heavy dashed, the solid, the dashed, the long dashed, and the dot-dashed lines correspond to the condition $`\phi <0`$, $`\phi >0`$, $`\phi <\sigma _\phi `$, $`\sigma _\phi <\phi <0`$, $`0<\phi <\sigma _\phi `$, and $`\phi >\sigma _\phi `$ respectively.
warning/0007/cond-mat0007501.html	ar5iv	text	# Ferromagnetic ordering in a generalized Hubbard model: weak intra-atomic interaction limit ## 1 Introduction In spite of great attention payed to investigation of ferromagnetism in narrow energy bands -, where the same system of electrons is responsible both for conductivity and for magnetic ordering (the transition $`3d`$-metals, ferromagnetic sulphides, oxides and selenides with metallic type of conductivity) the problem still remains open . There are some experimental results for the disulphides of transition metals which can be interpreted on the basis of itinerant mechanism of ferromagnetism. In these compounds the dependence of Curie temperature on the number of magnetic moments is non-typical from point of view of exchange mechanism of feromagnetism. In $`Fe_xCo_{1x}S_2`$ where electron concentration $`n`$ changes from 0 to 1 in doubly degenerate $`e_g`$ -subband the Curie temperature increases with decrease of electron concentration at $`0.75<n<1`$. Although the consistent theory of ferromagnetic ordering in transition metal compounds can be constructed only in the model including the orbital degeneracy of the band, the qualitative character of observed properties can be interpreted in the framework of generalized model of non-degenerate band . ## 2 The Hamiltonian The simplest model for description of the magnetic properties of narrow-band materials is the Hubbard model , but this model contains only diagonal matrix elements of Coulomb interaction in site representation, that not give rise to metallic ferromagnetism except in special situations such as a single hole in a half-filled band or a special lattice geometry . The importance of non-diagonal matrix elements of Coulomb interaction was pointed out in papers . We write the Hamiltonian of system of $`s`$-electrons in the following form $`H=`$ $``$ $`\mu {\displaystyle \underset{i\sigma }{}}a_{i\sigma }^+a_{i\sigma }+{\displaystyle \underset{ij\sigma }{}^{}}t_{ij}(n)a_{i\sigma }^+a_{j\sigma }++{\displaystyle \underset{ij\sigma }{}^{}}(T_2(ij)a_{i\sigma }^+a_{j\sigma }n_{i\overline{\sigma }}+h.c.)`$ $`+`$ $`U{\displaystyle \underset{i}{}}n_in_i+{\displaystyle \frac{1}{2}}{\displaystyle \underset{ij\sigma \sigma ^{^{}}}{}^{}}Ja_{i\sigma }^+a_{j\sigma ^{}}^+a_{i\sigma ^{}}a_{j\sigma },`$ where $`a_{i\sigma }^+`$, $`a_{i\sigma }`$ – creation and destruction operators of electron on site $`i`$, $`\sigma =,`$, $`n_{i\sigma }=a_{i\sigma }^+a_{i\sigma }`$, $`n=n_i+n_i`$, $`\mu `$ – chemical potential, $`t_{ij}(n)=t(ij)+T_1(ij)`$ is the effective hopping integral between nearest neighbours, $`t_{ij}`$ is the hopping integral of an electron from site $`j`$ to site $`i`$, $`T_1(ij)`$ and $`T_2(ij)`$ are the parameters of correlated hopping of electron, $`U`$ is the intraatomic Coulomb repulsion and $`J`$ is the exchange integral for the nearest neighbours. The prime by second sum in Eq. (2) signifies that $`ij`$. The peculiarities of the model described by the Hamiltonian (2) is taking into consideration the influence of sites occupation on the electron hoppings (correlated hopping), and the exchange integral. The correlated hopping, firstly, renormalize the initial hopping integral (it becomes concentration- and spin-dependent) and, secondly, lead to an independent on quasiimpulse shift of the band center, dependent on magnetic ordering. Taking into account the quantity of second order $`J`$ is on principle necessary to describe ferromagnetism in this model . In this paper we do not take into account the inter-atomic Coulomb interaction, which is important in the charge ordering state (this ordering type go beyond this article). To characterize the value of correlated hopping we introduce dimensionless parameters $`\tau _1=\frac{T_1(ij)}{\|t_{ij}\|}`$, $`\tau _2=\frac{T_2(ij)}{\|t_{ij}\|}`$. ## 3 Weak intra-atomic interaction To simplify the consideration we use model Hamiltonian (2). If intra-atomic Coulomb interaction is weak ($`U<\|t_{ij}(n)\|`$) then we can take into account the electron-electron interaction in the Hartree-Fock approximation: $`n_in_i=n_{}n_i+n_{}n_i,`$ (2) $`a_{i\sigma }^+n_{i\overline{\sigma }}a_{j\sigma }=n_{\overline{\sigma }}a_{i\sigma }^+a_{j\sigma }+a_{i\sigma }^+a_{j\sigma }n_{i\overline{\sigma }},`$ where the average values $`n_{i\sigma }=n_\sigma `$ are independant of site number (we suppose that distributions of electron charge and magnetic momentum are homogenous). Taking into account (2) we can write Hamiltonian (2) in the following form: $`H={\displaystyle \underset{ij\sigma }{\overset{}{}}}ϵ_\sigma (ij)a_{i\sigma }^+a_{j\sigma },`$ (3) where $`ϵ_\sigma (ij)=\mu +\beta _\sigma +n_{\overline{\sigma }}U+zn_\sigma J+t_{ij}(n\sigma );`$ (4) $`\beta _\sigma ={\displaystyle \frac{2}{N}}{\displaystyle \underset{ij}{}}T_2(ij)a_{i\overline{\sigma }}^+a_{j\overline{\sigma }},`$ (5) $`t_{ij}(n\sigma )=t_{ij}(n)+2n_{\overline{\sigma }}T_2(ij).`$ (6) The dependences of effective hopping integral on electron concentration and magnetization, a being of the spin-dependent displasement of band center are the essential distinction of single-particle energy spectrum in the model described by Hamiltonian (3) from the spectrum in the Hubbard model for weak interaction. An use of (3) allows, in particular, to explain the peculiarities of dependence of binding energy on atomic number in transition metals and also essentially modifies theory of ferromagnetism in a collective electron model. ## 4 Ferromagnetic ordering in weak intra-atomic interaction limit After the transition to Fourier representation we obtain for the Green function $`a_{p\sigma }\|a_{p^{}\sigma }^+_𝐤={\displaystyle \frac{1}{2\pi }}{\displaystyle \frac{1}{Eϵ_\sigma (𝐤)}}`$ (7) where the energy spectrum is $`ϵ_\sigma (𝐤)=\mu +\beta _\sigma +n_{\overline{\sigma }}Uzn_\sigma J+t(n\sigma )\gamma (𝐤);`$ (8) here the spin-dependent shift of the band center is $`\beta _\sigma ={\displaystyle \frac{2}{N}}{\displaystyle \underset{ij}{}}T_2(ij)a_{i\overline{\sigma }}^+a_{j\overline{\sigma }},`$ (9) $`\gamma (𝐤)=\underset{𝐑}{}e^{i\mathrm{𝐤𝐑}}`$, the spin and concentration dependent hopping integral is $`t(n\sigma )=t(n)+2n_{\overline{\sigma }}T_2.`$ (10) The concentration of electrons with spin $`\sigma `$ is $`n_\sigma ={\displaystyle \underset{\mathrm{}}{\overset{+\mathrm{}}{}}}\rho (ϵ)f(ϵ)𝑑ϵ.`$ (11) Here $`\rho (ϵ)`$ is the density of states, $`f(ϵ)`$ is the Fermi distribution function. Let us assume the rectangular density of states: $`\rho (ϵ)={\displaystyle \frac{1}{N}}{\displaystyle \underset{𝐤}{}}\delta (ϵϵ(𝐤))={\displaystyle \frac{1}{2w}}\theta (ϵ^2w^2)`$ (12) Then in the case of zero temperature we obtain: $`n_\sigma ={\displaystyle \underset{\mathrm{}}{\overset{+\mathrm{}}{}}}\rho (ϵ)\theta (E_\sigma (ϵ))𝑑ϵ={\displaystyle \frac{1}{2w}}{\displaystyle \underset{w}{\overset{w}{}}}\theta (E_\sigma (ϵ))𝑑ϵ={\displaystyle \frac{1}{2w}}{\displaystyle \underset{w}{\overset{ϵ_\sigma }{}}}𝑑ϵ,`$ (13) where $`ϵ_\sigma =w(2n_\sigma 1)`$. The value $`ϵ_\sigma `$ is the solution of the equation $`E_\sigma (ϵ)=0`$ from which we obtain $`ϵ_\sigma =\frac{\mu _\sigma }{\alpha _\sigma }`$, where $`\mu _\sigma =\mu \beta _\sigma +zn_\sigma Jn_{\overline{\sigma }}U`$ and $`\alpha _\sigma =12\tau _2n_{\overline{\sigma }}`$. The system parameters are related by the equation $`m(zJ+u)+\beta _{}\beta _{}=2m(1\tau _2).`$ (14) The shift of band center is obtained from $`\beta _\sigma ={\displaystyle \frac{2}{N}}{\displaystyle \underset{ij}{}}T_2(ij)a_{i\overline{\sigma }}^+a_{j\overline{\sigma }}={\displaystyle \frac{\tau _2}{2w}}{\displaystyle \underset{w}{\overset{ϵ_\sigma }{}}}ϵ𝑑ϵ=\tau _2wn_{\overline{\sigma }}(n_{\overline{\sigma }}1).`$ (15) One can see that $`\beta _{}\beta _{}=2\tau _2mw(1n).`$ (16) From the equation (14) one can see that in less than half-filled band correlated hopping leads to the stabilisation of ferromagnetism as well as inter-atomic exchange interaction and intra-atomic Coulomb interaction; the larger is electron concentration $`n`$ the smaller is the critical value of exchange integral for occurence of ferromagnetism. These our results are in accordance with the results of paper . From (14) and (16) we obtain the condition of feromagnetic ordering realization: $`{\displaystyle \frac{zJ+U}{2w}}+\tau _2(2n)>1.`$ (17) Here we use the notation $`w=w(n)=z\|t_0\|(1\tau _1n)`$ (18) where $`t_0`$ is the band hopping integral. In the case of non-zero temperatures the equation for magnetization is $`\mathrm{exp}\left({\displaystyle \frac{mJ_{eff}}{\theta }}\right)={\displaystyle \frac{\mathrm{sinh}\left((1n_{})\alpha _{}w/\theta \right)\mathrm{sinh}\left(n_{}\alpha _{}w/\theta \right)}{\mathrm{sinh}((1n_{})\alpha _{}w/\theta )\mathrm{sinh}(n_{}\alpha _{}w/\theta )}},`$ (19) where the effective exchange integral is $`J_{eff}=zJ+U+2zT_2(1n)`$ (20) The expression for Curie temperature can be obtained from equation (19) at $`m0`$. If $`\tau _1=\tau _2=0`$ then we obtain $`\theta _C={\displaystyle \frac{w_0}{2}}\mathrm{arcth}^1\left({\displaystyle \frac{2w_0}{zJ+U}}\right)`$ (21) If $`\tau _1=\tau _2=1/2`$ then $`\theta _C={\displaystyle \frac{w_0n(2n)(1\frac{n}{2})}{4}}\mathrm{arcth}^1\left({\displaystyle \frac{w_0(1\frac{n}{2})}{zJ+U+(1n)(1\frac{n}{2})w_0}}\right)`$ (22) It is important that at some values of parameters the Curie temperature can increase with the decrease of the electron concentration. ## 5 Discussion and Conclusions The analisys of obtained in this paper energy spectrum (8) shows that both intraatomic Coulomb repulsion $`U`$ and exchange integral $`J`$ favor spin polarization. Futhermore, taking into account the correlated hopping leads to the spin-dependent shift of the band center and to the band narrowing, what also give rise to ferromagnetism. These our results are in accordance with the results of paper . The peculiarities of the energy spectrum (8) lead to concentration dependence of the Curie temperature. In particular, the concentration dependent shift (16) of the band centers at $`n<1`$ is positive, at $`n>1`$ is negative. According to this fact the Curie temperature at decreasing $`n`$ from 1 can increase and at $`n>1`$ can decrease. The obtained dependence qualitaively agrees with the experimental data on Fe<sub>x</sub>Co<sub>1-x</sub>S<sub>2</sub> . The proposed approach can be extended to all values of $`n`$, and the peculiarities of ferromagnetic ordering in Co<sub>x</sub>Ni<sub>1-x</sub>S<sub>2</sub> where this concentration changes from 1 to 2 can be explained. In conclusion, taking into consideration both correlated hopping and inter-atomic exchange interaction essentially modify the $`s`$band model and favours the ferromagnetic ordering.
warning/0007/cond-mat0007014.html	ar5iv	text	# Entropy Burst from Parabolic Potentials ## Abstract The change of the energy of ground state is investigated in a thermodynamical process by using the model described by one-dimensional harmonic oscillator + two-dimensional isotropic parabolic potential barrier such as $`V(x,y,z)=m\omega ^2x^2/2m\gamma ^2(y^2+z^2)/2`$. In the process where two independent many-particle systems suddenly touch with each other, it is shown that the lowest energy after the interaction can possibly be smaller than that before the interaction and then the entropy burst can occur. In the previous paper it has been shown that even-dimensional parabolic potential barriers such as $`V(y,z)=m\gamma ^2(y^2+z^2)/2`$ have infinitely degenerate stationary-states with the zero real energy eigenvalue, which are described by stationary flows round the center of the potential. This result is really surprising, because in general the parabolic potential barriers have been considered as the model for unstable states \[2-7\]. In fact it is shown that one has only unstable states like resonances in one-dimensional parabolic potential barriers like $`V(x)=m\gamma ^2x^2/2`$, which are represented by the solutions of the conjugate space of Gel’fand triplet with imaginary energy eigenvalues $`i(n+1/2)\mathrm{}\gamma (n=0,1,2,\mathrm{})`$ \[2-7\]. The existence of such stationary states in even-dimensional is not only surprising fact but also suggests the existence of the stable many-body systems constructed from the stationary states. Following this line of consideration, we have shown that there exist stable many-body systems composed of the stationary flows in statistical mechanics including all states derived in the Gel’fand triplet formalism . An important difference appears in the fact that in the new statistical mechanics a new freedom arises from that of the imaginary part of energy eigenvalues which is only allowed in the extended spaces of the Gel’fand triplet. Actually the thermodynamical probability $`W`$ is expressed by the product of the thermodynamical probability $`W^{\mathrm{}}`$ with respect to the real part of the energy and that $`W^{\mathrm{}}`$ with respect to the imaginary part of the energy such that $`W=W^{\mathrm{}}W^{\mathrm{}}`$. Therefore the entropy which is defined by $`S=k_\mathrm{B}\mathrm{log}W`$ becomes the sum of the usual Boltzmann entropy $`S^{\mathrm{}}=k_\mathrm{B}\mathrm{log}W^{\mathrm{}}`$ induced from the freedom of the real part of the energy eigenvalues and a new one $`S^{\mathrm{}}=k_\mathrm{B}\mathrm{log}W^{\mathrm{}}`$ originating from that of the imaginary part as $`S=S^{\mathrm{}}+S^{\mathrm{}}`$. (In details, see ref. 9.) It is natural to expect that there exists some kind of entropy transfer between the two entropies $`S^{\mathrm{}}`$ and $`S^{\mathrm{}}`$. In this paper we shall study the entropy transfer in an explicit example by using a simple model which is described by one-dimensional harmonic oscillator (HO) \+ two-dimensional isotropic parabolic potential barrier (PPB) such as $$V(x,y,z)=\frac{1}{2}m\omega ^2x^2\frac{1}{2}m\gamma ^2(y^2+z^2),$$ (1) where $`m`$ is the mass of particle. Note that the HO part is necessary for confirming the entropy transfer, because the PPB cannot have any freedom of the real energy eigenvalues, that is, the real energy eigenvalue is fixed at 0. The energy eigenvalues of the potential (1) are given by $$E_{n_xn_yn_z}=\left(n_x+\frac{1}{2}\right)\mathrm{}\omega i(n_yn_z)\mathrm{}\gamma $$ (2) for the solutions involving the stationary states, where the first and the second terms, respectively, denote the contributions from the HO and the PPB and $`n_x,n_y,n_z=0,1,2,\mathrm{}`$. The stationary states appear for $`n_y=n_z`$, of which relation washes out the imaginary part from the eigenvalues . It is now transparent that the infinite degeneracy of the all states having the energy eigenvalues of (2) originates from the freedom of the imaginary eigenvalues. Note that all stationary states are described by stationary flows round the center of the potential . Actually the difference between the infinitely degenerate states is represented by the difference of the stationary flows. Note that complex velocity potentials $`W`$ which are well-known in hydrodynamics can be introduced and $`W=\pm \gamma (y+iz)^2/2`$ is obtained for the first few stationary eigenstates in the two-dimensional PPB $`V(y,z)=m\gamma ^2(y^2+z^2)/2`$ . Let us consider a system being in a thermal equilibrium, which is composed of $`N_1`$ particles trapped by the potential $$V_1=\underset{i=1}{\overset{N_1}{}}\left[\frac{1}{2}m\omega ^2x_i^2\frac{1}{2}m\gamma ^2(y_i^2+z_i^2)\right].$$ (3) The energy of the stationary ground state is uniquely determined by $$E_1=\frac{1}{2}N_1\mathrm{}\omega .$$ (4) This state, of course, has no degeneracy arising from the HO and then the Boltzmann entropy $`S^{\mathrm{}}=0`$. It, however, has a diverging entropy with respect to the new entropy $`S^{\mathrm{}}`$, since the infinite degeneracy arises from the PPB. It is noted here that the stable many-body systems in the PPB are described as the nets of stationary circular-flows, the two joints of which are not connected by the wavefunctions of the stationary eigenstates but they are connected by the stationary flows of the PPB . This fact means that one can directly look in the net-structure (texture) composed of the stationary flows, because the flows are basically observable in quantum mechanics. It is, therefore, understood that the systems have some kind of classical property that the components inside of the systems are observable. Thus one sees that the appearance of the entropy $`S^{\mathrm{}}`$ originates from this semiclassical property of the stable many-body systems in the even-dimensional PPB. Now let us study a thermal process which is described by the mixing of two independent many-particle systems written by the potential given in (3). We consider the following situation; the two independent systems are suddenly put on their interactive region, where one has the center at $`x=a`$ and $`y=z=b`$ and the other at $`x=a`$ and $`y=z=b`$, and then the mixing starts and they finally make a new stable system in a thermal equilibrium. Before the mixing the ground-state energy is uniquely determined by $$E_{\mathrm{bef}}=\frac{1}{2}(N_1+N_2)\mathrm{}\omega ,$$ (5) where $`N_1`$ and $`N_2`$ are, respectively, the number of the constituent particles of the system one and that of the other. After the mixing the two potentials being at the two centers have the effects on the all constituent particles and then the potential is written by $`V_\mathrm{T}={\displaystyle \underset{i=1}{\overset{N}{}}}\{`$ $`{\displaystyle \frac{1}{2}}m\omega ^2\left[(x_ia)^2+(x_i+a)^2\right]`$ $`{\displaystyle \frac{1}{2}}m\gamma ^2[(y_ib)^2+(z_ib)^2+(y_i+b)^2+(z_i+b)^2]\},`$ (6) where $`N=N_1+N_2`$. This potential is rewritten as $$V_\mathrm{T}=\underset{i=1}{\overset{N}{}}\left[m\omega ^2x_i^2m\gamma ^2(y_i^2+z_i^2)\right]+Nm(\omega ^2a^22\gamma ^2b^2).$$ (7) The differences between the potentials before and after the mixing appear in the following two points: One is the fact that the potential curvatures after the mixing become twice as large as those before the mixing and the other is the appearance of a real constant term $`Nm(\omega ^2a^22\gamma ^2b^2)`$. These effects can be seen in the ground-state energy such that $$E_{\mathrm{aft}}=\frac{1}{\sqrt{2}}N\mathrm{}\omega +Nm(\omega ^2a^22\gamma ^2b^2).$$ (8) The difference between the ground-state energies before and after the mixing is evaluated as $$\mathrm{\Delta }E_0=\left(\frac{1}{\sqrt{2}}\frac{1}{2}\right)N\mathrm{}\omega +Nm(\omega ^2a^22\gamma ^2b^2),$$ (9) where $`\mathrm{\Delta }E_0=E_{\mathrm{aft}}E_{\mathrm{bef}}`$. The difference can be negative, even if the first term of the difference is definitely positive. Namely, when the relation $$\left(\frac{1}{\sqrt{2}}\frac{1}{2}\right)\mathrm{}\omega +m\omega ^2a^2<2m\gamma ^2b^2$$ (10) is satisfied, the difference become negative. It should be noted that the negative contribution originates only from the two-dimensional PPB. One easily see that the difference always becomes positve, if the potential is described only by HO. The existence of the PPB is essential to derive the negative difference. The change of the potentials producing real energy eigenvalues before and after the mixing are illustrated in fig. 1. This real and negative energy $`\mathrm{\Delta }E_0`$ must be absorbed in the energy of the HO, because the PPB can only have the freedom of pure imaginary eigenvalues. This means that the ground-state energy before the mixing becomes an excited-state energy after the mixing. The excitation energy is evaluated by $`E_{\mathrm{exc}}=\|\mathrm{\Delta }E_0\|.`$ Note here that $`E_{\mathrm{exc}}`$ is generally a macroscopic order because it is proportional to the total particle number $`N`$. Thus the system have the freedom arising from the real energy, which is given by the Boltzmann entropy $$S^{\mathrm{}}=k_\mathrm{B}\mathrm{log}W^{\mathrm{}}(E_{\mathrm{exc}}),$$ (11) where $$W^{\mathrm{}}(E_{\mathrm{exc}})=\frac{(M+N1)!}{M!(N1)!}$$ and $`M`$ stands for the excitation number of the HO after the mixing, which is defined by the maximum integer being smaller than $`E_{\mathrm{exc}}/\sqrt{2}\mathrm{}\omega `$. It is well-known that one can express the entropy as $$S^{\mathrm{}}=k_\mathrm{B}\left[(M+N)\mathrm{log}(M+N)M\mathrm{log}MN\mathrm{log}N\right],$$ (12) where $`M,N1`$ are postulated. The temperature is now given by $$T=\frac{\sqrt{2}\mathrm{}\omega }{k_\mathrm{B}}\left\{\mathrm{log}\left[\frac{m(2\gamma ^2b^2\omega ^2a^2)+(\sqrt{2}+1)\mathrm{}\omega /2}{m(2\gamma ^2b^2\omega ^2a^2)(\sqrt{2}1)\mathrm{}\omega /2}\right]\right\}^1.$$ (13) We have shown that the entropy production from the PPB is possible. Namely, the mixing of two systems changes the ground state and then the energy of the ground state after the mixing can be much smaller than that before the mixing. The essential point of this process is the existence of stable states in the repulsive potentials like the PPB . PPB’s can possibly be good approximations to repulsive forces which are very week at the center of the force as same as the fact that HO’s are known as good approximations to attractive forces being very week at the center. We can, therefore, expect that entropy productions will occur in real thermodynamical processes. One may understand that the entropy production is caused by the entropy transfer between $`S^{\mathrm{}}`$ and $`S^{\mathrm{}}`$, though one cannot directly show it in the present model because the entropies $`S^{\mathrm{}}`$’s in the initial and the final states are infinity. It will happen in realistic processes that the number $`E_{\mathrm{exc}}/\sqrt{2}\mathrm{}\omega (M)`$ is not integer. In such processes some part of energy must be emitted from the systems to make a stable system after the mixing. Sometimes the most of the energy $`E_{\mathrm{exc}}`$ will possibly be emitted and the system goes to the fatal ground state, since the change of the curvature of the PPB makes all stationary states in the initial state unstable in the final state. In such cases the observers will see the burst of energy, which will be seen as the burst of entropy also. This fact indicates that one can make macroscopic energy bursts, provided that one can prepare the thermally stable systems composed of stationary flows in PPB’s. That is to say, there is a possibility of producing energies in purely thermal processes without any nuclear fusions at ordinary temperatures. Finally we would also like to note that it will be a very interesting trial to describe the birth of the Universe in the present scheme.
warning/0007/astro-ph0007078.html	ar5iv	text	# Metal Abundances in the Magellanic StreamBased upon observations with the NASA/ESA Hubble Space Telescope, obtained at the Space Telescope Science Institute, which is operated by the Association of Universities for Research in Astronomy, Inc., under NASA contract NAS5-26555. ## 1 Introduction Subtending $``$100 from the Magellanic Clouds through the south Galactic pole and beyond, the Magellanic Stream is the most striking extragalactic feature in the neutral hydrogen sky (Wannier & Wrixon 1972). Despite its prominence and apparent association with the Clouds, the Stream’s origin remains somewhat controversial. Tidally-disrupted (Putman et al. 1998) or ram pressure-stripped (Moore & Davis 1994) gas from the LMC or SMC remain the most likely scenarios, although ram pressure-stripped gas from the inter-Cloud region has also been suggested (Mathewson et al. 1987). A related question pertains to the still uncertain relationship between the Stream and the general population of Galactic High-Velocity Clouds (HVCs).<sup>1</sup><sup>1</sup>1The Magellanic Stream corresponds to HVC #493 in the Wakker & van Woerden (1991) catalog of High-Velocity Clouds. In our study, we use the classical definition of the Stream - i.e., the stream of gas trailing the Magellanic Clouds. Other HVCs leading the Clouds have been ascribed a similar origin to that of the Stream (Lu et al. 1998; Putman et al. 1998; Sahu 1998), but for our purposes the conventional definition is assumed throughout. As reviewed by Wakker & van Woerden (1997, § 6.3), several models ascribe a Magellanic Cloud origin to some of the most prominent HVC complexes (e.g., Complex C, VHVCs, Population EP), similar to that of the Stream. The chemical composition of gas comprising the Stream and HVCs has long been recognized as a key discriminant between competing models for their respective origins. Unfortunately, the dearth of suitable background probes against which to determine abundances of intervening HVC gas has made progress in this area difficult. To date, accurate metallicities have been derived for only a few HVCs, including \[S II/H I\]$`=0.60\pm 0.13`$, using the background probe NGC 3783, for HVC 287+22+240 (Lu et al. 1998), and \[S/H\]$`=1.03\pm 0.10`$, using Mrk 290, for Complex C (Wakker et al. 1999). A detailed study of the metallicity distribution within Complex C, using five different background probes, will be reported in a future paper (Gibson et al. 2000). A comparison of the metallicity of the Stream with these and other HVCs is hampered by this same lack of suitable background probes. All efforts to date have concentrated upon the background Seyfert galaxy Fairall 9, which intersects the MS I concentration of the Stream (see Figure 1 of Mathewson 1985 for nomenclature), $``$10 down-Stream from the SMC. Songaila (1981) detected the Stream in Ca II optical absorption, which when combined with the lack of a detection at Na I D, and subsequent ionization equilibrium arguments, led to the loose constraint that $`2\mathrm{}<[\mathrm{Ca}/\mathrm{H}]\mathrm{}<+0.25`$. This limit was not stringent enough to differentiate between a primordial or Magellanic Cloud origin for the Stream. More recently, Lu, Savage & Sembach (1994) examined two Stream clouds along the Fairall 9 sightline (one at $`+210`$ km s<sup>-1</sup>, the other at $`+170`$ km s<sup>-1</sup>). They derived S II $`\lambda \lambda `$1250,1253 and Si II $`\lambda \lambda `$1260,1526 abundance limits of \[S II/H I\]$`\mathrm{}<0.52`$ and \[Si II/H I\]$`\mathrm{}>1.15`$, for the $`+210`$ km s<sup>-1</sup> cloud, and \[S II/H I\]$`\mathrm{}<0.05`$ and \[Si II/H I\]$`\mathrm{}>0.70`$, for the $`+170`$ km s<sup>-1</sup> cloud. Under the assumption that the two clouds have the same metallicity, these constraints, for the full sightline, reduce to \[S II/H I\]$`\mathrm{}<0.52`$ and \[Si II/H I\]$`\mathrm{}>0.70`$. The low signal-to-noise ratio S/N in the vicinity of the sulfur lines only allowed Lu et al. to set upper limits on \[S/H\]; the saturated silicon lines allowed for only lower limits on \[Si II/H I\]. As part of our ongoing HST Guest Observer program on the origin and physical conditions in the local Ly$`\alpha `$ forest (cf. Penton, Stocke, & Shull 2000; Penton, Shull, & Stocke 2000), we revisited Fairall 9 with a similar GHRS set-up to that employed by Lu et al. (1994). Supplemented with the archival Lu et al. dataset, this has led to a four-fold increase in the effective integration time spent on this one target; the resultant increase in S/N has allowed us, for the first time, to derive an unequivocal metallicity determination for the Magellanic Stream. In § 2.1, we describe both the GHRS and neutral hydrogen data employed in our analysis of the Fairall 9 sightline. The analysis of the Stream absorption features seen in this dataset, and the corresponding implications for \[S II/H I\] and \[Si II/H I\], are described in § 3.1. We report serendipitous detections of the Stream (but $``$80 down-Stream of Fairall 9, near the tip and adjacent to the MS V concentration) in § 2.2, 3.2, 2.3, and 3.3. These detections of the Stream, seen in Mg II absorption in the spectra of III Zw 2 and NGC 7469, represent the only other probes, to date, toward which the Stream has been detected. We discuss the implications of our detections in § 4. Our results are summarized in § 5. ## 2 Data ### 2.1 Fairall 9 The first three entries of Table 1 list the HST GHRS data employed in our analysis of the Fairall 9 sightline. A single, merged, spectrum was created by first scaling the flux levels of the two 1996 exposures – z3e70404m and z3e70406t – to be consistent with the initial Lu et al. (1994) exposure (z26o0208n),<sup>2</sup><sup>2</sup>2This scaling was applied linearly across the region overlapping that of the initial exposure. and then averaging the three, weighting by the flux uncertainties. A detailed description of the spectrum preparation can be found in Penton, Stocke & Shull (2000). The effective integration time of the merged Fairall 9 spectrum totals $``$7.5 hrs, with a 3$`\sigma `$ equivalent width detection limit of 6 mÅ at $`\lambda =1250`$ Å. The S/N at $`\lambda =1250`$ Å is $``$30. The resulting merged GHRS spectrum is shown in Figure 1. The Galactic and Stream lines of S II $`\lambda \lambda `$1250, 1253 and Si II $`\lambda `$1260 are labeled accordingly. The three Galactic lines shown are redshifted (in the mean) by 0.028 Å, with respect to their rest frame wavelengths (taken from Verner et al. 1994, and listed in Table 3), corresponding to $``$$`+`$7 km s<sup>-1</sup> at $`\lambda =1250`$ Å. In contrast, the Galactic H I along this sightline (upper panel of Figure 4) peaks at $`v_{\mathrm{LSR}}1`$ km s<sup>-1</sup>. This suggests that an a posteriori velocity shift of $`\mathrm{\Delta }v8`$ km s<sup>-1</sup> might be appropriate for our merged GHRS Fairall 9 spectrum,<sup>3</sup><sup>3</sup>3Such a shift was applied by Penton, Stocke & Shull (2000) in the Ly$`\alpha `$ absorber analysis. although we have not done so here since it has no impact upon the Stream abundance analysis that follows. The neutral hydrogen column density of the Stream in this direction is N(H I)=$`9.35\times 10^{19}`$ cm<sup>-2</sup>. The two-component velocity structure of the high-velocity gas (components centered at $`+`$155 km s<sup>-1</sup> and $`+`$195 km s<sup>-1</sup>) has been noted before, although the spectrum shown in Figure 4 is the first to clearly show it. We will return to a discussion of the H I properties of this sightline in § 3.1. ### 2.2 III Zw 2 The second Stream background probe discussed here is III Zw 2. As evidenced by its entry in Table 1 (row 4), only a single 25 min HST GHRS G270M exposure is available. The corresponding 3$`\sigma `$ equivalent width detection limit and S/N, at $`\lambda =2800`$ Å, are 27 mÅ and 8, respectively. The spectrum itself is shown in Figure 2, with Galactic and Stream Mg II $`\lambda \lambda `$2796, 2803 features identified. The Galactic Mg II lines are blueshifted (in the mean) by 0.236 Å (i.e., $``$25 km s<sup>-1</sup> at $`\lambda =2800`$Å), with respect to their rest frame wavelengths (Pickering, Thorne & Webb 1998). This offset is to be expected, based upon the H I distribution along this line of sight, which shows a two-component structure (Galactic and intermediate-velocity gas) centered on $`v_{\mathrm{LSR}}30`$ km s<sup>-1</sup>. As was the case for our Fairall 9 spectrum (§ 2.1), we have not applied any a posteriori velocity shift to the GHRS data. ### 2.3 NGC 7469 As part of an HST FOS snapshot program, NGC 7469 was observed with the G270H grating for 5 min during 1996. The 3$`\sigma `$ equivalent width detection limit and S/N at $`\lambda =2800`$ Å are 162 mÅ and 29, respectively. The FOS spectrum is shown in Figure 3, with Galactic and Stream Mg II $`\lambda \lambda `$2796, 2803 features identified. An a posteriori velocity shift of $``$20 km s<sup>-1</sup> has been applied to the FOS wavelength calibration, in order to reconcile a systematic offset between the Galactic and Stream Mg II $`\lambda \lambda `$2796,2803 and H I features. The NRAO 140 H I spectrum employed in our analysis (upper panel of Figure 8) was kindly provided by Ken Sembach prior to publication (Murphy, Sembach & Lockman 2000). ## 3 Analysis ### 3.1 Fairall 9 Figure 4 shows the velocity stack for the Galactic (G) and Stream (MS) lines (S II $`\lambda `$1250, S II $`\lambda `$1253, and Si II $`\lambda `$1260) detected in our Fairall 9 GHRS spectrum, as well as an H I spectrum taken with the Parkes Narrowband Receiver (Haynes et al. 1999 - spectrum kindly provided by Lister Staveley-Smith). The S II $`\lambda `$1259 Stream line is blended with the Galactic Si II $`\lambda `$1260 line in the lower panel. In Table 2, we list the centroids of the Stream and Galactic lines (column 1), velocity range over which equivalent widths were determined (column 2), and the line equivalent widths (column 3).<sup>4</sup><sup>4</sup>4Equivalent widths were derived by integrating over the line profile; Gaussian fits were also undertaken (cf. Penton, Stocke & Shull 2000), but the results here are insensitive to the measurement technique employed, and thus we only report the former. The inferred ionic column densities are listed in columns 4 and 5 – the former (N<sub>τ=0</sub>) corresponds to the optically-thin assumption, while the latter (N$`_{\tau _v}`$) corresponds to the apparent optical depth method of Sembach & Savage (1992). For the Fairall 9 S II Stream features, N<sub>τ=0</sub> and N$`_{\tau _v}`$ are consistent within the uncertainties. Since the Si II $`\lambda `$1260 lines are clearly saturated, it is not surprising that N<sub>τ=0</sub> and N$`_{\tau _v}`$ are discrepant. This saturation allows us only to set lower limits on the Si II column densities. In what follows, we adopt N$`_{\tau _v}`$ in deriving the Stream’s metallicity, and we make no attempt to treat the two components seen in the H I spectrum separately in our analysis. In the weak-line (optically-thin) limit, the column density N<sub>τ=0</sub> (in cm<sup>-2</sup>) and equivalent width of a line W<sub>λ</sub> (in mÅ) are related through $$\mathrm{N}_{\tau =0}=1.13\times 10^{17}\frac{\mathrm{W}_\lambda }{f\lambda _0^2},$$ (1) where $`\lambda _0`$ (in Å) is the rest wavelength of the line and $`f`$ is its oscillator strength (e.g., Savage & Sembach 1996; equation 3). For the lines considered in this paper, the relevant values of $`\lambda _0`$ and $`f`$ are provided in columns 2 and 3 of Table 3. Under the apparent optical depth method of Sembach & Savage (1992), the apparent column density N$`_{\tau _v}`$, in the limit of a finite number of data points, is instead given by $$\mathrm{N}_{\tau _v}=\frac{3.767\times 10^{14}}{f\lambda _0}\underset{i=1}{\overset{n}{}}\mathrm{ln}\left[\frac{I_c(v_i)}{I(v_i)}\right]\mathrm{d}v_i,$$ (2) where $`I(v_i)`$ and $`I_c(v_i)`$ are the observed and estimated continuum intensities at velocity $`v_i`$ (equations A21 and A29 of Sembach & Savage). The statistical noise uncertainty associated with the apparent column density of equation 2 is $$\sigma _{\mathrm{N}_{\tau _v}}=\mathrm{N}_{\tau _v}\left(\sqrt{\underset{i=1}{\overset{n}{}}\sigma ^2(v_i)/I(v_i)^2(\mathrm{d}v_i)^2}/\underset{i=1}{\overset{n}{}}\mathrm{ln}\left[\frac{I_c(v_i)}{I(v_i)}\right]\mathrm{d}v_i\right),$$ (3) based upon equations A27 and A30 of Sembach & Savage. Uncertainties quoted in this analysis correspond to those of equation 3; those associated with continuum placement have not been considered here. For the high S/N Fairall 9 data, the uncertainty associated with the continuum placement is only an additional $``$5% beyond that of equation 3 (Penton, Stocke & Shull 2000) and is neglected here. The equivalent width ratio of the S II Stream features (from Table 2) is W<sub>λ</sub>(S II $`\lambda `$1253)/W<sub>λ</sub>(S II $`\lambda `$1250)= 1.72$`\pm `$0.24, further evidence that any saturation effects are mild (the expected theoretical ratio is 2.0). Regardless, allowing for these marginal optical depth effects, the apparent column densities of the Stream S II and Si II features seen in the Fairall 9 GHRS spectrum are N$`_{\tau _v}`$(S II $`\lambda `$1250)=$`(5.39\pm 0.61)\times 10^{14}`$ cm<sup>-2</sup>, N$`_{\tau _v}`$(S II $`\lambda `$1253)=$`(4.75\pm 0.41)\times 10^{14}`$ cm<sup>-2</sup>, and N$`_{\tau _v}`$(Si II)$``$N(Si II)$`>`$$`9.29\times 10^{13}`$ cm<sup>-2</sup>. The S II-weighted average is N$`_{\tau _v}`$(S II)$``$N(S II)= $`(4.95\pm 0.34)\times 10^{14}`$ cm<sup>-2</sup>. Expressing the sulfur abundance in terms of the logarithmic abundance A<sub>S</sub> with respect to that of the Sun A$`_\mathrm{S}_{}`$, we can write $$[\mathrm{S}/\mathrm{H}]=\mathrm{log}\frac{\mathrm{A}_\mathrm{S}}{\mathrm{A}_\mathrm{S}_{}}=\mathrm{log}\mathrm{N}(\mathrm{S})\mathrm{log}\mathrm{N}(\mathrm{H})\mathrm{log}\mathrm{A}_\mathrm{S}_{},$$ (4) since A$`{}_{\mathrm{S}}{}^{}`$N(S)/N(H), by definition. The nondetection of S I $`\lambda `$1262 in our Fairall 9 spectrum (W<sub>λ</sub> $`<6`$ mÅ) implies N(S I) $`<2.1\times 10^{14}`$ cm<sup>-2</sup>. With an ionization potential of only 10.36 eV, S I is easily ionized by the ambient halo radiation field and likely makes a negligible contribution to N(S). We currently have no observational limits on the contribution of S III to N(S) but, after Wakker et al. (1999) who consider a Complex C cloud of similar H I column density, we assume that N(S III)$``$N(S II). (Also, see our discussion of corrections for H II below). With A$`{}_{\mathrm{S}_{}}{}^{}=(1.862\pm 0.215)\times 10^5`$ (column 4 of Table 3), and explicitly restricting ourselves to S II and H I (as opposed to the global S and H), equation 4 reduces to $$[\mathrm{SII}/\mathrm{HI}]=(19.425\pm 0.058)\mathrm{log}\mathrm{N}(\mathrm{HI}).$$ (5) During a 02/17/99-02/22/99 observing run at Parkes Observatory, the Multibeam Narrowband Facility (Haynes et al. 1999) was employed to obtain the Fairall 9 sightline H I spectrum shown in the upper panel of Figure 4. The quality of this H I data is a vast improvement over that available to Lu et al. (1994) - i.e., the McGee, Newton & Morton (1983) and Morras (1983) datasets. The H I column density for the Stream, along this sightline, is N(H I)$`=(9.35\pm 0.47)\times 10^{19}`$ cm<sup>-2</sup>. In combination with equation 5, we can use this new determination of N(H I) to derive the sulfur abundance for the Stream in the Fairall 9 sightline, resulting in \[S II/H I\]=$``$0.55$`\pm `$0.06. If we had chosen to use N<sub>τ=0</sub>, as opposed to N$`_{\tau _v}`$, the derived value of \[S II/H I\] would have been reduced by only 0.04 dex. If 25% of the total hydrogen column is in the form of ionized hydrogen, as appears to be the case for the Mrk 290 sightline through Complex C (Wakker et al. 1999), the derived \[S II/H I\] could be 0.12 dex greater than \[S II/H\]. This is (perhaps) a conservative overestimate since, in the ionized gas, additional sulfur may be in the form of S III. A further source of systematic error in our estimate of N(H) is far more challenging to quantify. Our H I column was derived using the 14 beam provided by the Parkes dish, while the HST absorption measurements toward Fairall 9 sample the gas at sub-arcsecond resolution. Comparisons have been made of N(H I) derived from 21cm mapping with measurements of N(H I) derived from Ly$`\alpha `$ absorption along the lines of sight to early type stars in the Galactic halo (Lockman, Hobbs & Shull 1986; Savage et al. 2000; Wakker & Savage 2000). The ratio N(H I)<sub>Lyα</sub> / N(H I)<sub>21cm</sub> is somewhat less than unity, with a dispersion less than $`\pm `$50%. This low dispersion is in apparent contrast with the results of 21cm surveys of different angular resolution, which show variations in N(H I) as large as a factor of five over arcminute spatial scales (Wakker & Schwarz 1991). A simple, yet elegant, argument put forth by Bart Wakker suggests that the two observations may not be in contradiction with one another. Radio surveys with higher angular resolution (1) reveal that the number of fields with a given N(H I) is a steep inverse power law function of N(H I). This implies that higher resolution probes of a large beam will be more likely to intersect low N(H I) sightlines.<sup>5</sup><sup>5</sup>5Of the seven sightlines for which both $`<`$1 and 10$``$12 H I data exists (Table A1 of Wakker 2000), in only 4/7 of these cases is N(H I)$`_1^{}`$$`<`$N(H I)$`_{10^{}}`$. The ratio of N(H I)$`_1^{}`$$`/`$N(H I)$`_{10^{}}`$ for all seven is 0.90$`\pm `$0.13, with extrema of 0.75 and 1.24. This might lead one to conclude that the importance of the predicted “resolution effect” has been overestimated, but it should be stressed that it is based upon only seven data points. To be conservative we have retained the assumed $`\pm `$50% uncertainty in N(H I). In principle, this same argument could be extended to much smaller scales, as the gulf between high resolution 1 21cm maps and the scales probed by absorption lines remains large. To be conservative, we have incorporated a factor $`\pm `$0.17 dex ($`\pm `$50%) in our systematic error budget, to reflect the systematic uncertainty in our estimate of N(H I). It remains possible that we have been unfortunate enough to encounter a line of sight with $``$N(H I)$``$ that differs from the “pencil beam” N(H I) by a factor of five or more. Higher resolution radio mapping will better address the likelihood of this possibility. It is reassuring to note that, within the measurement errors, the S II and H I velocity profiles are similar (Figure 4). This lends some credence to the suggestion that the H I revealed by our Parkes data is representative of the “pencil beam” H I along the Fairall 9 sightline. Adding this second factor in quadrature with the systematic uncertainty in converting H I to H, we write our final derived value for the sulfur abundance of the Stream in the Fairall 9 sightline as $$[\mathrm{SII}/\mathrm{H}]=0.55\pm 0.06(\mathrm{r})_{0.21}^{+0.17}(\mathrm{s}),$$ (6) where ‘r’ and ‘s’ correspond to the associated random and systematic uncertainties, respectively. This result can now be compared with that derived in the earlier Lu et al. (1994) study. Their best S II constraint was set by their 3$`\sigma `$ upper limit on the S II $`\lambda `$1253 absorption, W$`{}_{\lambda }{}^{}87`$ mÅ (consistent with our result of W$`{}_{\lambda }{}^{}=62\pm 5`$ mÅ). Using our H I column density for this sightline, as opposed to the Morras (1983) value used by Lu et al., results in an upper limit to the sulfur abundance \[S II/H I\]$`0.46`$, consistent with our result of \[S II/H I\]=$``$0.55$`\pm `$0.06. For comparison, the gas-phase sulfur abundances of the Magellanic Clouds are \[S/H\]$`{}_{\mathrm{LMC}}{}^{}=0.57\pm 0.10`$ and \[S/H\]$`{}_{\mathrm{SMC}}{}^{}=0.68\pm 0.16`$ (Russell & Dopita 1992). Because of the magnitude of the residual uncertainties (particularly those of a systematic nature), associating the MS I gas with only one of the Clouds remains difficult. Equation 6 has a somewhat ad hoc provision for ionization corrections, and no provision for potential dust depletion effects. Fortunately, there is some confidence that S II is effectively free of both of these complicating effects (Wakker et al. 1999). In analogy to the derivation of equation 6, we can use the Si II $`\lambda `$1260 column density (Table 2) to derive a lower limit on the MS I silicon metallicity of $$[\mathrm{SiII}/\mathrm{H}]>1.55\pm 0.03(\mathrm{r})_{0.21}^{+0.17}(\mathrm{s}).$$ (7) Because we are limited to the saturated Si II $`\lambda `$1260 line, we could derive only the above lower limit with the current data. Lu et al. (1994) find W$`{}_{\lambda }{}^{}=449\pm 53`$ mÅ for the Si II $`\lambda `$1526 Stream absorption. For the total line of sight H I column density N(H I)=$`9.35\times 10^{19}`$ cm<sup>-2</sup>, this corresponds to a lower limit on the silicon abundance of \[Si II/H I\]$`\mathrm{}>`$$``$1.1. This particular limit is set by the much lower S/N Si II $`\lambda `$1526 line, which is not included in our G160M GHRS spectra. Therefore, we retain the more conservative limit set by Si II $`\lambda `$1260 (i.e., equation 7). ### 3.2 III Zw 2 Figure 5 shows the velocity profiles of the detected Galactic (G) and Magellanic Stream (MS) lines of Mg II $`\lambda \lambda `$2796,2803 in our G270M GHRS spectrum of III Zw 2. The centroids, equivalent widths, and inferred column densities (both optically thin and following the apparent optical depth methodology described in § 3.1) are listed in Table 4. While our GHRS spectrum is clearly of low signal-to-noise (S/N=8 at $`\lambda =2800`$ Å), both Stream Mg II lines appear unsaturated, supported by the fact that the ratio of their equivalent widths, W<sub>λ</sub>(Mg II $`\lambda `$2796)/W<sub>λ</sub>(Mg II $`\lambda `$2803)= 2.11$`\pm `$0.94, is consistent with the expected theoretical ratio of 2.0. N<sub>τ=0</sub> and N$`_{\tau _v}`$ agree to within the uncertainties, further evidence that only a mild degree of saturation is present. Employing the apparent optical depth technique, we derive Mg II column densities of N(Mg II $`\lambda `$2796)=$`(1.08\pm 0.10)\times 10^{13}`$ cm<sup>-2</sup> and N(Mg II $`\lambda `$2803)=$`(0.91\pm 0.14)\times 10^{13}`$ cm<sup>-2</sup>. We adopt the weighted average in the analysis which follows, N(Mg II)=$`(1.02\pm 0.08)\times 10^{13}`$ cm<sup>-2</sup>. Analogous to the derivation of equation 5, we can use N(Mg II) along the line of sight to III Zw 2 to write $$[\mathrm{Mg}\mathrm{II}/\mathrm{H}\mathrm{I}]=(17.429\pm 0.039)\mathrm{log}\mathrm{N}(\mathrm{H}\mathrm{I}).$$ (8) Unlike the case for Fairall 9, we have little in the way of useful observational constraints on the $`\mathrm{log}\mathrm{N}(\mathrm{H}\mathrm{I})`$ term of equation 8. The upper panel of Figure 5 shows the Hanning-smoothed H I spectrum from Hartmann & Burton (1997); the non-detection of H I and the position and velocity of our Mg II detection sets an upper limit to the H I column density of $``$5$`\times `$10<sup>18</sup> cm<sup>-2</sup>.<sup>6</sup><sup>6</sup>6The 5$`\sigma `$ detection limit for the Leiden-Dwingeloo Survey (Hartmann & Burton 1997) is T<sub>B</sub>=0.35 K (unsmoothed). For an HVC with a linewidth of 8 km s<sup>-1</sup>, the 5$`\sigma `$ H I column density detection limit corresponds to 5$`\times `$10<sup>18</sup> cm<sup>-2</sup>; for a 20 km s<sup>-1</sup> linewidth, N(H I)=5$`\times `$10<sup>18</sup> cm<sup>-2</sup> corresponds to a 2$`\sigma `$ detection limit. The MS V concentration of the Stream is centered on $`(\mathrm{},b)(92^{},51^{})`$, with $`v_{\mathrm{LSR}}350`$ km s<sup>-1</sup> (Mathewson 1985). The upper panel of Figure 6 shows an $``$2000 deg<sup>2</sup> region encompassing all of MS V, including the III Zw 2 sightline shown in the lower panel and the neighboring NGC 7469 sightline discussed in § 3.3. The lower panel of shows contours of neutral hydrogen, derived from the Leiden-Dwingeloo Survey (Hartmann & Burton 1997), in a $`14^{}\times 7^{}`$ region centered on the III Zw 2 sightline. The H I contour levels are $`3\times 10^{18}`$ cm<sup>-2</sup>, with the outermost level being N(H I)=$`1\times 10^{18}`$ cm<sup>-2</sup>. Only the velocity range $`400<v_{\mathrm{LSR}}<250`$ km s<sup>-1</sup> was included in this figure. The “edge” of the Stream proper coincides with the increased number of H I clouds between $`\mathrm{}98^{}`$ and $`\mathrm{}104^{}`$. In the lower panel of Figure 6, we have highlighted three nearby ($`\mathrm{}<`$4 away from the line of sight) HVCs to the III Zw 2 sightline. While these HVCs are at the threshold of the Leiden-Dwingeloo Survey detection limit, each has been detected previously (Hulsbosch & Wakker 1988; Wakker & van Woerden 1991): HVC 104.0-51.0-337 was detected by Hulsbosch & Wakker (1988), but is subsumed into HVC#493 by Wakker & van Woerden (1991); HVC 110.0-50.0-290 corresponds to HVC#520 in Wakker & van Woerden; HVC 108.0-53.0-325 corresponds to HVC#523. Much of the H I detected with N(H I)$`\mathrm{}<`$10<sup>18</sup> cm<sup>-2</sup> is probably noise (and therefore the lower panel of Figure 6 should be used cautiously), but at least these three nearby HVCs have been independently detected by Hulsbosch & Wakker (1988). More importantly, these neighboring HVCs demonstrate that gas presumably associated with the Stream exists $`\mathrm{}>15^{}`$ away from the peak of the H I emission. Because the “width” of the MS V concentration seen in the upper panel of Figure 6 is $``$10, the existence of MS “froth” beyond the outermost contour is not surprising. Equipped only with our measurement of the Mg II column density and the upper limit on the neutral hydrogen column density, it is difficult to constrain the metallicity of the absorbing gas. As an exploration of the possible range in inferred metallicity, we have constructed a grid of photoionization models using the photoionization code CLOUDY (Version 90.04 - Ferland 1996). We hope to constrain this range of physical parameters more tightly with future observations.<sup>7</sup><sup>7</sup>7Being radio-loud, III Zw 2 may lend itself to potential future 21cm absorption analyses. We model the absorbing gas as a plane-parallel slab illuminated on one side with a normally incident ionizing photon flux, $`\mathrm{log}\varphi =5.5`$ photons cm<sup>-2</sup> s<sup>-1</sup>. This is consistent with the estimates of Bland-Hawthorn & Maloney (1999a,b) for this region of the Magellanic Stream, based on measurements of the H$`\alpha `$ emission near our line of sight. As Bland-Hawthorn & Maloney discuss, this level of radiation is much higher than that due to the general metagalactic background, $`\mathrm{log}\varphi =4`$ (Shull et al. 1999), and is likely due to escaping stellar radiation from the Milky Way. As a result, we model the spectrum as a power law with spectral index $`\alpha _s=2`$, assuming F$`{}_{\nu }{}^{}\nu ^2`$ between 1 and 4 ryd, with a dropoff above 4 ryd of a factor of 100. This represents a rough approximation to the integrated radiation from Galactic OB associations, as well as a harder extragalactic component. However, adopting a T<sub>eff</sub> = 35000 K stellar atmosphere for the incident spectrum yields a grid of models which are indistinguishable from the following results. For a given density n(H) and total hydrogen column density N(H), the assumed metallicity of our cloud model was allowed to vary until N(Mg II)=$`(1.0\pm 0.05)\times 10^{13}`$ cm<sup>-2</sup> was achieved. The results of the grid of models are shown in Figure 7. Also shown is the area that can be currently excluded by the requirement that N(H I)$`<5\times 10^{18}`$ cm<sup>-2</sup> (shaded region). Our Mg II column density determination, coupled with the upper limit to N(H I), allows us to set a lower limit \[Mg II/H I\]$`\mathrm{}>`$$``$1.3. Even this has only limited use as a constraint on the metallicity \[Mg/H\], however, as for N(H I)$`<`$$`5\times 10^{18}`$ cm<sup>-2</sup> and $`\mathrm{log}\varphi =5.5`$, the ionization correction for H I is substantially more sensitive to changes in density than that for Mg II. ### 3.3 NGC 7469 As already noted in § 2.3, the Stream was detected in Mg II $`\lambda `$2796 in the FOS G270H spectrum of NGC 7469 (bottom panel of Figure 8). The Mg II $`\lambda `$2803 line lies below the 3$`\sigma `$ detection threshold and is not discussed further. The Mg II column density inferred from the apparent optical depth method (Table 5) is N(Mg II)$`>(0.47\pm 0.11)\times 10^{13}`$ cm<sup>-2</sup>. While the line does not appear saturated, the inferior spectral resolution of FOS ($``$230 km s<sup>-1</sup>), in comparison with GHRS ($``$19 km s<sup>-1</sup>), does not allow us to unequivocally rule out the presence of saturation effects. To be conservative, we assume the measured Mg II column is a lower limit and will revisit this sightline with our FUSE Science Team Cycle 1 observations of NGC 7469, the O VI properties for which have already been discussed by Sembach et al. (2000). An NRAO 140 H I spectrum (upper panel of Figure 8), kindly provided by Ken Sembach (to be published in Murphy et al. 2000), shows that the H I column density along this MS V sightline is N(H I)=$`(0.40\pm 0.04)\times 10^{19}`$ cm<sup>-2</sup>. These constraints lead to $$[\mathrm{Mg}\mathrm{II}/\mathrm{H}\mathrm{I}]>1.51\pm 0.11(\mathrm{r}).$$ (9) Errors in converting this measurement to an abundance \[Mg/H\] are potentially large, since as discussed, for N(H I) $`10^{18}`$ cm<sup>-2</sup>, the ionization correction for H I is substantially more sensitive to density than the correction for Mg II. For an assumed density n(H)=0.1 cm<sup>-3</sup>, \[Mg II/H I\] may overestimate \[Mg/H\] by an order of magnitude or more. In addition, the large uncertainty in N(H I) about this $``$10<sup>18</sup> cm<sup>-2</sup> range has the secondary effect of enhancing the density sensitivity of the ionization correction to H I. In general, however, ionization corrections imply that \[Mg II/H I\] is an overestimate of \[Mg/H\], ranging from a factor of 2 to 10 for n(H) ranging from (1$``$0.1) cm<sup>-3</sup>. ## 4 Discussion ### 4.1 Fairall 9 Based on its proximity to the SMC and the disrupted appearance of the SMC, it is likely that the MS I (and the Stream as a whole) originated in the SMC. If the Stream was drawn out of the SMC $``$1.5 Gyrs ago, as the best tidal models predict (Gardiner 1999), one would expect the metallicity to reflect that of the SMC at that epoch. According to the best chemical evolution models (Pagel & Tautvaisiene 1998; Figure 5), that should actually be $``$0.2 dex smaller than the present-day value of \[S/H\]=$``$0.68 (Russell & Dopita 1992). Note that this argument also holds for an LMC origin as well; the LMC metallicity 1.5 Gyrs ago, in the mean, is also predicted to be $``$0.2 dex lower than the present-day LMC value of \[S/H\]=$``$0.57 (Pagel & Tautvaisiene; Figure 4). In either scenario, the predicted \[S/H\] would appear to be mildly inconsistent with our observations, \[S II/H I\]=$``$0.55$`\pm `$0.06(r)$`\pm `$$`{}_{0.21}{}^{}{}_{}{}^{+0.17}`$(s). On the other hand, the stochastic nature of star formation and the magnitude of the observational scatter (see Figures 4 and 5 of Pagel & Tautvaisiene), severely limits the predictive power of even the best Magellanic Cloud chemical evolution models, applied to the most recent $`12`$ Gyrs. Thus, this apparent discrepancy should not be overinterpreted. An inherent assumption in the above picture is that the MS I gas can be linked directly to disk gas-phase abundances. The present-day population of disk H II regions, in both the LMC and SMC, shows no evidence for any substantial abundance gradients (Dufour 1975), nor are the gaseous disks of the Clouds (currently) substantially larger than their optical disks. In other words, based on the present state of the Magellanic Clouds, there is no reason to suspect that the Stream abundances should not reflect the disk gas-phase abundances. A caveat to this statement is that the present-day conditions may not necessarily reflect those when the Stream was purported to form $``$1.5 Gyrs ago. Even if an abundance gradient or extended gaseous disk existed when the Stream formed, unless the gradient was inverted (highly unusual), one would then expect the Stream abundances to be even lower than the disk gas-phase abundances. This would worsen the comparison of equation 6 with either the present-day Russell & Dopita (1992) LMC/SMC abundances, or those corrected downward by 0.2 dex based upon conservative chemical evolution model predictions. Higher spatial resolution H I data will be required to reduce the remaining substantial systematic uncertainties discussed in § 3.1, and to better ascertain whether a legitimate discrepancy exists between the MS I metallicity and that expected for an SMC disk origin. Impending synthesis data from our Australia Telescope Compact Array program (Putman & Gibson 2000) should aid in this regard. Complicating the above interpretation is the recent work of Rolleston et al. (1999), who found that, for three B-stars in the Magellanic Bridge, \[Mg/H\]=$`0.94\pm 0.14`$ and \[Si/H\]=$`1.23\pm 0.25`$. These Bridge stars appear to be $``$0.5 dex deficient in metals compared to a similar B-star (AV 304) in the SMC itself. These stars do not appear to have formed from gas of the present-day SMC or LMC composition. Some unenriched component most likely mixed with the SMC gas to yield abundances low enough to produce the Rolleston et al. results. Evidence is presented therein for the SMC not being as well-mixed as Dufour (1975) had claimed earlier, potentially providing a source for metal-poor contaminating gas. Invoking this metal-poor contamination is counter to that implied by our Fairall 9 MS I sightline, which appears to be enriched (by 0.1$``$0.2 dex), in comparison with the present-day mean gas-phase abundance of the SMC. Equations 6 and 7 imply \[Si/S\]$`\mathrm{}>`$$``$1. While this lower limit marginally excludes dust depletion of the magnitude seen in cool diffuse Galactic disk clouds (Savage & Sembach 1996; Figure 6), it is not sufficiently restrictive to discriminate between depletion of the magnitude seen in warm Galactic halo clouds (as might be expected for an SMC origin for the Stream gas – Welty et al. 1997) and warm Galactic disk clouds (as might be expected for an LMC origin – Welty et al. 1999), which typically show silicon depletion of order \[Si/S\]$``$$``$0.4$`\pm `$0.1 (Savage & Sembach 1996; Figure 6). Further, it is known that the gas-phase in both the LMC and SMC is relatively overabundant in silicon, showing \[Si/S\]$``$$`+`$0.08$`\pm `$0.04 (Welty et al. 1997,1999). In other words, our limit limit of \[Si/S\]$`\mathrm{}>`$$``$1 for the Magellanic Stream cannot be used as a constraint to differentiate between an LMC versus SMC origin. Impending FUSE Cycle 1 Science Team observations of Fairall 9 (scheduled for July 2000), when coupled with our GHRS data, should shed further light on the dust depletion pattern of the Magellanic Stream. ### 4.2 III Zw 2 Assuming that the gas near MS V, $``$80 down-Stream from Fairall 9, originates from within the Magellanic Clouds, it is tempting to assign it a magnesium abundance that reflects that of either the LMC or SMC. Unfortunately, the gas-phase magnesium abundance for both Clouds is highly uncertain. Only for the LMC is there a direct measure of the interstellar medium (ISM) gas-phase magnesium, through the IUE absorption line analysis of R136 (de Boer et al. 1985; Welty et al. 1999), the result for which was \[Mg/H\]$`=`$$``$1.1$`\pm `$0.3. In contrast, spectral synthesis of B-star and F-type supergiant atmospheres implies the much higher abundance \[Mg/H\]$``$$``$0.5$``$$``$0.1 (Welty et al. 1999; Table 11). For the SMC, no direct ISM absorption line abundances have been published yet, but B-star and supergiant spectral synthesis studies lead to \[Mg/H\]$``$$``$0.7$``$$``$0.4 (Welty et al. 1997; Rolleston et al. 1999). As the model curves of Figure 7 show, neither an LMC nor an SMC origin can be excluded based upon our knowledge of N(Mg II) and N(H I), as even the most extreme allowable values for the LMC and SMC magnesium abundances (see above) reside in the parameter space not yet excluded by the upper limit to N(H I) – i.e., the shaded region. In either case, the total hydrogen column along this sightline would be dominated by N(H II) and not N(H I), and the total magnesium column would be dominated by N(Mg III). This latter result should not be surprising, since the ionization potentials of H I and Mg II lie within 10% of each other. Where hydrogen is predominantly ionized, magnesium follows suit, although as discussed above, more magnesium remains in the singly ionized state. From the Leiden-Dwingeloo Survey (Hartmann & Burton 1997), we can derive the H I volume density for both the marginally resolved HVC 100.0-48.5-390 (Figure 6) and the clearly resolved MS V concentration at $`(\mathrm{},b)(92^{},51^{})`$. In both cases, H I column densities $`2\times 10^{19}`$ cm<sup>-2</sup> are found. If the radial extent of the clouds is assumed to be similar to their lateral extent, the angular width times the distance to the clouds must be of order N(H I)/n(H I). At the assumed distance of the LMC ($``$50 kpc), this corresponds to an H I volume density of n(H I)$`10^{1.7}`$cm<sup>-3</sup>. If one were to make the assumption that this H I volume density truly reflected the total hydrogen volume density in those higher column density clouds, and that this represented the typical total hydrogen density of this part of the Magellanic Stream, this would imply that the total hydrogen column density for the absorber toward III Zw 2 was N(H)=$`(14)\times 10^{18}`$ cm<sup>-2</sup> with N(H II)$``$N(H I). Because we have no definitive evidence of the nature (if any) of the dust depletion pattern in the Magellanic Stream, the above CLOUDY analysis necessarily adopted the simplifying assumption of zero dust – i.e., pure photoionization effects were employed, in an attempt to reconcile the observed Mg II and H I Stream constraints with the observed \[Mg/H\] of the LMC and SMC. Such an analysis is a useful exercise, but is no doubt an over-simplification, due to the potential complicating effects of dust depletion. If we subscribe to the conclusion of § 4.1, that the Stream dust depletion may range anywhere from nonexistent, to that seen in warm diffuse Galactic clouds, then our measured gas-phase Mg II abundance may underestimate the true magnesium abundance by up to a factor of ten. More sensitive H I observations, coupled with additional absorption measurements for other ions, are necessary. The LMC and SMC possess magnesium abundances of \[Mg/H\]$``$$``$1.1$``$$``$0.1 and \[Mg/H\]$``$$``$0.7$``$$``$0.4, respectively. Regardless of the origin of the high-velocity Stream, it is likely that some degree of dust depletion will be present, since both the LMC and SMC exhibit distinct depletion patterns; the former displays a pattern similar to that seen in warm Galactic disk clouds (Welty et al. 1999), while the pattern seen in the latter resembles that of warm Galactic halo clouds (Welty et al. 1997). Our silicon analysis of the Fairall 9 sightline (§ 3.1) showed that we could not exclude dust depletion patterns similar to those of warm diffuse Galactic gas clouds. This implies a potential correction of $``$0$``$$`+`$1 dex to the measured magnesium abundance (Savage & Sembach 1996; Figure 6). This would retain the consistency of the inferred \[Mg II/H I\] limits of both III Zw 2 (\[Mg II/H I\]$`\mathrm{}>`$$``$1.3) and NGC 7469 (\[Mg II/H I\]$`\mathrm{}>`$$``$1.5) with the present-day SMC stellar abundance value of \[Mg/H\]$``$$``$0.7$``$$``$0.4. ## 5 Summary We summarize in Table 6 the results of our column density measurements and limits for three sightlines through the Magellanic Stream. We have also included our inferences about the abundances of the gas seen along these sightlines, many of which remain uncertain. For the line of sight toward Fairall 9, where ionization corrections are expected to be small, measurements of \[S II/H I\] are consistent with the present-day abundances of both the LMC and SMC. If this gas was tidally-stripped from the SMC 1.5 Gyrs ago, its abundance is approximately a factor of two greater than that expected. Stochastic star formation effects and remaining systematic uncertainties (particularly in the H I column density) could weaken this statement. New HST/STIS and FUSE observations of Fairall 9 at S III will test the inherent assumption that S III ionization corrections are negligible. Similarly, both III Zw 2 and NGC 7469 need to be revisited with HST/STIS, in order to derive S II and S III column densities. NGC 7469 has been observed by FUSE and its O VI properties discussed by Sembach et al. (2000); further analysis of this FUSE dataset is currently underway. Fairall 9 will likewise be observed by the FUSE Science Team in July 2000. A fourth potential probe of the Magellanic Stream, in a high H I column density region of MS III, is NGC 7714. It was observed with GHRS by González-Delgado et al. (1999), but with the low-resolution G140L grating. The resulting resolution ($`>`$100 km s<sup>-1</sup>) was insufficient to accurately separate MS III absorption features from the saturated Galactic lines, as the expected separation (based upon the LDS H I spectrum for this sightline) is only $``$50 km s<sup>-1</sup>. Scheduled FUSE observations of NGC 7714 should resolve any far-UV MS III lines. This work was supported by the NASA HST General Observer Grants GO-06586.01-95A and GO-06593.01-95A, NASA Long-Term Space Astrophysics Program (NAG5-7262), NSF grant AST 96-17073, and the FUSE Science Team (NAS5-32985). It was based in part on observations with the NASA/ESA Hubble Space Telescope obtained at the Space Telescope Science Institute, which is operated by AURA, Inc., under NASA contract NAS5-26555. We wish to thank Phil Maloney and Bart Wakker for helpful discussions, Ken Sembach and Lister Staveley-Smith for providing the NGC 7469 and Fairall 9, respectively, neutral hydrogen data prior to publication, and Gary Ferland for the use of CLOUDY. A special thanks is due the referee (Chris Howk), whose detailed report improved greatly the quality of the final manuscript. Fig. 1. – HST GHRS spectrum of Fairall 9 taken with the G160M grating. Fairall 9 samples Magellanic Stream gas associated with the MS I concentration (see Figure 1 of Mathewson 1985). The raw spectrum has been smoothed by the inverse of the post-COSTAR, GHRS, large science aperture, line spread function (LSF – Gilliland 1994), in order to improve the S/N (at the admitted expense of some resolution) and eliminate any LSF-induced asymmetries in the line profiles (Penton, Stocke & Shull 2000). Absorption lines from both Galactic (G) and Magellanic Stream (MS) S II (at 1250 and 1253 Å) and Si II (at 1260 Å) are seen and labeled accordingly. Galactic S II $`\lambda `$1259 is also seen, but since the corresponding Stream line is blended with Galactic Si II $`\lambda `$1260, we exclude it from our analysis. Similarly, we will not discuss further the $`z=0.032`$ Ly$`\alpha `$ line seen at $`\lambda =1254.15`$ Å. Fig. 2. – HST GHRS spectrum of III Zw 2 taken with the G270M grating. III Zw 2 lies $``$15 from the peak of the H I emission associated with the MS V concentration of the Magellanic Stream (see Figure 1 of Mathewson 1985). The raw spectrum has been smoothed by the inverse of the post-COSTAR LSF (Gilliland 1994), as noted in the caption to Figure 1. Both Galactic (G) and Magellanic Stream (MS) Mg II $`\lambda \lambda `$ 2796,2803Å is seen in absorption. The centroid of the Galactic lines are blueshifted 0.24 Å with respect to their rest wavelengths (Pickering, Thorne & Webb 1998), reflecting the blended two-component structure of the Galactic gas along this sightline. Fig. 3. – HST FOS spectrum of NGC 7469 taken with the G270H grating. Galactic (G) and Magellanic Stream (MS) Mg II $`\lambda \lambda `$ 2796,2803Å is seen in absorption, although the Stream is clearly detected in Mg II $`\lambda `$2796 only. Fig. 4. – Velocity stack showing H I (upper panel), S II (middle panels) and Si II (lower panel) Galactic (G) and Magellanic Stream (MS) features along the Fairall 9 sightline. The Hanning-smoothed H I spectrum was collected with the Parkes Multibeam Narrowband System (Haynes et al. 1999) and was transformed from the heliocentric to the local standard of rest frame via $`v_{\mathrm{LSR}}=v_{}11.6`$ km s<sup>-1</sup>. A mild fourth-order polynomial baseline was subtracted from the raw spectrum, and the native Jy/beam converted to T<sub>B</sub> via a multiplicative scale factor of 0.82 (Staveley-Smith 1999). The corresponding H I column densities for both the Galactic and Magellanic Stream components are labeled. Fig. 5. – Velocity stack showing H I (upper panel) and Mg II (middle and lower panels) Galactic (G) and Magellanic Stream (MS) features along the III Zw 2 sightline. The Hanning-smoothed H I spectrum was taken from the Leiden-Dwingeloo Survey (Hartmann & Burton 1997). Fig. 6. – Upper Panel: Contours of neutral hydrogen column density N(H I) in the vicinity of the MS V component of the Magellanic Stream. Contour levels are $`5\times 10^{18}`$ cm<sup>-2</sup>, with the outermost level N(H I)=$`5\times 10^{18}`$ cm<sup>-2</sup> over the velocity range $`400<v_{\mathrm{LSR}}<250`$ km s<sup>-1</sup>. The H I data were taken from the Leiden-Dwingeloo Survey (Hartmann & Burton 1997); the spatial resolution is restricted by the 1/2-degree sampling grid (with a 1/2-degree beam) employed. Both the III Zw 2 and NGC 7469 sightlines are noted. The peak of the H I emission associated with MS V lies near ($`\mathrm{}`$,$`b`$)$``$($`92^{}`$,$`51^{}`$). Lower Panel: Neutral hydrogen in the immediate vicinity (i.e., the area covered by the marked box in the upper panel) of our III Zw 2 sightline. Contour levels are now $`3\times 10^{18}`$ cm<sup>-2</sup>, with the outermost level N(H I)$`=1\times 10^{18}`$ cm<sup>-2</sup>, over the velocity range $`400<v_{\mathrm{LSR}}<250`$ km s<sup>-1</sup>. The local standard of rest velocities of three independently-confirmed (Hulsbosch & Wakker 1988) High-Velocity Clouds within 4 of our III Zw 2 sightline are labeled. The lack of detectable H I at the position and velocity of our Mg II detection in the spectrum of III Zw 2 sets an upper limit of N(H I)=$`5\times 10^{18}`$ cm<sup>-2</sup>. Fig. 7. – Grid of CLOUDY (Version 90.04 - Ferland 1996) models showing contours of total magnesium to total hydrogen (\[Mg/H\]) as a function of total hydrogen column N(H) and volume n(H) densities. The models satisfy the observational constraints that N(Mg II)$`1\times 10^{13}`$ cm<sup>-2</sup> and N(H I)$`<5\times 10^{18}`$ cm<sup>-2</sup>, representative constraints imposed by our III Zw 2 analysis. A normally-incident ionizing photon flux of $`\varphi =3\times 10^5`$ photons cm<sup>-2</sup> s<sup>-1</sup> was assumed, consistent with that expected at MS V, according to the models of Bland-Hawthorn & Maloney (1999a,b). The shaded region corresponds to the (unfortunately) small part of parameter space excluded by the H I column density constraint. Fig. 8. – Velocity stack showing H I (upper panel) and Mg II $`\lambda `$2796 (middle and lower panels) Galactic (G) and Magellanic Stream (MS) features along the III Zw 2 sightline. The Hanning-smoothed H I spectrum was collected at the NRAO 140 and kindly provided by Ken Sembach prior to publication (Murphy, Sembach & Lockman 2000).
warning/0007/math0007076.html	ar5iv	text	# Invariant Rings and Quasiaffine Quotients. ## 1. Introduction The fourteenth of Hilbert’s famous problems () is the following. Let $`K/L`$ and $`L/k`$ be field extensions, and $`AK`$ be a finitely generated $`k`$-algebra. Does this imply that $`AL`$ is also a finitely generated $`k`$-algebra? This problem was motivated by the following special case: Let $`k`$ be a field and $`GGL(n,k)`$ a subgroup. Is the ring of invariants $`k[x_1,\mathrm{},x_n]^G`$ a finitely generated $`k`$-algebra? (This is a special case: Take $`K=k(x_1,\mathrm{},x_n)`$ and $`L=K^G`$.) For reductive groups this is indeed the case. This was already shown by Hilbert. However, for non-reductive groups there is the celebrated counterexample of Nagata ). Popov deduced from Nagata’s example that for every non-reductive algebraic group $`G`$ there exists an affine $`G`$ variety such that the ring of invariants is not finitely generated . In 1990, a new counterexample was found by Roberts . Later, further counterexamples were obtained by Deveney and Finston and by A’Campo-Neuen . Recently, Daigle and Freudenburg constructed examples in dimension $`6`$ and $`5`$ (,). Reformulated in a more geometric fashion, Hilbert’s 14th problem ask whether the ring of invariant functions is necessarily isomorphic to the ring of regular functions on some affine variety. From this point of view it is maybe not too surprising that the answer is negative in general. Quotients of affine varieties by actions of (non-reductive) algebraic groups are often quasi-affine without being affine, and for arbitrary quasi-affine varieties the ring of regular functions is not necessarily finitely generated (see e.g. ,,). Thus even if the ring of invariants is not isomorphic to the ring of regular functions on an affine variety it nevertheless may be isomorphic to the ring of regular functions on a quasi-affine variety. The purpose of this note is to demonstrate that this is indeed always the case. Actually we show that the a $`k`$-algebra occurs as the ring of invariants for some affine $`G`$-variety if and only if it is isomorphic to the algebra of regular functions on some quasi affine variety. ###### Theorem 1. Let $`k`$ be a field and $`R`$ an integrally closed $`k`$-algebra. Then the following properties are equivalent: 1. There exists an irreducible, reduced $`k`$-variety $`V`$ and a subgroup $`GAut_k(V)`$ such that $`Rk[V]^G`$. 2. There exists a quasi-affine irreducible, reduced $`k`$-variety $`V`$ such that $`Rk[V]`$. 3. There exists an affine irreducible, reduced $`k`$-variety $`V`$ and a regular action of $`G_a=(k,+)`$ on $`V`$ defined over $`k`$ such that $`Rk[V]^{G_a}`$. If $`char(k)=0`$, these properties are futhermore equivalent to the following: 1. There exists a finitely generated, integrally closed $`k`$-algebra $`A`$ and a locally nilpotent derivation $`D`$ on $`A`$ such that $`R\mathrm{ker}D`$. This result is based on the following more general theorem. ###### Theorem 2. Let $`k`$ be a field, $`V`$ an irreducible, reduced, normal $`k`$-variety, and $`L`$ a subfield of the function field $`k(V)`$, containing $`k`$. Let $`R=k[V]L`$. Then there exists a finitely generated $`k`$-subalgebra $`R_0`$ of $`R`$ such that 1. The quotient fields of $`R`$ and $`R_0`$ coincide. 2. For every prime ideal $`p`$ of height one in $`R`$ the prime ideal $`pR_0`$ of $`R_0`$ also has height one. 3. There is an open $`k`$-subvariety $`\mathrm{\Omega }Spec(R_0)`$ such that $`R=k[\mathrm{\Omega }]`$ (as subsets of $`Q(R)`$). These results can be used to construct some “quasi-affine” quotient for a group action on an algebraic variety. ###### Theorem 3. Let $`k`$ be a field, $`V`$ an irreducible, reduced, normal $`k`$-variety and $`GAut(V)`$. Then there exists a quasi-affine $`k`$-variety $`Z`$ and a rational map $`\pi :VZ`$ such that 1. The rational map $`\pi `$ induces an inclusion $`\pi ^{}:k[Z]k[V]`$. 2. The image of the pull-back $`\pi ^{}(k[Z])`$ coincides with the ring of invariant functions $`k[V]^G`$. 3. For every affine $`k`$-variety $`W`$ and every $`G`$-invariant morphism $`f:VW`$ there exists a morphism $`F:ZW`$ such that $`F\pi `$ is a morphism and $`f=F\pi `$. We may also translate our results in the language of category theory (also know asa “general nonsense”) and deduce the following. ###### Theorem 4. For a field $`k`$ let $`𝒱_k`$ denote the category whose objects are irreducible reduced normal $`k`$-varieties and whose morphisms are those dominant rational maps for which the pull-back of every regular function is again regular. Let $`𝒬_k`$ denote the full sub-category whose objects consist of all quasiaffine such varieties. Then for every object $`V𝒱_k`$ and every subgroup $`GAut_{𝒱_k}(V)`$ the functor $`Mor_{𝒱_k}(V,)^G`$ is representable in the category $`𝒬_k`$. ## 2. Preparations Following ideas of Nagata we employ the notions of Krull rings and ideal transforms as algebraic tools for our constructions. We give proofs for some basic facts although they are well known. This is for the benefit of being self-contained and because the proofs are so short. ### 2.1. Krull rings ###### Definition 1. An integral domain $`R`$ is called a Krull ring if there is a family $`F`$ of discrete valuations on the quotient field $`K`$ of $`R`$ such that $`R=\{fK:v(f)0vF\}`$ and $`\{vF:v(f)0\}`$ is finite for every $`fK`$. For basic facts on Krull rings, see . Noetherian integral domains integrally closed in their quotient fields are Krull rings. Intersections of Krull rings inside a fixed field are again Krull rings. For any Krull ring the family of valuations $`F`$ necessarily contains valuations associated to all prime ideals of height one, on the other hand for a Krull ring $`F`$ can be choosen as the set of all discrete valuations associated to prime ideals of height one. Next, we recall that for every $`k`$-variety $`V`$ the ring of functions $`k[V]`$ is a Krull ring and for every Krull ring $`R`$ and every group $`G`$ acting on $`R`$ by ring automorphisms $`R^G`$ is again a Krull ring. ###### Lemma 1. Let $`k`$ be a field, $`V`$ an irreducible, reduced and normal $`k`$-variety. Then $`k[V]`$ is a Krull ring. ###### Proof. Let $`(U_i)_{iI}`$ be a collecting of affine $`k`$-varieties covering $`V`$. For every $`iI`$ the ring $`k[U_i]`$ is noetherian and integrally closed and therefore a Krull ring. Now $`k[V]`$ equals the intersection of all the $`k[U_i]`$ considered as subrings of the function field $`k(V)`$. Hence $`k[V]`$ is a Krull ring. ∎ ###### Lemma 2. Let $`R`$ be a Krull ring and $`GAut(R)`$. Then $`R^G`$ is a Krull ring. ###### Proof. This is immediate, because $`R^G`$ is the intersection of two Krull rings, namely $`R`$ and $`Q(R)^G`$ where $`Q(R)`$ denotes the quotient field of $`R`$. ∎ ### 2.2. The Ideal transform ###### Definition 2. Let $`R`$ be an integral domain, $`I`$ an ideal. Then the $`I`$-transform of $`R`$ is defined as $`S(I,R)=\{xK:n:x(I^n)R\}`$ where $`K`$ denotes the quotient field of $`R`$. For Krull rings the ideal transform has particular nice descriptions in geometric as well as in algebraic form. First we explain the geometric description. ###### Lemma 3. Let $`R`$ be a Krull ring, $`I`$ an ideal. Then $`S(I,R)`$ equals the ring of regular functions on $`SpecRV(I)`$ (both considered as subsets of the quotient field of R). ###### Proof. Let $`K`$ denote the quotient field of $`R`$. Let $`fS(I,R)`$ and $`pX=SpecRV(I)`$. Since $`pV(I)`$, the prime ideal $`p`$ of $`R`$ does not contain $`I`$. Hence there is an element $`hIp`$. By definition $`fh^nR`$ for some natural number $`n`$. Since $`hp`$, this implies $`fR_p`$. We will now show the converse. Let $`f`$ be a regular function on $`X`$, i.e., $`fR_ppI`$, and let $`\mathrm{\Lambda }=\{vF:v(f)>0\}`$. For $`v\mathrm{\Lambda }`$ consider the associate prime ideal $`p_v=\{gR:v(g)<0\}`$. Now $`v(f)>0`$ implies that $`fR_{p_v}`$, hence $`Ip_v`$. Thus $`v(x)<0`$ for all $`xI`$ and $`v\mathrm{\Lambda }`$. Since $`\mathrm{\Lambda }`$ is finite, it follows that $`fI^nR`$ for $`n`$ sufficiently large, i.e. $`fS(I,R)`$. ∎ Next we come to the algebraic description of the ideal transform. ###### Lemma 4. Let $`R`$ be a Krull ring, $`K`$ its quotient field, $`F`$ the set of all discrete valuations corresponding to prime ideals of height one in $`R`$ and $`I`$ an ideal in $`R`$. Then $`S(I,R)=\{xK:v(x)0vF^{}\}`$ with $`F^{}=\{vF:Ip_v\}`$. ###### Proof. Let $`xS(I,R)`$ and $`vF^{}`$. Then there is an element $`fI`$ such that $`v(f)=0`$. Now $`xI^nR`$ implies $`v(x)0`$. Conversely, assume $`x`$ is an element in $`K`$ such that $`v(x)0`$ for all $`vF^{}`$. Let $`\mathrm{\Lambda }=\{vF:v(x)>0`$. Then $`\mathrm{\Lambda }`$ is finite and $`Ip_v`$ for all $`v\mathrm{\Lambda }`$. This implies $`xI^nR`$ for $`n`$ sufficiently large. ∎ We will need the fact that for prime ideals of height one the ideal transform is non-trivial. ###### Lemma 5. Let $`R`$ be a Krull ring, $`p`$ a prime ideal of height one. Then $`S(p,R)R`$. ###### Proof. Let $`K`$ be the quotient field of $`R`$, $`R=\{xK:v(x)0vF\}`$, $`p=\{xR:v_0(x)<0\}`$. Choose $`xP`$ and define $`\mathrm{\Lambda }=\{vF:v(1/x)>0\}\{v_0\}`$. Now $`\mathrm{\Lambda }`$ is finite, and $`p_v`$ being a prime ideal of height one for every $`v`$ implies that $`p_v`$ is not contained in $`p`$ for any $`v\mathrm{\Lambda }`$. Hence there is an element $`yR`$ with $`yp`$ but $`v(y)<0`$ for all $`v\mathrm{\Lambda }`$. Then $`y^nxS(p,R)R`$ for $`n`$ sufficiently large. ∎ ### 2.3. In the situation in which we are interested prime ideals of height one are determined by their zero sets. ###### Lemma 6. Let $`k`$ be a field, $`V`$ an integral, normal $`k`$-variety, $`R`$ a $`k`$-sub algebra of $`k[V]`$ such that $`R`$ is a Krull ring and $`R=Q(R)k[V]`$. Let $`p`$ be a prime ideal of height one in $`R`$. Then $`p=I(Z(p))R`$, where $`I(Z(p))`$ denotes the set of all $`fk[V]`$ vanishing on the zero set of $`p`$. Moreover, $`Z(p)`$ contains an irreducible component $`Z_0`$ such that $`I(Z_0)=p`$. ###### Proof. Let $`v`$ denote the discrete valuation corresponding to $`p`$. By $`()`$ there is an element $`xQ(R)R`$ with $`v(x)>0`$ and a number $`N`$ such that $`xp^NR`$. Thus $`x`$ defines a rational function on $`V`$ whose poles are contained in $`Z(p)`$. Now assume $`gI(Z(p))`$ but $`gp`$. This would imply that $`v(g)=0`$ and $`g`$ vanishes on $`Z(p)`$. Since the poles of $`x`$ are contained in $`Z(p)`$, it follows that $`g^mxk[V]`$ for $`m`$ large enough. On the other hand $`v(g^mx)=v(x)>0`$ implies $`g^mxR`$. Thus the assumption $`I(Z(p))p`$ yields a contradiction to the requirement that $`R=Q(R)k[V]`$. For the final statement, let $`(Z_i)_{iI}`$ denote the irreducible components of $`Z(p)`$. Then $`I(Z(p))=_iI(Z_i)`$. Since both $`I(Z(p))`$ and all of the $`I(Z_i)`$ are prime ideals, it follows that there exists an $`iI`$ such that $`I(Z_i)=I(Z(p))=p`$. ∎ ###### Lemma 7. Let $`k`$ be a field and $`V`$ a $`k`$-variety. For every $`k`$-algebra $`Rk[V]`$ let $`_R`$ denote the equivalence relation on the geometric ($`\overline{k}`$-rational) points of $`V`$ given by $`x_Ry`$ iff $`f(x)=f(y)`$ for all $`fR`$. Then for every $`k`$-sub algebra $`Rk[V]`$ there exists a finitely generated $`k`$-algebra $`R_0R`$ such that $`_R`$ coincides with $`_{R_0}`$ Furthermore $`R_0`$ can be choosen as a Krull ring, if $`R`$ is a Krull ring. ###### Proof. The equivalence relation $`_R`$ defines a subset $`E_RV\times V`$ via ($``$) $$E_R=\{(x,y)V\times V:x_Ry\}=\{(x,y):f(x)=f(y)fR\}$$ Then $`E_R`$ is the $`k`$-sub variety defined by the radical of the ideal of $`k[V\times V]`$ generated by $`\pi _1^{}f\pi _2^{}f`$ with $`f`$ running through $`R`$. If $`R^{}R`$ are two $`k`$-sub algebras of $`k[V]`$, then $`E_RE_R^{}`$. Since $`V\times V`$ is noetherien it follows that for any $`k`$-sub algebra $`Rk[V]`$ there exists a finitely generated $`k`$-sub algebra $`R_0R`$ such that $`E_R=E_{R_0}`$. Finally, if $`R`$ is a Krull ring, then $`R`$ is integrally closed in its quotient field. Hence the integral closure of $`R_0`$ is again contained in $`R`$. Furthermore the integral closure of $`R_0`$ is again finitely generated as a $`k`$-algebra. Thus we may assume that $`R_0`$ is integrally closed. But integrally closed finitely generated $`k`$-algebras are Krull rings. ∎ Let us now fix some key assumptions. ###### Key Assumptions. In the sequel, $`k`$ is a field, $`V`$ an irreducible, reduced and normal $`k`$-variety, $`Rk[V]`$ a $`k`$-sub algebra such that $`Q(R)k[V]=R`$, $`R_0`$ is a finitely generated $`k`$-algebra such that $`R_0R`$, $`Q(R_0)=Q(R)`$, $`R_0`$ is integrally closed in $`Q(R_0)=Q(R)`$ and $`E_R=E_{R_0}`$ where $`E_R`$ is defined as in $`()`$ above. We will prove the statements of theorem 2 hold for such a choice of $`R_0`$. ###### Lemma 8. Under the key assumptions $`height(pR_0)=1`$ for every prime ideal $`p`$ of height one in $`R`$. ###### Proof. Let $`W=Spec(R_0)`$, this is an affine $`k`$-variety. The prime ideals of height one in $`R_0`$ correspond to the irreducible hypersurfaces in $`W`$. Let $`\tau :VW`$ denote the morphism induced by $`R_0k[V]`$. Let $`p`$ be a prime ideal of height one in $`R`$. We have seen above that there is an irreducible subvariety $`ZV`$ such that $`I(Z)R=p`$. From $`E_R=E_{R_0}`$ we infer that there exists an irreducible subvariety $`YZ`$ such that $`\tau \|_Y:YW`$ is generically quasi-finite and $`Z`$ is the smallest $`E_R`$-saturated subvariety containing $`Y`$. Now let $`X`$ be an irreducible subvariety of $`V`$ such that $`dim(X)=dim(Y)+1`$ and $`XZ(p)`$. Let $`I=I(X)R`$. Then $`pI`$. On the other hand $`II(Y)`$ and $`I(Y)R=I(Z)R=p`$. Thus $`height(p)=1`$ implies $`I=\{0\}`$. It follows that $`\tau (X)`$ must be Zariski-dense in $`W=Spec(R_0)`$. Since $`dim(X)=dim(Y)+1`$ it now follows from $`\tau \|_Y`$ being generically quasi-finite that $`\overline{\tau (Y)}`$ either is a hypersurface or equals $`W`$. The latter is excluded since $`p\{0\}`$ and $`YZ(p)`$. Thus $`\overline{\tau (Y)}=\overline{\tau (Z)}`$ has to be a hypersurface implying that $`height(pR_0)=1`$. ∎ ###### Lemma 9. Under the key assumptions for any two distinct prime ideals of height one $`p_1,p_2R`$ we have $`p_1R_0p_2R_0`$. ###### Proof. This is an immediate consequence of $`E_R=E_{R_0}`$ and $`I(Z(p_i))R=p_i`$. ∎ ###### Corollary 1. Let $`F`$ resp. $`F_0`$ denote the set of discrete valuations of $`K=Q(R)`$ corresponding to prime ideals of height one in $`R`$ resp. $`R_0`$. Then $`FF_0`$. ###### Lemma 10. The set $`F_0F`$ is finite. ###### Proof. Let $`vF_0F`$, $`H_0W=Spec(R_0)`$ the corresponding hypersurface and $`H=\tau ^1(H_0)`$. We claim that $`\tau (H)`$ is not Zariski-dense in $`H_0`$. Indeed, if it were Zariski-dense, there would exist an irreducible subvariety $`H^{}`$ with $`\overline{\tau (H^{})}=H`$. But this implies $`I(H^{})R_0=I(H_0)`$. Now let $`p`$ be a prime ideal of height one contained in $`I(H^{})`$ (such an ideal exists, because $`I(H^{})`$ is a prime ideal and $`R`$ is a Krull ring). Then $`pR_0`$ is a prime ideal of height one. Since $`I(H_0)`$ is of height one, it follows that $`I(H_0)=pR_0`$ contrary to our assumption $`vF`$. Thus $`H`$ is a hypersurface in $`W`$ with $`\tau (\tau ^1(H))`$ not being dense in $`H`$. Since $`\tau :VW`$ is dominant, there are only finitely many such hypersurfaces in $`W`$. ∎ ## 3. Proofs of the theorems ### 3.1. Proof of theorem 2. ###### Proof. Let $`k`$, $`V`$, $`L`$ and $`R`$ be as in the theorem. Note that $`L`$ is a finitely generated field extension of $`k`$, because $`k(V)/k`$ is finitely generated and $`Lk(V)`$. By lemma 7 there is a finitely generated $`k`$-algebra $`R_1`$ with $`R_1R`$, $`E_R=E_{R_1}`$ and $`R_1`$ being a Krull ring. We may adjoin finitely many further elements of $`R`$ to $`R_1`$ and thereby assume that the quotient fields of $`R`$ and $`R_1`$ coincide. Then we choose $`R_0`$ as the integral closure of $`R_1`$ in $`L`$. Since $`R`$ is integrally closed, we have $`R_0R`$. Furthermore $`E_R=E_{R_0}`$ and $`R_1`$ is again finitely generated, because it is the integral closure of a finitely generated $`k`$-algebra. Thus $`R_0`$ fulfills the “Key Assumptions”. Statement $`(1)`$ of the theorem is clear, and statement $`(2)`$ follows from lemma 8. Next we define $`F`$ and $`F_0`$ as in the corollary above. By lemma 10 the difference set $`F_0F`$ is finite. Hence we may define an ideal of $`R_0`$ by $$I=\mathrm{\Pi }_{v\{F_0F\}}p_v$$ with $`p_v=\{xR_0:v(x)<0\}`$. Since $`p_\mu `$ is of height one for every $`\mu F_0`$, it is clear that $`Ip_v`$ for $`vF`$. Therefore $`R=S(I,R_0)`$ by lemma 4. Finally, statement $`(3)`$ of the theorem follows with the aid of lemma 3. ∎ Before starting the proof of theorem 3 we need the subsequent lemma. ###### Lemma 11. Let $`k`$ be a field, $`V`$ a $`k`$-variety and $`W`$ a quasi-affine $`k`$-variety. Then there is a one-to-one correspondence between $`k`$-algebra homomorphisms $`\varphi :k[W]k[V]`$ and rational maps $`f:VW`$ with $`f^{}k[W]k[V]`$. ###### Proof. Given a $`k`$-algebra homomorphism $`\varphi :k[W]k[V]`$, choose a finitely generated $`k`$-sub algebra $`Ak[W]`$ such that the quotient fields of $`k[W]`$ and $`A`$ coincide. The restriction of $`\varphi `$ yields a $`k`$-morphism $`F`$ from $`V`$ to $`Z=Spec(A)`$ and the inclusion $`Ak[W]`$ yields a birational morphism $`\tau :WZ`$. Now $`\tau ^1F`$ is the desired rational map. ∎ ### 3.2. Proof of theorem 3. ###### Proof. We apply theorem 2 with $`L=k(V)^G`$ and set $`Z=\mathrm{\Omega }`$. There is an inclusion $`k[Z]=k[\mathrm{\Omega }]=Rk[V]`$. By the preceding lemma this induces a rational map $`\pi :VZ`$ with $`\pi ^{}k[Z]k[V]`$. Since $`R=k[V]L=k[V]^G`$, it follows that $`\pi ^{}k[Z]=k[V]^G`$. Finally note that for every affine $`k`$-variety $`W`$ and every $`G`$-invariant morphism $`f:VW`$ we obtain an inclusion $`f^{}k[W]k[V]^Gk[Z]`$ which implies that there exists a morphism $`F:ZW`$ such that $`f=F\pi `$. ∎ ### 3.3. Proof of theorem 1. ###### Proof. The implication $`(1)(2)`$ follows from theorem 3. $`(3)(1)`$ is trivial and $`(2)(3)`$ follows from the proposition below. Finally, $`(3)(4)`$ for the case of characteristic zero is implied by the well-known correspondence between locally nilpotent derivations and $`G_a`$-actions on affine varieties in characteristic zero. ∎ ###### Proposition 1. Let $`k`$ be a field and let $`V`$ be a normal quasi-affine $`k`$-variety. Then there exists a normal affine $`k`$-variety $`W`$ and a regular action of the additive group $`G=G_a`$ defined over $`k`$ on $`W`$ such that $`k[V]k[W]^G`$. ###### Proof. Let $`VY`$ be an open embedding in a normal affine $`k`$-variety $`Y`$, and let $`S=YV`$. Let $`D`$ denote the union of codimension $`1`$-components of $`S`$ and choose a regular function $`f_1`$ on $`Y`$ such that $`D`$ is contained in the zero set of $`f_1`$. Then choose a regular function $`f_2`$ on $`Y`$ such that $`f_2`$ vanishes on $`D`$, but does not vanish on any irreducible component of $`Z(f_1)D`$. If $`char(k)=p>0`$, we replace $`f_i`$ by $`f_i^{p^N}`$ for a sufficiently large $`N`$. In this way we may assume that both functions $`f_i`$ are defined over a finite Galois extension $`k^{}/k`$ with Galois group $`\mathrm{\Gamma }`$. Now we may replace $`f_i`$ by $`\mathrm{\Pi }_{\sigma \mathrm{\Gamma }}{}_{}{}^{\sigma }f_{i}^{}`$. Therefore we may assume that both $`f_i`$ are defined over $`k`$. We obtain a $`k`$-morphism $`f:Y𝔸^2`$. By construction $`D`$ is the union of codimension $`1`$-components of $`E=f^1\{(0,0)\}`$. Since regular functions extend through subvarieties of codimension at least 2 on normal varieties, it follows that $$k[V]k[YD]k[\mathrm{\Omega }]$$ for $`\mathrm{\Omega }=YE`$. Next we consider $$S=\{\left(\begin{array}{cc}a& b\\ c& d\end{array}\right):adbc=1\}$$ and the natural projection $`\pi :S𝔸^2`$ given by $`\pi :(a,b,c,d)(a,b)`$. This realizes $`𝔸^2\{(0,0)\}`$ as the quotient of $`S`$ by the $`G_a`$-action given by $$t:\left(\begin{array}{cc}a& b\\ c& d\end{array}\right)\left(\begin{array}{cc}a& b+ta\\ c& d+tc\end{array}\right)$$ Now the fiber product $`W=Y\times _{𝔸^2}S`$ is an affine variety and $`W\mathrm{\Omega }\times _{𝔸^2\{(0,0)\}}S`$, because the image of $`S`$ in $`𝔸^2`$ is contained in $`𝔸^2\{(0,0)\}`$ and $`\mathrm{\Omega }=f^1(𝔸^2\{(0,0)\})`$. The $`G_a`$-action on $`S`$ induces a $`G_a`$-action on the fibered product $`W`$ and evidently $`k[\mathrm{\Omega }]k[W]^{G_a}`$. ∎ ### 3.4. Proof of theorem 4 ###### Proof. Let $`X𝒱_k`$ and $`GAut_{𝒱_k}(X)`$. Each element in $`G`$ is a birational self-map of $`X`$ such that the induced field automorphism of $`k(X)`$ stabilizes $`k[X]`$. In particular $`GAut(k[X])`$. Due to theorem 2 there exists an object $`\mathrm{\Omega }𝒬_k`$ and a $`𝒱_k`$-morphism $`\pi :X\mathrm{\Omega }`$ such that $`k[\mathrm{\Omega }]\stackrel{}{}k[X]^G`$. Consider now an object $`W𝒬_k`$ with a $`G`$-invariant $`𝒱_k`$-morphism $`f:XW`$. Then $`f`$ is dominant rational map with $`f^{}(k[W])k[X]^G`$. Thus we obtain an $`k`$-algebra homomorphism $`f^{}:k[W]k[\mathrm{\Omega }]=k[X]^G`$, which by lemma 11 induces a $`𝒬_k`$-morphism from $`\mathrm{\Omega }`$ to $`W`$. Therefore $`Mor_{𝒱_k}(X,W)^GMor_{𝒬_k}(\mathrm{\Omega },W)`$. ∎ ## 4. An example Let $`k`$ be a field of characteristic zero. In Daigle and Freudenburg gave an example of a locally nilpotent derivation $`D`$ of $`k[𝔸^5]`$ such that $`\mathrm{ker}D`$ is not finitely generated. This is the lowest-dimensional example known today. In coordinates $`x,s,t,u,v`$ the derivation $`D`$ can be written as $$D=x^3\frac{}{s}+s\frac{}{t}+t\frac{}{u}+x^2\frac{}{v}$$ The associated group action of the additive group $`G_a`$ is given by $$\mu (r):(x,s,t,u,v)(x,s+rx^3,t+rs+\frac{r^2}{2}x^3,u+rt+\frac{r^2}{2}s+\frac{r^3}{6}x^3,v+rx^2)$$ The action is free outside the set of fixed points $$(𝔸^5)^\mu =\{(0,0,0,v,u):u,vk\}$$ An explicit calculation shows that the following regular functions on $`𝔸^5`$ are invariant: $`\varphi _1`$ $`=x`$ $`\varphi _2`$ $`=2x^3ts^2`$ $`\varphi _3`$ $`=3x^6u3x^3ts+s^3`$ $`\varphi _4`$ $`=xvs`$ $`\varphi _5`$ $`=x^2tss^2v+2x^3tv3x^5u`$ $`=`$ $`(\varphi _2\varphi _4\varphi _3)/\varphi _1`$ $`\varphi _6`$ $`=18x^3tsu+9x^6u^2+8x^3t^3+6s^3u3x^6t^2s^2`$ $`=`$ $`(\varphi _2^3+\varphi _3^2)/\varphi _1^6`$ Further explicit calculations yields the following: ###### Lemma 12. Let $`V=\{w=(w_1,\mathrm{},w_6)𝔸^6:w_5w_1=w_2w_4w_3,w_6w_1^6=w_2^3+w_3^2\}`$. Then $`V`$ is an affine subvariety of $`𝔸^6`$ with $`SingV=\{w𝔸^6:w_1=w_2=w_3=0\}`$. $`SingV`$ is a Weil divisor of $`V`$. The regular functions $`(\varphi _i)_{i=\mathrm{1..6}}`$ defined above give an invariant morphism $`\varphi :𝔸^5V`$ such that 1. $`\varphi `$ has rank $`4`$ outside $`E=\{x=s=0\}`$. 2. $`\varphi `$ maps $`𝔸^5E`$ surjectively on $`VSingV`$. 3. For every $`pVSingV`$ the fiber $`\varphi ^1(p)`$ coincides with a $`G_a`$-orbit. It follows that $$k[𝔸^5]^\mu =k[𝔸^5E]^\mu k[VSingV]$$ Thus the ring of invariants is indeed isomorphic to the ring of regular functions of a quasi affine variety, namely $`VSingV`$. ## 5. Appendix Naturally one would prefer having a regular morphism in $`()`$ instead of having merely a rational map. One might hope for a general result implying that this map is automatically regular, and endeavor to prove a statement like the following. ###### Falsity. Let $`V`$ , $`W`$ be affine varieties, $`f:VW`$ be a regular map, $`HW`$ a hypersurface and assume that $`f^{}(k[WH])k[V]`$. Then $`f(V)WH`$ (implying that one has a regular morphism from $`V`$ to $`WH`$ and not merely a rational map.) But this is wrong: ###### Example. Let $`V=\{(x_1,x_2,x_3)𝔸^3:x_10\}`$, $`W=\{(z_1,z_2,z_3,z_4)𝔸^4:z_1z_4=z_2z_3\}`$, $`H=\{(z_1,z_2,z_3,z_4):z_1=z_2=0\}`$ and $`f:VW`$ given by $`f(x_1,\mathrm{},x_3)=(x_1x_2,x_1x_3,x_2,x_3)`$. Then $`H`$ is a hypersurface in $`W`$, but $`f^1(H)=\{(x_1,0,0):x_10\}`$ is a curve, and therefore has codimension two in $`V`$. As a consequence $`f^{}k[WH]k[Vf^1(H)]=k[V]`$, although $`f(V)H=\{(0,0,0,0)\}\mathrm{}`$.
warning/0007/nucl-th0007048.html	ar5iv	text	# 1 Introduction ## 1 Introduction The subject of this talk can be introduced by commenting on the four conceptual items mentioned in the title, but in reverse order. First, high-energy nuclear collisions. The goal is to create quark-gluon plasma. So far the conventional signatures are aimed at detecting the primordial quark matter. Suppose that is found. Then what? If we can learn from the condensed-matter physics, then we know that critical phenomena form a vibrant area of research. That leads to the second item in the title: critical behavior. Our concern will be the quark-to-hadron phase transition (PT). The hadron-to-quark PT is, in my view, too turbulent in heavy-ion collisions to lend itself to systematic studies, despite some recent attempt to use percolation as a model to characterize it. The third item in the title is “measures.” For a thermal system the heat capacity can be a measure that exhibits the critical behavior $`C\left\|TT_c\right\|^\alpha `$, while for a magnetic system the magnetic susceptibility behaves as $`\chi \left\|TT_c\right\|^\gamma `$. What is the measure for quark-hadron PT? Is there a measure, call it $`\mu `$, which behaves as $`\mu \left\|TT_c\right\|^\kappa `$? If so, how do we find it? That brings us to the last item: observable. Temperature is not directly observable. (The transverse momentum $`p_T`$ is an unreliable, gross indicator of the average phenomenon.) Moreover, the system is not static and the emitted hadrons may undergo final-state interactions that can smear the signature that is to be observed. The aim of this talk is to show how all these issues are addressed and to indicate where the answers are to be found. The measure that I believe will convey the information about quark-hadron PT is about fluctuations, not only bin-to-bin fluctuations, but also event-to-event fluctuations. A way to illustrate this belief is to consider the annual precipitation (AP) of two places with widely different climatic characteristics. The AP of certain regions in Oregon and Florida are very similar, about 40<sup>′′</sup> per year, but their rainfall patterns are nearly opposite: drizzly in Oregon, stormy in Florida. The problem is that in determining the PA the daily rainfall is integrated over the entire year, and then averaged over all years. That is exactly how the average multiplicity of heavy-ion collisions is determined: integrated over the whole hadronization time of an event (which can take over 30 fm/c), and then averaged over all events. Clearly, to learn about the nature of rainfall or the properties of hadronization, it is necessary to study spatial and temporal fluctuations. ## 2 The Problem If quark-gluon plasma is indeed formed in a heavy-ion collision, it is reasonable to assume that the plasma occupies a cylindrical volume, expanding rapidly in the longitudinal direction, and more slowly radially. The medium most likely has a temperature profile that is high at the center of the transverse plane and decreases with increasing radius due to the radial expansion and cooling. When $`T`$ reaches $`T_c`$ at the surface, the quark matter undergoes a $`PT`$ into hadrons. Thus the characteristics of critical behavior are to be found on the cylindrical surface at any given time of the evolutionary history of the system. That is like raining at a particular instant of observation. If we collect all the hadrons emitted throughout the long hadronization processes, the features of critical behavior are likely to overlap and get averaged out. It is therefore important to select a narrow $`\mathrm{\Delta }t`$ cut in the hadronization time. If the emission time and the transverse momentum $`p_T`$ of the hadrons are correlated in a one-to-one relationship, then a narrow $`\mathrm{\Delta }p_T`$ cut would accomplish what we need. Unfortunately, that relationship is not one-to-one. On the contrary, at each emission time the emitted particles have a very wide $`p_T`$ distribution. That wide distribution turns out to provide us with a way to proceed. Suppose we make a narrow $`\mathrm{\Delta }p_T`$ cut and study the hadronic patterns in the $`\eta `$-$`\varphi `$ 2D space orthogonal to $`\stackrel{}{p}_T`$. Then at any emission time only a small fraction of the emitted hadrons would enter the narrow $`\mathrm{\Delta }p_T`$ window. This random selection process at each time frame thus decorrelates the hadrons collected over the entire hadronization process. The measured hadronic patterns in the $`\eta `$-$`\varphi `$ space is then the superposition of small parts of many configurations (which we call configuration mixing), each of which we shall stimulate before mixing. The object of our analysis is to find the scaling property of each configuration that exhibits the critical behavior, and then to see whether it survives the mixing process. It is well known that a system at the critical point forms clusters of all sizes that exhibit Kadanoff scaling. The details of the dynamical system are not important, since the critical behavior depends mainly on the dimension and symmetry of the system. On the basis of that universality we shall use the Ising model as a simple device to simulate the configuration of hadronic clusters formed on the surface of the plasma cylinder at each time frame of hadronization. There are reasons to believe that the QCD dynamics for the chiral PT suggests a cross-over near the second-order PT that belongs to the same universality class as the Ising system . We shall define hadron formation on the Ising lattice in such a way as to be able to simulate a cross-over. From the simulated configurations we can see the hadronic clusters randomly scattered over a background of voids. In our first-attempt analysis of the patterns formed, we shall ignore the complications that arise from the $`p_T`$ distribution and consider directly the scaling properties of the clusters and voids in uncorrelated configurations that we shall simulate. ## 3 Cluster Analysis The use of the Ising model on a 2D lattice to simulate second-order PT is standard . To have a cross-over in quark-hadron PT without the use of an external field in the Ising Hamiltonian is not standard . The method uses a cell on the lattice, of dimension $`ϵ\times ϵ`$ where $`ϵ`$ may be taken to be 4, to define hadron density $`\rho _c`$ at the cth cell $$\rho _c=\lambda \varphi _c^2\theta \left(\varphi _c\right)$$ (1) where $$\varphi _c=\underset{jc}{}\text{sgn}\left(m_L\right)\sigma _j,m_L=\underset{jL^2}{}\sigma _j.$$ (2) The first sum in (2) is over all sites in a cell having site-spins $`\sigma _j`$; the second sum is over all lattice sites, giving the total magnetization $`m_L`$. We use the direction of $`m_L`$ of each configuration to serve as the direction of an external field and define the order parameter $`\varphi _c`$ as the cell spin at $`c`$ relative to that direction. $`\lambda `$ in (1) is a parameter relating hadron density to Ising spins; all measurable quantities should be insensitive to $`\lambda `$. Because of the $`\theta `$-function in (1), $`\rho _c`$ is either positive where hadrons are emitted, or zero where a void exists. There are two quantities we have found that can effectively serve as measures of fluctuations of hadronic clusters from bin to bin and from configuration to configuration. One is $`J(M,T)`$ defined $$J(M,T)=\frac{\rho _k}{\rho }\mathrm{}n\frac{\rho _k}{\rho }$$ (3) and the other is $`K(M,T)`$ whose definition is omitted here, but can be found in Ref. . In (3) $`\mathrm{}`$ refers to averaging over all bins \[$`M`$ of them, $`M=(L/\delta )^2`$\] and over all configurations. We discuss only the properties of $`J(M,T)`$ here; those of $`K(M,T)`$ are similar. What we have found about $`J(M,T)`$ is that it is factorizable , i.e., $$J(M,T)=\alpha (T)J_c(M),$$ (4) where $`J_c(M)=J(M,T_c)`$. Thus, by definition, $`\alpha \left(T_c\right)=1`$. Yet, when $`T`$ is away from $`T_c`$, but not too far in $`T<T_c,\alpha (T)`$ satisfies the power-law behavior $$\alpha (T)\left(T_cT\right)^\zeta ,\zeta =1.88.$$ (5) This may be regarded as the critical behavior that we have been searching for. However, $`T`$ is not directly measurable. We can use $`\overline{\rho }`$ instead, where $`\overline{\rho }`$ is the cell density averaged over all cells in one configuration. Then we find $$J(M,\overline{\rho })=\alpha \left(\overline{\rho }\right)J_0(M)$$ (6) with $$\alpha \left(\overline{\rho }\right)\left(\overline{\rho }\overline{\rho }_0\right)^{\overline{\zeta }},\overline{\zeta }=0.97,$$ (7) where $`\overline{\rho }_0`$ is the average density where $`J(M,\overline{\rho }_0)=J_0(M)`$ is the highest. The behavior in (7) is measurable and can be used to signal the fluctuation property of a quark-hadron PT. ## 4 Void Analysis In addition to finding the critical behavior of the hadronic clusters, we can also analyze the scaling properties of the voids where no hadrons appear . First, it is necessary to define precisely what a void is, and how to avoid a single hadron from changing its quantification. To that end we define a bin to be empty if the average density $`\overline{\rho }_k`$ of the kth bin is less than $`\rho _0`$, where $`\rho _0`$ is a parameter that can be varied by the analyst of the data. We do not set $`\rho _0`$ to zero in order to eliminate the small fluctuations due to irregular occurrences of hadrons in or at the edges of a bin. A void is then the area of contiguous empty bins, and its area is therefore $$V_n=\underset{k}{}\theta \left(\rho _0\overline{\rho }_k\right),$$ (8) where the sum is over all connected empty-bins of the nth void. $`V_n`$ is in the unit of number of bins. A quantity less dependent on the size of the bins is the fraction $`x_n=V_n/M`$. For every configuration, there exists a set of void fractions $`\{x_1,x_2,x_3,\mathrm{}x_m\}`$ where $`m`$ is the total number of voids. The comparison of void patterns of different configurations would be difficult if they are characterized by those sets. To ease that comparison, we define the $`G`$ moments $$G_q=\frac{x^q}{x^q},$$ (9) where $`x^q=\left(_{n=1}^mx_n^q\right)/m`$, an average performed for each configuration. Thus $`G_q`$ is a number (for each $`q`$) that characterizes the void pattern of a configuration. At the critical point this $`G_q`$ fluctuates widely from configuration to configuration. It depends on the bin size. At $`T`$ in the vicinity of $`T_c`$ we find that $`G_q`$, averaged over all configurations, satisfies scaling law for a variety of $`\rho _0`$ values, i.e., $$G_qM^{\gamma _q}.$$ (10) That means voids of all sizes occur. Furthermore, the scaling exponents depend on $`q`$ linearly $$\gamma _q=c_0+cq$$ (11) with $`c=0.8`$ at $`T_c`$ and $`\rho _0=20`$. This value of $`\rho _0`$ is less than 8% of the maximum $`\overline{\rho }`$. The value of $`c`$ depends on $`T`$ and $`\rho _0`$ mildly, and can be used to quantify the nature of void patterns. Since $`G_q`$ fluctuates very widely, i.e., the distribution $`P(G_q)`$ is very broad, the moments of $`G_q`$ can reveal that fluctuation. We have studied the derivative of the first moment, which is $$S_q=G_q\mathrm{}nG_q$$ (12) and found that it also satisfies a scaling law $$S_qM^{\sigma _q}$$ (13) whose exponent depends on $`q`$ linearly $$\sigma _q=s_0+sq$$ (14) with $`s=0.76`$ at $`T=T_c`$ and $`\rho _0=20`$. The value of $`s`$ also depends on $`T`$ and $`\rho _0`$ mildly, and therefore can further be used to quantify the nature of the fluctuation of void patterns from configuration to configuration. The dependences of $`c`$ and $`s`$ on $`\rho _0`$ and $`T`$, though mild, can be used to diagnose the temperature range of quark-hadron $`PT`$. For $`\rho _0>20`$ the values of $`c`$ and $`s`$ determined in our analysis vary in a range as much as 20% for $`T`$ varying in the range of $`\pm 2\%`$ around $`T_c`$, but they converge to the values quoted as $`\rho _020`$. If $`c`$ and $`s`$ are not unique for a range of $`T`$ around $`T_c`$, it means that $`\gamma _q`$ and $`\sigma _q`$ are not fixed (for fixed $`q`$) if PT takes place over a range of $`T`$, not exactly at $`T_c`$. Then there would be no scaling behavior. Thus, if the experimental data reveal the loss of scaling as $`\rho _0`$ is increased, it implies that the $`PT`$ takes place over a range of $`T`$. If, however, scaling persists, then $`PT`$ occurs only at one $`T`$, presumably $`T_c`$. ## 5 Complications in Hadronization There are two aspects about hadronization in heavy-ion collisions that one should be concerned about for the type of measures we are proposing. One is the effect due to final-state interaction as the hadrons traverse the hadron gas. The other is the effect of configuration mixing. The former is studied by requiring the produced hadron in each configuration to take $`\nu `$ steps of random walk, since scattering by a hadron gas implies the random deflections of the initial $`\stackrel{}{p}_T`$ that result in random wandering in the $`\eta `$-$`\varphi `$ plane. We have studied the effect of such random walks (for $`\nu `$ up to 6) on $`J(M,T)`$ and found that the factorizability (4) persists and that $`\alpha (T)`$ retains the essential character of (5) . The other complication of configuration mixing due to the non-correspondence between $`\mathrm{\Delta }p_T`$ and $`\mathrm{\Delta }t`$ is more difficult to treat. We have considered only the simple case of no correlation between successive time frames. We mix four independently simulated configurations by taking one quadrant from each and then apply our void analysis to the mixed configurations . We find that the effects on $`\gamma _q`$ and $`\sigma _q`$ are essentially negligible. This problem will be followed by a more realistic treatment of the configuration mixing where the wide $`p_T`$ distribution of the emitted particles at any given time will be taken into account together with the hadron-void correlation in successive time frames. ## 6 Multifragmants Finally, let me make some remarks about “Common Aspects” that is a part of the title of this session. Despite the valiant attempts made by the organizers of this conference to promote interaction between the two subfields encompassed here, cross-fertilization has not been a vibrant phenomenon that has taken place. However, it seems that the ideas about the void analysis discussed here can be translated into the study of multifragmentation in the intermediate mass range. Let me therefore venture a suggestion on an analysis of the fragments with the hope that it may yield a feature of the experimental results more challenging to the model builders than what they have hitherto dealt with. Suppose that for each event an experiment can measure a certain number of fragments whose total mass may be less than the total mass of the nuclei in collision and may fluctuate from event to event. Nevertheless, for all the fragments $`M_i`$ that have been detected, one can calculate the fraction $`x_i=M_i/_iM_i`$. If it is the charges of the fragments that are measured, then let $`x_i=Z_i/_iZ_i`$. Each event is then characterized by a set of fractions $`\{x_1,x_2,x_3,\mathrm{}x_n\}`$. The idea is to go beyond studying the averages and examine the nature of the fluctuation of the fragments detected. Let us define $$H_q(n)=\frac{1}{n}\underset{i=1}{\overset{n}{}}x_i^q$$ (15) where $`n`$ is the total number of fragments measure in an event. Clearly, $`H_0=1`$ and $`H_1=1/n`$. At higher $`q`$, $`H_q(n)`$ are progressively smaller, but are increasingly more dominated by the large $`x_i`$ components. If $`P\left(H_q\right)`$ is the distribution of $`H_q`$ after all events are sampled, then the average is $$H_q=\frac{1}{𝒩}\underset{e=1}{\overset{𝒩}{}}H_q^e(n)=𝑑H_qH_qP(H_q)$$ (16) where $`e`$ labels an event and $`𝒩`$ is the total number of events. $`H_q`$ provides some average information about the multifragments, but to study the fluctuation from event to event one should calculate the normalized moments of moments $$C_{p,q}=H_q^p/H_q^p$$ (17) This $`C_{p,q}`$ may contain too much information. The derivative at $`p=1`$ may already be sufficient to provide some interesting insight into the erratic nature of the fluctuation $$E_q=\frac{d}{dp}C_{p,q}\|_{p=1}=\frac{H_q}{H_q}\mathrm{}n\frac{H_q}{H_q},$$ (18) the study of which may be called erraticity , for convenience. The experimental determination of $`E_q`$ as a function of $`q`$ may pose a serious challenge to the theoretical models, some of which may not contain enough dynamical fluctuations to generate the erraticity revealed by the data. ## 7 Acknowledgments The work described here was carried out at different stages in collaboration with Z. Cao, Y. F. Wu and Q. H. Zhang. It was supported, in part, by the U. S. Department of Energy under Grant No. DE-FG03-96ER40972.
warning/0007/astro-ph0007192.html	ar5iv	text	# Feedback from Protostellar Outflows in Star and Star Cluster Formation ## 1. Introduction Stars, even low-mass stars, are not born quietly. In the process of accreting gas through an accretion disk, each star returns a fraction of the inspiralling material to the interstellar medium at speeds comparable to the escape velocity from the stellar surface. These powerful winds, observable as jets of ionized gas, sweep the ambient material into opposed streams known as bipolar molecular outflows. Because of their intensity, protostellar winds pose severe challenges to those who wish to model star forming regions. On a small scale, the wind from a forming star is capable of stripping gas from its parent molecular core even as that core feeds material to the star. On a larger scale, the outflow driven by a star’s wind can punch through a massive molecular clump that is forming a stellar cluster. There is now plentiful observational evidence that these interactions are common, even ubiquitous, within star-forming clouds. Moreover, there is evidence to support the argument (Norman & Silk 1980) that the feedback from protostellar winds is the primary driver of turbulence within molecular clouds, and may be responsible for establishing Larson’s (1981) line width-size relation in these clouds (McKee 1989). Because these phenomena have profound consequences for the formation of stars and star clusters and for the evolution of molecular clouds, it is important to quantify the ways in which protostellar winds impact their molecular environment. In the past, models of feedback in star formation have either ignored the effects of low-mass stars (e.g., Elmegreen 1983) or considered only spherically symmetric protostellar winds (e.g., Nakano, Hasegawa, & Norman 1995), despite the fact that these winds are clearly jet-like. More sophisticated models of protostellar outflows (such as those of Shu et al. 1991; Masson & Chernin 1993; Li & Shu 1996) have not been applied to the basic questions of feedback. Partly, this has been caused by a lack of consensus (see Richer et al. 2000) as to which of the several classes of models for outflows best describes their dynamics. We shall argue that there is a single model for protostellar outflows that is both motivated by theories of protostellar winds and validated by its ability to explain some of the common, otherwise mysterious features of protostellar outflows. Although this model cannot explain all the details of these flows, it provides a stable platform from which to launch an exploration wind-cloud interactions. We shall then use this model to address basic questions of star formation: the stellar initial mass function (IMF); the efficiency of star formation in clusters; and the dynamical evolution of star-forming gas during stellar cluster formation. The work presented here has been conducted in collaboration with Chris McKee and appears or is intended to appear in: Matzner (1999), Matzner & McKee (1999a), Matzner & McKee (2000), and Matzner, Bertoldi, & McKee (2000). ## 2. Common Properties of Bipolar Molecular Outflows Protostellar molecular outflows typically share several common features, as discussed by Masson & Chernin (1993) and Lada & Fich (1996). These are: a roughly linear position-velocity (PV) diagram (a “Hubble law”); a lack of receding material in the approaching lobe, and vice versa; and a power law distribution of mass with velocity: $`dM/v_{\mathrm{obs}}v_{\mathrm{obs}}^\mathrm{\Gamma }`$, where $`v_{\mathrm{obs}}`$ is the line-of-sight velocity relative to the systemic velocity and $`\mathrm{\Gamma }`$ is typically about -1.8. Of these, the power law distribution of mass with velocity is the most difficult to explain, as it often holds quite well over a factor of five or ten in $`v_{\mathrm{obs}}`$ (e.g., Lada & Fich 1996). In order to construct a theory for the feedback effects of protostellar winds it is necessary to develop a model for these outflows, which represent the interaction between winds and the ambient gas. To be tenable, such a model must confront the common properties of outflows listed above. Moreover, it must also be compatible with the ways in which protostellar winds are launched and collimated, a topic we now address. ## 3. Structure and Intensity of Protostellar Winds Protostellar winds are launched centrifugally from the inner regions of protostellar accretion disks (Blandford & Payne 1982). This is made possible by the presence of intense poloidal magnetic fields, which act as rigid tubes anchored in the disk, along which gas is flung away. Once the wind has traveled sufficiently far, the field weakens and gas begins to travel ballistically rather than being forced to corotate with the disk. Since it is still stuck to the same field line as when it was launched, and since that field line stretches back to the rotating disk, the magnetic field is wound into a tight spiral. By the Biot-Savart law, the tightly-wound field implies that the wind encloses a current. Magnetic stresses associated with current within the wind can generate forces that push the wind toward or away from its axis. These forces are important for the collimation of the wind on scales ($`\stackrel{<}{}1`$ AU) comparable to the launching region of the accretion disk, but on the larger scales of molecular outflows ($`0.1`$ pc), the wind relaxes into transverse pressure balance. Therefore, the wind’s magnetic flux lines tend toward a force-free state in which there is no current within the wind: the enclosed current lies entirely along the axis, rather than threading the wind itself. The field thus settles into the force-free state $`B_\varphi =2I(r)/(c\varpi )`$, where $`\varpi =r\mathrm{sin}\theta `$ is the cylindrical radius and the current $`I(r)`$ is a very slowly varying function of $`r`$. (There must be currents, and hence forces, in the wind for $`I(r)`$ to vary at all: these effect a continuing gradual collimation of the wind.) This force-free distribution of toroidal magnetic field is sufficient to determine the ram pressure of the wind at large distances from its source. For, the magnetic field lines must corotate with the disk in a steady state. This, along with the fact that the wind will expand to fill the entire solid angle available to it, allows one to trace the field lines from their origin on the disk to their destination on the sky (Ostriker 1997; Matzner 1999). Then, the conservation of mass and energy along each streamline determines the ram pressure of the wind in each direction. In the case of a wind whose density just above the disk varies with disk radius $`\varpi _0`$ as $`\rho _{w0}\varpi _0^q`$, and whose field lines expand significantly ($`\varpi \varpi _0`$) according to the force-free distribution of the field, the wind ram pressure varies as $`\rho _wv_w^2C(r)\varpi _0^{(1q)/2}\varpi ^2`$ (Ostriker 1997). Now, the disk radius $`\varpi _0`$ is generally restricted within a very narrow range compared to the wind radius $`\varpi `$: for instance, in the model of Shu et al. (1995), $`\varpi _0`$ takes only a single value whereas $`\varpi `$ extends from a few AU to $`0.1`$ pc. For this reason, the ram pressure of the wind is very close to $`\rho _wv_w\varpi ^2`$ in practice (Matzner & McKee 1999a). This force distribution, which is characteristic of the x-wind model (Shu et al. 1995), is thus common among hydromagnetic winds. In reality, variability and instabilities are likely to broaden the wind force inside some angle $`\theta _0`$; therefore, we consider the smoothed distribution $$\rho _wv_w=\frac{\dot{m}_wv_w}{4\pi r^2}\frac{1}{\mathrm{ln}(2/\theta _0)}\frac{1}{(\mathrm{sin}\theta )^2+(\mathrm{sin}\theta _0)^2}.$$ (1) In the model for outflows presented below, observations imply $`\theta _00.01`$. The fact that hydromagnetic winds tend to $`\rho _wv_w(\mathrm{sin}\theta )^2`$ is consistent with the theoretical expectation (Matzner & McKee 1999a) that the fast Alfvénic Mach number, $`B^2/(4\pi \rho _wv_w^2)`$, is roughly constant. This expectation follows from the fact that this ratio is unity near the source, and varies logarithmically with distance along each streamline. The $`(\mathrm{sin}\theta )^2`$ force distribution implies that the cumulative wind momentum within some angle, $`p_w(<\theta )^\theta 2\pi r^2\rho _wv_w\mathrm{sin}\theta ^{}d\theta ^{}dt`$, varies as $`\mathrm{ln}(\theta /\theta _0)`$ for a broad range of angles, and this property will be useful for understanding outflows and their effects. A protostar’s wind lasts as long as it accretes material through its disk, which is essentially the free-fall time of the unstable core from which it formed, typically $`10^5`$ yr. If the wind’s mass is a fraction $`f_w`$ of the star’s mass, then the wind momentum is $`f_wv_wm_{}`$. We estimate that $`f_wv_w40`$ km s<sup>-1</sup>, roughly consistent with the observational estimate of Richer et al. (2000); for a wind velocity $`v_w200`$ km s<sup>-1</sup>, this implies $`f_w1/5`$, between the theoretical predictions of Shu et al. (1988) and Pelletier and Pudritz (1992). We expect $`f_wv_w`$ to be approximately independent of stellar mass, because $`v_w`$ scales roughly with the stellar escape velocity, which itself is regulated by deuterium burning during accretion (Stahler 1988). ## 4. Bipolar Molecular Outflows from Hydromagnetic Protostellar Winds The protostellar winds described above are essentially radial, both because they expand to fill the available angle, and because the wind coasts after achieving force balance. At wind speeds typical of low-mass protostars, the shocks that separate the wind from the surrounding gas are radiative (Koo & McKee 1992), so the swept-up shell conserves momentum in each direction. For a steady wind and an ambient medium whose density varies as $`\rho _aQ(\theta )r^2`$, Shu et al. (1991) showed that these ingredients reproduce the Hubble law for outflows. However, Matzner & McKee (1999a) demonstrated that neither a steady wind nor $`\rho _ar^2`$ is necessary for this conclusion: as long as the ambient medium is relatively featureless (e.g., a power law) on the scale of the outflow, the outflow motion is self-similar even if the wind intensity varies. For radial motion this implies the Hubble law, since $`𝐯(\theta ,t)𝐫(\theta ,t)/t`$. The collimated wind force distribution $`\rho _wv_w(\mathrm{sin}\theta )^2`$ leads to outflows that are highly elongated and have bowshock-shaped tips of width $`r\theta _0`$. This morphology is thus consistent with a “wide-angle” driving wind. A collimated outflow with $`𝐯𝐫`$ lacks blue material in its red lobe and vice versa, consistent with observations. That leaves only the power-law distribution of mass with velocity, $`dm/dv_{\mathrm{obs}}v_{\mathrm{obs}}^{1.8}`$, to be explained. Masson & Chernin (1992) argued that the radial outflow model of Shu et al. (1991) was inconsistent with this relation, but they only considered angular distributions $`\rho _wv_w^2(\mathrm{cos}\theta )^\beta `$ for some $`\beta `$. Their results do not apply to hydromagnetic winds, for which $`\rho _wv_w^2(\mathrm{sin}\theta )^2`$. Indeed, Li & Shu (1996) demonstrated that the outflow driven by a steady x-wind into a magnetically-flattened core, with $`\rho _aQ(\theta )r^2`$, could produce $`dm/dv_{\mathrm{obs}}v_{\mathrm{obs}}^2`$. I shall now show that this conclusion is much more general (see also Matzner & McKee 1999a), and is essentially independent of the ambient medium. The outflow shell expands at the rate dictated by momentum conservation; since $`\rho _wv_w^2\theta ^2`$ for small $`\theta `$, and since one expects the ambient medium also to be a power law of $`\theta `$, the rate of expansion obeys a power-law relation $`v(\theta )\theta ^x`$ quite generally. The power $`x`$ can be computed, but it is not important for the present argument. Now, hydromagnetic winds satisfy the cumulative momentum distribution $`p_w(<\theta )\mathrm{ln}(\theta /\theta _0)`$ for a wide range of angles, as described in §3. But since $`v\theta ^x`$, momentum conservation \[$`p_{\mathrm{shell}}(<\theta )=p_w(<\theta )`$\] implies a cumulative momentum distribution with outflow velocity $`p_{\mathrm{shell}}(>v)\mathrm{ln}(v/v_{\mathrm{max}})`$, where $`v_{\mathrm{max}}`$ is the rate of expansion along the wind axis. This, in turn, implies that $`dm/dv=d^2p_{\mathrm{shell}}(>v)/dv^2v^2`$; and, since $`v_{\mathrm{obs}}v`$ for an elongated outflow, $`dm/dv_{\mathrm{obs}}v_{\mathrm{obs}}^2`$, i.e., $`\mathrm{\Gamma }=2`$. Notice that essentially nothing about the ambient medium entered into this argument; therefore, $`\mathrm{\Gamma }2`$ is generic. Figure 1 demonstrates that models with this scaling can reproduce the mass-velocity distributions of real outflows. The commonly-observed features of molecular outflows are thus the natural products of momentum-conserving shells driven by hydromagnetic protostellar winds: they are insensitive to the details of the wind’s launching region, of variations in the wind’s intensity, and of the ambient density distribution. Although this model remains to be tested in detail, it matches observation well enough to deserve attention as a mechanism for feedback. ## 5. The efficiency of star and star cluster formation and the IMF The intensity of a protostellar wind along its axis ensures that within some angle, gas will be blown away. Because nearby gas is typically accreting onto the star itself or in the process of forming other stars, this process limits the efficiency with which individual stars and multiple stellar systems can form (see Matzner & McKee 2000, for a more detailed discussion). Gas dispersal by protostellar winds is more ubiquitous and less violent than disruption by massive stars. To estimate the amount of gas ($`M_{\mathrm{ej}}`$) lost per outflow, first note that outflows propagate by conserving momentum in each direction. To be ejected, material must travel faster than the escape velocity $`v_{\mathrm{esc}}`$ of the system. The momentum of the escaping gas is at least $`M_{\mathrm{ej}}v_{\mathrm{esc}}`$; therefore, $`M_{\mathrm{ej}}v_{\mathrm{esc}}<f_wv_wm_{}`$. In reality, momentum is lost in those directions where the wind is too weak to eject anything, and wasted in those where the ejected gas travels faster than $`v_{\mathrm{esc}}`$. Since the wind momentum is distributed logarithmically with angle, the fraction available to eject material at around $`v_{\mathrm{esc}}`$ is roughly $`\mathrm{ln}(1/\theta _0)`$; in fact, $`\mathrm{ln}(2/\theta _0)`$ turns out to be a better estimate. Moreover, an outflow’s momentum is reduced by a factor $`c_g`$ (typically of order unity) due to the action of gravity as it traverses the cloud. Defining the efficiency parameter $`X`$ and the efficiency per star $`\epsilon `$, $$Xc_g\mathrm{ln}(2/\theta _0)v_{\mathrm{esc}}/(f_wv_w),\epsilon m_{}/(m_{}+M_{\mathrm{ej}}),$$ (2) Matzner & McKee (2000) find that for the formation of an individual star from a collapsing molecular core, $$\epsilon ^1(2X)^1+[(2X)^2+(1+f_w)^2]^{1/2};$$ (3) and for the formation of a star within a larger cluster-forming clump, $$\epsilon ^11+(2X)^1.$$ (4) In both cases, $`\epsilon X`$ for $`X1`$ (strong winds) and $`\epsilon 1`$ for $`X1`$ (weak winds). The core and clump cases differ for a number of reasons. In the formation of a single star, unlike in multiple star formation, we can safely assume that all of the core mass is either accreted or ejected. Clumps have higher values of $`v_{\mathrm{esc}}`$ (which raises $`X`$), but their density profiles are shallower (which lowers $`X`$ through $`c_g`$); these effects roughly cancel. Finally, mass is ejected at angles much further from the wind axis in the core case; for this reason, the ejected wind mass is included in (3) but neglected in (4). For the formation of an individual star, (3) assumes that the collapsing core is spherically symmetric. However, protostellar cores are likely to be partially supported by magnetic fields, which will cause them to flatten along the magnetic axis. Since the protostellar wind is likely to share this axis, magnetic flattening segregates material away from the wind axis, reducing the amount that will be ejected and increasing the efficiency. Employing the models of flattened, magnetized cores presented by Li & Shu (1996), Matzner & McKee (2000) find that equation (3) remains valid if $`X`$ is replaced by $`(1+H_0)^{3/2}X`$, where $`1(1+H_0)\stackrel{<}{}2`$ is the factor by which magnetic support increases the core’s mean density. In the limit of a completely disk-like core ($`H_01`$), only the wind mass would be lost \[$`\epsilon 1/(1+f_w)`$\]. Individual stars form at about $`30\%`$ efficiency if cores are spherical ($`H_0=0`$), or $`75\%`$ if they are significantly flattened ($`H_0=1`$), as shown in the left panel of figure 2. For the formation of a star within a stellar cluster-forming clump, the anisotropy of the region is not likely to correlate with the wind axis. But, (4) assumes that each star forms at the center of a spherical distribution. Comparing again simulations in which stars form in a distributed manner throughout a cloud, Matzner & McKee (2000) find that (4) is an excellent approximation nevertheless, as shown in the right panel of figure 2. Equation (4) predicts efficiencies $`30\%\stackrel{<}{}\epsilon \stackrel{<}{}50\%`$ for the conditions typical of low-mass stellar cluster-forming regions, like those in Orion B (Lada 1992): each star expels one to two times its own mass from the clump. It is difficult to compare $`\epsilon `$ with its observational analogue, the ratio of stellar to total mass, but the two quantities are similar in a few well-studied cluster-forming regions (Matzner & McKee 2000). Using the efficiency of individual star formation, we can estimate the degree to which the stellar IMF differs from the mass function of pre-stellar cores (Nakano et al. 1995), a topic that has recently come under observational scrutiny (Motte, André, & Neri 1998). This comparison, discussed in detail by Matzner & McKee (2000), involves a consideration of how much $`X`$, and thus $`\epsilon `$, differs among cores of different masses. We find that the variation of $`\epsilon `$ is quite subtle, and that the slope of the IMF is thus only slightly flatter than the slope of the core mass function. It should be noted that equations (3) and (4) predict efficiencies significantly higher than those found by Nakano et al. (1995), and that this is a direct result of the collimation of protostellar winds. By this process (in contrast to the effect of a massive star), mass is lost continuously as stars form. The efficiencies found above are consistent with the formation of bound clusters (Mathieu 1983), which raises the question of why so few actually do remain bound (Lada & Lada 1991). This could be due to a violent event associated with massive star formation; alternatively, clusters could be unbound by motions of their parent clumps, a topic we address below. ## 6. Dynamical Evolution of Clumps During Cluster Formation Embedded stellar clusters are found exclusively within the most massive clumps inside molecular clouds. Since these clumps are massive enough to be confined by their own self-gravity despite the weight of the molecular cloud above them (McKee 1999), we can gain insight into the production of a stellar cluster by approximating a clump as a distinct reservoir of gas, and examining its dynamical evolution while a stellar cluster forms within it. This is the subject of an ongoing investigation (Matzner 1999; Matzner & McKee 1999b; Matzner et al. 2000), and the results reported here should be considered preliminary. Like any self-gravitating gaseous system, a cluster-forming clump must satisfy the virial theorem for self-gravitating gas (here in Eulerian form; McKee & Zweibel 1992): $$\frac{1}{2}\ddot{I}=2(𝒯𝒯_0)+𝒲+\frac{1}{2}\frac{d}{dt}_S(\rho 𝐯r^2)𝑑S,$$ (5) where $`IM_{\mathrm{cl}}R_{\mathrm{cl}}^2/2`$ is the trace of the clumps’s moment of inertia tensor; $`𝒲`$ and $``$ are its gravitational and magnetic energies; the kinetic energy of its internal motions is $`𝒯`$; $`𝒯_0`$ is an energy associated with its confining pressure; and the last term accounts for gas that may leave the system. Equation (5) has been used to identify equilibrium states of star-forming molecular clouds (McKee 1989; Bertoldi & McKee 1996), but it can also be taken as a dynamical equation if not taken to equal zero. In order to accomplish this, each of the energetic terms must be approximated as a function of four dynamical variables: the clump’s mass $`M_{\mathrm{cl}}`$ and the mass $`M_{}`$ of its embedded cluster, its radius $`R_{\mathrm{cl}}`$, and its internal velocity dispersion $`v_{\mathrm{rms}}`$. To simplify matters, the results presented here assume a spherical, magnetically supercritical ($`M_{\mathrm{cl}}=2M_\mathrm{\Phi }`$) clump with a $`\rho r^1`$ density profile and no rotation; any stars are assumed to move along with the gas. To complete our description of the system, it is necessary to specify the evolution of $`M_{\mathrm{cl}}`$, $`M_{}`$, and $`v_{\mathrm{rms}}`$. All three are affected by star formation and the effects of protostellar outflows on the gas. For this, we use the theory of McKee (1989): the local rate of gas conversion into stars ($`t_g^1`$) is identical to the local ambipolar-diffusion rate ($`t_{\mathrm{AD}}^1`$), because the formation of protostellar cores requires the assembly of magnetically supercritical condensations. Because of ionization by external FUV photons, star formation is inhibited within four visual magnitudes of the surface; interior to that, it proceeds at about a twelfth of the local free-fall rate. As stars form, their mass is added to $`M_{}`$ and subtracted from $`M_{\mathrm{cl}}`$; also subtracted is the mass of gas that is ejected according to equation (4). For a massive molecular clump, $`v_{\mathrm{rms}}`$ is dominated by non-thermal turbulent motions. The turbulent energy decays at a characteristic rate that can be parameterized by $`(dv_{}^{2}{}_{\mathrm{rms}}{}^{}/dt)_{\mathrm{decay}}=v_{\mathrm{rms}}^2/(\eta t_{\mathrm{ff}})`$, where $`t_{\mathrm{ff}}`$ is the free-fall time of the clump. Therefore, turbulence must be replenished every $`\eta `$ free-fall times; we expect $`1\stackrel{<}{}\eta <10`$, where numerical experiments (e.g., Stone, Ostriker, & Gammie 1998) favor $`\eta \stackrel{<}{}1`$. Turbulent energy can derive from two sources. For one, it can be exchanged for other energies of the system, through work done during compressions or expansions of the cloud; for this, we use the analytical results of Zweibel & McKee (1995) and McKee & Zweibel (1995) to estimate the energetics of turbulent, magnetized gas. This process does not, however, prevent the cloud from collapsing. Alternatively, turbulence may be stirred up by protostellar outflows. Since these conserve momentum, it is most convenient to consider the clump’s turbulent momentum ($`M_{\mathrm{cl}}v_{\mathrm{rms}}`$) enhanced an amount $`\varphi _wp_w`$ by a wind of momentum $`p_w`$; here, $`\varphi _w0.6`$ is the fraction of the wind momentum that does not escape the clump (Matzner 1999). Putting these elements together, it is useful to consider two approximations: one in which energy is added continuously according to the mean star formation rate within the clump, and another where stars form randomly at the current star formation rate, are chosen individually from the IMF, and affect the cloud instantaneously. These approaches differ for clumps of about five hundred solar masses or less, in which only ten to twenty stars form per $`t_{\mathrm{ff}}`$. In the continuous approximation it is possible to identify equilibria (Norman & Silk 1980; McKee 1989) in which turbulent decay is just offset by injection from outflows. These are stable, in the sense that a compressed cloud’s turbulence is enhanced faster than it decays. But they can also be overstable, as energy in an oscillation can grow during compressive phases. Perturbations due to individual forming stars will stimulate oscillations of the system, whose response depends critically on the rate of turbulent decay. Figure 3 shows the effect of changing $`\eta `$ for a clump of $`500M_{}`$: when turbulence decays slowly ($`\eta =10`$), a cluster-forming cloud will remain very close to its equilibrium state, whether or not star formation is taken to be continuous. For moderate turbulent decay ($`\eta =3`$), individual star formation seeds overstable oscillations about equilibrium: stars form in bursts. For fast turbulent decay ($`\eta =1`$), a continuous simulation can remain near its equilibrium, but one with individual star formation cannot consistently replenish the lost turbulent energy every dynamical time; the system collapses. Although the approximations used in these simulations (e.g., spherical symmetry and no rotation) affect these outcomes, it appears that the duration and constancy of star formation are sensitive to the rate of turbulent decay. ## 7. Conclusions Observations of star-forming regions and of individual accreting protostars have made it increasingly clear that stars interact violently with their parent clouds as they form. To model this phenomenon requires an understanding of protostellar winds, and of the bipolar molecular outflows they drive into the surrounding gas. We have seen that the typically observed features of molecular outflows can be explained as the natural product of magnetically-collimated winds and radiative shocks. This yields a model for outflows which can be used to estimate how much mass is ejected, either from an individual collapsing core or from a massive clump in the process of forming a stellar cluster. The estimates of star formation efficiency ($`30\%50\%`$) presented in §5 are much higher than those of Nakano et al. (1995), because they account for the collimation of protostellar winds. Lastly, the feedback effects of protostellar outflows can be incorporated into simple dynamical models for the process of stellar cluster formation within massive molecular clumps. The investigations described in §6 show that the dynamics of stellar cluster formation are sensitive to the rate at which turbulence decays, and this will have observational consequences. In the future, large-scale numerical simulations of star formation must include protostellar outflows, for without them such simulations are incomplete. ### Acknowledgments. I am grateful to Lee Hartmann, Zhi-Yun Li, Scott Kenyon, Chris McKee and Frank Shu for many useful and insightful comments. Motivation for the models of §6 came from a conversation with Frank Bertoldi. I very much appreciate the gracious hospitality of Roger Blandford and Sterl Phinney during my visits to Caltech, and of Lars Bildsten during a visit to the ITP. 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warning/0007/gr-qc0007008.html	ar5iv	text	# Gravitational collapse in 2+1 dimensional AdS spacetime ## 1 Introduction The past several years has seen growing interest in the properties and dynamics of asymptotically anti de-Sitter (AdS) spacetimes, predominantly due to the discovery of black hole solutions in 2+1 dimensional AdS spacetime and the AdS/CFT conjecture . The existence of vacuum <sup>3</sup><sup>3</sup>3with a negative cosmological constant black holes (also called BTZ black holes) is surprising because the local solution to the field equations is isometric to AdS, and hence has constant curvature. What makes a BTZ spacetime different from AdS is its global structure, which can be obtained by making appropriate identifications within the universal covering space of AdS . The natural question that such a construction poses is: how similar are these black holes to their more familiar 3+1 dimensional (4D) counter-parts? In particular, do these black holes have thermodynamic properties when considered within the framework of a quantum theory, and can they form through dynamical collapse processes? It turns out that BTZ black holes do bear striking resemblance to 4D black holes in many respects (see for review articles). In this paper we present the results of a numerical study of the collapse and formation of non-rotating BTZ black holes from a massless scalar field in 2+1D AdS spacetime. Of particular interest is whether critical phenomena are present at the threshold of black hole formation—namely if by fine-tuning of initial data, we can make the system asymptote (at “intermediate times”) to a solution which is universal in the sense of not depending on details of the initial data. Furthermore, if the black hole transition is “Type II”, so that there is no smallest mass of black hole which can be formed, then we expect their to be a scaling relation for the black hole mass of the form $`M=K(pp^{})^{2\gamma }`$. Here $`p`$ is a parameter in a family of initial data such that $`p=p^{}`$ is the critical solution, $`K`$ is a family dependent constant and $`\gamma `$ is a universal exponent (see for a recent review). The ‘extra’ factor of 2 in the exponent is expected for BTZ black holes—see section 3.2. As we will show, it turns out that the system does exhibit a continuously self-similar (CSS) solution in the critical limit, with a scaling exponent $`\gamma =1.2\pm 0.05`$. Earlier works on black hole formation in AdS considered disks of dust , null radiation , thin dust rings , and the collision of point-particles . In the case of dust-ring collapse Peleg and Steif found a scaling exponent of $`1/2`$ at the transition between black hole and naked singularity formation. Birmingham and Sen found the same exponent at the threshold of formation in the case of colliding particles. Husain and Olivier have also studied the massless scalar field in 2+1 dimensions using a double null formalism, and have formed black holes with their code . Our paper is organized as follows. In section 2 we describe the system of coordinates and numerical scheme we have chosen to use, and the resultant field equations and boundary conditions. An interesting consequence of our analysis is that we are unable to derive boundary conditions for the scalar field at the edge of the universe that are analogous to the out-going radiation conditions often employed in numerical relativity. In AdS spacetime the scalar field reaches time-like infinity $``$ in finite proper time as measured by a central observer, and the only consistent boundary conditions we can place on the scalar field confine it to the universe. This is reassuring from the standpoint of global energy conservation, but complicates the search for the universal scaling relation between black hole mass and parameter-space distance to the critical solution. The system behaves as if the scalar field is within a finite sized box, and so when a black hole forms all of the scalar field initially present eventually falls into the hole. $`M(p)`$ is therefore trivially a function of how the initial energy distribution scales with $`p`$. In section 3 we present results from the evolution of several families of initial data, focusing on critical behavior. To obtain $`\gamma `$, we follow the work of Garfinkle and Duncan, and examine the scaling of the maximum value attained by the curvature scalar $`R`$ in the sub-critical regime. We also study the effect that a central point particle (characterized by the angle deficit of the spacetime) has on the critical solution. As expected, we find that the more massive the point particle, the smaller the initial amplitude of the scalar field that gives rise to the critical solution. One might thus expect to have a one-parameter family of critical solutions with an overall scale related to the particle mass. It is surprising, therefore, that the scalar field always grows to the same amplitude in a near-critical evolution. A phase shift in central proper time is the only qualitative difference attributable to the mass of the particle. At the end of section 3 we study the interior structure of black holes that form, giving evidence that a ‘crushing’ spacelike curvature singularity forms within the event horizon. Thus the interior structure is significantly different from the BTZ solution, which has constant curvature (though the BTZ singularity is still crushing for extended objects falling into it). ## 2 The Einstein Klein-Gordon system in AdS spacetime We solve the Einstein field equations in 3 spacetime dimensions with cosmological constant $`\mathrm{\Lambda }1/\mathrm{}^2`$, coupled to a massless Klein-Gordon (KG) field $$R_{ab}\frac{1}{2}Rg_{ab}+\mathrm{\Lambda }g_{ab}=\kappa T_{ab},$$ (1) where the stress-energy-momentum tensor for the KG field $`\varphi `$ is $$T_{ab}=\varphi _{;a}\varphi _{;b}\frac{1}{2}g_{ab}\varphi _{;c}\varphi ^{;c}.$$ (2) Covariant differentiation is denoted by a semi-colon, while a comma denotes partial differentiation. We only consider circularly symmetric configurations of a minimally-coupled scalar field in this paper. Hence, $`\varphi `$ satisfies the wave equation $$\mathrm{}\varphi =\varphi _{;a}^a=0,$$ (3) and in coordinates ($`t,r,\theta `$) adapted to the symmetry, characterized by the Killing vector $`/\theta `$, $`\varphi (r,t)`$ is only a function of the radial coordinate, $`r`$, and time coordinate, $`t`$. One of the many peculiar features of AdS spacetime is its causal structure. In particular, null infinity $``$ is time-like, and any observer living in AdS spacetime can send and receive light-like signals to and from $``$ in finite proper time . These properties of AdS make it challenging to deal with numerically, as the scalar field traverses the entire universe on a local dynamical time-scale. Also, as we will show in section 2.1, the only regular boundary conditions on the field $`\varphi `$ at $``$ are Dirichlet conditions, so we cannot ignore the unusual causal structure of the spacetime by, for instance, placing out-going radiation boundary conditions on $`\varphi `$ at a finite proper distance from the origin. For these reasons, we adopt a coordinate system in which the metric takes the form: $$ds^2=\frac{e^{2A(r,t)}}{\mathrm{cos}^2(r/\mathrm{})}\left(dr^2dt^2\right)+\mathrm{}^2\mathrm{tan}^2(r/\mathrm{})e^{2B(r,t)}d\theta ^2.$$ (4) $`A(r,t)`$ and $`B(r,t)`$ are arbitrary functions of $`(r,t)`$, and it is straight-forward to show that when $`A=B=0`$ the above metric describes AdS spacetime; i.e. it is a solution to (1) with $`T_{ab}=0`$. Notice that, in this metric, radial null geodesics travel with constant coordinate speed $`dr/dt=\pm 1`$, and $``$ is at $`r=\pi \mathrm{}/2`$. The metric is singular at $``$, but we can place regular boundary conditions on $`A`$ and $`B`$ there, so that the spacetime is asymptotically AdS. Also, if we interpret $`\theta `$ as a periodic angular variable then the above metric has the correct topology to represent a BTZ black hole, as the topological censorship theorems require that the boundary at infinity share the topology of any event horizon that may exist in the interior of the spacetime . However, for the non-rotating collapse described in this paper, $`\theta `$ has no dynamical significance. Defining $$\mathrm{\Phi }(r,t)=\varphi _{,r},\mathrm{\Pi }(r,t)=\varphi _{,t}$$ (5) and using units where $`\kappa =4\pi `$, we get the following set of equations upon expanding (1)–(3) with the metric (4): $$A_{,rr}A_{,tt}+\frac{(1e^{2A})}{\mathrm{}^2\mathrm{cos}^2(r/\mathrm{})}+2\pi (\mathrm{\Phi }^2\mathrm{\Pi }^2)=0,$$ (6) $$B_{,rr}B_{,tt}+B_{,r}\left(B_{,r}+\frac{2}{\mathrm{}\mathrm{cos}(r/\mathrm{})\mathrm{sin}(r/\mathrm{})}\right)(B_{,t})^2+\frac{2(1e^{2A})}{\mathrm{}^2cos^2(r/\mathrm{})}=0,$$ (7) $`B_{,rr}+B_{,r}\left(B_{,r}A_{,r}+{\displaystyle \frac{1+\mathrm{cos}^2(r/\mathrm{})}{\mathrm{}\mathrm{cos}(r/\mathrm{})\mathrm{sin}(r/\mathrm{})}}\right)`$ $`{\displaystyle \frac{A_{,r}}{\mathrm{}\mathrm{cos}(r/\mathrm{})\mathrm{sin}(r/\mathrm{})}}A_{,t}B_{,t}+{\displaystyle \frac{(1e^{2A})}{\mathrm{}^2cos^2(r/\mathrm{})}}+2\pi (\mathrm{\Phi }^2+\mathrm{\Pi }^2)=0,`$ (8) $$B_{,rt}+B_{,t}\left(B_{,r}A_{,r}+\frac{cot(r/\mathrm{})}{\mathrm{}}\right)A_{,t}\left(B_{,r}+\frac{1}{\mathrm{}\mathrm{sin}(r/\mathrm{})\mathrm{cos}(r/\mathrm{})}\right)+4\pi \mathrm{\Phi }\mathrm{\Pi }=0$$ (9) and $$\left[\mathrm{tan}(r/\mathrm{})e^B\mathrm{\Phi }\right]_{,r}\mathrm{tan}(r/\mathrm{})\left[e^B\mathrm{\Pi }\right]_t=0.$$ (10) Within the context of the 3+1, or ADM, formalism, equations (2) and (9) are the Hamiltonian and momentum constraints respectively, while equations (6) and (7) are combinations of the evolution and constraint equations. Equation (10) is the wave equation for the scalar field. There are two unknown geometric variables—$`A(r,t)`$ and $`B(r,t)`$; hence one needs to use at least two of the four equations (6) - (9) to dynamically determine the geometry. In this work, we have chosen to use equations (6) and (7) to update $`A`$ and $`B`$. As is common practice in such a “free evolution scheme”, we can then use residuals of the constraints (2) and (9) as one way of estimating the level of error in our solution. With regards to initial conditions, we choose to freely specify $`\mathrm{\Phi }(r,0)`$ and $`\mathrm{\Pi }(r,0)`$ (we have to specify two scalar-field degrees of freedom at each $`r`$), as well as $`B(r,0)`$ and $`B_{,t}(r,0)`$. $`A(r,0)`$ and $`A_{,t}(r,0)`$ are then fixed from the constraint equations (see sec. 2.2 for more details). This procedure is clearly somewhat ad hoc, but has worked very well in our study. The Ricci scalar of this spacetime is $$R=\frac{4\pi \mathrm{cos}(r/\mathrm{})^2}{e^{2A}\mathrm{}^2}\left(\mathrm{\Phi }^2\mathrm{\Pi }^2\right)\frac{6}{\mathrm{}^2}.$$ (11) The Weyl tensor is zero, and other non-zero curvature scalars can be expressed as polynomial functions of $`R`$. ### 2.1 Regularity conditions We require that the solution for our dynamical variables $`A(r,t),B(r,t),\mathrm{\Phi }(r,t)`$ and $`\mathrm{\Pi }(r,t)`$ be regular at the origin, $`r=0`$, and at $``$, $`r=\pi \mathrm{}/2`$. The field equations then essentially dictate the allowed boundary conditions on these variables. By inspection of (6)–(10) we obtain the following conditions. At $`r=0`$ $`A_{,t}(0,t)`$ $`=`$ $`B_{,t}(0,t)`$ (12) $`A_{,r}(0,t)`$ $`=`$ $`0`$ (13) $`B_{,r}(0,t)`$ $`=`$ $`0`$ (14) $`\mathrm{\Phi }(0,t)`$ $`=`$ $`0`$ (15) $`\mathrm{\Pi }_{,r}(0,t)`$ $`=`$ $`0`$ (16) and at $`r=\pi \mathrm{}/2`$ $`A(\pi \mathrm{}/2,t)`$ $`=`$ $`A_{,r}(\pi \mathrm{}/2,t)=A_{,t}(\pi \mathrm{}/2,t)=0`$ (17) $`B_{,r}(\pi \mathrm{}/2,t)`$ $`=`$ $`0`$ (18) $`\mathrm{\Phi }(\pi \mathrm{}/2,t)`$ $`=`$ $`0`$ (19) $`\mathrm{\Pi }(\pi \mathrm{}/2,t)`$ $`=`$ $`0.`$ (20) Note that condition (16) on $`\mathrm{\Pi }(0,t)`$ is a direct consequence of the defining relation for $`\mathrm{\Pi }(r,t)`$ (5), and the regularity condition for $`\mathrm{\Phi }(0,t)`$ (15). Also note that we have multiple conditions for $`B`$ at the outer boundary, and for $`A`$ and $`B`$ at the origin. We have chosen to implement the Neumann conditions for $`A`$ and $`B`$ at the origin and the Dirichlet condition for $`A`$ at $``$, and then to monitor the other conditions as a consistency check during evolution. Conditions (17)–(20) ensure that the spacetime is asymptotically AdS. It is interesting that the field equations enforce Dirichlet boundary conditions on $`\mathrm{\Phi }`$ and $`\mathrm{\Pi }`$, effectively preventing us from implementing out-going radiation boundary conditions at $``$ (if we wanted to let the field “leak out of the universe” when it reaches $``$). To see this more clearly, consider the energy fluxes $`T_{ab}\eta ^a\eta ^b`$ and $`T_{ab}\mathrm{}^a\mathrm{}^b`$ along outgoing and ingoing null vectors, $`\mathrm{}^a`$ and $`\eta ^a`$, respectively, normalized so that $`\mathrm{}^a\eta _a=1`$ $$\mathrm{}^a=\frac{\mathrm{cos}(r/\mathrm{})}{\sqrt{2}e^A}\left[\frac{}{t}+\frac{}{r}\right]^a$$ (21) $$\eta ^a=\frac{\mathrm{cos}(r/\mathrm{})}{\sqrt{2}e^A}\left[\frac{}{t}\frac{}{r}\right]^a$$ (22) A straight forward calculation using (2) gives $$E_\pm =\frac{\mathrm{cos}(r/\mathrm{})^2(\mathrm{\Phi }\pm \mathrm{\Pi })^2}{2e^{2A}},$$ (23) where $`E_+`$ is the influx and $`E_{}`$ the outflux. Thus no- outflux/influx boundary conditions can be obtained in the usual way by differentiating $`\mathrm{\Phi }\pm \mathrm{\Pi }`$ with respect to $`r`$ and $`t`$ in turn, and utilizing the fact that, from (5), $`\mathrm{\Phi }_{,t}=\mathrm{\Pi }_{,r}`$: $`\mathrm{\Phi }_{,r}\pm \mathrm{\Phi }_{,t}=0`$ (24) $`\mathrm{\Pi }_{,r}\pm \mathrm{\Pi }_{,t}=0.`$ (25) Here, the plus sign corresponds to no-influx, and the minus sign to no-outflux. However, at the outer boundary, regularity forces $`\mathrm{\Phi }(\pi \mathrm{}/2,t)=\mathrm{\Pi }(\pi \mathrm{}/2,t)=0`$, and hence $`\mathrm{\Phi }_{,t}(\pi \mathrm{}/2,t)=\mathrm{\Pi }_{,t}(\pi \mathrm{}/2,t)=0`$, so there is no distinction between the no-influx and no-outflux condition. The only situation consistent with both conditions is that no flux crosses the outer boundary in either direction. Even when we try to derive no-outflux/influx conditions with the asymptotic behavior of $`\varphi `$ factored out, namely letting $`\varphi =\mathrm{cos}^2(r/\mathrm{})\widehat{\varphi }`$ and placing boundary conditions on $`\widehat{\varphi }`$, we find that the wave equation on $``$ cannot distinguish between no-outflux and no-influx conditions. Also, in early experiments we were unable to obtain stable numerical evolution with the no-influx boundary conditions (24) and (25) applied at a finite proper circumference, corresponding to $`r<\pi \mathrm{}/2`$. The Dirichlet boundary condition at $``$ is also consistent with the behavior of a massive scalar field in an AdS background, where an infinite effective-potential barrier prevents any of the field from reaching $``$, regardless of how small the mass is. Of course, all of this does not mean that an effective outgoing radiation condition can not be implemented for the massless field in asympotically AdS spacetimes. In any case, in the context of the current study, we would be apt to view such a condition as a numerical convenience, rather than being of any intrinsic physical interest. ### 2.2 Initial conditions For initial conditions at $`t=0`$, we are free to specify the scalar field gradients $`\mathrm{\Phi }(r,0)`$ and $`\mathrm{\Pi }(r,0)`$, the metric function $`B(r,0)`$ and its time derivative $`B_{,t}(r,0)`$. We then numerically solve for $`A(r,0)`$ and $`A_{,t}(r,0)`$ using the hamiltonian and momentum constraints (2) and (9). The freedom that we have to specify $`B(r,0)`$ amounts to a choice of what the proper circumference, ($`\mathrm{}\mathrm{tan}(r/\mathrm{})e^B`$), and its initial time derivative are, as a function of the radial coordinate $`r`$. That we do not have the freedom to choose $`B`$ for all time is a consequence of the gauge condition that radial light-like signals travel with unit coordinate velocity. For simplicity we set $`B(r,0)=B_{,t}(r,0)=0`$. We believe (though are unable to prove so), that the set of conditions just described is capable of generating all possible initial data, which is regular and free of trapped surfaces, for the minimally-coupled scalar field in asymptotically AdS spacetime (in 2+1 dimensions). The presence of trapped surfaces at $`t=0`$ is incompatible with the conditions on $`B(r,0)`$ and $`B_{,t}(r,0)`$— in our coordinate system $`dr/dt=1`$ along an outgoing null curve, and hence a non-zero $`B(r,0)`$ and/or $`B_{,t}(r,0)`$ is required to describe non-positive outward null-expansion. However, in this study we are only interested in initial data that is free of trapped surfaces, so the conditions on $`B(r,0)`$ and $`B_{,t}(r,0)`$ are not restrictive. For the initial scalar field profile, $`\varphi (r,0)`$, we consider three families of functions—a gaussian curve raised to the $`n^{th}`$ power $$\varphi (r,0)=Pe^{\left((rr_0)/0pt\right)^{2n}},$$ (27) a ‘kink’ (based on an $`\mathrm{arctan}`$ function) for which $`\mathrm{\Phi }=\varphi /r`$ is $$\mathrm{\Phi }(r,0)=\frac{2P\sqrt{0pt}\mathrm{cos}(r/\mathrm{})\mathrm{sin}(r/\mathrm{})[\mathrm{}\mathrm{sin}(r/\mathrm{})\mathrm{cos}(r/\mathrm{})+2(rr_0)(12\mathrm{sin}(r/\mathrm{})^2)]e^{(rr_0)^2/0pt^2}}{\pi \mathrm{}[0pt\mathrm{sin}(r/\mathrm{})^4\mathrm{cos}(r/\mathrm{})^4+(rr_0)^2]}$$ (28) and a family of harmonic functions <sup>4</sup><sup>4</sup>4we call these functions ‘harmonic’ because without back-reaction and for initially static configurations ($`\mathrm{\Pi }(r,0)=0`$) the exact solution to the wave equation is periodic in time. $$\varphi (r,0)=P\mathrm{cos}^2(rn/\mathrm{}),$$ (29) where $`P,r_0,0pt`$ and $`n`$ are constant parameters. Then, depending upon whether we want to model initially ingoing, outgoing or static fields, we set $`\mathrm{\Pi }(r,0)=\mathrm{\Phi }(r,0)`$, $`\mathrm{\Pi }(r,0)=\mathrm{\Phi }(r,0)`$ or $`\mathrm{\Pi }(r,0)=0`$ respectively. Note that this method cannot give purely ingoing or outgoing pulses—$`\mathrm{\Pi }(r,t)=\pm \mathrm{\Phi }(r,t)`$ is not an exact solution to the wave equation, and a little bit of energy always propagates in the opposite direction to that desired. As noted previously, we set $`B(r,0)=B_{,t}(r,0)=0`$. The remaining geometric variables, $`A(r,0)`$ and $`A_{,t}(r,0)`$ are then computed from the Hamiltonian and momentum constraints (2) and (9). We integrate the constraints outwards from $`r=0`$, setting $`A_{,t}(0,0)=0`$. For the most part we will consider the collapse of a scalar field initially exterior to empty AdS space. This corresponds to setting $`A(0,0)=0`$. However, in section 3.2.2, we will briefly consider the effect of collapsing the field in the presence of a point particle at the origin, the calculation of which involves introducing an angle deficit into the spacetime. From the metric (or by examining the parallel transport of a vector about $`r=0`$ in an infinitesimal loop), the angle deficit $`\omega `$ at $`t=0`$ is related to $`A(0,0)`$ as follows: $$\omega =2\pi (1e^{A(0,0)}).$$ (30) Of more interest is the relationship between $`A(0,0)`$ and the mass of the point-particle, $`M_{pp}`$. The remainder of this section is devoted to finding this relationship, and in the process we will define a general mass aspect function $`M(r,t)`$ for the spacetime. When the scalar field gradients identically vanish (which they do at $``$, and, to an excellent approximation, at $`r=0`$ for the initial data that we consider), the Hamiltonian constraint has the simple solution $$e^{2A}=\frac{k}{k\mathrm{cos}^2(r/\mathrm{})},$$ (31) where $`k`$ is a constant of integration. We can relate $`k`$ to the BTZ mass parameter $`M`$ of the spacetime by appealing to the usual form in which the BTZ solution is expressed: $$ds^2=(M+\overline{r}^2/\mathrm{}^2)d\overline{t}^2+\frac{1}{M+\overline{r}^2/\mathrm{}^2}d\overline{r}^2+\overline{r}^2d\theta ^2.$$ (32) $`M=1`$ is AdS spacetime, $`M0`$ are black hole solutions and $`M<0,M1`$ are spacetimes with conical singularities, or point particles at the origin (the range of $`\overline{r}`$ is from $`0`$ to $`\mathrm{}`$). For general (non-vacuum) solutions let us define the mass aspect $`M(\overline{r},\overline{t})`$ as follows $$\|\overline{r}\|^2M(\overline{r},\overline{t})+\overline{r}^2/\mathrm{}^2.$$ (33) Then, in our coordinate system (4), $`M(r,t)`$ takes the following form $$M(r,t)=e^{2(BA)}\left[e^{2A}\mathrm{tan}^2(r/\mathrm{})+\mathrm{}^2\mathrm{sin}^2(r/\mathrm{})((B_{,t})^2(B_{,r})^2)2\mathrm{}\mathrm{tan}(r/\mathrm{})B_{,r}\mathrm{sec}^2(r/\mathrm{})\right].$$ (34) Using the field equations (6)-(9) it is straightforward to show that $`M`$ is a conserved quantity in regions of the spacetime where $`\mathrm{\Phi }`$ and $`\mathrm{\Pi }`$ are zero (in particular at $``$). At $`t=0`$, where $`B=B_{,t}=0`$, we can substitute (31) into (34) to find $`k`$: $$k=\frac{1}{1+M}.$$ (35) When $`M0`$ our metric (4) with the chosen initial conditions is singular at the horizon of an empty BTZ spacetime, but the metric turns out to be well behaved at $`t=0`$ for initial data that does not contain trapped surfaces <sup>5</sup><sup>5</sup>5however, because of our choice of gauge, we know that the coordinate system must become singular within one light-crossing time (LCT) of the formation of a black hole. The event horizon is a null hypersurface travelling outward with unit coordinate velocity, thus the coordinate distance between the event horizon and $``$ will go to zero within a time $`t=\pi \mathrm{}/2`$.. Finally, from (31) and (35) the contribution, $`M_{\mathrm{PP}}`$, of the point particle to the mass of the spacetime is $$M_{\mathrm{PP}}=1e^{2A(0,0)}$$ (36) ### 2.3 Numerical Scheme We solve the set of equations (6),(7) and (10) by converting them to a system of finite difference equations on a uniform coordinate grid using a two-time level Crank-Nicholson scheme. We also add Kreiss-Oliger style dissipation to control high-frequency solution components; this is crucial for the stability of our method. At first, we used standard 2nd order accurate 3-point finite difference stencils for the spatial derivatives at each time level—centered-difference operators at interior points, a forward-difference operator at the inner boundary and a backward-difference operator at the outer boundary. However, we found that these operators excited a small instability in the metric variables in the vicinity of the outer boundary. The resultant ripples would propagate inwards and cause problems in situations where black hole formation was imminent. The primary source of these ripples was truncation error in the solution $`A`$ exciting small oscillations in $`B`$. Specifically, $`A`$ acts as a source term in the evolution equation for $`B`$ (7), and $`B`$ happens to be very sensitive to small errors in $`A`$ near the outer boundary (essentially since the leading term of $`A`$, when considered as a power series in $`\mathrm{cos}^2(r/\mathrm{})`$, cancels with the spatial derivatives of $`B`$ in (7) initially, and so higher order, less accurately known terms of $`A`$ are responsible for $`B`$’s “acceleration”). To reduce these problems we now use a 5-point, 4th order accurate spatial derivative operator at interior grid points, and 6-point 4th order backward and forward operators near boundaries that have the same truncation error as the interior operator. Also, we find that using the momentum constraint (9) to solve for $`A`$ at the next to last grid point is necessary to obtain convergence of the solution as we go to finer spatial resolution (for some as yet unknown reason the evolution equation was exciting a growing mode on finer grids at that point). The program to perform the evolution was written in Fortran 77 and RNPL (Rapid Numerical Prototyping Language ); animations and pictures from several evolutions can be obtained from our website . ### 2.4 Detecting black holes and excising singularities To detect black hole formation we search for trapped surfaces, defined to be surfaces where the expansion of outgoing null curves normal to the surface is negative. If cosmic censorship holds, then trapped surfaces are always found within the event horizon of a black hole, though at the end of the simulation we can trace null rays backwards from $``$ to confirm this. In our coordinate system the condition for a surface to be trapped is $$1+\mathrm{}\mathrm{cos}(r/\mathrm{})\mathrm{sin}(r/\mathrm{})(B_{,r}+B_{,t})<0.$$ (37) We estimate the mass of the black hole by monitoring the proper circumference $`2\pi \mathrm{}\mathrm{tan}(r_{\mathrm{AH}}/\mathrm{})e^{B(r_{\mathrm{AH}},t)}`$ of the apparent horizon (the outer-most trapped surface), and use the relationship between BTZ black hole mass and event horizon circumference ((32)—the horizon is at $`\overline{r}=\sqrt{M}\mathrm{}`$)): $$M\mathrm{tan}^2(r_{\mathrm{AH}}/\mathrm{})e^{2B(r_{\mathrm{AH}},t)}.$$ (38) If all of the scalar field is absorbed by the black hole during evolution, then the estimated mass should eventually become equal to the initial, asymptotic mass of the spacetime as given by (34) in the limit $`r\pi \mathrm{}/2`$. As we will show in section 3.3, shortly after an apparent horizon(AH) forms, we find what appears to be a spacelike curvature singularity forming within the AH. If we use a straightforward evolution scheme, the metric and scalar field variables quickly diverge, and any given simulation just as quickly breaks down. At the same time, we would like to probe the structure of the spacetime approaching the singularity, as well as to continue to following the evolution outside the AH as long as our coordinate system allows (approximately 1 light-crossing time). To accomplish this, we have implemented singularity excision, a technique fundamentally motivated by the black-hole-excision strategy first proposed by Unruh . Our excision strategy is as follows. We monitor the magnitude of the metric variables, and when they grow beyond a certain threshold <sup>6</sup><sup>6</sup>6for a threshold we choose a number that is sufficiently large so we are fairly certain (from past experiments) that if any variable grows beyond the threshold then a crash is imminent at any point we excise that point plus a small buffer zone (of 4 to 6 grid points) on either side of it. (Note that the non-excised region of the grid will no longer be contiguous if the excised point is further away from the original grid boundaries than the size of the buffer zone). At the new grid boundaries exterior to the excised region, we continue to solve for the metric and field variables using the evolution equations, but replace all centered-difference operators with forward and backward-difference operators, as appropriate, so that the solution is not “numerically influenced” by the excised grid points. Physically, the solution that one would obtain within the causal future of the excised zone is meaningless, so we also remove this region of the grid during subsequent evolution. In our coordinate system this is easy to implement, as radial null curves travel at constant, unit coordinate velocity. Thus, if our grid-spacing is $`\mathrm{\Delta }r`$, after an amount of time $`\mathrm{\Delta }t=\mathrm{\Delta }r`$ we expand the excised region by 1 grid point on either side. Also, we continue to monitor the metric variables on the remainder of the grid, and when they grow beyond the threshold at any other points we expand the excised region to include those points (and a buffer). Thus the excised piece of the grid is always contiguous. In principle, it would not be difficult to keep track of multiple excised zones, though we did not find it necessary to do so for the interior solution shown in sec. 3.3—a single zone is sufficient to obtain a good view of all of the interior up to the putative spacetime singularity. We have tested the singularity excision scheme by excising a light cone from a solution that remains regular, and verifying that the excised solution does converge to the regular solution as $`\mathrm{\Delta }r`$ decreases. In summary, we briefly clarify the difference between singularity and black hole excision. First, notice that we never use trapped surfaces to trigger the excision of a region of the grid. Thus, our code could, without modification, excise naked and coordinate singularities. The boundary of the excised region is always null or spacelike, so the scheme might not be able to distinguish between timelike and null curvature singularities. However, if a timelike singularity was encountered, it may still be possible to deduce its nature by examining the curvature invariants just exterior to the excised surface. For example, suppose during evolution a light-like region was excised, and curvature invariants started diverging as one approached the initial excised point, yet remained relatively “small” and finite just outside the future light-cone of the excised point, then one would have reasonable evidence for a timelike singularity. Second, with the singularity excision scheme, we excise only the region of the grid to the causal future of the singularity. In the case of a black hole spacetime, this results in a more complete view of the spacetime than what one would obtain with the standard black hole excision strategy (which would have in Fig. 19, for example, excised the region of the spacetime labeled “region of trapped surfaces”, and everything to the left of it). ## 3 Results In this section we discuss results from the evolution of several sets of initial data, focusing on the threshold of black hole formation. For convenience we set $`\mathrm{}=2/\pi `$ so that $``$ is at $`r=1`$, though the results presented here are valid for any non-zero, finite $`\mathrm{}`$, through an appropriate rescaling of the metric variables and scalar field gradients. Specifically, consider the following coordinate transformation $$\stackrel{~}{r}=\frac{r}{\mathrm{}},\stackrel{~}{t}=\frac{t}{\mathrm{}},$$ (39) with $`\stackrel{~}{r}`$ defined on the range $`[0,\pi /2]`$. Then it is easy to see that the $`\mathrm{}`$ dependence cancels from all equations (6)-(10) when expressed in terms of $`\stackrel{~}{r}`$ and $`\stackrel{~}{t}`$. So, given a solution $`A(\stackrel{~}{r},\stackrel{~}{t})`$, $`B(\stackrel{~}{r},\stackrel{~}{t})`$, $`\mathrm{\Phi }(\stackrel{~}{r},\stackrel{~}{t})`$ and $`\mathrm{\Pi }(\stackrel{~}{r},\stackrel{~}{t})`$ to the rescaled field equations we can find a corresponding solution for any $`\mathrm{}`$ by inverting the transformation (39) (see also (5)): $`A(\stackrel{~}{r},\stackrel{~}{t})`$ $``$ $`A(r/\mathrm{},t/\mathrm{})`$ (40) $`B(\stackrel{~}{r},\stackrel{~}{t})`$ $``$ $`B(r/\mathrm{},t/\mathrm{})`$ $`\mathrm{\Phi }(\stackrel{~}{r},\stackrel{~}{t})`$ $``$ $`\mathrm{}\mathrm{\Phi }(r/\mathrm{},t/\mathrm{})`$ $`\mathrm{\Pi }(\stackrel{~}{r},\stackrel{~}{t})`$ $``$ $`\mathrm{}\mathrm{\Pi }(r/\mathrm{},t/\mathrm{}),`$ with $`r`$ ranging from $`0`$ to $`\pi \mathrm{}/2`$. Notice that the initial energy density, being proportional to ($`\mathrm{\Phi }^2+\mathrm{\Pi }^2`$), scales like $`\mathrm{}^2`$, so there is no straight-forward method to extrapolate a solution to the limit of zero cosmological constant, where $`\mathrm{}\mathrm{}`$. We present results from 4 families of initial data: an ingoing gaussian ((27) with $`n=1`$), an ingoing squared gaussian ((27) with $`n=2`$), an ingoing kink (28), and a time-symmetric, $`n=1`$, harmonic function (29). In each case we vary the amplitude $`P`$ when tuning to the black hole threshold<sup>7</sup><sup>7</sup>7though we did check (for the gaussian) that we get the same critical solution when tuning the width, keeping the amplitude fixed, and for the first three families we have chosen $`0pt=0.05`$ and $`r_0=0.2`$. Except in section 3.2.2, where we briefly study collapse onto a point-particle, we have set $`A(0,0)=0`$ in all cases, corresponding to angle deficit-free spacetimes. The 3 ingoing families were simulated using a finest numerical grid of size 4096 points, with a Courant factor of 0.1; thus 40960 time steps are required per light-crossing time (for some of the critical solutions presented in the next section a grid size of 8192 points was used with a Courant factor of 0.2) For the time-symmetric $`\mathrm{cos}^2`$ function we do not need as many points to get good convergence results (because of the milder field gradients), so that the highest resolution required for that family was a 1024-point grid. In fact, we get acceptable results even after 50 light-crossing times with 1024 points for the $`\mathrm{cos}^2`$ data, whereas the more compact ingoing families start having noticeable errors (estimated from convergence tests) in near-critical evolution after 3-4 LCT’s with 4096 points. Fig. 1 shows the initial scalar field gradient, $`\mathrm{\Phi }(r,0)`$, of typical amplitude for each of the families. Fig. 2 shows the metric function $`A(r,0)`$ for a gaussian (the other families have similar shapes for $`A`$), and for later reference we show how $`A(r,t)`$ and $`B(r,t)`$ have evolved at $`t=0.6`$. In order to provide the reader with some feeling for the dynamics of a “typical” evolution, Fig. 3 shows a “space-time” plot of the evolution of a sample gaussian with $`P=0.1302`$ that does not form a black hole within 4 LCT’s (and it should not, as the asymptotic mass of the spacetime is $`1.062\text{x}10^2`$). ### 3.1 Parameter space survey, varying $`P`$ Figs. 4 and 5 show plots of the asymptotic mass, $`M(P)`$, of the spacetime, as a function of the amplitude $`P`$, about the region $`M=0`$ of parameter space, for the guassian and harmonic families. The second curve on each plot shows the initial mass estimate of a black hole (if one formed during the 2 LCTs of the gaussian evolution, or 50 LCTs of the harmonic evolution) at the time an apparent horizon is first detected. For these amplitudes, Figs. 6 and 7 show the time $`t`$ and coordinate position $`r`$ of apparent horizon formation. Qualitatively, the features of corresponding plots for the kink and squared guassian (also evolved for 2 LCT’s) are very similar to those for the gaussian, so for brevity we do not show them. To within the resolution of our simulation, the final black hole mass always approaches the asymptotic mass—in other words, we do not detect any remnant scalar field (black hole ‘hair’). See Fig. 8 for typical examples. Due to the “reflecting” boundary conditions at time-like $``$, this is not too surprising, although one might have expected something like a low amplitude, long wave-length, periodic scalar remnant. The scalar field also tends to zero at late times along the event horizon, though in that region of the spacetime our results are not good enough to obtain useful decay exponents. Fig. 7 for the harmonic family shows almost chaotic dependence of the time of AH formation as a function of amplitude, as $`M(P)`$ decreases towards $`M=0`$. There is evidence that this behavior is also present for the other families of initial data, but we have not run those simulations at the necessary resolution to give convincing evidence. What appears to be happening is the following. First of all, it is more “difficult” for a distribution of the scalar field corresponding to $`M0`$ to form a black hole—the distribution needs to be compact and centrally condensed. Thus, when we implode a relatively “space-filling” distribution with $`M`$ small (and positive) a black hole will not form on the first bounce through the origin. However, because of the boundary conditions at infinity, the scalar field will reflect off $``$, and, as the field has evolved through a strong field (non-linear) regime in the interior, the distribution of energy will be different on the subsequent implosion. Moreover, because of the strong gravitational field, the scalar field has a tendency to spend more time in the vicinity of the origin on average, preventing it from dispersing throughout the spacetime . So, one may expect that if the asymptotic mass $`M`$ is positive, a region of phase space will eventually be traversed during evolution, where it is favorable for a black hole to form, no matter how near-zero is $`M`$. However, due to the chaotic nature of the curve in Fig. 7, we cannot extrapolate $`t_0(M)`$ to $`t_0=\mathrm{}`$ in order to directly test this conjecture. ### 3.2 The critical regime To search for critical behavior in the gravitational collapse of the four families of initial data introduced in the previous section, we vary the amplitude $`P`$ in each case to find the threshold of black hole formation. Ideally, we would simply seek the amplitude $`P^{}`$ where a black hole forms for $`P>P^{}`$, while for $`P<P^{}`$ the scalar field bounces around forever without collapse. Unfortunately, such a search is not practical; as mentioned in the previous section we do not have the computational resources to follow compact initial data for numerous LCT’s, and, even with the $`\mathrm{cos}^2`$ data, we do not see any trends that would allow us to conclude that if a black hole has not formed after, say, $`n`$ LCT’s, then it probably will not form at all. Thus, what we do instead is tune to the threshold of black hole formation on the initial implosion; i.e. we base our search on whether or not a black hole forms before any initially out-going radiation reflects off $``$ and then falls in, contributing to the collapse. This point of parameter space is labeled as $`P^{}`$ in Figs. 4 and 5, and coincides with the place where the initial mass estimate dips to near zero (though for the harmonic data—as mentioned in the caption of Fig. 5—for amplitudes a little larger than $`P^{}`$ an apparent horizon first forms further out, engulfing the one that is about to form at the smaller radius; see also Fig. 7). Near this threshold, it turns out that shortly after the initial implosion, the scalar field and geometry close to the origin evolve towards a universal, continuously self-similar (CSS) form. We remind the reader that a function which is CSS depends only on a single scale-invariant variable $`x`$. Now, the coordinates $`(r,t)`$ in which we solve the equations of motion are not well-adapted to self-similarity. However, after some experimentation we found that a natural scale-invariant independent variable in our system is $$x=\frac{\overline{r}}{t_c},$$ (41) where $`\overline{r}=\mathrm{}\mathrm{tan}(r/\mathrm{})e^B`$ is proportional to the proper circumference of an $`r=\text{constant}`$ ring, and $`t_c`$ is proper time as measured by the central ($`r=0`$) observer. By convention, $`t_c`$ is negative and increases to the accumulation point $`t_c^{}0`$. To better visualize the CSS behavior, we also transform to logarithmic coordinates: $$Z\mathrm{ln}(\overline{r}),T\mathrm{ln}(t_c).$$ (42) A CSS function, $`f(x)=f(e^{Z+T})`$, then looks like a wave propagating to the left with unit velocity as $`T`$ increases to $`\mathrm{}`$. Figs. 912 show scale-invariant functions $`\varphi _{,Z}(Z,T)`$, $`\varphi _{,ZZ}(Z,T)`$ (the second derivative better demonstrates the “wave nature” of the critical solution), the mass aspect $`M(Z,T)`$, and $`(\overline{r}^2R)(Z,T)`$, for a gaussian evolution with $`P=0.133059219`$, which is ‘close’ to the critical solution ($`\mathrm{ln}(PP^{})=17.5`$; see sec. 3.2.1). In principle, the closer to criticality we tune the initial pulse, the longer the scale-invariant behavior should persist in logarithmic space. In practice, of course, finite computational precision and grid resolution prohibits fine-tuning to arbitrary accuracy—the figures plotted here show data which is about as close to criticality as we can get with 8192 grid points. In terms of the mass aspect in Fig. 11, one can surmise that the critical solution is (locally) a kink-like transition from the AdS value $`M=1`$ to a zero mass state; though, interestingly enough, the value of the curvature scalar $`R`$ at the origin diverges like $`1/t_c^2`$ as one approaches the accumulation point (we will discuss this in more detail below; also, bear in mind that in Fig. 12 we are plotting $`\overline{r}^2R`$, not $`R`$ itself). This behavior of the mass aspect suggests that the the transition at the critical point is Type II—in other words, there is no lower, positive bound on the mass of black holes that can be formed by the scalar field. Fig. 13 demonstrates the universality of the solution in the critical regime. Here we plot $`\varphi _{,ZZ}`$ (as in Fig. 10 for the gaussian) at the same time $`T`$ for each family in a near-critical evolution. The harmonic function appears to have a slightly larger amplitude, but, as we shall now argue, this is apparently just a slicing effect. As mentioned in sec. 2.2, because of the gauge that we use, and since we choose to solve for $`A(r,0)`$ and $`A_{,t}(r,0)`$ using the constraint equations (2) and (9), the only slicing freedom we have remaining is in the initial conditions for $`B(r,0)`$ and $`B_{,t}(r,0)`$. Once $`B(r,0)`$ and $`B_{,t}(r,0)`$ are specified, we have no control over the manner in which the slice evolves. For the three compact, ingoing families, the critical behavior develops at times ranging from $`t=0.25`$ to $`t=0.30`$, and because of the similar initial spatial distribution of the energy densities, the evolution has proceeded along very similar slices. On the other hand, the harmonic data approaches the critical solution at about $`t=1.25`$, at which point the slice has evolved quite differently from the other three families near their respective critical times (we note, however, that by plotting as a function of $`T`$ we do “match” the slices at the origin). To demonstrate that the slices evolve differently, we plot in Fig. 14 the normalized inner product between $`/t_c`$ (in an $`\overline{r},t_c`$ coordinate basis) and $`t_c`$ for the 4 families, at the same time used in Fig. 13. This inner product is the Lorentz gamma factor, $`W`$ (assuming the vectors are time-like), between $`\overline{r}=\text{constant}`$ observers, and those moving normal to the hypersurface $`t_c=\text{constant}`$: $$W=\frac{\|(/t_c)^\alpha _\alpha t_c\|}{\|/t_c\|\|t_c\|}.$$ (43) This quantity will be the same along identical slices of a spacetime (since such slices will have the same normal vectors); thus the harmonic solution slice is clearly different as one moves away from the origin. Another interesting feature of this plot for the harmonic data is that it shows gravitational collapse occurring a short distance away from the unfolding critical behavior, since, to the right of the peak, the vector $`/t_c`$ has become space-like (equivalently the surface $`\overline{r}=\text{constant}`$ has become space-like—see the discussion on the singularity structure in sec. 3.3, and in particular Fig. 20). At this point in parameter space for the harmonic function there is a lot more mass in the spacetime than that involved in the critical evolution, and this is causing an apparent horizon to form at a larger radius (see Figs. 5 and 7). Presumably, for smaller amplitudes one could tune to a threshold solution after several light-crossing times, and perhaps then one would more cleanly uncover the critical solution. To give more evidence that all the solutions are indeed approaching a universal one in the critical regime, we need to compare them on a common spacetime slice. In Fig. 15 we show the same function of the scalar field as in Fig. 13 transformed to a Christodoulou type coordinate system ($`\overline{r},v`$), where a $`v=\text{constant}`$ curve is an ingoing null geodesic . We normalized $`v`$ so that $`dv=dt_c`$ at the origin; i.e. $`v`$ also measures central proper time. Thus comparing solutions on the same $`v=\text{constant}`$ surface removes any slicing ambiguity <sup>8</sup><sup>8</sup>8we are grateful to David Garfinkle for suggesting this procedure to us. As can be seen from the figure, the transformed solutions are all quite similar, though we lose some accuracy in the transformation (which is why we have elected not to use these coordinates in all of the plots in Figs. 9 \- 13). #### 3.2.1 The scaling exponent $`\gamma `$ Another characteristic feature of Type II critical behavior in gravitational collapse is the universal scaling exponent $`\gamma `$ in the relation $`M=K(pp^{})^{2\gamma }`$. To measure this relationship in the current context, one needs to wait for the system to settle down to a steady-state to ensure that the apparent horizon is coincident with the event horizon, and hence that the mass estimate (38) gives the correct mass. In AdS, the boundary conditions at $``$ prevent us from performing this measurement—initially outgoing radiation that did not contribute to the near-critical black hole formation will eventually reflect off $``$ and pollute our measurement. However, as discussed by Garfinkle and Duncan , in the near-critical regime (above or below $`p^{}`$) any quantity with dimension $`L^q`$, where $`L`$ is a length scale, should exhibit a scaling relation with an exponent of $`q\gamma `$. Thus, following those authors, we find the maximum value attained by the Ricci scalar $`R`$ at $`r=0`$ in sub-critical evolution for $`t_c<0`$. Plots of $`\mathrm{max}_{t<t_c}\mathrm{ln}\|R(0,t)\|`$ vs. $`\mathrm{ln}(P^{}P)`$ for the four families studied is shown in Fig. 16. Since $`RL^2`$, these figures show that the scaling exponent $`\gamma `$ of the 2+1D AdS Klein-Gordon system is about $`1.2\pm 0.05`$. Notice that the mass aspect $`M`$ as defined in (33) is dimensionless (which is consistent with the scale-invariance of $`M`$ as plotted in Fig. 11). On the other hand, when we keep $`\mathrm{}`$ fixed and vary $`P`$, the resulting black hole mass (being proportional to $`r_{ah}^2`$) has a length scale of $`2`$, so one would expect the mass-parameter scaling relationship for BTZ black holes to go like $`M=K(PP^{})^{2\gamma }`$, where $`\gamma `$ is the same value $`1.151.25`$ found above for the scaling of $`R`$. The initial-mass estimate curves as shown in Figs. 4 and 5 do roughly exhibit this scaling behavior for $`P>P^{}`$. #### 3.2.2 Critical behavior in the presence of a point particle Here we briefly show how the presence of a point particle (angle deficit) alters the critical solution. The particle contributes to the mass of the spacetime (36), so the more massive the particle (up to the maximum $`M_{\mathrm{PP}}=1`$ in our units) the less scalar field energy is needed to form a black hole, and consequently we have smaller amplitudes, $`P^{}`$, at threshold. Interestingly, we find the same critical solution in all cases (see Fig. 17 for 3 examples), the only noticeable differences being a systematic phase shift in $`T`$ related to the mass of the particle. The kink-like transition in the mass aspect has the same shape as well, but it ranges from the particle mass at $`r=0`$ to $`M=0`$. To within the resolution of our simulations (which was at 2048 grid-points in this case) the critical exponent is also the same, namely within the range $`\gamma =1.15`$ to $`1.25`$. #### 3.2.3 The critical solution from a CSS ansatz? Given that we have self-similar behavior in the critical regime, it would be useful to find the exact solution assuming a CSS ansatz. Traditionally this is done by assuming the existence of a homothetic Killing vector. $`\xi `$ (see ) $$_\xi g_{ab}=2g_{ab}.$$ (44) This implies that in coordinates adapted to the homotheticity, so that $`\xi =/\tau `$, each component of $`g_{ab}`$ has the form $`e^{2\tau }f`$, for some function $`f`$ independent of $`\tau `$. Furthermore, $`_\xi R_{ab}=0`$, so that $`R_{ab}`$ and hence the Einstein tensor $`G_{ab}`$ are independent of $`\tau `$. This ansatz is not consistent with the field equations (1) in the presence of the cosmological constant if we assume that the scalar field is self-similar (see ), as we observe in the collapse simulations. Essentially, the scalar field stress-energy tensor (2) would need to decouple into a piece that exactly cancels the cosmological constant term plus a scale-invariant term, but we do not think that this is possible for a minimally-coupled scalar field. It may be that in the 2+1D AdS system a different symmetry, such as a conformal Killing vector, would be needed to generate the critical solution. Or perhaps the critical solution is only approximately homothetic over a limited region of the spacetime. Nevertheless, we have not yet found a symmetry-reduced system that reproduces the observed critical behavior.<sup>9</sup><sup>9</sup>9Note added in preparation: David Garfinkle has very recently found a CSS solution in the limit where the cosmological constant vanishes that appears to quite accurately describe the critical solution that we have found . His result is quite intriguing—the cosmological constant is essential for black holes to form, yet apparently it plays very little role in the solution at the threshold of formation! ### 3.3 Singularity structure In all of the solutions that we have studied so far we find that after an apparent horizon forms what appears to be a spacelike curvature singularity develops within the horizon. Specifically, the surface of excision along which the metric variables $`A`$ and $`B`$ and, consequently, the curvature invariants begin to diverge, is spacelike. By itself, demonstrating a spacelike surface of arbitrarily large curvature is not sufficient to prove that the singularity is spacelike—a counter-example would be the mass-inflation null singularity <sup>10</sup><sup>10</sup>10We are grateful to Lior Burko for pointing this out to us. However, if we extrapolate to the surface of infinite curvature, based upon the growth of the Ricci scalar prior to excision, we still find a spacelike surface (in fact, $`R`$ grows so rapidly prior to excision —roughly like $`1/t^4`$ along an $`r=constant`$ surface if we translate $`t`$ to zero at the singularity—that the surface of infinite curvature essentially coincides with the excision surface at the resolution of Fig 19 below). In addition, $`B(t,r)\mathrm{}`$ along this surface, indicating that the proper circumference measure $`\mathrm{}\mathrm{tan}(r/\mathrm{})e^B`$ goes to zero there (see Fig. 20 below). Thus, as with vacuum BTZ black holes, this singularity is crushing<sup>11</sup><sup>11</sup>11or deformationally strong, see . It is straight-forward (though tedious) to see that $`r=0`$ in the non-rotating $`BTZ`$ black hole is a strong singularity as defined by Tipler (though it is not a curvature singularity!). We have not repeated the formal calculations in terms of Jacobi fields in our collapse simulations, but because of the central, space-like nature of the singularity back-reaction is not likely to weakening it. Note added in revision: shortly after this paper was first published, Lior Burko studied the structure of the singularity in 2+1D AdS spacetime using a ‘qausi-homogenious’ approximation, and did find the singularity to be strong and spacelike .: any extended object reaching the singularity is forced to zero proper circumference, regardless of any angular momentum or internal pressures that the object might have. Figs. 18 and 19 are spacetime plots (essentially Penrose diagrams) of $`\mathrm{\Phi }(r,t)`$ and the Ricci scalar $`R(r,t)`$, respectively, for a gaussian initial pulse with $`P=0.133051`$. On the pictures we have superimposed the region of trapped surfaces and the inferred event horizon of the space time, found by tracing a null ray backwards in time from the place where the AH meets $``$ on the coordinate grid. Fig. 20 show contours of proper circumference for the same solution. The point $`P=0.133051`$ in parameter space is slightly sub-critical (as we have defined criticality, see Sec. 3.2)—a black hole forms because the bit of outgoing energy present at $`t=0`$ bounces off $``$ and falls back onto the nearly collapsed scalar field, pushing it over the limit. This gives us a very clear view of the interior structure; for a more massive pulse the singularity forms shortly after the initial implosion, resulting in a thin sliver of an interior in $`(r,t)`$ coordinates. From Fig. 19 one can see a striking peak that forms in $`R`$ after the scalar field has bounced through the origin and is travelling outwards. In this particular case $`R`$ has a value of order $`10^{10}`$ in the interior, it then grows to order $`+10^8`$ over a very short distance before decreasing to the AdS value of $`6/\mathrm{}^215`$. This near-discontinuous behavior in $`R`$ is characteristic of sub-critical evolutions, and becomes more extreme as one nears the critical solution. As one approaches the excised space-like surface in Fig. 19, $`R`$ starts to grow very rapidly, reaching values up to $`\|10^{28}\|`$ before excision (this may not be clear on the figure—we chose the gray scale to highlight the near-discontinuous behavior in $`R`$). $`R`$ actually oscillates between large positive and negative values along this surface, but our calculations are not sufficiently accurate to conclude that the oscillation is genuine. In particular, $`R`$ is extremely sensitive to the difference $`\mathrm{\Pi }^2\mathrm{\Phi }^2`$ (see (11)), and $`\mathrm{\Pi }^2`$ is usually around the same order of magnitude as $`\mathrm{\Phi }^2`$ there. We also note that the maximum value attained by $`R`$ along the excised surface becomes smaller towards $``$. This is to be expected, since in the 2+1D system, some scalar field is necessary to produce a value of $`R`$ differing from the AdS value (again, see (11)), and as we move towards $``$ along the excised surface there is progressively less scalar field energy remaining. ## 4 Concluding remarks We have studied black hole formation from the collapse of a minimally-coupled massless scalar field in 2+1 dimensional AdS spacetime. Outside of the event horizon the spacetime settles down to a BTZ form; in the interior a central, spacelike curvature singularity develops. At the threshold of black hole formation we find that the scalar field and spacetime geometry evolve towards a universal, continuously self-similar form. When a point particle is present at the origin the critical solution is shifted in central proper time by an amount related to the mass of the particle. By examining the behavior of the curvature scalar during sub-critical evolution we deduced that the universal scaling exponent $`\gamma `$ for this system is roughly $`1.2\pm 0.05`$. This value is quite different from the scaling exponent $`1/2`$ derived by Peleg and Steif for the collapse of thin rings of dust and by Birmingham and Sen for particle collisions. However, those works considered different forms of matter, and the phase transition was between black hole and naked singularity formation. Thus one would not expect the same exponent. Also, the local spacetime geometry about a dust ring or point-particles is necessarily (empty) AdS, hence such systems cannot exhibit any of the features, other than mass scaling, that are characteristic of critical gravitational collapse. Some questions remain unanswered in this work. First, what is the exact nature of the critical solution? In other words, what is the character of the symmetry (if any) responsible for the self-similar behaviour, as the system does not appear to admit a global homothetic Killing vector <sup>12</sup><sup>12</sup>12though, as mentioned in the footnote of sec. 3.2.3, David Garfinkle has found a CSS solution that is apparently relevant to the AdS critical solution . Second, will any distribution of energy that could conceivably form a black hole (i.e. with asymptotic mass $`M>0`$) eventually do so if one waits long enough (because of the Dirichlet boundary conditions imposed on the scalar field at $``$)? A third question, related to the first two, is whether the critical solution we have found is a true black-hole-formation threshold solution. In other words, that we have a found a universal, CSS solution via a fine-tuning process indicates that this critical solution is one-mode unstable; so, does perturbing the critical solution “one way” result in a black hole, and perturbing it the “other way” cause the scalar field to remain regular, never forming a black hole? The asymptotic nature of AdS spacetime, which is ultimately responsible for the boundary conditions of the scalar field at $``$, prevent us from answering this question in our collapse simulations. With regards to future work, it would be useful to extend these results to different scalar-field/geometry couplings, include a mass and potential terms in the Lagrangian, and to add angular momentum to the initial data to study the formation of rotating black holes. It would also be interesting to understand the critical behavior in light of the AdS/CFT correspondence. Even though our calculation is purely classical, there should be a regime where the classical evolution is a good approximation to the full bulk theory, and consequently there should be a dual CFT description of the critical phenomena. Acknowledgements We would like to thank David Garfinkle, Viqar Husain, Lior Burko, Iñaki Olabarrieta, Michel Olivier, Bill Unruh, Jason Ventrella, and Don Witt for many stimulating discussions. We are grateful to David Garfinkle for suggesting to us the method we used to obtain $`\gamma `$, as well as the use of the ingoing null coordinate system to compare near-critical solutions. MWC would particularly like to thank Robert Mann for many early discussions about this problem during the 1999 Classical and Quantum Physics of Strong Gravitational Fields program held at the Institute for Theoretical Physics, UC Santa Barbara. This work was supported by NSERC and by NSF PHY97-22068 and PHY94-07194. Most calculations were carried out on the vn.physics.ubc.ca Beofwulf cluster which was funded by the Canadian Foundation for Innovation.
warning/0007/hep-th0007184.html	ar5iv	text	# 1 Introduction ## 1 Introduction In recent years, there has been much investigation of domain walls which appear in many areas of physics. These domain walls interpolate between degenerate discrete minima of a scalar potential, with dependence on one spatial coordinate. They can occur naturally when a discrete symmetry is spontaneously broken. Domain walls can also appear in supersymmetric field theories when the superpotential has more than two critical points corresponding to degenerate minima of the scalar potential. In particular, it has been found that domain walls in supersymmetric theories can saturate the Bogomol’nyi bound. Such domain walls are called BPS domain walls and preserve half of the original supersymmetry. The existence of BPS domain walls corresponds to the central extension of $`𝒩=1`$ superalgebra, and the topological charge of the walls becomes the central charge $`Z`$ of the superalgebra. - The BPS bound and supercharges are determined by this central charge $`Z`$. BPS domain walls in supersymmetric theories have been extensively studied in models with degenerate isolated vacua. - Moreover, it has been found that such BPS domain walls can form a junction when three or more different isolated vacua occur in separate regions of space. The BPS state of the junction preserves $`1/4`$ supersymmetry, and the BPS bound is determined by two kinds of central charges, $`Z`$ and $`Y`$, appearing in the $`𝒩=1`$ superalgebra. There has been progress in the study of the general properties of such BPS junctions, - for example, the negative contribution of the charge $`Y`$ to the junction mass and the non-normalizability of zero modes on the BPS junction. It has also been argued that BPS junctions can create a network and that they can play a role in our world in higher-dimensional spacetime with a negative cosmological constant. In this way, it has been found that BPS domain walls have many interesting properties using models with several isolated vacua. It is essential in these models that isolated vacua have different values of the superpotential, since their differences are related to the energy densities saturating the BPS bound. In many supersymmetric theories, however, the vacuum manifold consists of a continuously degenerated moduli space. Since supersymmetric vacua are the extrema of the superpotential ($`W^{}=0`$), the connected parts of the moduli space have the same values of the superpotential, and thus each connected part is mapped to a single point in the superpotential space. Hence we can expect the existence of domain walls in the models, with the moduli space composed of several disjoint parts, rather than isolated points, because these disjoint vacua (in the field space) are, in general, mapped to different points in the superpotential space. In short, the moduli spaces of disjoint supersymmetric vacua appear the same as the isolated vacua in the superpotential space. In this work, we investigate BPS domain walls in $`𝒩=1`$ four-dimensional supersymmetric field theories with continuous global symmetry. If the models have vacua with spontaneously broken global symmetry, there exists a flat direction along the broken symmetry, or the moduli space of the vacua. Domain walls in such theories can be expected to connect pairs of the vacua in the disjoint moduli spaces if the Homotopy group $`\pi _0`$ is nontrivial.<sup>1</sup><sup>1</sup>1 Investigation of the Homotopy group gives the necessary conditions for the existence of domain walls, but there is not always a solution of the equation of motion, in particular the BPS equation for the BPS domain walls. If a BPS domain wall connects such disjoint moduli spaces for broken global symmetry, the configuration itself breaks the symmetry. Hence there can be a family of BPS walls interpolating between two disjoint moduli spaces; the BPS bound for walls is given by the difference between the superpotential values corresponding to two vacua, and this never changes under the symmetry transformation. In fact, we show that applying a symmetry transformation to a BPS domain wall solution produces another solution of the BPS equation. Therefore we can expect additional moduli of BPS walls, in addition to the location of the wall’s center. There is another reason why we study the BPS walls in models with continuous global symmetry. It is known that when a supersymmetric model possesses global symmetry, the superpotential has a larger symmetry, or the complexification of the original global symmetry, owing to the holomorphy of the superpotential. The vacuum manifold has non-compact flat directions, corresponding to the imaginary parts of the vacuum expectation values of the fields. The Nambu-Goldstone theorem for supersymmetric models has been proven. From this, it is known that when a global symmetry is spontaneously broken in supersymmetric vacua, there appear NG supermultiplets as many as the number of the broken generators of the complexified group. Since the complexified group is the symmetry of the superpotential, not that of the whole model nor of the BPS equation, it is a highly nontrivial problem to determine whether there can exist BPS walls interpolating between two vacua along disjoint non-compact flat directions. We examine this problem by using two supersymmetric models with global $`O(2)`$ symmetry, consisting of two chiral superfields. Unlike the global $`O(2)`$ symmetry, $`O(2)^𝐂`$ transformations of a BPS wall solution are not solutions of the BPS equation. However, we show that there can exist moduli of BPS walls corresponding to the shift of vacua along the non-compact flat direction. This moduli is different from the imaginary part of the parameter of the $`O(2)^𝐂`$ transformation. In sect. 2, we discuss the general properties of BPS domain walls in the model with continuous global symmetry. In sect. 3, we introduce our two models with $`O(2)`$ symmetry. We examine the existence of complex BPS walls interpolating between non-compact flat directions in both the models. In sect. 4, we reach conclusions for both models and discuss the features of BPS domain walls in general models with global symmetry. We also discuss a possible extension of the supersymmetric Nambu-Goldstone theorem. ## 2 BPS walls and continuous symmetry We consider supersymmetric field theories with only chiral superfields, and the Kähler potential is assumed to be linear: $`K=\varphi ^{}\varphi `$. The supersymmetric vacua are given as the extrema of the superpotential $`W(\varphi _k)`$, given by $`{\displaystyle \frac{W}{\varphi _k}}=0,k=1,\mathrm{}K,`$ (2.1) where the $`\varphi _k`$ are the scalar components of the chiral superfields. It is known that, denoting two solutions of Eq. (2.1) by $`\{\varphi _k\}_I`$ and $`\{\varphi _k\}_J`$, and the corresponding values of the superpotential by $`W_I`$ and $`W_J`$, there exists the lower bound of the surface energy density, or tension, for walls connecting these two vacua expressed by $`{\displaystyle \frac{\text{Energy}}{\text{Area}}}2\|W_JW_I\|.`$ (2.2) The BPS wall for which the equality in Eq. (2.2) holds satisfies the equation $`_z\varphi _k=e^{i\alpha }{\displaystyle \frac{W^{}}{\varphi _k^{}}},`$ (2.3) where $`\alpha =\mathrm{arg}(W_JW_I)`$. Here we have considered the wall depending on the coordinate $`z`$. Equation (2.3) is called the “BPS equation“. If the superpotential $`W`$ is invariant under the global symmetry $`G`$, $`W(\varphi )W(g\varphi )=W(\varphi ),\varphi \stackrel{g}{}g\varphi ,gG,`$ (2.4) where $`\varphi `$ belongs to unitary representation of $`G`$, Eq. (2.1) is also invariant under $`G`$: $`{\displaystyle \frac{W(\varphi )}{\varphi _i}}\stackrel{g}{}g_{ij}^{1T}{\displaystyle \frac{W(\varphi )}{\varphi _j}}.`$ (2.5) Since the superpotential includes only chiral superfields, the invariant group $`G`$ of the superpotential is enlarged to its complexification, $`G^𝐂`$. It is known that, in addition to the ordinary Nambu-Goldstone bosons corresponding to broken $`G`$ symmetry, there appear so-called quasi-Nambu-Goldstone bosons corresponding to broken $`G^𝐂`$ symmetry. With the fermions of their superpartner, they constitute massless chiral superfields. The vacuum manifold as a $`G^𝐂`$-orbit is parameterized by these massless bosons, and the quasi-Nambu-Goldstone bosons just parameterize the non-compact flat directions.<sup>2</sup><sup>2</sup>2 It is known that, in the case of the F-term breaking, there must exist at least one quasi-Nambu-Goldstone boson. Then the vacuum manifold inevitably becomes non-compact. Therefore, in the moduli space of its supersymmetric vacua, there exists a non-compact flat direction along the direction of the imaginary part of the scalar fields. We can see that the BPS equation (2.3) is covariant under transformation of the global symmetry $`G`$, but it is not covariant under the transformation of $`G^𝐂`$, since the BPS equation includes both holomorphic and anti-holomorphic fields. Then, if we can find a solution of Eq. (2.3), configurations obtained through transformation of this solution by elements of $`G`$ are also solutions of the BPS equation. However, configurations obtained through transformations of a solution by elements of $`G^𝐂`$ are not generally solutions of the BPS equation. Therefore, if the model has more than two disjoint flat directions, it is a nontrivial problem to determine whether there exist BPS walls interpolating between them. We examine this problem in two supersymmetric models. ## 3 BPS walls in models with flat directions ### 3.1 Moduli spaces of our models with flat directions In this paper, we consider the following two supersymmetric models with flat directions.<sup>3</sup><sup>3</sup>3 The two models that we consider in this paper are not renormalizable. Therefore these models must be interpreted as effective theories. First we consider a model with one flat direction. Its superpotential is $`W(\varphi )={\displaystyle \frac{1}{4}}(\stackrel{}{\varphi }^{\mathrm{\hspace{0.17em}2}}a^2)^2,\stackrel{}{\varphi }=\left(\begin{array}{c}\varphi ^1\\ \varphi ^2\end{array}\right),`$ (3.1) where $`\varphi ^1`$ and $`\varphi ^2`$ are chiral superfields composing the doublet of $`O(2)`$, $`\stackrel{}{\varphi }`$, and $`a`$ is a constant parameter. By a field redefinition, we can take this parameter $`a`$ to be real and positive without loss of generality. This model has two disjoint vacua: Vac. I $`\stackrel{}{\varphi }=0,W={\displaystyle \frac{a^4}{4}},`$ Vac. II $`\stackrel{}{\varphi }^{\mathrm{\hspace{0.17em}2}}=a^2,W=0.`$ (3.2) Let us note that the $`\varphi ^i`$ are the scalar components of chiral superfields here. (We denote the chiral superfields and their scalar components by the same letter.) Vac. I is $`O(2)`$ symmetric, but Vac. II spontaneously breaks $`O(2)`$ symmetry. The expectation value for Vac. II can be labeled as $`\varphi ^1=a\mathrm{cos}\theta ,\varphi ^2=a\mathrm{sin}\theta .`$ (3.3) Now the fields $`\varphi ^1`$ and $`\varphi ^2`$ can take complex values, and we can regard $`\theta `$ as a complex parameter. Therefore the vacuum manifold of this model is enlarged to an $`O(2)^𝐂`$-orbit: If we set $`\stackrel{}{\varphi }=\stackrel{}{x}+i\stackrel{}{y}`$, the two disjoint vacua in the Eq. (3.2) become Vac. I $`\stackrel{}{x}=\stackrel{}{y}=\stackrel{}{0},`$ Vac. II $`\stackrel{}{x}^2\stackrel{}{y}^2=a^2,\text{and}\stackrel{}{x}\stackrel{}{y}=0.`$ (3.4) Hence Vac. II can be rewritten as a two-dimensional surface in the three-dimensional linear space $`(x^1,x^2,\|y\|)`$, where $`\|y\|=\sqrt{\stackrel{}{y}2}`$ (see Fig. 1). Vac. II breaks this $`O(2)^𝐂`$ symmetry spontaneously. We consider the BPS wall connecting $`O(2)^𝐂`$ symmetric and $`O(2)^𝐂`$ broken vacua, and show that no BPS wall can connect the complex vacuum \- the vacuum with a complex value of the fields shifting along the flat direction in this model (see Fig. 1). Next we consider the model with two flat directions. Its superpotential is $`W(\varphi )={\displaystyle \frac{1}{6}}\stackrel{}{\varphi }^{\mathrm{\hspace{0.17em}2}}(\stackrel{}{\varphi }^{\mathrm{\hspace{0.17em}2}}a^2)^2,`$ (3.6) where $`\stackrel{}{\varphi }`$ is an $`O(2)`$ doublet composed of the chiral superfields $`\varphi ^1`$ and $`\varphi ^2`$, and the parameter $`a`$ is assumed to be a positive real constant for simplicity. This model has three disjoint vacua: Vac. I $`\stackrel{}{\varphi }=0,W=0,`$ Vac. II $`\stackrel{}{\varphi }^{\mathrm{\hspace{0.17em}2}}={\displaystyle \frac{a^2}{3}},W={\displaystyle \frac{2}{81}}a^6,`$ Vac. III $`\stackrel{}{\varphi }^{\mathrm{\hspace{0.17em}2}}=a^2,W=0.`$ (3.7) Setting $`\stackrel{}{\varphi }=\stackrel{}{x}+i\stackrel{}{y}`$, as in the previous model, Vac. II and Vac. III can be rewritten as two hyperboloids with different sizes and Vac. I as the origin in the space $`(x^1,x^2,\|y\|)`$ (see Fig. 1). We see that Vac. I is $`O(2)^𝐂`$ symmetric, but Vac. II and Vac. III break $`O(2)^𝐂`$ symmetry spontaneously. We consider the two kinds of BPS walls, connecting Vac. I and Vac. II, and connecting Vac. II and Vac. III. Then we show that the BPS walls can connect the complex vacua of Vac. II and Vac. III, but cannot connect Vac. I and complex vacua of Vac. II. ### 3.2 BPS walls in model I Here, we construct BPS saturated walls in the model with one flat direction (Model I). The BPS equation (2.3) for this wall is $`{\displaystyle \frac{\varphi ^i}{z}}=\varphi ^i(\stackrel{}{\varphi }^2a^2).`$ (3.8) First we show that there is no complex solution of this BPS equation. When we map the field space to the superpotential space, two disjoint vacua are mapped to two points. It is known that the configuration of the BPS wall can be mapped to a line segment connecting these two points in the superpotential space. Now, the difference between the values of the superpotentials for the two vacua, $`\mathrm{\Delta }W=a^4/4`$, is real. This means that the configuration of the BPS wall in the superpotential space is also real. If we set $`\stackrel{}{\varphi }=\stackrel{}{x}+i\stackrel{}{y}`$, the imaginary part of the superpotential is $`\mathrm{}W=4(\stackrel{}{x}\stackrel{}{y})(\stackrel{}{x}^2\stackrel{}{y}^2a^2)`$, so we find that BPS solution must satisfy the constraint $`\stackrel{}{x}\stackrel{}{y}=0`$. Using this constraint, the BPS equation of Eq. (3.8) can be rewritten as $`{\displaystyle \frac{d}{dz}}(\stackrel{}{x}+i\stackrel{}{y})=(\stackrel{}{x}i\stackrel{}{y})(\stackrel{}{x}^2\stackrel{}{y}^2a^2).`$ (3.9) From this equation we can derive the following equations: $`{\displaystyle \frac{d}{dz}}\left({\displaystyle \frac{x^2}{x^1}}\right)={\displaystyle \frac{d}{dz}}\left({\displaystyle \frac{y^1}{y^2}}\right)=0,{\displaystyle \frac{d}{dz}}(x^iy^j)=0,\text{for}i,j=1,2.`$ (3.10) The first of these two equations implies that the $`O(2)`$ rotation parameter $`\theta `$ is constant for the BPS wall. Combining these with the constraint $`\stackrel{}{x}\stackrel{}{y}=0`$, we can parameterize the BPS wall as $`\stackrel{}{\varphi }(z)=v(z)\left(\begin{array}{c}\mathrm{cos}\theta \\ \mathrm{sin}\theta \end{array}\right)+iu(z)\left(\begin{array}{c}\mathrm{cos}(\theta +\pi /2)\\ \mathrm{sin}(\theta +\pi /2)\end{array}\right).`$ (3.15) In Fig. 2 (a), we plot the moduli space of this model in the $`(u,v)`$-plane. Substituting Eq. (3.15) into the second equation of Eq. (3.10), we can immediately find $`{\displaystyle \frac{d(uv)}{dz}}=0,uv=\text{const}\sqrt{c},`$ (3.17) where $`c`$ is a real integral constant. From Fig. 2(a) we find that there is no complex BPS solution connecting Vac. I and vacua along the flat direction of Vac. II: In order for a BPS wall to reach Vac. I, we need to set $`uv=\sqrt{c}=0`$, and this is reduced to a real solution \[$`u(z)=0`$\] for the boundary condition of Vac. II on the other side. Hence we consider this solution of Eq. (3.9). This solution can be found as $`v=\varphi _Wa\sqrt{{\displaystyle \frac{1}{1+\mathrm{exp}\left[2a^2(zz_0)\right]}}},u=0,`$ (3.18) where $`z_0`$ is an integral constant, representing the position of the center of the domain wall. We plot this real solution in Fig. 2(b). Using an $`O(2)`$ transformation, the general real solutions can be written as $`\varphi ^1=\varphi _W\mathrm{cos}\theta ,\varphi ^2=\varphi _W\mathrm{sin}\theta ,`$ (3.19) where $`\theta `$ is a real parameter. The wall separates the two vacua in the broken phase and the unbroken phase. The wall interpolating between the broken and unbroken phase of the discrete symmetry $`Z_2`$ is discussed in Ref. . ### 3.3 BPS wall in the model II In this section, we construct BPS walls in the model with two flat directions (model II). This model has three disjoint vacua as in the case of Eq. (3.7). The difference between the values of the superpotentials for each pair of the three vacua is real, as in the previous model. There exists no BPS wall connecting Vac. I and Vac. III, because the two values of the superpotential corresponding to these two vacua are the same, and the BPS bound (2.2) becomes zero. For this reason we consider two kinds of walls: walls interpolating between Vac. II and Vac. III (“outer walls”), and walls interpolating between Vac. I and Vac. II (“inner walls”). The BPS equations (2.3) for these walls are $`{\displaystyle \frac{\varphi ^i}{z}}=\varphi ^i\left(\stackrel{}{\varphi }^2{\displaystyle \frac{a^2}{3}}\right)(\stackrel{}{\varphi }^2a^2),`$ (3.20) where the boundary conditions are $`\stackrel{}{\varphi }(\mathrm{})=a^2`$ \[$`\stackrel{}{\varphi }(\mathrm{})=0`$\] and $`\stackrel{}{\varphi }(\mathrm{})=a^2/3`$ for the outer (inner) walls. The map of the BPS walls into the superpotential space must be real, as in the previous model: If we set $`\stackrel{}{\varphi }=\stackrel{}{x}+i\stackrel{}{y}`$, the imaginary part of the superpotential in this model becomes $`\mathrm{}W={\displaystyle \frac{1}{3}}(\stackrel{}{x}\stackrel{}{y})[3(\stackrel{}{x}^{\mathrm{\hspace{0.17em}2}}\stackrel{}{y}^{\mathrm{\hspace{0.17em}2}}a^2)(\stackrel{}{x}^{\mathrm{\hspace{0.17em}2}}\stackrel{}{y}^{\mathrm{\hspace{0.17em}2}}a^2/3)4(\stackrel{}{x}\stackrel{}{y})^2].`$ (3.21) Thus $`\stackrel{}{x}\stackrel{}{y}=0`$ is a sufficient condition. <sup>4</sup><sup>4</sup>4 We can show that this is also a necessary condition using the continuity of the solution. With this condition, Eq. (3.10) is again valid, and we can set $`\stackrel{}{\varphi }`$ as in Eq. (3.15). Hence we can set $`\theta =0`$ in Eq. (3.15) by using the $`O(2)`$ transformation, without loss of generality, yielding $`\stackrel{}{\varphi }=\left(\begin{array}{c}v\\ iu\end{array}\right)`$, where $`v`$ and $`u`$ are real scalar fields. In Fig. 3, we illustrate the moduli space of this model in the $`(u,v)`$-plane. Equation (3.20) becomes $`{\displaystyle \frac{dv}{dz}}`$ $`=`$ $`v\left(v^2u^2{\displaystyle \frac{a^2}{3}}\right)(v^2u^2a^2),`$ $`{\displaystyle \frac{du}{dz}}`$ $`=`$ $`u\left(v^2u^2{\displaystyle \frac{a^2}{3}}\right)(v^2u^2a^2).`$ (3.23) We can then find $`{\displaystyle \frac{d(uv)}{dz}}=0.`$ (3.24) Hence, we can set $`uv=\text{ const}=\sqrt{c}`$. We find, from Fig. 3 (a), that there can exist a complex BPS wall solution connecting Vac. II and Vac. III, but no complex BPS wall can connect Vac. I and Vac. II, for the same reason as in model I. The first equation in Eq. (3.23) becomes $`{\displaystyle \frac{dv^2}{dz}}=2{\displaystyle \frac{1}{v^2}}\left((v^2)^2{\displaystyle \frac{a^2}{3}}v^2c\right)((v^2)^2a^2v^2c).`$ (3.25) This can be integrated to give $`e^{\frac{4a^2}{3}(zz_0)}=\left\|{\displaystyle \frac{v^2\frac{1}{2}(\frac{a^2}{3}+\sqrt{\frac{a^4}{9}+4c})}{v^2\frac{1}{2}(\frac{a^2}{3}\sqrt{\frac{a^4}{9}+4c})}}\right\|^{\frac{1}{\sqrt{\frac{a^4}{9}+4c}}}\left\|{\displaystyle \frac{v^2\frac{1}{2}(a^2\sqrt{a^4+4c})}{v^2\frac{1}{2}(a^2+\sqrt{a^4+4c})}}\right\|^{\frac{1}{\sqrt{a^4+4c}}},`$ (3.26) where $`z_0`$ is the center of the wall. For the complex solution interpolating between Vac. II and Vac. III, (3.26) can be rewritten as $`e^{\frac{4a^2}{3}(zz_0)}=\left[{\displaystyle \frac{v^2\frac{1}{2}(\frac{a^2}{3}+\sqrt{\frac{a^4}{9}+4c})}{v^2\frac{1}{2}(\frac{a^2}{3}\sqrt{\frac{a^4}{9}+4c})}}\right]^{\frac{1}{\sqrt{\frac{a^4}{9}+4c}}}\left[{\displaystyle \frac{v^2\frac{1}{2}(a^2\sqrt{a^4+4c})}{\frac{1}{2}(a^2+\sqrt{a^4+4c})v^2}}\right]^{\frac{1}{\sqrt{a^4+4c}}}.`$ (3.27) Since we cannot obtain an explicit solution $`v(z)`$ of this equation, we plot $`v(z)`$ in the Fig. 3(b) as an implicit solution of a complex BPS wall. We must note that the complex solution of $`uv=\sqrt{c}`$ is not the $`O(2)^𝐂`$ transformation of the solution of $`uv=0`$. Let us consider a vacuum transformed by a $`O(2)^𝐂`$ parameter from a real expectation value in Vac. II. The complex BPS wall solution connects this Vac. II to the Vac. III transformed by a different $`O(2)^𝐂`$ parameter from the corresponding real expectation value in Vac. III. Therefore the $`O(2)^𝐂`$ transformation of a BPS solution does not become a solution of the BPS equation; the parameter $`c`$ which labels the imaginary direction is not associated with the $`O(2)^𝐂`$ symmetry. We can find an explicit solution for real BPS walls. For the real solution, the integrated BPS equation can be obtained by setting $`c=0`$ in Eq. (3.26). We have $`X\stackrel{\mathrm{def}}{=}\mathrm{exp}\left[{\displaystyle \frac{4a^4}{3}}(zz_0)\right]={\displaystyle \frac{\|v^2a^2\|v^4}{\|v^2\frac{a^2}{3}\|^3}}={\displaystyle \frac{\|\mathrm{\Phi }a^2\|\mathrm{\Phi }^2}{\|\mathrm{\Phi }\frac{a^2}{3}\|^3}},`$ (3.28) where we have defined $`\mathrm{\Phi }\stackrel{\mathrm{def}}{=}v^2=(\varphi ^1)^2`$. We solve this equation in the outer region, $`\frac{a^2}{3}(\varphi ^1)^2a^2`$, and the inner region, $`0(\varphi ^1)^2\frac{a^2}{3}`$, separately. In the case of the outer solutions, $`\frac{a^2}{3}(\varphi ^1)^2a^2`$, Eq. (3.28) can be rewritten as the third order equation $`(X+1)\mathrm{\Phi }^3a^2(X+1)\mathrm{\Phi }^2+{\displaystyle \frac{a^4}{3}}X\mathrm{\Phi }{\displaystyle \frac{a^6}{27}}X=0.`$ (3.29) Thus the third order equation can be solved to yield $`(\varphi ^1)^2={\displaystyle \frac{a^2}{3}}\left[1+\left({\displaystyle \frac{1}{1+X}}+\sqrt{{\displaystyle \frac{X}{(X+1)^3}}}\right)^{\frac{1}{3}}+\left({\displaystyle \frac{1}{1+X}}\sqrt{{\displaystyle \frac{X}{(X+1)^3}}}\right)^{\frac{1}{3}}\right]`$ (3.30) for a real solution. (The two other solutions are complex and thus inappropriate.) In the case of the inner solutions, $`0(\varphi ^1)^2\frac{a^2}{3}`$, Eq. (3.28) can be rewritten as $`(X1)\mathrm{\Phi }^3a^2(X1)\mathrm{\Phi }^2+{\displaystyle \frac{a^4}{3}}X\mathrm{\Phi }{\displaystyle \frac{a^6}{27}}X=0.`$ (3.31) In this case, we must solve this equation for each case $`X=1`$ and $`X1`$ separately. When $`X=1`$, the solution of this equation is $`(\varphi ^1)^2=a^2/9`$, and this corresponds to the expectation value at the center of the wall ($`z=z_0`$). When $`X1`$ ($`zz_0`$), there are three candidates for the solution of the outer wall: $`(\varphi ^1)^2={\displaystyle \frac{a^2}{3}}[1+({\displaystyle \frac{1}{1X}}+\sqrt{{\displaystyle \frac{X}{(X1)^3}}})^{\frac{1}{3}}\left(\begin{array}{c}1\\ e^{\frac{2\pi }{3}i}\\ e^{\frac{2\pi }{3}i}\end{array}\right)`$ $`+({\displaystyle \frac{1}{1X}}\sqrt{{\displaystyle \frac{X}{(X1)^3}}})^{\frac{1}{3}}\left(\begin{array}{c}1\\ e^{\frac{2\pi }{3}i}\\ e^{\frac{2\pi }{3}i}\end{array}\right)].`$ (3.32) These solutions are not real and positive, so we must choose the correct one for the regions $`z<z_0`$ ($`X<1`$) and $`z>z_0`$ ($`X>1`$). In the region $`z>z_0`$, the first solution is appropriate for the real solution. In the region of $`z<z_0`$, the third solution is appropriate. (In the latter case, the first solution cannot satisfy the correct boundary conditions, $`(\varphi ^1)^2(\mathrm{})=0`$, and the second solution tends to infinity in the limit $`zz_0`$.) In summary, we obtain the inner wall solution by using the third solution in the left ($`z<z_0`$) and the first solution in the right ($`z>z_0`$): $`(\varphi ^1)^2=\{\begin{array}{c}\frac{a^2}{3}\left[1+\left(\frac{1}{1X}+\sqrt{\frac{X}{(X1)^3}}\right)^{\frac{1}{3}}\left(\frac{1}{1X}+\sqrt{\frac{X}{(X1)^3}}\right)^{\frac{1}{3}}\right](z>z_0)\hfill \\ \frac{a^2}{3}\left[1+\left(\frac{1}{1X}+\sqrt{\frac{X}{(X1)^3}}\right)^{\frac{1}{3}}e^{\frac{2\pi }{3}i}+\left(\frac{1}{1X}\sqrt{\frac{X}{(X1)^3}}\right)^{\frac{1}{3}}e^{\frac{2\pi }{3}i}\right](z<z_0)\hfill \end{array}`$ (3.35) The profiles of the outer and inner wall solutions are plotted in Fig. 4. ## 4 Conclusions and discussion We considered BPS domain walls in models with continuously degenerate moduli spaces. We discussed only two $`O(2)`$ symmetric models explicitly, but many results can be straightforwardly generalized to other models with a global symmetry $`G`$. When a model has a continuous symmetry, $`O(2)`$ in our models, the BPS equation of the wall becomes covariant under this symmetry, so the BPS wall also has this symmetry. If we can find a BPS solution, configurations obtained through transformations of this solution by elements of $`G`$ are also BPS solutions, so they constitute a family of BPS walls. Although the boundary conditions change under these transformations, the tensions of the walls never change. In supersymmetric field theories, the symmetry $`G`$ of the superpotential is enlarged to its complexification, $`G^𝐂`$, due to the holomorphy of the superpotential. Therefore the vacuum manifold includes non-compact flat directions corresponding to the directions of imaginary parts of the vacuum expectation values. As the BPS equation is not covariant under $`G^𝐂`$, it is a highly nontrivial problem to determine whether there can exist complex BPS walls interpolating between two disjoint non-compact flat directions. To examine this problem, we considered two models with flat directions. We found that there is no complex BPS wall in the first model, while there can exist a family of complex BPS walls in the second model. We have learned from the examination of these two models that we must consider complex configurations for BPS walls in models with continuous symmetry. This is an important lesson, since only real configurations of BPS walls have been considered in the literature. We have not yet found a criterion to determine whether or not a complex BPS wall exists in general models. Let us now examine general structures for the existence of complex BPS walls by counting the number of degrees of freedom in these two models. Since the BPS equation is a first order differential equation, it can be expected that the general solution has the same number of integral constants as the number of BPS equations, unless we enforce the boundary conditions. Since we considered supersymmetric models with two chiral superfields, there are four BPS equations corresponding to the four real scalar degrees of freedom. However, we have been able to eliminate one degree of freedom, since it must be the case of that any BPS solution maps to a straight line in the superpotential space. We thus can expect that BPS solutions can maximally include three free parameters as the integral constants. In fact, three parameters, $`z_0`$, $`\theta `$ and $`c`$, have appeared as the integral constants in the BPS solutions in the second model. However, the third parameter $`c`$ is not contained in the BPS wall solution of the first model: It was eliminated by the boundary condition. We can interpret the parameters $`z_0`$ and $`\theta `$ (and $`c`$), labeling the solutions of the BPS walls, as the “moduli” of the BPS wall solutions, since the tension of the wall does not change when we continuously vary the values of these parameters. The configurations obtained under such variations are all solutions of a BPS equation, and their tension realize the same BPS bound. These parameters, however, have slightly different meanings: since $`z_0`$ represents the location of the center of the wall, we can vary this parameter without changing the boundary conditions. Contrastingly, we cannot vary $`\theta `$ and $`c`$ without changing the boundary condition. Next we discuss the nature of these parameters in terms of symmetry. Two of the three parameters represent the Nambu-Goldstone modes corresponding to the symmetries broken by the existence of the wall configuration; $`z_0`$ corresponds to translation along the $`z`$-axis in the spacetime, and $`\theta `$ to the continuous internal symmetry, $`O(2)`$. Therefore the BPS wall solution apparently contains these free parameters. The parameter $`c`$ can be considered to represent the deformation of the BPS wall along the non-compact flat direction, which originates from the complexified symmetry of the superpotential. This is, however, not the symmetry of the whole model (the Kähler potential is invariant under $`G`$ but not $`G^𝐂`$), and therefore the additional parameter $`c`$ does not directly correspond to the complex symmetry. This is why the BPS wall solutions do not always contain $`c`$. Concerning this fact, we must comment on the similarity with the results in Ref. . As discussed above, the parameter $`c`$ in our model is the additional integral constant, which depends on the details of the model. This quantity is similar to the additional integrals of motion in Ref. ,<sup>5</sup><sup>5</sup>5 Similar additional constants are discussed in the context of non-supersymmetric models in Ref. . in the sense that in both models the additional constants do not correspond directly to the symmetry of the theory. However, we must emphasize that these quantities have essentially different origins: The additional integrals of motion in Ref. represent the spatial distance (in the spacetime) of two separated BPS walls, while the quantity $`c`$ in our model controls the shift of the BPS walls along the flat direction in the internal space. Let us discuss an interesting problem regarding the Nambu-Goldstone theorem suggested by our models. The moduli space of supersymmetric vacua is parameterized by the NG and the quasi-NG bosons associated with the spontaneously broken $`G^𝐂`$ symmetry of the superpotential, and with their superpartners they constitute massless NG chiral multiplets as described by the supersymmetric extension of the Nambu-Goldstone theorem. However, the configuration of the BPS domain wall spontaneously breakes half of the supersymmetry (and the translational symmetry along the z-axis). Therefore, in the entire four-dimensional spacetime, $`𝒩=1`$ massless NG supermultiplets are justified only at infinite distance from the wall. The supersymmetric Nambu-Goldstone theorem must be deformed around the wall. This fact may be a reason why the complex parameter of the $`O(2)^𝐂`$ transformation does not appear as the moduli of BPS wall solutions. It would be interesting to examine the extension of the Nambu-Goldstone theorem to the case of a BPS wall background, or the case in which half of the supersymmetry is spontaneously broken. Before ending this conclusion, we point out some interesting features of our models. We found that there can exist BPS walls connecting $`O(2)`$ symmetric and $`O(2)`$ broken vacua. (For conventional BPS walls, the broken symmetry is usually discrete, and vacua separated by the wall are both in the broken phase.) Mass spectra are different on opposite sides of the walls in our models: We can expect massless (quasi-)NG bosons and their superpartners only in the broken phase. It is a future problem to examine the wave functions of these massless modes in order to determine this difference. Our second model has three disjoint vacua, and the maps from two of them to the superpotential space coincide accidentally. By modifying the model slightly, we can construct a model with three disjoint vacua mapped to three distinct points in the superpotential space. Hence our examinations can be extended to the case of the BPS domain wall junction. We expect that the new types of BPS domains wall found in this paper will play an important role in the further understanding of non-perturbative aspects of supersymmetric quantum field theories. ## Acknowledgements We would like to thank N. Sakai for valuable discussions. We are also grateful to T. Kugo, H. Kawai, H. Kunitomo and N. Sasakura for useful comments. The work of M. Nitta is supported in part by JSPS Research Fellowships.
warning/0007/nlin0007041.html	ar5iv	text	# 7.1 The inversion of the Bäcklund transformation ### 7.1 The inversion of the Bäcklund transformation Let us call the direct map by $`B_P`$. The inverse map acts from $`\stackrel{~}{L}(u)`$ to $`L(u)`$. We can rewrite the equations for the $`B_P`$ in the inverse form $$M^{}(u)\stackrel{~}{L}(u)=L(u)M^{}(u),M^{}(u)=\left(\begin{array}{cc}\gamma & p\\ q& u\lambda +pq/\gamma \end{array}\right).$$ (7.1) To define the inverse map we must find expressions for the co-factor matrix $`M^{}(u)`$, or for the variables $`p`$ and $`q`$, in terms of $`\stackrel{~}{}`$-variables, i.e. in terms of the entries of $`\stackrel{~}{L}(u)`$. We have already the expressions (3.10), $$p=\frac{J_{}}{2\alpha },q=\frac{\stackrel{~}{J}_+}{2\alpha },$$ (7.2) which define $`q`$. To obtain the formula for the variable $`p`$ we will use again the spectrality property. The matrix $`M^{}(\lambda )`$ has a one-dimensional kernel $`\stackrel{~}{\mathrm{\Omega }}`$, $$M^{}(\lambda )\stackrel{~}{\mathrm{\Omega }}=\left(\begin{array}{cc}\gamma & p\\ q& pq/\gamma \end{array}\right)\stackrel{~}{\mathrm{\Omega }}=0,\stackrel{~}{\mathrm{\Omega }}=\left(\begin{array}{c}p\\ \gamma \end{array}\right).$$ (7.3) The main difference comparing to the formulas of the direct map is that the inverse map will be parametrized by the point $`Q=(\lambda ,\mu )\mathrm{\Gamma }`$, not the $`P=(\lambda ,\mu )\mathrm{\Gamma }`$. Therefore, $`\stackrel{~}{\mathrm{\Omega }}`$ is an eigenvector of the matrix $`\stackrel{~}{L}(u)`$ with the eigenvalue $`Q=(\lambda ,\mu )`$: $$\left(\begin{array}{cc}\stackrel{~}{A}(\lambda )+\mu & \stackrel{~}{B}(\lambda )\\ \stackrel{~}{C}(\lambda )& \stackrel{~}{A}(\lambda )+\mu \end{array}\right)\left(\begin{array}{c}p\\ \gamma \end{array}\right)=0.$$ (7.4) This gives us the needed formula for the variable $`p`$, $$p=\gamma \frac{\stackrel{~}{A}(\lambda )\mu }{\stackrel{~}{C}(\lambda )}=\gamma \frac{\stackrel{~}{B}(\lambda )}{\stackrel{~}{A}(\lambda )+\mu }.$$ (7.5) To prove that this does indeed give the inverse map, we have to show that the two formulas, (3.17) and (7.5), define in fact the same variable $`\mu `$: $$\mu =\stackrel{~}{A}(\lambda )\frac{p}{\gamma }\stackrel{~}{C}(\lambda )\stackrel{\mathrm{?}}{=}A(\lambda )\frac{q}{\gamma }B(\lambda ).$$ (7.6) It is easy to see that this equation is the $`(11)`$-element of the matrix identity: $$M^{}(\lambda )\stackrel{~}{L}(\lambda )=L(\lambda )M^{}(\lambda ).$$ (7.7) We will denote as $`_Q`$ the map which is inverse to the map $`B_P`$. Generally speaking, we have constructed 4 different maps, $`B_P`$, $`B_Q`$, $`_Q`$, and $`_P`$, with two pairs of maps which are inverse to each other: $$_QB_P=B_P_Q=_PB_Q=B_Q_P=\text{Id}.$$ (7.8) ### 7.2 The two-point map $`B_{P_1,Q_2}`$ We now construct a composite map which is a product of the map $`B_{P_1}B_{(\lambda _1,\mu _1)}`$ and $`_{Q_2}_{(\lambda _2,\mu _2)}`$: $$B_{P_1,Q_2}=_{Q_2}B_{P_1}.$$ (7.9) The second parameter of the basic map, namely the $`\gamma `$, is taken the same in both maps, so $`\gamma _1=\gamma _2`$. Obviously, when $`\lambda _1=\lambda _2`$ (and $`\mu _1=\mu _2`$) this composite map will turn into an identity map. The first map $`B_{P_1}`$ reads as follows: $$M_1(u)L(u)=\stackrel{~}{L}(u)M_1(u),M_1(u)=\left(\begin{array}{cc}u\lambda _1+p_1q_1/\gamma & p_1\\ q_1& \gamma \end{array}\right),$$ (7.10) where the formulas for the variables $`p_1`$ and $`q_1`$ are $$p_1=\frac{J_{}}{2\alpha }=\gamma \frac{\stackrel{~}{A}(\lambda _1)\mu _1}{\stackrel{~}{C}(\lambda _1)}=\gamma \frac{\stackrel{~}{B}(\lambda _1)}{\stackrel{~}{A}(\lambda _1)+\mu _1},$$ (7.11) $$q_1=\frac{\stackrel{~}{J}_+}{2\alpha }=\gamma \frac{A(\lambda _1)\mu _1}{B(\lambda _1)}=\gamma \frac{C(\lambda _1)}{A(\lambda _1)+\mu _1}.$$ (7.12) The second map $`B_{Q_2}`$ reads as follows: $$M_2(u)\stackrel{~}{L}(u)=\stackrel{~}{\stackrel{~}{L}}(u)M_2(u),M_2(u)=\left(\begin{array}{cc}\gamma & p_2\\ q_2& u\lambda _2+p_2q_2/\gamma \end{array}\right),$$ (7.13) where the formulas for the variables $`p_2`$ and $`q_2`$ are $$p_2=\frac{\stackrel{~}{\stackrel{~}{J}}_{}}{2\alpha }=\gamma \frac{\stackrel{~}{A}(\lambda _2)\mu _2}{\stackrel{~}{C}(\lambda _2)}=\gamma \frac{\stackrel{~}{B}(\lambda _2)}{\stackrel{~}{A}(\lambda _2)+\mu _2},$$ (7.14) $$q_2=\frac{\stackrel{~}{J}_+}{2\alpha }=\gamma \frac{\stackrel{~}{\stackrel{~}{A}}(\lambda _2)\mu _2}{\stackrel{~}{\stackrel{~}{B}}(\lambda _2)}=\gamma \frac{\stackrel{~}{\stackrel{~}{C}}(\lambda _2)}{\stackrel{~}{\stackrel{~}{A}}(\lambda _2)+\mu _2}.$$ (7.15) Notice that $`q_1`$ is equal to $`q_2`$, hence we omit the sub-index, $`q_1=q_2=q`$. The composite map $`B_{P_1,Q_2}`$ acts from $`L(u)`$ to $`\stackrel{~}{\stackrel{~}{L}}(u)`$, $$M(u)L(u)=\stackrel{~}{\stackrel{~}{L}}(u)M(u),$$ (7.16) $$M(u)=\frac{1}{\gamma }M_2(u)M_1(u)=\left(\begin{array}{cc}u\lambda _1+\frac{q}{\gamma }(p_1p_2)& p_1p_2\\ \frac{q}{\gamma }(\lambda _1\lambda _2\frac{q}{\gamma }(p_1p_2))& u\lambda _2\frac{q}{\gamma }(p_1p_2)\end{array}\right).$$ (7.17) In order to get rid of the intermediate $`\stackrel{~}{}`$-variables, we will use the spectrality property with respect to two points, $`P_1=(\lambda _1,\mu _1)`$ and $`Q_2=(\lambda _2,\mu _2)`$. Obviously, both spectralities are still valid after composing the maps. For the point $`P_1`$ we get the following equations: $$M(\lambda _1)\mathrm{\Omega }_1=0,\mathrm{\Omega }_1=\left(\begin{array}{c}\gamma \\ q\end{array}\right),L(\lambda _1)\mathrm{\Omega }_1=\mu _1\mathrm{\Omega }_1$$ $$q=\gamma \frac{A(\lambda _1)\mu _1}{B(\lambda _1)}=\gamma \frac{C(\lambda _1)}{A(\lambda _1)+\mu _1};$$ (7.18) $$M^{}(\lambda _1)\stackrel{~}{\stackrel{~}{\mathrm{\Omega }}}_1=0,\stackrel{~}{\stackrel{~}{\mathrm{\Omega }}}_1=\left(\begin{array}{c}p_1p_2\\ \lambda _1\lambda _2\frac{q}{\gamma }(p_1p_2)\end{array}\right),\stackrel{~}{\stackrel{~}{L}}(\lambda _1)\stackrel{~}{\stackrel{~}{\mathrm{\Omega }}}_1=\mu _1\stackrel{~}{\stackrel{~}{\mathrm{\Omega }}}_1$$ $$p_1p_2=\frac{\gamma (\lambda _2\lambda _1)\stackrel{~}{\stackrel{~}{B}}(\lambda _1)}{\gamma \left(\stackrel{~}{\stackrel{~}{A}}(\lambda _1)+\mu _1\right)q\stackrel{~}{\stackrel{~}{B}}(\lambda _1)}=\frac{\gamma (\lambda _1\lambda _2)\left(\stackrel{~}{\stackrel{~}{A}}(\lambda _1)\mu _1\right)}{q\left(\stackrel{~}{\stackrel{~}{A}}(\lambda _1)\mu _1\right)+\gamma \stackrel{~}{\stackrel{~}{C}}(\lambda _1)}.$$ (7.19) For the point $`Q_2`$ we get the second set of equations: $$M(\lambda _2)\mathrm{\Omega }_2=0,\mathrm{\Omega }_2=\left(\begin{array}{c}p_1p_2\\ \lambda _1\lambda _2\frac{q}{\gamma }(p_1p_2)\end{array}\right),L(\lambda _2)\mathrm{\Omega }_2=\mu _2\mathrm{\Omega }_2$$ $$p_1p_2=\frac{\gamma (\lambda _2\lambda _1)B(\lambda _2)}{\gamma \left(A(\lambda _2)+\mu _2\right)qB(\lambda _2)}=\frac{\gamma (\lambda _1\lambda _2)\left(A(\lambda _2)\mu _2\right)}{q\left(A(\lambda _2)\mu _2\right)+\gamma C(\lambda _2)};$$ (7.20) $$M^{}(\lambda _2)\stackrel{~}{\stackrel{~}{\mathrm{\Omega }}}_2=0,\stackrel{~}{\stackrel{~}{\mathrm{\Omega }}}_2=\left(\begin{array}{c}\gamma \\ q\end{array}\right),\stackrel{~}{\stackrel{~}{L}}(\lambda _2)\stackrel{~}{\stackrel{~}{\mathrm{\Omega }}}_2=\mu _2\stackrel{~}{\stackrel{~}{\mathrm{\Omega }}}_2$$ $$q=\gamma \frac{\stackrel{~}{\stackrel{~}{A}}(\lambda _2)\mu _2}{\stackrel{~}{\stackrel{~}{B}}(\lambda _2)}=\gamma \frac{\stackrel{~}{\stackrel{~}{C}}(\lambda _2)}{\stackrel{~}{\stackrel{~}{A}}(\lambda _2)+\mu _2}.$$ (7.21) Equations (7.18) and (7.21) are already known to us (cf. (7.12) and (7.15)). The formulas (7.19) and (7.20) for the variable $`p_1p_2`$ are new. They are equivalent to the formulas (7.11) and (7.14) expressed in terms of entries of $`L(u)`$ and $`\stackrel{~}{\stackrel{~}{L}}(u)`$. Concluding, we have constructed a two-point Bäcklund transformation which is factorised to two one-point Bäcklund transformations and which is explicitly given, together with its inverse, by the formulas: $$M(u)L(u)=\stackrel{~}{\stackrel{~}{L}}(u)M(u),M(u)=\left(\begin{array}{cc}u\lambda _1+xX& X\\ x^2X+(\lambda _1\lambda _2)x& u\lambda _2xX\end{array}\right),$$ (7.22) where $$x:=\frac{A(\lambda _1)\mu _1}{B(\lambda _1)}=\frac{C(\lambda _1)}{A(\lambda _1)+\mu _1}=\frac{\stackrel{~}{\stackrel{~}{A}}(\lambda _2)\mu _2}{\stackrel{~}{\stackrel{~}{B}}(\lambda _2)}=\frac{\stackrel{~}{\stackrel{~}{C}}(\lambda _2)}{\stackrel{~}{\stackrel{~}{A}}(\lambda _2)+\mu _2},$$ (7.23) $$X:=\frac{(\lambda _2\lambda _1)B(\lambda _1)B(\lambda _2)}{B(\lambda _1)(A(\lambda _2)+\mu _2)B(\lambda _2)(A(\lambda _1)\mu _1)};$$ (7.24) $$=\frac{(\lambda _1\lambda _2)B(\lambda _1)(A(\lambda _2)\mu _2)}{(A(\lambda _1)\mu _1)(A(\lambda _2)\mu _2)+B(\lambda _1)C(\lambda _2)}$$ (7.25) $$=\frac{(\lambda _2\lambda _1)B(\lambda _2)(A(\lambda _1)+\mu _1)}{(A(\lambda _1)+\mu _1)(A(\lambda _2)+\mu _2)+B(\lambda _2)C(\lambda _1)}$$ (7.26) $$=\frac{(\lambda _1\lambda _2)(A(\lambda _1)+\mu _1)(A(\lambda _2)\mu _2)}{(A(\lambda _1)+\mu _1)C(\lambda _2)(A(\lambda _2)\mu _2)C(\lambda _1)}$$ (7.27) $$=\frac{(\lambda _2\lambda _1)\stackrel{~}{\stackrel{~}{B}}(\lambda _2)\stackrel{~}{\stackrel{~}{B}}(\lambda _1)}{\stackrel{~}{\stackrel{~}{B}}(\lambda _2)\left(\stackrel{~}{\stackrel{~}{A}}(\lambda _1)+\mu _1\right)\stackrel{~}{\stackrel{~}{B}}(\lambda _1)\left(\stackrel{~}{\stackrel{~}{A}}(\lambda _2)\mu _2\right)}$$ (7.28) $$=\frac{(\lambda _1\lambda _2)\stackrel{~}{\stackrel{~}{B}}(\lambda _2)\left(\stackrel{~}{\stackrel{~}{A}}(\lambda _1)\mu _1\right)}{\left(\stackrel{~}{\stackrel{~}{A}}(\lambda _2)\mu _2\right)\left(\stackrel{~}{\stackrel{~}{A}}(\lambda _1)\mu _1\right)+\stackrel{~}{\stackrel{~}{B}}(\lambda _2)\stackrel{~}{\stackrel{~}{C}}(\lambda _1)}$$ (7.29) $$=\frac{(\lambda _2\lambda _1)\stackrel{~}{\stackrel{~}{B}}(\lambda _1)\left(\stackrel{~}{\stackrel{~}{A}}(\lambda _2)+\mu _2\right)}{\left(\stackrel{~}{\stackrel{~}{A}}(\lambda _2)+\mu _2\right)\left(\stackrel{~}{\stackrel{~}{A}}(\lambda _1)+\mu _1\right)+\stackrel{~}{\stackrel{~}{B}}(\lambda _1)\stackrel{~}{\stackrel{~}{C}}(\lambda _2)}$$ (7.30) $$=\frac{(\lambda _1\lambda _2)\left(\stackrel{~}{\stackrel{~}{A}}(\lambda _2)+\mu _2\right)\left(\stackrel{~}{\stackrel{~}{A}}(\lambda _1)\mu _1\right)}{\left(\stackrel{~}{\stackrel{~}{A}}(\lambda _2)+\mu _2\right)\stackrel{~}{\stackrel{~}{C}}(\lambda _1)\left(\stackrel{~}{\stackrel{~}{A}}(\lambda _1)\mu _1\right)\stackrel{~}{\stackrel{~}{C}}(\lambda _2)}.$$ (7.31) The above formulas give several equivalent expressions for the variables $`x`$ and $`X`$ since the points $`(\lambda _1,\mu _1)`$ and $`(\lambda _2,\mu _2)`$ belong to the spectral curve $`\mathrm{\Gamma }`$, i.e. are bound by the following relations: $$\mu _1^2=A^2(\lambda _1)+B(\lambda _1)C(\lambda _1),\mu _2^2=A^2(\lambda _2)+B(\lambda _2)C(\lambda _2),$$ (7.32) $$\mu _1^2=\stackrel{~}{\stackrel{~}{A}}^2(\lambda _1)+\stackrel{~}{\stackrel{~}{B}}(\lambda _1)\stackrel{~}{\stackrel{~}{C}}(\lambda _1),\mu _2^2=\stackrel{~}{\stackrel{~}{A}}^2(\lambda _2)+\stackrel{~}{\stackrel{~}{B}}(\lambda _2)\stackrel{~}{\stackrel{~}{C}}(\lambda _2).$$ (7.33) ### 7.3 Two-point map as a discrete-time map We will see in this Section that the two-point map constructed above is a one-parameter, $`\lambda _1`$, time-discretization of a family of flows parameterised by the point $`Q_2=(\lambda _2,\mu _2)`$, with the difference $`\lambda _1\lambda _2`$ playing the role of the time-step. Indeed, consider the limit $`\lambda _1\lambda _2`$, $$\lambda _1=\lambda _2+\epsilon ,\epsilon 0.$$ (7.34) It is easy to see from the formulas of the previous subsection that $$x=x_0+O(\epsilon ),x_0=\frac{A(\lambda _2)\mu _2}{B(\lambda _2)}=\frac{C(\lambda _2)}{A(\lambda _2)+\mu _2}$$ (7.35) and $$X=\epsilon X_0+O(\epsilon ^2),X_0=\frac{B(\lambda _2)}{2\mu _2}.$$ (7.36) Then we derive that the matrix $`M(u)`$ has the following asymptotics: $$M(u)=(u\lambda _2)\left(1\frac{\epsilon }{2\mu _2(u\lambda _2)}\left(\begin{array}{cc}A(\lambda _2)+\mu _2& B(\lambda _2)\\ C(\lambda _2)& A(\lambda _2)+\mu _2\end{array}\right)\right)+O(\epsilon ^2).$$ (7.37) If we now define the time-derivative $`\dot{L}(u)`$ as $$\dot{L}(u)=\underset{\epsilon 0}{lim}\frac{\stackrel{~}{\stackrel{~}{L}}(u)L(u)}{\epsilon },$$ (7.38) then in the limit we obtain from the equation of the map, $`M(u)L(u)=\stackrel{~}{\stackrel{~}{L}}(u)M(u)`$, the Lax equation for a corresponding continuous flow that our Bäcklund transformation discretizes, namely: $$\dot{L}(u)=[L(u),\frac{L(\lambda _2)}{2\mu _2(u\lambda _2)}].$$ (7.39) This is a Hamiltonian flow with $`\mu _2`$, $$\mu _2=\sqrt{A^2(\lambda _2)+B(\lambda _2)C(\lambda _2)}=\sqrt{\alpha ^2+\underset{j=1}{\overset{n}{}}\left(\frac{H_j}{\lambda _2a_j}+\frac{s_j^2}{(\lambda _2a_j)^2}\right)},$$ as the Hamiltonian function, $$\dot{L}(u)=\text{i}\{\mu _2,L(u)\}.$$ (7.40) This means that the two-point map discretizes a one-parameter family of flows. Having chosen the parameter $`\lambda _2`$ to be equal to any of the poles of the Lax matrix (parameters of the model) $`a_j`$, $`j=1,\mathrm{},n`$, the map leads to $`n`$ different maps, each discretizing the flow with the corresponding Hamiltonian $`H_j`$, $`j=1,\mathrm{},n`$. Indeed, take the limit $`\lambda _2a_j`$, $$\lambda _2=a_j+\epsilon ,\epsilon 0.$$ (7.41) Then we have $$\mu _2=\frac{s_j}{\epsilon }+\frac{H_j}{2s_j}+O(\epsilon ),$$ (7.42) and in this limit the Lax equation (7.39)–(7.40) turns into $$\dot{L}(u)=\frac{\text{i}}{2s_j}\{H_j,L(u)\}=[L(u),\frac{1}{2s_j(ua_j)}\left(\begin{array}{cc}s_j^3& s_j^{}\\ s_j^+& s_j^3\end{array}\right)].$$ (7.43) Let us denote a collection of these maps by $`\{B_{P_1}^{H_j}\}_{j=1}^n`$. The map $`B_{P_1}^{H_k}`$ discretizes the flow governed by the Hamiltonian $`H_k`$ with $`\lambda _1a_k`$ playing the role of the discrete time-step parameter. The map (and its inverse) is defined by the two-point matrix $`M(u)`$ (7.22) with the following expressions for the variables $`x`$ and $`X`$: $$x=\frac{A(\lambda _1)\mu _1}{B(\lambda _1)}=\frac{\stackrel{~}{\stackrel{~}{s}}_k^{\mathrm{\hspace{0.33em}3}}s_k}{\stackrel{~}{\stackrel{~}{s}}_k^{}},$$ (7.44) $$X=\frac{(a_k\lambda _1)B(\lambda _1)s_k^{}}{B(\lambda _1)(s_k^3+s_k)s_k^{}(A(\lambda _1)\mu _1)}=\frac{(a_k\lambda _1)\stackrel{~}{\stackrel{~}{B}}(\lambda _1)\stackrel{~}{\stackrel{~}{s}}_k^{}}{\stackrel{~}{\stackrel{~}{s}}_k^{}\left(\stackrel{~}{\stackrel{~}{A}}(\lambda _1)+\mu _1\right)\stackrel{~}{\stackrel{~}{B}}(\lambda _1)\left(\stackrel{~}{\stackrel{~}{s}}_k^{\mathrm{\hspace{0.33em}3}}s_k\right)}.$$ (7.45) All these maps are explicit Poisson maps, preserving Hamiltonians and having the spectrality property with respect to the pair of variables $`(\lambda _1,\mu _1)`$. 8. Concluding remarks One of the very important branches of the theory of finite-dimensional integrable systems is the area of discrete-time integrable systems. The interest to this area was revived in the beginning of 90’s by Veselov in the series of works (see ). He defined integrable Lagrange correspondences as discrete-time analogs of integrable continuous flows, clarified their geometric meaning as finite shifts on Jacobians and gave several important examples. Since then the subject has got a boost and has been developed further by many authors. It has not been our intention to give a review of many important recent contributions made to the area, because it would require much more space. Instead, here we only mention in brief the main features of a new recent approach to constructing integrable maps which was introduced in , developed in , which has also been used in this paper and which will be referred to as Bäcklund transformations for finite-dimensional integrable systems. One of the new features of this approach to disrete-time integrability is the spectrality property which is a projection on the classical case of the famous (quantum) Baxter equation. It was discovered on the examples of Toda lattice and elliptic Ruijsenaars system in and was generalized to the integrable case of the DST model in . Later on, it was observed that the property is universal and that, in effect, it gives a canonical way of parameterising the corresponding shift on the Jacobian which is characterized by adding a point $`(\lambda ,\mu )`$ to a divisor of points on the spectral curve $`\mathrm{\Gamma }`$ (cf. ). A direct consequence of the spectrality property is the explicitness of the constructed maps. This new point, which is an obvious advantage because explicit iterative maps are much more useful than implicit maps (given as a system of non-linear equations), was clearly demonstrated in . This new aspect of constructing explicitly given maps has been also adopted and illustrated in detail in the present paper. There were several examples of explicit maps known before, like McMillan’s map, but all those cases were exceptional, for in the generic situation, according to Veselov’s approach, integrable Lagrange/Poisson correspondences are multi-valued maps, i.e. correspondences rather than maps. Using the spectrality property as extra data allows to overcome this drawback and to construct discrete-time integrable flows as genuine maps. Another new feature of the proposed construction of integrable time-discretizations is an identification of the most elementary, one-point, basic map and construction of composite maps, like the two-point map, as compositions of the one-point map and its inverse. The choice of the matrix $`M(u)`$ (3.8) generating the one-point map is dictated by the algebraic considerations explained in . In brief, the matrix $`M(u)`$ should be a simple $`L`$-operator of the quadratic algebra, $$\{L_1(u),L_2(v)\}=[r(uv),L_1(u)L_2(v)],$$ (8.1) with the same rational $`r`$-matrix (2.7) as in the linear algebra (2.6). The number of zeros of the $`detM(u)`$ is the number of essential Bäcklund parameters, so that the matrix $`M(u)`$ in (3.8) is one-point and the matrix $`M(u)`$ in (7.22) is two-point. The fact that the right ansatz for the matrix $`M(u)`$ obeys the algebra (8.1) usually garanties that the resulted map will be Poisson, see for details. In the present paper we have observed a new ‘spectrality’ property of the basic one-point map with respect to the parameter $`\gamma `$ in $$detM(u)=\gamma (u\lambda ).$$ (8.2) We have also shown that the two-point map factorises to two one-point maps. The two-point map constructed above is probably most general map for the considered sl(2) Gaudin model, meaning that it gives a discretization of continuous flows given by any Hamiltonian $`H_j`$, $`j=1,\mathrm{},n`$, from the spectral curve, $$v^2=A^2(u)+B(u)C(u)=\alpha ^2+\underset{j=1}{\overset{n}{}}\left(\frac{H_j}{ua_j}+\frac{s_j^2}{(ua_j)^2}\right).$$ (8.3) So, at least in principle, any other integrable map for this model should be a function of the $`n`$ maps constructed in this paper. There is no established name for integrable maps with all the qualities mentioned above, namely: i) spectrality, ii) explicitness, iii) Poissonicity, iv) limits to continuous flows, v) preservation of the same integrals as for the continuous flows which these maps discretize. We are using for them the same name, Bäcklund transformations, as was used in the references . The application of the constructed maps as exact numerical integrators of the continuous flows is considered in . ## Acknowledgements VBK wishes to acknowledge the support from the EPSRC Advanced Research Fellowship AF/100072. The hospitality and support of the Universita degli Studi ‘Roma Tre’ extended to VBK during his visit to Rome in 1998 where part of this work was done is kindly acknowledged too.
warning/0007/hep-th0007191.html	ar5iv	text	# Supergravity and a Confining Gauge Theory: Duality Cascades and 𝜒SB–Resolution of Naked Singularities ## 1 Introduction A fruitful extension of the basic AdS/CFT correspondence stems from studying branes at conical singularities . Consider, for instance, a stack of D3-branes placed at the apex of a Ricci-flat 6-d cone $`Y_6`$ whose base is a 5-d Einstein manifold $`X_5`$. Comparing the metric with the D-brane description leads one to conjecture that type IIB string theory on $`AdS_5\times X_5`$ is dual to the low-energy limit of the world volume theory on the D3-branes at the singularity. A useful example of this correspondence has been to study D3-branes on the conifold . When the branes are placed at the singularity, the resulting $`𝒩=1`$ superconformal field theory has gauge group $`SU(N)\times SU(N)`$. It contains chiral superfields $`A_1,A_2`$ transforming as $`(𝐍,\overline{𝐍})`$ and superfields $`B_1,B_2`$ transforming as $`(\overline{𝐍},𝐍)`$, with superpotential $`𝒲=\lambda ϵ^{ij}ϵ^{kl}\mathrm{Tr}A_iB_kA_jB_l`$. The two gauge couplings do not flow, and indeed can be varied continuously without ruining conformal invariance. For many singular spaces $`Y_6`$ there are also fractional D3-branes which can exist only within the singularity . These fractional D3-branes are D5-branes wrapped over (collapsed) 2-cycles at the singularity. In the case of the conifold, the singularity is a point. The addition of $`M`$ fractional branes at the singular point changes the gauge group to $`SU(N+M)\times SU(N)`$; the four chiral superfields remain, now in the representation $`(𝐍+𝐌,\overline{𝐍})`$ and its conjugate, as does the superpotential . The theory is no longer conformal. Instead, the relative gauge coupling $`g_1^2g_2^2`$ runs logarithmically, as pointed out in , where the supergravity equations corresponding to this situation were solved to leading order in $`M/N`$. In this solution was completed to all orders; the conifold suffers logarithmic warping, and the relative gauge coupling runs logarithmically at all scales. The D3-brane charge, i.e. the 5-form flux, decreases logarithmically as well. However, the logarithm in the solution is not cut off at small radius; the D3-brane charge eventually becomes negative and the metric becomes singular. In it was conjectured that this solution corresponds to a flow in which the gauge group factors repeatedly drop in size by $`M`$ units, until finally the gauge groups are perhaps $`SU(2M)\times SU(M)`$ or simply $`SU(M)`$. It was further suggested that the strong dynamics of this gauge theory would resolve the naked singularity in the metric. Here, we show that this conjecture is correct. The flow is in fact an infinite series of Seiberg duality transformations — a “duality cascade” — in which the number of colors repeatedly drops by $`M`$ units. Once the number of colors in the smaller gauge group is fewer than $`M`$, non-perturbative effects become essential. We will show that these gauge theories have an exact anomaly-free $`𝐙_{2M}`$ R-symmetry, which is broken dynamically, as in pure $`𝒩=1`$ Yang-Mills theory, to $`𝐙_2`$. In the supergravity, this occurs through the deformation of the conifold.<sup>1</sup><sup>1</sup>1For a five-dimensional supergravity approach to chiral symmetry breaking, see . In short, the resolution of the naked singularity found in occurs through the chiral symmetry breaking of the gauge theory. The resulting space, a warped deformed conifold, is completely nonsingular and without a horizon, leading to confinement. If the low-energy gauge theory has fundamental matter, a horizon appears and leads to screening. ## 2 Branes and Fractional Branes on the Conifold ### 2.1 The Conifold The conifold is described by the following equation in $`𝐂^4`$: $$\underset{n=1}{\overset{4}{}}z_n^2=0.$$ (1) Equivalently, using $`z_{ij}=\frac{1}{\sqrt{2}}_n\sigma _{ij}^nz_n`$, where $`\sigma ^n`$ are the Pauli matrices for $`n=1,2,3`$ and $`\sigma ^4`$ is $`i`$ times the unit matrix, it may be written as $$\underset{i,j}{det}z_{ij}=0.$$ (2) This is a cone whose base is a coset space $`T^{11}=(SU(2)\times SU(2))/U(1)`$, with topology $`S^2\times S^3`$ and symmetry group $`SU(2)\times SU(2)\times U(1)`$. As discussed in , fractional D3 branes at the singularity $`z_n=0`$ of the conifold are simply D5-branes which are wrapped on the $`S^2`$ of $`T^{11}`$. The Einstein metric of $`T^{11}`$ may be written down explicitly : $$ds_{T^{11}}^2=\frac{1}{9}\left(d\psi +\underset{i=1}{\overset{2}{}}\mathrm{cos}\theta _id\varphi _i\right)^2+\frac{1}{6}\underset{i=1}{\overset{2}{}}\left(d\theta _i^2+\mathrm{sin}^2\theta _id\varphi _i^2\right).$$ (3) It will be useful to employ the following basis of 1-forms on the compact space : $`g^1={\displaystyle \frac{e^1e^3}{\sqrt{2}}},g^2={\displaystyle \frac{e^2e^4}{\sqrt{2}}},`$ $`g^3={\displaystyle \frac{e^1+e^3}{\sqrt{2}}},g^4={\displaystyle \frac{e^2+e^4}{\sqrt{2}}},`$ $`g^5=e^5,`$ (4) where $`e^1\mathrm{sin}\theta _1d\varphi _1,e^2d\theta _1,`$ $`e^3\mathrm{cos}\psi \mathrm{sin}\theta _2d\varphi _2\mathrm{sin}\psi d\theta _2,`$ $`e^4\mathrm{sin}\psi \mathrm{sin}\theta _2d\varphi _2+\mathrm{cos}\psi d\theta _2,`$ $`e^5d\psi +\mathrm{cos}\theta _1d\varphi _1+\mathrm{cos}\theta _2d\varphi _2.`$ (5) In terms of this basis, the Einstein metric on $`T^{11}`$ assumes the form $$ds_{T^{11}}^2=\frac{1}{9}(g^5)^2+\frac{1}{6}\underset{i=1}{\overset{4}{}}(g^i)^2,$$ (6) and the metric on the conifold is $$ds_6^2=dr^2+r^2ds_{T^{11}}^2.$$ (7) ### 2.2 The Gauge Theory If we place $`N`$ D3-branes and $`M`$ fractional D3-branes on the conifold, we obtain an $`SU(N+M)\times SU(N)`$ gauge group. The two gauge group factors have holomorphic scales $`\mathrm{\Lambda }_1`$ and $`\stackrel{~}{\mathrm{\Lambda }}_1`$. The matter consists of two chiral superfields $`A_1,A_2`$ in the $`(𝐍+𝐌,\overline{𝐍})`$ representation and two fields $`B_1,B_2`$ in the $`(\overline{𝐍+𝐌},𝐍)`$ representation. The superpotential of the model is $$W=\lambda _1\mathrm{tr}(A_iB_jA_kB_{\mathrm{}})ϵ^{ik}ϵ^j\mathrm{}.$$ (8) The model has a $`SU(2)\times SU(2)\times U(1)`$ global symmetry; the first (second) factor rotates the flavor index of the $`A_i`$ $`(B_i)`$, while the “baryon” $`U(1)`$ sends $`A_iA_ie^{i\alpha }`$, $`B_iB_ie^{i\alpha }`$.<sup>2</sup><sup>2</sup>2The question of whether this $`U(1)`$ is actually gauged is subtle. We believe that it is global, and arguments for this were given in . There are also two spurious $`U(1)`$ transformations, one an R-symmetry and one a simple axial symmetry, under which $`\lambda _1`$, $`\mathrm{\Lambda }_1`$ and $`\stackrel{~}{\mathrm{\Lambda }}_1`$ are generally not invariant. The charges of the matter and the couplings under the symmetries (excepting the $`SU(2)`$ flavor symmetries) are given in Table 1. Although $`U(1)_A`$ and $`U(1)_R`$ are anomalous, there is a discrete $`𝐙_{2M}`$ R-symmetry under which the theory is invariant. In particular, if we let $$[A_i,B_j][A_i,B_j]e^{2\pi in/4M},n=1,2,\mathrm{},2M,$$ (9) and rotate the gluinos by $`e^{2\pi in/2M}`$, then the superpotential rotates by $`e^{2\pi in/M}`$ with $`\lambda _1`$, $`\mathrm{\Lambda }_1`$ and $`\stackrel{~}{\mathrm{\Lambda }}_1`$ unchanged. \| \| $`SU(N_+)`$ \| $`SU(N)`$ \| $`SU(2)`$ \| $`SU(2)`$ \| $`U(1)_B`$ \| $`U(1)_A`$ \| $`U(1)_R`$ \| \| --- \| --- \| --- \| --- \| --- \| --- \| --- \| --- \| \| $`A_1,A_2`$ \| $`𝐍_+`$ \| $`\overline{𝐍}`$ \| $`\mathrm{𝟐}`$ \| $`\mathrm{𝟏}`$ \| $`\frac{1}{2N_+N}`$ \| $`\frac{1}{2N_+N}`$ \| $`\frac{1}{2}`$ \| \| $`B_1,B_2`$ \| $`\overline{𝐍_+}`$ \| $`𝐍`$ \| $`\mathrm{𝟏}`$ \| $`\mathrm{𝟐}`$ \| $`\frac{1}{2N_+N}`$ \| $`\frac{1}{2N_+N}`$ \| $`\frac{1}{2}`$ \| \| $`\mathrm{\Lambda }_1^{3N_+2N}`$ \| \| \| \| \| $`0`$ \| $`\frac{2}{N_+}`$ \| $`2M`$ \| \| $`\stackrel{~}{\mathrm{\Lambda }}_1^{3N2N_+}`$ \| \| \| \| \| $`0`$ \| $`\frac{2}{N}`$ \| $`2M`$ \| \| $`\lambda _1`$ \| \| \| \| \| $`0`$ \| $`\frac{2}{N_+N}`$ \| $`0`$ \| Table 1. Quantum numbers in the $`SU(N+M)\times SU(N)`$ model; we have written $`N_+N+M`$ for concision. The classical field theory is well aware that it represents branes moving on a conifold . Let us consider the case where the $`A_i`$ and $`B_k`$ have diagonal expectation values, $`A_i=\mathrm{diag}(a_i^{(1)},\mathrm{},a_i^{(N)})`$, $`B_i=\mathrm{diag}(b_i^{(1)},\mathrm{},b_i^{(N)})`$. The F-term equations for a supersymmetric vacuum $$B_1A_iB_2B_2A_iB_1=0,A_1B_kA_2A_2B_kA_1=0$$ (10) are automatically satisfied in this case, while the D-term equations require $`\|a_1^{(r)}\|^2+\|a_2^{(r)}\|^2\|b_1^{(r)}\|^2\|b_2^{(r)}\|^2=0`$. Along with the phases removed by the maximum abelian subgroup of the gauge theory, the D-terms leave only $`3N`$ independent complex variables. Define $`n_{ik}^{(r)}=a_i^{(r)}b_k^{(r)}`$; then the D-term and gauge invariance conditions are satisfied by using the $`n_{ik}^{(r)}`$ as coordinates. These $`4N`$ complex coordinates satisfy the condition, for each $`r`$, $$\underset{i,k}{det}n_{ik}^{(r)}=0.$$ (11) This is the same as equation (2). Thus, for each $`r=1,\mathrm{},N`$, the coordinates $`n_{11}^{(r)}`$, $`n_{12}^{(r)}`$, $`n_{21}^{(r)}`$, $`n_{22}^{(r)}`$, are naturally thought of as the position of a D3-brane moving on a conifold. There are various combinations of the fields and parameters which are invariant under the global symmetries. One is $$I_1\lambda _1^{3M}\frac{\stackrel{~}{\mathrm{\Lambda }}_1^{3N2(N+M)}}{\mathrm{\Lambda }_1^{3(N+M)2N}}[\mathrm{tr}(A_iB_jA_kB_{\mathrm{}}ϵ^{ik}ϵ^j\mathrm{})]^{2M}$$ (12) In addition, there are simple invariants such as $$R_1^{(1)}=\frac{\mathrm{tr}[A_iB_j]\mathrm{tr}[A_kB_{\mathrm{}}]ϵ^{ik}ϵ^j\mathrm{}}{\mathrm{tr}(A_iB_jA_kB_{\mathrm{}}ϵ^{ik}ϵ^j\mathrm{})};$$ (13) there are many other similar invariants, in each of which the same number of $`A`$ and $`B`$ fields appear in numerator and denominator but with color and flavor indices contracted differently. Finally there is a constant invariant $$J_1\lambda _1^{(N+M)+N}\mathrm{\Lambda }_1^{3(N+M)2N}\stackrel{~}{\mathrm{\Lambda }}_1^{3N2(N+M)}$$ (14) which plays the role of a dimensionless complex coupling analogous to $`\tau `$ in $`𝒩=4`$ Yang-Mills. The superpotential of the model will be renormalized and takes the general form $$W=\lambda _1\mathrm{tr}(A_iB_jA_kB_{\mathrm{}})ϵ^{ik}ϵ^j\mathrm{}F_1(I_1,J_1,R_1^{(s)})$$ (15) where $`F_1`$ is a function which we will not fully determine. ### 2.3 The conformal case: $`M=0`$ If there are no fractional D3-branes, then the $`U(1)_R`$ is anomaly-free, and the theory is superconformal (or if the couplings $`g_1,g_2,\lambda `$ are chosen completely arbitrarily, it will flow until it becomes conformal in the infrared.) There are two dimensionless global invariants $`\lambda ^2\mathrm{\Lambda }_1\stackrel{~}{\mathrm{\Lambda }}_1`$, the overall coupling $`\tau _1+\tau _2`$, and $`\stackrel{~}{\mathrm{\Lambda }}_1/\mathrm{\Lambda }_1`$, the relative coupling $`\tau _1\tau _2`$, which are built purely from the parameters and may be chosen arbitrarily. Thus there are two exactly marginal operators in the theory which preserve the continuous global symmetries. (There are other marginal operators which partially preserve these symmetries.) This was the case studied in , where it was shown the supergravity dual of this field theory is simply $`AdS_5\times T^{11}`$. In order to match the two couplings to the moduli of the type IIB theory on $`AdS_5\times T^{11}`$, one notes that the integrals over the $`S^2`$ of $`T^{11}`$ of the NS-NS and R-R 2-form potentials, $`B_2`$ and $`C_2`$, are moduli. In particular, the two gauge couplings are determined as follows : $$\frac{1}{g_1^2}+\frac{1}{g_2^2}e^\varphi ,$$ (16) $$\frac{1}{g_1^2}\frac{1}{g_2^2}e^\varphi \left[\left(_{S^2}B_2\right)1/2\right],$$ (17) where $`(_{S^2}B_2)`$ is normalized in such a way that its period is equal to $`1`$.<sup>3</sup><sup>3</sup>3These equations are crucial for relating the SUGRA background to the field theory beta functions when the theory is generalized to $`SU(N+M)\times SU(N)`$ . The matching between the moduli is one of the simplest checks of the duality. It is further possible to build a detailed correspondence between various gauge invariant operators in the $`SU(N)\times SU(N)`$ gauge theory and modes of the type IIB theory on $`AdS_5\times T^{11}`$ . In , it was noted that there exists a type IIA construction which is T-dual to $`N`$ D3-branes at the conifold. It involves two NS5-branes: one oriented in the $`(12345)`$ plane, and the other in the $`(12389)`$ plane. The coordinate $`x^6`$ is compactified on a circle of circumference $`l_6`$, and there are $`N`$ $`(1236)`$ D4-branes wrapped around the circle. If the NS5-branes were parallel, then the low-energy field theory would be the $`𝒩=2`$ supersymmetric $`SU(N)\times SU(N)`$ gauge theory with bifundamental matter (this type IIA configuration is T-dual to $`N`$ D3-branes at the $`𝒩=2`$ $`𝐙_2`$ orbifold singularity). Turning on equal and opposite masses for the two adjoint chiral superfields corresponds to rotating one of the NS5-branes. Under this relevant deformation the $`𝐙_2`$ orbifold field theory flows to the $`𝒩=1`$ supersymmetric conifold field theory . In terms of the type IIA brane construction, the two gauge couplings are determined by the positions of the NS5-branes along the $`x^6`$ circle. If one of the NS5-branes is located at $`x_6=0`$ and the other at $`x_6=a`$, then $$\frac{1}{g_1^2}=\frac{l_6a}{g_s},\frac{1}{g_2^2}=\frac{a}{g_s}.$$ (18) The couplings are equal when the NS5-branes are located diametrically opposite each other (in the type IIB language this corresponds to $`_{S^2}B_2`$ being equal to half of its period). As the NS5-branes approach each other, one of the couplings becomes strong. This simple geometrical picture will be useful for analyzing the RG flows in the following sections. ### 2.4 The RG cascade: $`M>0`$ Now let us consider the effect of adding $`M`$ fractional D3-branes, which as shown in corresponds to wrapping $`M`$ D5-branes over the $`S^2`$ of $`T^{11}`$. The D5-branes serve as sources of the magnetic RR 3-form flux through the $`S^3`$ of $`T^{11}`$. Therefore, the supergravity dual of this field theory involves $`M`$ units of the 3-form flux, in addition to $`N`$ units of the 5-form flux: $$_{S^3}F_3=M,_{T^{11}}F_5=N.$$ (19) In the SUGRA description the 3-form flux is the source of conformal symmetry breaking. Indeed, now $`B_2`$ cannot be kept constant and acquires a radial dependence : $$_{S^2}B_2Me^\varphi \mathrm{ln}(r/r_0),$$ (20) while the dilaton stays constant at least to linear order in $`M`$. Since the $`AdS_5`$ radial coordinate $`r`$ is dual to the RG scale , (17) implies a logarithmic running of $`\frac{1}{g_1^2}\frac{1}{g_2^2}`$ in the $`SU(N+M)\times SU(N)`$ gauge theory. This is in accord with the exact $`\beta `$-functions: $`{\displaystyle \frac{d}{d\mathrm{log}(\mathrm{\Lambda }/\mu )}}{\displaystyle \frac{8\pi ^2}{g_1^2}}`$ $`3(N+M)2N(1\gamma ),`$ (21) $`{\displaystyle \frac{d}{d\mathrm{log}(\mathrm{\Lambda }/\mu )}}{\displaystyle \frac{8\pi ^2}{g_2^2}}`$ $`3N2(N+M)(1\gamma ),`$ (22) where $`\gamma `$ is the anomalous dimension of operators $`\mathrm{Tr}A_iB_j`$. A priori, the conformal invariance of the field theory for $`M=0`$ requires that $`\gamma =\frac{1}{2}+O(M/N)`$. Taking the difference of the two equations in (21) we then find $$\frac{8\pi ^2}{g_1^2}\frac{8\pi ^2}{g_2^2}M\mathrm{ln}(\mathrm{\Lambda }/\mu )[3+2(1\gamma )],$$ (23) in agreement with (20) found on the SUGRA side. The constancy of the dilaton $`\varphi `$ to order $`M`$ is consistent with the field theory only if $`\gamma =\frac{1}{2}+O[(M/N)^2]`$. Fortunately, the field theory in Table 1 has an obvious symmetry $`MM,NN+M`$, which to leading order in $`M/N`$ is $`MM`$ with $`N`$ fixed. Clearly $`\gamma `$ is even under this symmetry and so cannot depend on $`M/N`$ at first order. The SUGRA analysis of was carried out to the linear order in $`M/N`$. Luckily, it is possible to construct an exact solution taking into account the back-reaction of $`H_3`$ and $`F_3`$ on other fields . In this solution $`e^\varphi =g_s`$ is exactly constant, which translates into the vanishing of the $`\beta `$-function for $`\frac{1}{g_1^2}+\frac{1}{g_2^2}`$ in the dual field theory. As in , <sup>4</sup><sup>4</sup>4We are not keeping track of the overall factor multiplying $`M`$, which is determined by the flux quantization. $$F_3=M\omega _3,B_2=3g_sM\omega _2\mathrm{ln}(r/r_0),$$ (24) $$H_3=dB_2=3g_sM\frac{1}{r}dr\omega _2,$$ (25) where $$\omega _2=\frac{1}{2}(g^1g^2+g^3g^4)=\frac{1}{2}(\mathrm{sin}\theta _1d\theta _1d\varphi _1\mathrm{sin}\theta _2d\theta _2d\varphi _2),$$ (26) $$\omega _3=\frac{1}{2}g^5(g^1g^2+g^3g^4).$$ (27) The relative factor of $`3`$ in (24), which is related to the coefficients in the metric (3), appears to be related to the factor of $`3`$ in the $`𝒩=1`$ beta function (23). This gives the correct value of beta function from a purely geometrical point of view. Both $`\omega _2`$ and $`\omega _3`$ are closed. Note also that $$g_s_6F_3=H_3,g_sF_3=_6H_3,$$ (28) where $`_6`$ is the Hodge dual with respect to the metric $`ds_6^2`$. Thus, the complex 3-form $`G_3`$ satisfies the self-duality condition $$_6G_3=iG_3,G_3=F_3+\frac{i}{g_s}H_3.$$ (29) This is consistent with $`G_3`$ being either a $`(0,3)`$ form or a $`(2,1)`$ form on the conifold. The Calabi-Yau form carries $`U(1)_R`$ charge equal to 2, while $`G_3`$ does not transform under the $`U(1)_R`$. Hence, the only consistent possibility appears to be that $`G_3`$ is a harmonic $`(2,1)`$ form.<sup>5</sup><sup>5</sup>5We are grateful to S. Gubser and E. Witten for discussions on this issue. It follows from (28) that $$g_s^2F_3^2=H_3^2,$$ (30) which implies that the dilaton is constant, $`\varphi =0`$. Since $`F_{3\mu \nu \lambda }H_3^{\mu \nu \lambda }=0`$, the RR scalar vanishes as well. The 10-d metric is $$ds_{10}^2=h^{1/2}(r)dx_ndx_n+h^{1/2}(r)(dr^2+r^2ds_{T^{11}}^2),$$ (31) where $$h(r)=b_0+4\pi \frac{g_sN+a(g_sM)^2\mathrm{ln}(r/r_0)+a(g_sM)^2/4}{r^4}$$ (32) and $`a`$ is a constant of order 1. Note that, for the ansatz (31), the solution for $`h`$ may be determined from the trace of the Einstein equation:<sup>6</sup><sup>6</sup>6 We are grateful to A. Tseytlin for explaining this to us. $$h^{3/2}_6^2hg_s^2F_3^2+H_3^2=2g_s^2F_3^2,$$ (33) where $`_6^2`$ is the Laplacian on the conifold. Since $`F_3^2M^2r^6h^{3/2}`$, the solution (32) follows directly. An important feature of this background, which is not visible to linear order in $`M`$, is that $`F_5`$ acquires a radial dependence . This is because $$F_5=dC_4+B_2F_3,$$ (34) and $`\omega _2\omega _3\mathrm{vol}(T^{11})`$. Thus, we may write $$F_5=_5+_5,_5=𝒦(r)\mathrm{vol}(\mathrm{T}^{11}),$$ (35) and $$𝒦(r)=N+ag_sM^2\mathrm{ln}(r/r_0).$$ (36) The novel phenomenon in this solution is that the 5-form flux present at the UV scale $`r=r_0`$ may completely disappear by the time we reach a scale $`r=\stackrel{~}{r}`$ where $`𝒦(\stackrel{~}{r})=0`$. This is related to the fact that the flux $`_{S^2}B_2`$ is not a periodic variable in the SUGRA solution: as this flux goes through a period, $`𝒦(r)𝒦(r)M`$ which has the effect of decreasing the 5-form flux by $`M`$ units. We will shortly relate this decrease, which we refer to for now as the “RG cascade”, to Seiberg duality. In order to eliminate the asymptotically flat region for large $`r`$ we use the well-known device of setting $`b_0=0`$ (this corresponds to choosing the special solution of sec. 5 in ). In terms of the scale $`\stackrel{~}{r}`$, we then have $$𝒦(r)=ag_sM^2\mathrm{ln}(r/\stackrel{~}{r}),h(r)=\frac{4\pi g_s}{r^4}[𝒦(r)+ag_sM^2/4]$$ (37) This solution has a naked singularity at $`r=r_s`$ where $`h(r_s)=0`$. Writing $$h(r)=\frac{L^4}{r^4}\mathrm{ln}(r/r_s),L^2g_sM,$$ (38) we then have a purely logarithmic RG cascade: $$ds^2=\frac{r^2}{L^2\sqrt{\mathrm{ln}(r/r_s)}}dx_ndx_n+\frac{L^2\sqrt{\mathrm{ln}(r/r_s)}}{r^2}dr^2+L^2\sqrt{\mathrm{ln}(r/r_s)}ds_{T^{11}}^2.$$ (39) This is essentially the metric of sec. 5 in expressed in terms of a different radial coordinate. Since $`T^{11}`$ expands slowly toward large $`r`$, the curvatures decrease there so that corrections to the SUGRA are negligible. Therefore, there is no obstacle for using this solution as $`r\mathrm{}`$ where the 5-form flux diverges. The field theory explanation of the divergence is that the RG cascade goes on forever as the scale is increased, generating bigger and bigger $`N`$ in the UV. As the theory flows to the IR, the cascade must stop, however, because negative $`N`$ is physically nonsensical. Thus, we should not be able to continue the solution (39) to $`r<\stackrel{~}{r}`$ where $`𝒦(r)`$ is negative. The radius of $`T^{11}`$ at $`r=\stackrel{~}{r}`$ is of order $`\sqrt{g_sM}`$. The gauge group at this scale is essentially $`SU(M)`$, and it is satisfying to see the appearance of $`g_sM`$, which is the ‘t Hooft coupling. As usual, if the ‘t Hooft coupling is large then the SUGRA solution has small curvatures. Nevertheless, the fact that the solution of is singular tells us that it has to be modified, at least in the IR. After understanding the RG cascade, we will study the dynamics of the corresponding field theory, and will see how this singularity is removed. ## 3 The $`𝒩=1`$ RG Cascade is a Duality Cascade We now trace the jumps in the rank of the gauge group to a well-known phenomenon in the dual $`𝒩=1`$ field theory, namely, Seiberg duality . The essential observation is that $`1/g_1^2`$ and $`1/g_2^2`$ flow in opposite directions and, according to (21), there is a scale where the $`SU(N+M)`$ coupling, $`g_1`$, diverges. To continue past this infinite coupling, we perform a $`𝒩=1`$ duality transformation on this gauge group factor. The $`SU(N+M)`$ gauge factor has $`2N`$ flavors in the fundamental representation. Under a Seiberg duality transformation, this becomes an $`SU(2N[N+M])=SU(NM)`$ gauge group with $`2N`$ flavors, which we may call $`a_i`$ and $`b_i`$, along with “meson” bilinears $`M_{ij}=A_iB_j`$. The fields $`a_i`$ and $`b_i`$ are fundamentals and antifundamentals of $`SU(N)`$, while the mesons are in the adjoint-plus-singlet of $`SU(N)`$. The superpotential after the transformation $$W=\lambda _1\mathrm{tr}M_{ij}M_k\mathrm{}ϵ^{ik}ϵ^j\mathrm{}F_1(I_1,J_1,R_1^{(s)})+\frac{1}{\mu }\mathrm{tr}M_{ij}a_ib_j,$$ (40) where $`\mu `$ is the matching scale for the duality transformation , shows the $`M_{ij}`$ are actually massive. We may integrate them out $$0=2\lambda _1M_k\mathrm{}ϵ^{ik}ϵ^j\mathrm{}F_1(I_1,J_1,R_1^{(s)})\frac{1}{\mu }\mathrm{tr}a_ib_j$$ (41) leaving a superpotential $$W=\lambda _2\mathrm{tr}a_ib_ja_kb_{\mathrm{}}ϵ^{ik}ϵ^j\mathrm{}F_2(I_2,J_2,R_2^{(s)})$$ (42) Here $`F_2`$, $`\lambda _2`$, $`I_2`$, $`J_2`$ and $`R_2`$ are defined similarly as in the original theory. Thus we obtain an $`SU(N)\times SU(NM)`$ theory which resembles closely the theory we started with. <sup>7</sup><sup>7</sup>7 The fact that the quartic superpotential is left roughly invariant by the duality transformation in theories of this type has long been considered of interest. It was first noted in , where it was used to study duality in $`SO(3)`$ gauge theories, and in , where its wider significance in Seiberg duality transformations was established. Let us study the matching more carefully. We define, for reasons which will become clear in a moment, the strong coupling scale of the $`SU(NM)`$ factor to be $`\stackrel{~}{\mathrm{\Lambda }}_2`$. The strong coupling scale of the $`SU(N)`$ factor is not the same as it was before the duality (since the number of flavors in the $`SU(N)`$ gauge group has changed) and its old scale $`\stackrel{~}{\mathrm{\Lambda }}_1`$ must be replaced with a new strong coupling scale $`\mathrm{\Lambda }_2`$. The matching conditions relating these scales are of the form $$\lambda _2\frac{1}{\mu ^2\lambda _1}$$ (43) and $$\mathrm{\Lambda }_1^{3(N+M)2N}\stackrel{~}{\mathrm{\Lambda }}_2^{3(NM)2N}\mu ^{2N}\lambda _1^M\stackrel{~}{\mathrm{\Lambda }}_1^{3N2(N+M)}\lambda _2^M\mathrm{\Lambda }_2^{3N2(NM)}.$$ (44) It is easy to check that $$I_2I_1\mathrm{and}J_21/J_1.$$ (45) (Note that the inversion of $`J`$ is a sign of electric-magnetic duality, the generalization of $`\tau 1/\tau `$.) Matching of baryon numbers in the Seiberg duality assures that $`(A)^{(N+M)}(a)^{(NM)}`$. We will not attempt to match the $`R_i^{(s)}`$. With these matchings, the dual theory has the global charges given in Table 2. Remarkably, this theory has the same form as the previous one with $`NNM`$. Thus the renormalization group flow is self-similar: the next step is that the $`SU(N)`$ gauge group now becomes strongly coupled, and under a Seiberg duality transformation the full gauge group becomes $`SU(NM)\times SU(N2M)`$, and so forth. \| \| $`SU(N)`$ \| $`SU(N_{})`$ \| $`SU(2)`$ \| $`SU(2)`$ \| $`U(1)_B`$ \| $`U(1)_A`$ \| $`U(1)_R`$ \| \| --- \| --- \| --- \| --- \| --- \| --- \| --- \| --- \| \| $`a_1,a_2`$ \| $`\overline{𝐍}`$ \| $`𝐍_{}`$ \| $`\mathrm{𝟐}`$ \| $`\mathrm{𝟏}`$ \| $`\frac{1}{2NN_{}}`$ \| $`\frac{1}{2NN_{}}`$ \| $`\frac{1}{2}`$ \| \| $`b_1,b_2`$ \| $`𝐍`$ \| $`\overline{𝐍_{}}`$ \| $`\mathrm{𝟏}`$ \| $`\mathrm{𝟐}`$ \| $`\frac{1}{2NN_{}}`$ \| $`\frac{1}{2NN_{}}`$ \| $`\frac{1}{2}`$ \| \| $`\mathrm{\Lambda }_2^{3N2N_{}}`$ \| \| \| \| \| $`0`$ \| $`\frac{2}{N}`$ \| $`2M`$ \| \| $`\stackrel{~}{\mathrm{\Lambda }}_2^{3N_{}2N}`$ \| \| \| \| \| $`0`$ \| $`\frac{2}{N_{}}`$ \| $`2M`$ \| \| $`\lambda _2`$ \| \| \| \| \| $`0`$ \| $`\frac{2}{NN_{}}`$ \| $`0`$ \| Table 2. Quantum numbers of the dual $`SU(N)\times SU(NM)`$ theory; we have written $`N_{}=NM`$ for concision. This flow will stop, of course, at or before the point where $`NkM`$ becomes zero or negative. Note that the Seiberg duality transformation is the same in both the so-called conformal window ($`3N_c>N_f>\frac{3}{2}N_c`$) and in the free magnetic phase ($`\frac{3}{2}N_f>N_c+1`$). Even for $`N_f=N_c+1`$ the effect on the superpotential described in is not essential, since it is accounted for in the function $`F`$. The first significant changes occur when $`N_f=N_c`$, since for $`N_fN_c`$ the classical moduli space is drastically modified. Thus, the RG flow just described proceeds step by step until the gauge group has the form $`SU(M+p)\times SU(p)`$, where $`0<pM`$. At this point we should do a more careful analysis, which we will carry out in the next section. It is instructive to consider the type IIA brane picture of the duality cascade of $`SU(N+M)\times SU(N)`$ theories.<sup>8</sup><sup>8</sup>8We should remind the reader that the classical IIA picture of Seiberg duality, introduced by , has the feature that it is not a string duality but a motion which transforms one theory into its dual, one which a priori need not leave the infrared physics invariant. In particular, there is no direct relation between the classical brane motion, or semiclassical brane bending, and the actual dynamics of the field theories. To see that the IIA story gives the right answer, one should use its M theory generalization , but even there, the M5-brane, about which only holomorphic information is available, does not generally match the dynamics of the field theory, which is not holomorphic. These issues have been studied and explained in detail; see especially . We mention this only to warn our readers not to take this paragraph for more than a heuristic argument. To implement such theories we have to add $`M`$ D4-branes stretched only one way between the D4-branes, rather than all around the circle . Such D4-branes are T-duals of the “fractional” branes which are the D5-branes wrapped over the 2-cycle of $`T^{11}`$. These new D4-branes violate the balance of forces on the NS5-branes, and the latter undergo logarithmic bending , which is the RG flow in this picture. Although the NS5 and NS5’ brane bend along the circular $`x^6`$ direction, this does not force them to intersect, because they are oriented in perpendicular directions. However, their $`x^6`$ positions become equal somewhere away from the fractional branes, and it seems natural to interpret this as a divergence of the $`SU(N+M)`$ coupling. Under such circumstances it is natural to move the NS branes so as to eliminate this divergence, moving one of them once around the $`x^6`$ circle. When the NS and NS’ brane cross during this motion, the $`N+M`$ fractional D4-branes shrink to zero size and then re-grow. In doing so, they flip their orientation and become anti-D4-branes. Meanwhile, the other $`N`$ fractional branes stretch more than once around the circle, but where they are doubled they are partially cancelled by the $`N+M`$ anti-D4-branes. This leaves $`NM`$ D4-branes in one segment and $`N`$ in the other — exactly our starting point but with $`NNM`$. After the crossing, the NS5-branes are still bent in the same directions as before, so again their $`x^6`$ positions become equal and we are led to repeat the motion around the circle. Finally, the number $`N`$ becomes of order $`M`$ and something more drastic should happen . For this physics, the analysis of becomes essential. ## 4 Chiral Symmetry Breaking and the Deformation of the Conifold The solution of is well-behaved for large $`r`$ but becomes singular at sufficiently small $`r`$. The solution must be modified in such a way that this singularity is removed. In this section we argue that the conifold (2) should be replaced by the deformed conifold $$\underset{i=1}{\overset{4}{}}z_i^2=2\underset{i,j}{det}z_{ij}=ϵ^2,$$ (46) in which the singularity of the conifold is removed through the blowing-up of the $`S^3`$ of $`T^{11}`$. There are a number of arguments in favor of this idea. One suggestive observation is that in the solution of , the source of the singularity can be traced to the infinite energy in the $`F_3`$ field. At all radii there are $`M`$ units of flux of $`F_3`$ piercing the $`S^3`$ of $`T^{11}`$, and when the $`S^3`$ shrinks to zero size this causes $`F_3^2`$ to diverge. If instead the $`S^3`$ remained of finite size, as occurs in the deformed conifold, this problem would be evaded. However, the most powerful argument that the conifold is deformed comes from the field theory analysis, which shows clearly that the spacetime geometry is modified by the strong dynamics of the infrared field theory. We will see that the theory has a deformed moduli space, with $`M`$ independent branches, each of which has the shape of a deformed conifold. The branches are permuted by the $`𝐙_{2M}`$ R-symmetry, which is spontaneously broken down to $`𝐙_2`$. This breaking of the R-symmetry is exactly what we would expect in a pure $`SU(M)`$ $`𝒩=1`$ Yang-Mills theory, although here it proceeds through scalar as well as gluino expectation values. The theory will also have domain walls, confinement, magnetic screening, and other related phenomena. The complete analysis of the nonperturbative dynamics of the field theory in Table 1 is mathematically intensive, and we have not attempted it. In this section we present a simplified version of the analysis which captures the physics which we are interested in. In an appendix we present more general (although still partial) results that show our conclusions are robust. In particular, our goal is to discover what happens in the far infrared of the flow, where the D3 brane charge has cascaded (nearly) to zero and only the $`M`$ fractional D3 branes remain. If there are no D3 branes left, we expect we have pure $`𝒩=1`$ Yang-Mills in the far infrared, a theory which breaks its $`𝐙_{2M}`$ R-symmetry to $`𝐙_2`$ and has $`M`$ isolated vacua, domain walls, and confinement. However, while this may be correct, we have no access to the supergravity background through this analysis. What we need is a probe which can see if and how the fractional D3-branes have modified the conifold itself. The right choice, it turns out, is to probe the space with a single additional D3 brane. In this case the gauge group is $`SU(M+1)\times `$$`SU(1)`$” — in short, simply $`SU(M+1)`$ — with fields $`C_i`$ and $`D_j`$ in the $`𝐌+\mathrm{𝟏}`$ and $`\overline{𝐌+\mathrm{𝟏}}`$ representations, $`i,j=1,2`$, and with superpotential $`W=\lambda C_iD_jC_kD_lϵ^{ik}ϵ^{jl}`$. Define $`N_{ij}=C_iD_j`$, which is gauge invariant. As in the discussion surrounding equation (11), the expectation values of $`N_{ij}`$ specify the position of the probe brane; in the classical theory, we have $`det_{i,j}N_{ij}=0`$, indicating the probe is moving on the original, singular conifold. At low energy the theory can be written in terms of these invariants and develops the nonperturbative superpotential first written down by Affleck, Dine and Seiberg $$W_L=\lambda N_{ij}N_k\mathrm{}ϵ^{ik}ϵ^j\mathrm{}+(M1)\left[\frac{2\mathrm{\Lambda }^{3M+1}}{N_{ij}N_k\mathrm{}ϵ^{ik}ϵ^j\mathrm{}}\right]^{\frac{1}{M1}}.$$ (47) The equations for a supersymmetric vacuum are $$0=\left(\lambda \left[\frac{2\mathrm{\Lambda }^{3M+1}}{(N_{ij}N_k\mathrm{}ϵ^{ik}ϵ^j\mathrm{})^M}\right]^{\frac{1}{M1}}\right)N_{ij}.$$ (48) The apparent solution $`N_{ij}=0`$ for all $`i,j`$ actually gives infinity on the right-hand side. The only solutions are then $$(N_{ij}N_k\mathrm{}ϵ^{ik}ϵ^j\mathrm{})^M=\frac{2\mathrm{\Lambda }^{3M+1}}{\lambda ^{M1}}.$$ (49) As predicted, this equation has $`M`$ independent branches, in each of which $`N_{ij}N_k\mathrm{}ϵ^{ik}ϵ^j\mathrm{}`$ is a $`M^{th}`$ root of $`\mathrm{\Lambda }^{3M+1}/\lambda ^{M1}`$. The $`𝐙_{2M}`$ discrete non-anomalous R-symmetry rotates $`N_{ij}N_k\mathrm{}ϵ^{ik}ϵ^j\mathrm{}`$ by a phase $`e^{2\pi i/M}`$, and thus the $`M`$ branches transform into one another under the symmetry. In short, the $`𝐙_{2M}`$ is spontaneously broken down to $`𝐙_2`$. The low-energy effective superpotential is $$W=M\lambda N_{ij}N_k\mathrm{}ϵ^{ik}ϵ^j\mathrm{}M\left[2\lambda \mathrm{\Lambda }^{3M+1}\right]^{1/M}$$ (50) which reflects the $`M`$ branches. Most importantly, on each of these branches the classical condition on the $`N_{ij}`$ has been modified to read $$\underset{i,j}{det}N_{ij}=\frac{1}{2}N_{ij}N_k\mathrm{}ϵ^{ik}ϵ^j\mathrm{}=\left(\frac{\mathrm{\Lambda }^{3M+1}}{[2\lambda ]^{M1}}\right)^{1/M}$$ (51) Comparing with equation (46) we see that the probe brane in the quantum theory moves on the deformed conifold; the classical singularity at the origin of the moduli space has been resolved through chiral symmetry breaking. The above constraint on the expectation values for $`N_{ij}`$ implies that in the perturbative region (where semiclassical analysis is valid) they can break the gauge group only down to $`SU(M)`$, with no massless charged matter. This gauge theory is thus in the universality class of pure $`SU(M)`$ Yang-Mills, and will share many of its qualitative properties. However, the existence of massive matter $`C_i,D_j`$ in the fundamental representation of $`SU(M)`$ (note that if $`N_{11}`$ is large then $`C_2,D_2`$ have mass $`\lambda N_{11}`$) implies that confinement occurs only in an intermediate range of distances. As in QCD with heavy quarks, pair production of the massive quarks breaks the confining flux tubes, so a linear potential between external sources exists only between the length scales $`1/\sqrt{T}`$ and $`m_q/T`$, where $`T`$ is the string tension and $`m_q`$ is the dynamical quark mass. For $`N_{11}N_{22}\left(\mathrm{\Lambda }^{3M+1}/\lambda ^{M1}\right)^{1/2M}`$, their minimal values, we expect little sign of a linear potential at any length scale, as in physical QCD. Only for $`p=0`$ do we expect confinement at all scales. More generally, for $`1<p<M`$, one obtains a moduli space corresponding to $`p`$ probe branes moving on the deformed conifold. If $`pM`$, both the $`SU(p)`$ gauge coupling and its ’t Hooft coupling are small at the strong-dynamics scale of $`SU(M+p)`$. Furthermore, the $`SU(p)`$ factor has vanishing beta function in the far IR, where it has three adjoint chiral superfields (namely, three of the $`N_{ij}`$) and is essentially a copy of $`𝒩=4`$ Yang-Mills. Consequently, we expect no strong dynamics from the $`SU(p)`$ sector, and the theory is very close to $`SU(M+p)`$ with $`2p`$ light flavors. In this case a similar analysis to the above is essentially correct. At large expectation values, the gauge theory is broken to $`SU(M)`$ $`𝒩=1`$ Yang-Mills times $`SU(p)`$ $`𝒩=4`$ Yang-Mills, with massive states in the bifundamental representation of the group factors. Details of this analysis are given in the appendix. As before, pair production of these massive states eliminates confinement at large distances; electric sources are screened by massive states which leave them charged only under the nonconfining group $`SU(p)`$. The pattern of chiral symmetry breaking gives us another qualitative argument why the conifold must be deformed. The original conifold has a $`U(1)_R`$ symmetry under which the $`z_{ij}`$ in (2) rotate by a phase. In Table 1 we saw this was broken by instantons to $`𝐙_{2M}`$, but for large $`M`$ this is a $`1/M`$ effect and need not show up in the leading order supergravity. However, if we expect the infrared theory to behave similarly to pure $`𝒩=1`$ Yang-Mills, then we expect this symmetry to be spontaneously broken to $`𝐙_2`$. This breaking is a leading-order effect and most definitely should be visible in the supergravity. The only natural modifications of the conifold are its resolution and its deformation; only the latter breaks the classical $`U(1)_R`$ symmetry, and it indeed breaks it to $`𝐙_2`$, as is obvious from equation (46). As a final argument, we consider expectations from the IIA/M brane construction. Classically we have NS and NS’ branes filling four-dimensional space and extending in the $`v=x^4+ix^5`$ and $`w=x^8+ix^9`$ directions respectively. They are separated along the compact direction $`x^6`$ by a distance $`a`$, which along with $`l_6`$ sets the two classical gauge couplings, as explained in (18). In one $`x^6`$ segment between the NS and NS’ brane we suspend $`M+1`$ D4 branes; in the other there is only one D4 brane. A single complete wrapped D4-brane — our probe — is free to move anywhere in the $`v,w,x^7`$ space, independently of the other branes, while the other $`M`$ suspended D4-branes are pinned to $`v=w=0`$. To understand the quantum theory, we must move to M theory , where we combine $`x^6`$ with the new compact coordinate $`x^{10}`$ using $`t=e^{x^6+ix^{10}}`$. Classically the equations for the NS and NS’ brane are $`w=0,t=1`$ and $`v=0,t=e^a`$. The M theory expectation is that, in the quantum theory, the probe brane will become an independent M5 brane wrapped on the $`t`$ directions, while the suspended D4-branes join with the NS and NS’ branes to make a single M5-brane, which we will refer to as our MQCD brane. This type of behavior was first seen in $`𝒩=2`$ and $`𝒩=1`$ supersymmetric Yang-Mills . Indeed the MQCD brane which appears in our case should be very similar to that of $`𝒩=1`$ super-Yang-Mills, since in the limit the $`x^6`$ direction becomes large they should become equal. The brane for super-Yang-Mills fills the coordinates $`x^0,x^1,x^2,x^3`$ and is embedded in the coordinates $`v=x^4+ix^5,w=x^8+ix^9,t=e^{x^6+ix^{10}}`$ as a Riemann surface defined through the equations $$(vw)^M=\mathrm{\Lambda }_L^{2M},v^M=t.$$ (52) Notice classically the equations include $`vw=0`$, corresponding to the presence of the NS and NS’ brane. However, the quantum Yang-Mills M-brane has $`vw`$ equal to a nonzero constant, and has $`M`$ possible orientations, one for each possible phase of a condensate. What is the connection with the deformed conifold? As shown in , a type IIA NS-brane and NS’-brane satisfying the equation $`vw=0`$, that is, intersecting at a point, are T-dual to the conifold. This lifts without change to M theory. We saw this equation appears in the construction of classical Yang-Mills, and it will appear in our classical theory as well. Meanwhile, if the NS and NS’ branes are at the same $`t`$, that is, if they intersect, then they can be deformed into a single object with equation $`vw=`$ constant $`0`$. This object is T-dual to the deformed conifold. Again this also lifts without change to M theory. Now notice that the Yang-Mills M-brane has this as one of its defining equations (52). This shows the NS and NS’ brane have been glued together into a single object. Without the suspended D4-branes, this could only occur if the joined NS and NS’-brane had equal $`t`$ coordinates, but in the presence of the suspended D4-branes, which extend along the $`t`$ direction, the NS and NS’-branes can be separated in $`t`$, as in (52). Thus the Yang-Mills M-brane shows that the suspended D4-branes allow a quantum effect in M theory by which the conifold can be deformed even when the two gauge couplings (18) are both finite. In our case, we similarly expect the two Riemann surfaces — the probe and the MQCD brane — to have $`M`$ branches, with a continuous variable specifying the position of the probe brane in the space, and a discrete variable labeling the orientation of the MQCD brane. However, when the probe is far away and the $`x^6`$ direction is large, our MQCD brane should closely resemble that of Yang-Mills. We therefore expect the equations governing it to have the same qualitative form. In particular, we expect that the $`𝐙_{2M}`$ discrete symmetry rotating the phase of $`t`$ by $`2\pi `$ is broken to $`𝐙_2`$, through the modification of the equation $`vw=0`$ to $`(vw)^M=`$ constant. By T-duality this indicates that the classical conifold is quantum deformed by the fractional branes. ## 5 Back to Supergravity: The Deformed Conifold Ansatz The field theory analysis of the previous section shows that the naive $`U(1)`$ (really $`𝐙_{2M}`$) R-symmetry is actually broken to a $`𝐙_2`$. On the other hand, the SUGRA background (31) has an exact $`U(1)`$ symmetry realized as shifts of the angular coordinate $`\psi `$ on $`T^{11}`$. The presence of this unwanted symmetry in the IR may also be the reason for the appearance of the naked singularity. In this section we propose that the solution of this problem is to replace the conifold by its deformation (46) in the ansatz (31). This indeed breaks the $`U(1)`$ symmetry $`z_ke^{i\alpha }z_k`$, $`k=1,\mathrm{},4`$, down to its $`𝐙_2`$ subgroup $`z_kz_k`$. Another reason to focus on the deformed conifold is that it gives the correct moduli space for the field theory, as shown in the previous section. The metric of the deformed conifold was discussed in some detail in . It is diagonal in the basis (2.1): $`ds_6^2={\displaystyle \frac{1}{2}}ϵ^{4/3}K(\tau )[{\displaystyle \frac{1}{3K^3(\tau )}}(d\tau ^2+(g^5)^2)+\mathrm{cosh}^2\left({\displaystyle \frac{\tau }{2}}\right)[(g^3)^2+(g^4)^2]`$ $`+\mathrm{sinh}^2\left({\displaystyle \frac{\tau }{2}}\right)[(g^1)^2+(g^2)^2]],`$ (53) where $$K(\tau )=\frac{(\mathrm{sinh}(2\tau )2\tau )^{1/3}}{2^{1/3}\mathrm{sinh}\tau }.$$ (54) For large $`\tau `$ we may introduce another radial coordinate $`r`$ via $$r^3ϵ^2e^\tau ,$$ (55) and in terms of this radial coordinate $$ds_6^2dr^2+r^2ds_{T^{11}}^2.$$ (56) The determinant of the metric (5) is $$g_6ϵ^8\mathrm{sinh}^4\tau ,$$ (57) which vanishes at $`\tau =0`$. Indeed, at $`\tau =0`$ the angular metric degenerates into $$d\mathrm{\Omega }_3^2=\frac{1}{2}ϵ^{4/3}(2/3)^{1/3}[\frac{1}{2}(g^5)^2+(g^3)^2+(g^4)^2],$$ (58) which is the metric of a round $`S^3`$ . The additional two directions, corresponding to the $`S^2`$ fibered over the $`S^3`$, shrink as $$\frac{1}{8}ϵ^{4/3}(2/3)^{1/3}\tau ^2[(g^1)^2+(g^2)^2].$$ (59) In what follows we will set $`ϵ=12^{1/4}`$, so that $`\frac{1}{2}ϵ^{4/3}(2/3)^{1/3}=1`$. The collapse of the $`S^2`$ implies that at $`\tau =0`$ $`F_3`$ must lie within the remaining $`S^3`$, $$F_3(\tau =0)=Mg^5g^3g^4,$$ (60) which may be shown to be a closed 3-form. On the other hand, for large $`\tau `$, $`F_3`$ should approach its value $$\frac{M}{2}g^5(g^1g^2+g^3g^4)$$ (61) found in the UV ansatz (24). These two closed 3-forms differ by an exact one, $$g^5(g^1g^2g^3g^4)=d(g^1g^3+g^2g^4)$$ (62) Therefore, the simplest ansatz which interpolates smoothly between $`\tau =0`$ and large $`\tau `$ is $`F_3=M\left\{g^5g^3g^4+d[F(\tau )(g^1g^3+g^2g^4)]\right\}`$ $`=M\{g^5g^3g^4(1F))+g^5g^1g^2F+F^{}d\tau (g^1g^3+g^2g^4)\},`$ (63) with $`F(0)=0`$ and $`F(\mathrm{})=1/2`$. Note also that this ansatz preserves the $`𝐙_2`$ symmetry which interchanges $`(\theta _1,\varphi _1)`$ with $`(\theta _2,\varphi _2)`$. A similarly $`𝐙_2`$-symmetric ansatz for $`B_2`$ is $$B_2=g_sM[f(\tau )g^1g^2+k(\tau )g^3g^4].$$ (64) Using the identity $$g^5(g^1g^3+g^2g^4)=d(g^1g^2g^3g^4),$$ (65) we find that $$H_3=dB_2=g_sM[d\tau (f^{}g^1g^2+k^{}g^3g^4)+\frac{1}{2}(kf)g^5(g^1g^3+g^2g^4)].$$ (66) We further have $$_5=B_2F_3=g_sM^2\mathrm{}(\tau )g^1g^2g^3g^4g^5,$$ (67) where $$\mathrm{}=f(1F)+kF.$$ (68) The most general radial ansatz for the 10-d metric, consistent with the symmetries of the deformed conifold, is $`ds_{10}^2=A^2(\tau )dx_ndx_n+B^2(\tau )(d\tau ^2)+C^2(\tau )(g^5)^2+D^2(\tau )[(g^3)^2+(g^4)^2]`$ $`+E^2(\tau )[(g^1)^2+(g^2)^2].`$ (69) The reason we are allowed to assume that $`A,\mathrm{},E`$ depend only on $`\tau `$ is that before the introduction of the 3-form fields, the metric has the form (5), and our ansatz for $`F_3`$ and $`H_3`$ does not break this symmetry. The flux of $`F_3`$ is distributed uniformly over the $`S^3`$ near the apex of the deformed conifold; therefore, the $`M`$ D5 branes wrapped over the $`S^2`$ may be thought of as smeared over the $`S^3`$. It is not hard to check that $`F_{3\mu \nu \lambda }H_3^{\mu \nu \lambda }=0`$, which implies that the RR scalar vanishes. It is not a priori clear whether the dilaton is constant for the deformed solution, but in what follows we will assume that such a background does exist, i.e. that $$g_s^2F_3^2=H_3^2.$$ (70) Furthermore, guided by the simple form of the solution constructed in and reviewed in section 2, we will assume that the 10-d metric takes the following form: $$ds_{10}^2=h^{1/2}(\tau )dx_ndx_n+h^{1/2}(\tau )ds_6^2,$$ (71) where $`ds_6^2`$ is the metric of the deformed conifold (5). This is the same type of “D-brane” ansatz as (31), but with the conifold replaced by the deformed conifold as the transverse space. This form will also permit additional D3-brane probes to be directly included in the ansatz. The type IIB equations satisfied by the 3-form fields are $$d(e^\varphi F_3)=F_5H_3,d(e^\varphi H_3)=g_s^2F_5F_3.$$ (72) First, let us calculate $$_5g_sM^2dx^0dx^1dx^2dx^3d\tau \frac{\mathrm{}(\tau )}{K^2h^2\mathrm{sinh}^2(\tau )}.$$ (73) To write down the first equation we need $`F_3=Mh^1dx^0dx^1dx^2dx^3[(1F)\mathrm{tanh}^2\left({\displaystyle \frac{\tau }{2}}\right)d\tau g^1g^2`$ $`+F\mathrm{coth}^2\left({\displaystyle \frac{\tau }{2}}\right)d\tau g^3g^4+F^{}g^5(g^1g^3+`$ $`g^2g^4)].`$ (74) Assuming a constant $`\varphi `$ and using (65) we find $$(1F)\mathrm{tanh}^2(\tau /2)F\mathrm{coth}^2(\tau /2)+2h\frac{d}{d\tau }(h^1F^{})=\alpha (kf)\frac{\mathrm{}}{K^2h\mathrm{sinh}^2\tau },$$ (75) where $`\alpha `$ is a normalization factor proportional to $`(g_sM)^2`$. Let us turn to the second of the equations (72). Since $`H_3=g_sMh^1dx^0dx^1dx^2dx^3[g^5(k^{}\mathrm{tanh}^2\left({\displaystyle \frac{\tau }{2}}\right)g^1g^2`$ $`+f^{}\mathrm{coth}^2\left({\displaystyle \frac{\tau }{2}}\right)g^3g^4){\displaystyle \frac{1}{2}}(fk)d\tau (g^1g^3+g^2g^4)]`$ $`,`$ (76) we find $`h{\displaystyle \frac{d}{d\tau }}(h^1\mathrm{coth}^2(\tau /2)f^{}){\displaystyle \frac{1}{2}}(fk)=\alpha {\displaystyle \frac{\mathrm{}(1F)}{K^2h\mathrm{sinh}^2\tau }}`$ $`h{\displaystyle \frac{d}{d\tau }}(h^1\mathrm{tanh}^2(\tau /2)k^{})+{\displaystyle \frac{1}{2}}(fk)=\alpha {\displaystyle \frac{\mathrm{}F}{K^2h\mathrm{sinh}^2\tau }}.`$ (77) where $`\alpha `$ is the same normalization factor as in (75). We have been assuming that the dilaton is constant. The equation that guarantees this is (70). Writing it out with our ansatz gives $`{\displaystyle \frac{(k^{})^2}{\mathrm{cosh}^4(\tau /2)}}+{\displaystyle \frac{(f^{})^2}{\mathrm{sinh}^4(\tau /2)}}+{\displaystyle \frac{2(fk)^2}{\mathrm{sinh}^2\tau }}`$ $`={\displaystyle \frac{(1F)^2}{\mathrm{cosh}^4(\tau /2)}}+{\displaystyle \frac{F^2}{\mathrm{sinh}^4(\tau /2)}}+{\displaystyle \frac{8(F^{})^2}{\mathrm{sinh}^2\tau }}.`$ (78) In order to complete the system of equations we need the Einstein equations for the metric. In view of the simplified ansatz (71) for the metric it is sufficient to use the trace of the Einstein equation: $$h^{3/2}_6^2hg_s^2F_3^2+H_3^2=2g_s^2F_3^2,$$ (79) where now $`_6^2`$ is the Laplacian on the deformed conifold. Using (5), we find that the explicit form of this equation is $$\frac{1}{\mathrm{sinh}^2\tau }\frac{d}{d\tau }(h^{}K^2(\tau )\mathrm{sinh}^2\tau )=\frac{\alpha }{4}\left[\frac{(1F)^2}{\mathrm{cosh}^4(\tau /2)}+\frac{F^2}{\mathrm{sinh}^4(\tau /2)}+\frac{8(F^{})^2}{\mathrm{sinh}^2\tau }\right].$$ (80) ### 5.1 The First-Order Equations and Their Solution In searching for BPS saturated supergravity backgrounds, the second order equations should be replaced by a system of first-order ones (see, for instance, ). Luckily, this is possible for our ansatz. We have been able to find a system of simple first-order equations, from which (75), (5), (5) and (79) follow: $`f^{}`$ $`=`$ $`(1F)\mathrm{tanh}^2(\tau /2),`$ $`k^{}`$ $`=`$ $`F\mathrm{coth}^2(\tau /2),`$ $`F^{}`$ $`=`$ $`{\displaystyle \frac{1}{2}}(kf),`$ (81) and $$h^{}=\alpha \frac{f(1F)+kF}{K^2(\tau )\mathrm{sinh}^2\tau }.$$ (82) Note that the first three of these equations, (5.1), form a closed system and need to be solved first. In fact, these equations imply the self-duality of the complex 3-form with respect to the metric of the deformed conifold: $`_6G_3=iG_3`$.<sup>9</sup><sup>9</sup>9 We believe, according to a discussion in section 2.4, that $`G_3`$ is a harmonic $`(2,1)`$ form on the deformed conifold. Inspection of these equations shows that the small $`\tau `$ behavior is<sup>10</sup><sup>10</sup>10It is also possible to shift $`f`$ and $`k`$ by the same constant. The effect of this shift will be considered in section 5.2. $$f\tau ^3,k\tau ,F\tau ^2.$$ (83) On the other hand, for large $`\tau `$ the 3-forms have to match onto the conifold solution , $$f\frac{\tau }{2},k\frac{\tau }{2},F\frac{1}{2}.$$ (84) Remarkably, it is possible to find the solution with these boundary conditions in closed form. Combining (5.1) we find the following second-order equation for $`F`$: $$F^{\prime \prime }=\frac{1}{2}[F\mathrm{coth}^2(\tau /2)+(F1)\mathrm{tanh}^2(\tau /2)].$$ (85) The solution is $$F(\tau )=\frac{\mathrm{sinh}\tau \tau }{2\mathrm{sinh}\tau },$$ (86) from which we obtain $`f(\tau )`$ $`=`$ $`{\displaystyle \frac{\tau \mathrm{coth}\tau 1}{2\mathrm{sinh}\tau }}(\mathrm{cosh}\tau 1),`$ $`k(\tau )`$ $`=`$ $`{\displaystyle \frac{\tau \mathrm{coth}\tau 1}{2\mathrm{sinh}\tau }}(\mathrm{cosh}\tau +1).`$ (87) Now that we have solved for the 3-forms on the deformed conifold, the warp factor may be determined by integrating (82). First we note that $$\mathrm{}(\tau )=f(1F)+kF=\frac{\tau \mathrm{coth}\tau 1}{4\mathrm{sinh}^2\tau }(\mathrm{sinh}2\tau 2\tau ).$$ (88) This behaves as $`\tau ^3`$ for small $`\tau `$. It follows that, for small $`\tau `$, $$h=a_0+a_1\tau ^2+\mathrm{}.$$ (89) For large $`\tau `$ we impose, as usual, the boundary condition that $`h`$ vanishes. The resulting integral expression for $`h`$ is $$h(\tau )=\alpha \frac{2^{2/3}}{4}_\tau ^{\mathrm{}}𝑑x\frac{x\mathrm{coth}x1}{\mathrm{sinh}^2x}(\mathrm{sinh}(2x)2x)^{1/3}.$$ (90) We have not succeeded in evaluating this in terms of elementary or well-known special functions. For our purposes it is enough to show that $$h(\tau 0)a_0;h(\tau \mathrm{})\frac{3}{4}2^{1/3}\alpha \tau e^{4\tau /3}.$$ (91) This is nonsingular at the tip of the deformed conifold and, from (55), matches the form of the large-$`\tau `$ solution (38). The small $`\tau `$ behavior follows from the convergence of the integral (90), while at large $`\tau `$ the integrand becomes $`xe^{4x/3}`$. Thus, for small $`\tau `$ the ten-dimensional geometry is approximately $`R^{3,1}`$ times the deformed conifold: $$ds_{10}^2a_0^{1/2}dx_ndx_n+a_0^{1/2}\left(\frac{1}{2}d\tau ^2+d\mathrm{\Omega }_3^2+\frac{1}{4}\tau ^2[(g^1)^2+(g^2)^2]\right).$$ (92) Very importantly, for large $`g_sM`$ the curvatures found in our solution are small everywhere. This is true even far in the IR. Indeed, since the integral (90) converges, $$a_0\alpha (g_sM)^2.$$ (93) Therefore, the radius-squared of the $`S^3`$ at $`\tau =0`$ is of order $`g_sM`$, which is the ‘t Hooft coupling of the gauge theory found far in the IR. As long as this is large, the curvatures are small and the SUGRA approximation is reliable. We have now seen that the deformation of the conifold allows the solution to be non-singular. Qualitatively, this is because the conserved $`F_3`$ flux prevents the 3-cycle from collapsing. This is why we expect to find a metric with a collapsing 2-cycle but finite 3-cycle, and these are the properties of the deformed conifold. It may be of further interest to consider more general metrics of the form (5), and to allow the dilaton to vary. In that event it still seems likely that the qualitative properties of the solution near the apex will not change. ### 5.2 Correspondence with the Gauge Theory In this section we point out some interesting features of the SUGRA background we have found and show how they realize the expected phenomena in the dual field theory. In particular, we will now demonstrate that there is confinement and magnetic screening, and argue that there are domain walls and baryon vertices with a definite mass scale. In many ways our results resemble those found in the $`𝒩=1^{}`$ theory , but the specific details are quite different; the confining vacua of $`𝒩=1^{}`$ involve a spacetime with a spherical 5-brane sitting in it, while our present spacetime is purely given by supergravity. First we should ask the question: how does the dimensional transmutation manifest itself in supergravity? The answer is related to the presence of parameter $`ϵ`$ in the deformed conifold metric (5). Reinstating this parameter is accomplished through $$ds_6^2ϵ^{4/3}ds_6^2.$$ (94) We are then free to redefine $`hhϵ^{8/3}`$ to remove the $`ϵ`$ dependence from the transverse part of the metric. Very importantly, the dependence then appears in the longitudinal part, and the metric assumes the form $$ds_{10}^2=h^{1/2}(\tau )m^2dx_ndx_n+h^{1/2}(\tau )ds_6^2,$$ (95) so that $`mϵ^{2/3}`$ sets the 4-d mass scale. This scale then appears in all 4-d dimensionful quantities. Now let us see the theory has confining flux tubes. The key point is that in the metric for small $`\tau `$ (92) the function multiplying $`dx_ndx_n`$ approaches a constant. This should be contrasted with the $`AdS_5`$ metric where this function vanishes at the horizon, or with the singular metric of where it blows up. Consider a Wilson contour positioned at fixed $`\tau `$, and calculate the expectation value of the Wilson loop using the prescription . The minimal area surface bounded by the contour bends towards smaller $`\tau `$. If the contour has a very large area $`A`$, then most of the minimal surface will drift down into the region near $`\tau =0`$. From the fact that the coefficient of $`dx_ndx_n`$ is finite at $`\tau =0`$, we find that a fundamental string with this surface will have a finite tension, and so the resulting Wilson loop satisfies the area law. Since for large $`g_sM`$ the SUGRA description is reliable for all $`\tau `$, we seem to have found a “pure supergravity proof” of confinement in $`𝒩=1`$ gauge theory. A similar result was found previously in but involved a spacetime containing an NS5-brane with D3-brane charge. A simple estimate shows that the string tension scales as $$T_s\frac{m^2}{g_sM}.$$ (96) To see that magnetic charge is screened, we must identify the correct massive magnetically-charged source. The correct choice is a fractional D1-brane, that is, a D3-brane wrapped on the $`S^2`$ of $`T^{11}`$, attached to the boundary of the space at $`\tau =\mathrm{}`$. On the six-dimensional deformed conifold the $`S^2`$ is fibered over $`\tau `$ such that the resulting three-dimensional bundle has only one boundary, at $`\tau =\mathrm{}`$; near $`\tau =0`$ the $`S^2`$ shrinks to zero size and the bundle locally has topology $`R^3`$. Therefore, a D3-brane with a single boundary can be wrapped on this bundle, corresponding to a fractional D1-brane attached at $`\tau =\mathrm{}`$ which quietly ends at $`\tau =0`$. Strictly speaking, this only shows monopole charge is not confined; to show it is screened one must go further and show this object does not couple to any massless modes. As we showed in section 2, the field theory has an anomaly-free $`𝐙_{2M}`$ R-symmetry at all scales. The UV limit of our background, which coincides with the solution found in , has a $`U(1)`$ R-symmetry associated with the rotations of the angular coordinate $`\psi `$. For large $`M`$ it is is somewhat difficult to distinguish between the $`U(1)`$ and its discrete subgroup $`𝐙_{2M}`$. In fact, the anomaly in the $`U(1)`$, which breaks it down to $`𝐙_{2M}`$, is an effect of fractional D-instantons, the euclidean D-string world sheets propagating inside $`T^{11}`$. The Wess-Zumino term present in the D-string action, which is associated with the topologically non-trivial $`F_3`$, has to be quantized (this is simply the $`F_3`$ flux quantization). As a result, the phase in the D-string path integral assumes $`𝐙_{2M}`$ rather than $`U(1)`$ values. Our metric provides a geometrical realization for the phenomenon of chiral symmetry breaking found in the field theory; the dynamical breaking of the $`𝐙_{2M}`$ down to $`𝐙_2`$ occurs via the deformation of the conifold. In the pure supergravity limit we have discussed, the spontaneous chiral symmetry breaking generates an $`\eta ^{}`$-like Goldstone boson (the zero mode in our solution corresponding to rotation of the coordinate $`\psi `$), which must get a mass of order $`1/M`$ from these fractional instantons. To see how this mass arises, and how it relates to the domain walls which we discuss in a moment, would be very interesting. It is by now clear why the conifold ansatz adopted in and reviewed in section 2 is too restrictive: it has the $`U(1)`$ symmetry everywhere. On the other hand, our deformed conifold ansatz breaks it down to $`𝐙_2`$, with the $`U(1)`$ symmetry becoming asymptotically restored at large radius. Thus, the deformation of the conifold ties together several crucial IR effects: resolution of the naked singularity found in , breaking of the chiral symmetry down to $`𝐙_2`$, and quark confinement. At the same time, the deformation does not destroy the logarithmic running of the couplings found in because it does not affect the geometry far in the UV. Due to the deformation, the full SUGRA background has a finite 3-cycle. We now interpret various branes wrapped over this 3-cycle in terms of the gauge theory. Note that the 3-cycle has the minimal volume near $`\tau =0`$, hence all the wrapped branes will be localized there. This should be contrasted with wrapped branes in $`AdS_5\times X_5`$ where they are allowed to have an arbitrary radial coordinate. A wrapped D3-brane plays the role of a baryon vertex which ties together $`M`$ fundamental strings. Note that for $`M=0`$ the D3-brane wrapped on the $`S^3`$ gave a dibaryon ; the connection between these two objects becomes clearer when one notes that for $`M>0`$ the dibaryon has $`M`$ uncontracted indices, and therefore joins $`M`$ external charges. Meanwhile, a D5-brane wrapped over the $`S^3`$ appears to play the role of the domain wall separating two inequivalent vacua of the gauge theory. As we expect, flux tubes can end on this object , and baryons can dissolve in it; as in , we may also build the domain walls from the baryons. Indeed, D3 and D5-branes play the roles of baryon vertices and domain walls in $`𝒩=1^{}`$; however in that case they do not wrap a cycle but instead have a boundary on the NS5-brane in the space . Calculations using the metric (95) show that the baryon mass is $$M_bmM,$$ (97) while the D5-brane domain wall tension is $$T_{wall}\frac{1}{g_s}m^3.$$ (98) Additionally, one can obtain the glueball spectrum in this theory. To do so requires finding the spectrum of eigenmodes of various supergravity fields in the metric background we have constructed. Since the background is known explicitly as a function of $`\tau `$, the calculation should be no more difficult than in . Unlike the case of $`𝒩=1^{}`$, where the presence of a narrow throat near a single NS5-brane could make the computation potentially unreliable for the lowest modes , there is no possible subtlety here, as the bulk space is large and everywhere nonsingular. Of course, there will be Kaluza-Klein modes on the $`S^3`$ which are not present in the pure $`𝒩=1`$ Yang-Mills theory. These are analogous to the extra modes which appear in both and ; their presence is expected, since they are necessary whenever pure $`𝒩=1`$ Yang-Mills is embedded into a theory that is fully in the supergravity regime. Only in the limit of pure $`𝒩=1`$ Yang-Mills, which we discuss below, can they be removed. A simple estimate of the glueball and KK modes masses shows that, in the SUGRA limit both scale as $`m/(g_sM)`$. Comparing with the string tension, we see that $$T_sg_sM(m_{glueball})^2.$$ (99) Thus, there is a large separation of scales between string tension and glueball mass in supergravity (a similar problem was observed in ) which goes away at small $`g_sM`$. Finally, we should address the possibility that $`N`$ is not a multiple of $`M`$. Note that in our solution the 5-form flux vanishes for $`\tau =0`$: $$F_5=\mathrm{}(\tau )\tau ^3.$$ (100) This suggests that the IR solution given above describes a large number $`M`$ of wrapped D5-branes without any D3-branes. Therefore, for small $`\tau `$ the background should be dual to $`SU(M)`$ gauge theory (the SUGRA is reliable only if both $`M`$ and $`g_sM`$ are large). More generally, however, the field theory analysis tells us that theories that may arise in the IR have gauge groups $`SU(M+p)\times SU(p)`$, with $`M>p0`$. If $`M`$ is large and $`p`$ is of order $`1`$, then the dual supergravity background should be the same as for $`p=0`$, to leading order in $`M`$. The extra $`p`$ colors should come from $`p`$ actual D3-branes, placed at various points in our background. Then the moduli space for each D3-brane is essentially the deformed conifold, in agreement with the field theory analysis. When far from $`\tau =0`$, the D3-branes represent the IR $`𝒩=4`$ $`SU(p)`$ factor in the theory. The ’t Hooft coupling on these branes is $`g_sp1`$, so when they are brought to $`\tau =0`$ the theory represented is essentially $`SU(M+p)`$ with $`2p`$ classically massless flavors and a quartic superpotential. The nonperturbative analysis of this theory, given in section 4 and in the appendix, then applies, giving chiral symmetry breaking and a moduli space with $`M`$ branches. Note that confinement is lost in the presence of the D3-branes, in agreement with the field theory. Strings hanging from the boundary can simply end on the D3-branes, corresponding to the statement that external sources are screened by massive dynamical quarks and end up carrying only $`SU(p)`$ charge.<sup>11</sup><sup>11</sup>11Similar findings were also obtained in $`𝒩=1^{}`$ . Many of the $`𝒩=1^{}`$ vacua have dynamical massive $`W`$-bosons, whose pair production eliminates confinement. The representation of this gauge theory physics in the string theory is closely related to the representation presented here and in the last paragraph of this section. The corresponding Wilson loop will have a perimeter law. Of course if the quarks are heavy (i.e., if the D3-branes are at large $`\tau `$) then relatively short flux tubes should be stable. It would be interesting to actually demonstrate this fact, which follows not from topology but from quantum dynamics. On the other hand, if $`p`$ is of the same order as $`M`$, then the flux due to the D3-branes is large and should be included in the SUGRA solution. First, let us try to change the boundary condition on $`F_5`$ so that $`F_5`$ no longer vanishes at $`\tau =0`$ but is $`p`$. We find a consistent solution for the 5-form by replacing $`\mathrm{}(\tau )\mathrm{}(\tau )+C`$, where $`C`$ is a constant of order $`p/(g_sM^2)`$. From (82) we find that the effect of this on the warp factor is $`hh+\stackrel{~}{h}`$ where $$\stackrel{~}{h}(\tau )=\alpha C_\tau ^{\mathrm{}}𝑑x\frac{1}{K^2(x)\mathrm{sinh}^2x}.$$ (101) This yields a singular behavior of $`\stackrel{~}{h}`$ for small $`\tau `$: $$\stackrel{~}{h}\frac{\alpha C}{\tau }.$$ (102) The new behavior of $`h`$ does change significantly the physical interpretation of the solution. Now the coefficient of the $`dx_ndx_n`$ term scales as $`\tau ^{1/2}`$ for small $`\tau `$; hence, the Wilson loop no longer satisfies the area law. Again, we find agreement with the field theory. This gravity background corresponds to making the charged matter as light as possible (that is, making the expectation values of the scalar fields all as small as possible.) In this regime we expect no metastable flux tubes; the dynamical charges in the fundamental representation of $`SU(M+p)`$ will screen external electric sources, until the sources are charged only under $`SU(p)`$, which does not confine. Thus, the new behavior of the metric (102) incorporates the loss of confinement found upon addition of dynamical quarks. However, supergravity may receive large corrections in the small $`\tau `$ region because curvatures blow up at $`\tau =0`$ where we find a singular horizon.<sup>12</sup><sup>12</sup>12 We are grateful to A. Tseytlin for useful discussions of this point. Thus, requiring that $`F_5`$ does not vanish at $`\tau =0`$ actually causes a singularity. Can we construct a non-singular SUGRA solution which incorporates screening? We believe that the correct approach may be to add D3-brane sources with total charge $`p`$ (this way $`F_5`$ may smoothly turn on from zero at $`\tau =0`$ to $`p`$ at some finite value of $`\tau `$). This idea also agrees with the incorporation of small $`p`$ via actual D3-branes. We postpone construction of such non-singular ‘Coulomb branch’ solutions until a later publication. ### 5.3 The Dual of Pure $`𝒩=1`$ Yang-Mills Theory As we have shown above, supergravity serves as a reliable dual of a cascading $`SU(N+M)\times SU(M)`$ gauge theory, provided that $`g_sM`$ is very large. We have also shown that, under appropriate circumstances, at the bottom of the cascade, we essentially have an $`SU(M)`$ theory, with the other gauge group disappearing. An immediate question that arises is: can our results be used to learn something about the pure glue $`𝒩=1`$ theory? To start answering this question, let us note that the field $`B_2`$ is multiplied by $`g_sM`$, while the jumps in the cascade occur after $`B_2`$ has changed by an amount of order $`1`$. Thus, the range of $`\tau `$ which describes any particular gauge group in the cascade is of order $`1/(g_sM)`$. This implies the supergravity regime is not sufficient for constructing such a dual, because for large $`g_sM`$ the cascade jumps occur very frequently, and we find the pure glue theory only for small $`\tau `$. There, at the tip of the deformed conifold, both $`B_2`$ and $`F_5`$ are very small, $`F_3`$ is of order $`M`$, and the metric is approximately given by (92). To have a reliable dual of the pure glue theory, valid for a large range of $`\tau `$, we need to take the limit of small $`g_sM`$ (and thus small $`B_2`$, holding $`M`$ fixed) which is the opposite of the limit where supergravity has small corrections. In this limit the $`S_3`$ at the apex of the conifold becomes small and the space acquires large curvature. This situation is familiar from previous studies aimed at finding a string theory dual of a pure glue gauge theory . Nevertheless, our work does constitute progress towards formulating a stringy dual because our SUGRA background captures the correct topology of the resulting string background. Indeed we are led to conjecture that the type IIB string dual of the pure glue $`𝒩=1`$ $`SU(M)`$ theory is given by a $`g_sM0`$ limit of a warped deformed conifold background, with $`M`$ units of the $`F_3`$ flux piercing its 3-cycle, and with $`B_2`$ and $`F_5`$ approaching zero at the apex. This would be the space generated by the fractional D3-branes alone, with no admixture of regular D3-branes. Hence it is relevant to pure $`SU(M)`$ theory with no quark flavors. Of course, studying such a theory for small $`g_sM`$ is difficult due to the well-known problems with RR flux and large curvature. However, the self-dual 5-form flux, which brings in some additional problems, is small, which raises hopes of a novel sigma model formulation. We note also that the addition of a small number of D3-branes to this story will permit the study of the $`SU(M+p)\times SU(p)`$ theory, which essentially reduces, for small $`g_s`$ and $`pM`$, to $`SU(M+p)`$ with $`2p`$ flavors and an all-important quartic superpotential. It is far from certain that this construction can give any insight into QCD, since the light charged scalars play such a central role in the dynamics. However, if these scalars can easily be removed (along with the gauginos) through explicit supersymmetry breaking, there might be additional interest in this approach. ## 6 Discussion We have not addressed the question of how to compute field theory correlation functions in this context, where our space does not approach Anti-de-Sitter space at large $`r`$. However, it is easy to see this space still has a boundary, and from the behavior of $`h(\tau )`$ it is clear that the logarithm is a subleading effect at large $`r`$. Correspondingly, at large $`NM`$, there is a sense in which the operators in the field theory have the same spectrum that they have for $`M=0`$, since $`\gamma \frac{1}{2}`$. We therefore believe that for low-lying supergravity modes, corresponding to operators of dimension much less than $`M`$, the story will not be modified in a significant way from that discussed in . For operators of dimension $`\mathrm{\Delta }>\frac{3}{2}M`$, we expect more interesting effects. These operators appear to exist at scales where $`\frac{3}{2}N>\mathrm{\Delta }`$, but should be eliminated when $`\frac{3}{2}N<\mathrm{\Delta }.`$ In the gauge theory, it is known what should occur ; operators of high dimension present classically are actually removed by quantum effects, which in the low-energy dual theory appear as simple group theory. On the gravity side we may speculate that high-lying bulk modes propagate in from the boundary until the region where $`N\frac{3}{2}\mathrm{\Delta }`$; there $`T^{11}`$ has shrunk down such that these modes blow up into the “giant gravitons” of . Only modes with dimension $`\mathrm{\Delta }<\frac{3}{2}M`$ can propagate all the way to $`\tau =0`$. It is easy to see that our story of the duality cascade can be orientifolded. This is obvious from the type IIA string theory brane construction. It is also clear from the corresponding $`SO\times Sp`$ gauge theory, although we have not analyzed the field theory dynamics to see how the orientifolded conifold is deformed. A number of other modifications, including theories whose IIA version involves multiple NS and NS’ branes, could potentially be interesting. This might be especially true for theories which are qualitatively different in the infrared from pure Yang-Mills, such as those studied in . Another interesting choice would be to orbifold the theory along the lines of , so that the low energy theory is non-supersymmetric $`SU(M)^2`$ with a Dirac fermion in the bifundamental representation. In contrast to the case studied in , the masslessness of this fermion would be exact, as it is guaranteed by the $`𝐙_{2M}`$ R-symmetry, and therefore chiral symmetry breaking and confinement in this QCD-like theory could be exhibited in the supergravity regime. Finally, it is interesting to resurrect a scenario discarded five years ago for its apparent absurdity. Namely, it is conceivable that the standard model — a small gauge group — itself lies at the base of a duality cascade. This is certainly possible, since the addition of supersymmetry and some appropriately chosen massive matter at the TeV scale easily could make the theory into one which could emerge from such an RG flow. In it was in fact pointed out that this was the natural scenario if the standard model, with its very small gauge groups, is a low-energy Seiberg-dual description of some other theory; every natural choice for an ultraviolet theory has a larger gauge group than the standard model, and typically hits a Landau pole below the Planck scale, requiring additional duality transformations, still larger gauge groups, more Landau poles, and continuation ad nauseum. This was termed the “duality wall” (since in some cases the duality transformations piled up so fast that an infinite number were required in a finite energy range.) But now we see this continuous generation of larger and larger gauge groups — ugly and unmotivated within field theory, and driving the field theory into highly non-perturbative regimes — can correspond to a perfectly natural spacetime background on which strings may propagate. If we imagine that the ultraviolet of the duality cascade is cut off in a compact space (along the lines of , following ) we may conjecture that the standard model coupled to gravity is best described, at high energy, by a compactified string theory on a space with a logarithmic (or otherwise warped) throat, with the weakly coupled standard model emerging as a good description only at energies below, say, 1–100 TeV. Such a model provides another possible way, somewhat related to ideas of , to explain the hierarchy between the gravitational and electroweak scales: it is perhaps given by TeV$`=m_{Pl}\times e^{cN/M}`$, where $`M`$ is of order 2 to 5, $`c`$ is a number of order one, and $`N`$ is the number of colors of the gauge group at the Planck scale. ## Acknowledgements We are grateful to K. Dasgupta, S. Frolov, S. Gubser, S. Gukov, J. Maldacena, J. Polchinski, A. Tseytlin and E. Witten for useful discussions. The work of I.K. was supported in part by the NSF grant PHY-9802484 and by the James S. McDonnell Foundation Grant No. 91-48; that of M.J.S. was supported by NSF grant PHY95-13835 and by the W.M. Keck Foundation. ## 7 Appendix In this appendix we analyze the field theory in somewhat greater detail, confirming and extending the results of section 3. First, we may check the results of section 3 in another region of moduli space. Consider first $`SU(M+1)`$ with two flavors. Suppose we permit $`C_1`$ and $`D_1`$ to have equal expectation values $`v`$, so that $`N_{11}=v^2`$. This breaks the $`SU(M+1)`$ to $`SU(M)`$. If $`\lambda `$ were zero, this would leave $`SU(M)`$ with one flavor $`C_2`$ and $`D_2`$, plus two gauge singlets $`C_1D_2=N_{12}`$ and $`C_2D_1=N_{21}`$; the corresponding strong coupling scale would be $`\mathrm{\Lambda }^{M+1}/v^2`$. However, the presence of nonzero $`\lambda `$ gives mass to these fields, leaving the $`SU(M)`$ gauge theory with a flavor of mass $`\lambda v^2`$. The effective Lagrangian is then $$W=2\lambda v^2N_{22}+\left[\frac{\mathrm{\Lambda }^{M+1}/v^2}{N_{22}}\right]^{\frac{1}{M1}}$$ (103) which again leads to $`M`$ branches with the correct values of $`N_{ij}N_k\mathrm{}ϵ^{ik}ϵ^j\mathrm{}`$. That our discussion of the $`SU(M+1)`$ theory in section 3 was only part of the story can be seen by starting one step higher, with $`SU(2M+1)\times SU(M+1)`$, which reduces after one duality transformation to the $`SU(M+1)`$ case. The $`SU(2M+1)`$ gauge group has one more flavor than color, and therefore, as $`\lambda _10`$, the theory is governed by the results of . For $`\lambda _1=0`$ the superpotential must go over to $$W\frac{detP_{ijb}^a}{\mathrm{\Lambda }_1^{4M+1}}C_{ia}P_{ijb}^aD_j^b,$$ (104) where $`a,b`$ are color indices of $`SU(M+1)`$, and $`PAB`$, $`CA^{2M+1}`$, $`DB^{2M+1}`$. From this we learn the function $`F_1(I_1,J_1,R_1^{(1)})`$ is not equal to one, and in fact, in the limit $`\lambda _10`$, that is, $`I_1,J_10`$, we have $`F_1(I_1,J_1,R_1^{(1)})\sqrt{I_1/J_1}f(R_1^{(1)})`$. The low energy theory is then $`SU(M+1)`$ with two flavors $`C_i,D_i`$ but with superpotential $$W=\lambda _2C_iD_jC_kD_lϵ^{ik}ϵ^{jl}F_2(I_2,J_2)$$ (105) Here $`(C_iD_jC_kD_lϵ^{ik}ϵ^{jl})`$ is the only invariant involving $`C`$ and $`D`$; there are no $`R_2`$ ratios. The low energy effective superpotential is now $$W_L=\lambda N_{ij}N_k\mathrm{}ϵ^{ik}ϵ^j\mathrm{}F_2(I_2,J_2).$$ (106) Note that $`F_2(I_20,J_20)=1+\sqrt{I_2/J_2}`$; some other limits can be studied but will not be needed here. The vacuum equations are $$0=\lambda \left[F(I_2,J_2)+I_2\frac{F(I_2,J_2)}{I_2}\right]N_{ij}.$$ (107) This gives an equation for $`I_2(N_{ij}N_k\mathrm{}ϵ^{ik}ϵ^j\mathrm{})^{2M}`$, whose solution must be $$I_2=G(J_2).$$ (108) The holomorphic function $`G(J_2)`$ is not zero everywhere (since for $`I_20`$, $`J_20`$ it is not zero) so it can only be zero at special points. Consequently $`N_{ij}N_k\mathrm{}ϵ^{ik}ϵ^j\mathrm{}`$ is generally nonzero. Since the $`𝐙_{2M}`$ symmetry rotates $`N_{ij}N_k\mathrm{}ϵ^{ik}ϵ^j\mathrm{}`$ by $`e^{2\pi i/M}`$, we again find $`M`$ separate branches. Again there are no restrictions on the individual values of the $`N_{ij}`$, so each branch takes the form of a deformed conifold, with a nonzero superpotential. Thus we obtain the same result as before; only the magnitude of the deformation is modified from our previous analysis. This analysis is too weak to rule out the possibility that there might be several independent solutions for $`I_2`$ given a single value of $`J_2`$. This would lead to several sets of branches, each set consisting of $`M`$ copies of a deformed conifold; the different sets would have deformations of different magnitudes. In the limit $`\mathrm{\Lambda }_1\mathrm{}`$ only one set would remain, as in our earlier analysis of the $`SU(M+1)`$ theory. Next, we consider the case of $`SU(M+p)\times SU(p)`$, $`0<p<M`$. We will first perform the analysis by taking the $`SU(p)`$ coupling small. We define $`(N_{ij})_\beta ^\alpha =(C_i)_a^\alpha (D_i)_\beta ^a`$, where $`\alpha ,\beta `$ are $`SU(p)`$ indices and $`a,b`$ are $`SU(M+p)`$ indices. If the $`SU(p)`$ coupling were set to zero, then we would have an $`SU(M+p)`$ gauge theory with $`2p`$ flavors. An Affleck-Dine-Seiberg superpotential would be generated, giving $$W=\lambda (N_{ij})_\beta ^\alpha (N_k\mathrm{})_\alpha ^\beta ϵ^{ik}ϵ^j\mathrm{}+(Mp)\left(\frac{\mathrm{\Lambda }_1^{3M+p}}{det_{ij,\alpha \beta }N}\right)^{\frac{1}{Mp}}$$ (109) where in the determinant we treat $`N`$ as a $`2p\times 2p`$ matrix. A little algebra gives the equations for a supersymmetric vacuum as $$det[(N_{ij})_\beta ^\alpha ]\left(\frac{\mathrm{\Lambda }_1^{3M+p}}{\lambda ^{Mp}}\right)^{\frac{2}{M}}$$ (110) and $$\lambda (N_{ij})_\beta ^\alpha (N_k\mathrm{})_\alpha ^\beta ϵ^{ik}ϵ^j\mathrm{}(\lambda ^p\mathrm{\Lambda }^{3M+p})^{1/M}$$ (111) It is possible to show that these equations represent $`M`$ branches, each of which is $`p`$ copies of the deformed conifold — in other words, the moduli space of $`p`$ probe branes moving on the deformed conifold. First, note that $`N_{ij}^0(N_{ij})_\alpha ^\alpha `$ is a gauge invariant operator. If we demand that the $`SU(p)`$-adjoint fields $`(N_{ij})_\beta ^\alpha \frac{1}{p}\delta _\beta ^\alpha (N_{ij})_\gamma ^\gamma `$ vanish, then the equations above become $$det[(N_{ij}^0)]\left(\frac{\mathrm{\Lambda }_1^{3M+p}}{\lambda ^{Mp}}\right)^{\frac{2}{M}}$$ (112) and $$\lambda N_{ij}^0N_k\mathrm{}^0ϵ^{ik}ϵ^j\mathrm{}(\lambda ^p\mathrm{\Lambda }^{3M+p})^{1/M}$$ (113) which gives $`M`$ branches, each of which is a single copy of the deformed conifold. This region of moduli space corresponds to taking all $`p`$ probe branes to have the same positions on the conifold. As before the $`𝐙_{2M}`$ global symmetry is broken to $`𝐙_2`$; it is easy to see that the superpotential on the $`M`$ branches rotates by a phase under the broken $`𝐙_M`$. Expectation values for elements of $`(N_{ij})_\beta ^\alpha \frac{1}{p}\delta _\beta ^\alpha (N_{ij})_\gamma ^\gamma `$ correspond to moving the $`p`$ probe branes apart; taking the special cases where these fields are diagonal, it is easy to show that each set of eigenvalues of $`(N_{ij})_\beta ^\alpha `$, $`i,j=1,2`$, sweeps out its own copy of the deformed conifold. When the $`SU(p)`$ gauge coupling is turned back on, the superpotential will include unknown functions of the invariants $`I`$, $`J`$, and $`R`$. These can be generated by a number of different physical phenomena, including instantons in regions where the $`SU(p)`$ group is partially broken. However, as before, these functions change the quantitative features of the deformation of the conifold without altering the basic picture we have obtained. Furthermore, we expect no additional significant infrared dynamics. Above the strong-dynamics scale for $`SU(M+p)`$, the $`SU(p)`$ gauge group is infrared free. Below it, the $`SU(p)`$ group contains three adjoint fields $`(N_{ij})_\beta ^\alpha `$ which have a trilinear superpotential — in short, a copy of $`𝒩=4`$ Yang-Mills. The $`SU(p)`$ sector is therefore scale-invariant and nonconfining at low energy. Lastly, we expect that the $`SU(p)`$ dynamics plays no role in the supergravity regime for $`pM`$. Supergravity requires we work at small gauge coupling and large ’t Hooft coupling for $`SU(M)`$, but in this regime $`SU(p)`$ will have small ’t Hooft coupling and will be described by weakly-coupled field theory. In the end, then, we again expect $`M`$ branches, given by equations of the same qualitative form as above. The case $`p=M`$ is the most subtle. For $`SU(2M)\times SU(M)`$, the $`SU(2M)`$ theory has equal numbers of flavors and colors, and consequently its moduli space is modified quantum mechanically . If we turn off the $`SU(M)`$ coupling, the superpotential becomes $$W=\lambda (N_{ij})_\beta ^\alpha (N_k\mathrm{})_\alpha ^\beta ϵ^{ik}ϵ^j\mathrm{}F_1(I_1/J_1)+X(det[(N_{ij})_\beta ^\alpha ]\overline{}\mathrm{\Lambda }_{2M}^{4M}),$$ (114) where the “baryon” $``$ is the gauge invariant operator $`A_1^MA_2^M`$, and the anti-baryon is similarly constructed from $`B_i`$. Here the equations seem to have multiple solutions. One solution is $$X=0;N=0;=\overline{}=i\mathrm{\Lambda }_{2M}^{2M}.$$ (115) In this case, the $`SU(M)`$ gauge group is unbroken and, when its coupling is restored, it generates $`M`$ distinct and isolated vacua via usual gaugino condensation. Alternatively, we may have $$=\overline{}=0;det[(N_{ij})_\beta ^\alpha ]=\mathrm{\Lambda }_{2M}^{4M};[(N_{ij})_\beta ^\alpha (N_k\mathrm{})_\alpha ^\beta ϵ^{ik}ϵ^j\mathrm{}G_1(I_1/J_1)]^M=\mathrm{\Lambda }_{2M}^{4M},$$ (116) where we have not determined $`G_1`$. As before this leads to $`M`$ branches, each of which has $`M`$ probe branes moving on a deformed conifold. This suggests that the complete solution to a theory with gauge group $`SU(N+M)\times SU(N)`$ might involve not one set of $`M`$ branches but many. The smallest set would consist of $`pN\mathrm{mod}M`$ D3-branes moving on the deformed conifold. The next smallest set would consist of $`p+M`$ D3-branes. Next would follow a branch with $`p+2M`$ D3-branes, and so forth, growing in size without limit. To see whether this is the case requires a more thorough and complete field theory analysis, which we have not performed. In any case, these partial results all support the main claims of the paper: that all branches which appear are consistent with probe branes moving on a deformed conifold, and that each branch is one of $`M`$ identical branches which are rotated by the spontaneously broken $`𝐙_{2M}`$ R-symmetry.
warning/0007/hep-ph0007219.html	ar5iv	text	# Backreaction and the Parametric Resonance of Cosmological Fluctuations ## I Introduction It has long been realized that reheating is a crucial part of the inflationary scenario. During reheating the large energy density contained within the coherently oscillating inflaton field is converted into particle excitations of whatever fields are coupled to the inflaton, vastly increasing the temperature and entropy density and setting the stage for the standard big bang phase. If inflation is ever to be a useful picture for describing the early universe, then it is essential to understand the details of how the vacuum energy is transformed into familiar particles. Reheating can occur very efficiently through the process of parametric resonance . Field modes within certain resonance bands in $`k`$-space grow exponentially with time, defining the “preheating” era. The possibility of resonant growth of linear scalar metric perturbations was first studied in . Recently it has been argued that the resonance of scalar metric perturbations can extend to $`kaH`$, i.e. that super-Hubble perturbations can be amplified . This opens up the possibility of new observational consequences, since the scales relevant to the cosmic microwave background and large-scale structure are much larger than the Hubble radius during preheating. The importance of the gauge-invariant formalism for cosmological perturbations and the study of the “traditionally conserved” Bardeen parameter $`\zeta `$ was emphasized in , where it was found that simple single-field chaotic inflation models do not exhibit super-Hubble growth beyond what is expected in the absence of parametric resonance. The absence of parametric amplification of super-Hubble modes in these single field models was shown to hold in a full nonlinear treatment , and a general no-go theorem in these models was suggested in . For the first model which was claimed to exhibit growth of super-Hubble metric perturbations beyond that of the usual theory of reheating (see also ), it was soon realized that the growth was unimportant since it followed a period of exponential damping during inflation . This damping of super-Hubble modes arises because the field perturbations which are amplified during preheating have an effective mass greater than the Hubble parameter during inflation. This results in a very “blue” power spectrum at the end of inflation, with a severe deficit at the largest scales . The relatively plentiful small-scale modes can also grow resonantly during preheating. Thus the end of parametric resonance occurs when the backreaction of the dominant small-scale modes becomes important, and the cosmological-scale modes are still negligible. An obvious class of models to study, then, consists of those with small masses during inflation and strong super-Hubble resonance . A simple example was provided by Bassett and Viniegra , namely that of a massless self-coupled inflaton $`\varphi `$ coupled to another scalar field $`\chi `$, i.e. a model with potential $$V(\varphi ,\chi )=\frac{\lambda }{4}\varphi ^4+\frac{g^2}{2}\varphi ^2\chi ^2.$$ (1) This model has been studied in detail, but in the absence of metric perturbations, by Greene et al., who found that the model contains a strong resonance band for $`\chi `$ fluctuations which extends to $`k=0`$ for the choice $`g^2=2\lambda `$. Bassett and Viniegra found that super-Hubble metric perturbations are resonantly amplified as well in this model (see also ). To date, however, none of the linearized analyses of parametric amplification of super-Hubble-scale metric fluctuations in the model (1) has included the effects of backreaction on the evolution of the fluctuations. The backreaction of the growing modes on the background fields is expected to shut the growth down at some point, but exactly when? Backreaction is also the only hope to make models which exhibit parametric amplification of super-Hubble cosmological perturbations compatible with the Cosmic Background Explorer (COBE) normalization . In this paper we will investigate the effects of backreaction on the growth of matter and metric fluctuations using the Bassett and Viniegra model (1) as our toy model. We will study the growth of scalar field and scalar metric perturbations, including the effect of backreaction in the Hartree approximation. We carefully treat the evolution during inflation, which can be very important for super-Hubble scales. We will compare the large-scale normalization predicted for this model with the COBE value, and find that, although backreaction is crucial in limiting the growth of the fluctuations, the final amplitude is larger than allowed by the COBE normalization (for supersymmetry-motivated coupling constant values), unless the period of inflation is very long. Note that the final amplitude of fluctuations in our model is independent of the scalar field coupling constant (unlike what happens without parametric resonance effects), but that it may depend on the duration of the inflationary period. We also extend the model to study the effect of $`\chi `$-field self-coupling, which can be important in limiting the growth of fluctuations. ## II Model and Linearized Dynamics ### A Equations of motion and analytical theory Our model is the two-real-scalar-field, gravitationally minimally coupled model specified by the Lagrangian density $$=\sqrt{g}\left(\frac{1}{2}_\mu \varphi ^\mu \varphi +\frac{1}{2}_\mu \chi ^\mu \chi \frac{\lambda }{4}\varphi ^4\frac{g^2}{2}\varphi ^2\chi ^2\right).$$ (2) The field $`\varphi `$ drives inflation, while $`\chi `$ is significant only after parametric resonance begins, so the inflationary dynamics is essentially that of $`(\lambda /4)\varphi ^4`$ chaotic inflation. Note that the behaviour of this system is expected to be robust under the addition of a small mass term $`m_\varphi ^2\varphi ^2`$ with $`m_\varphi \sqrt{\lambda }m_{\text{Pl}}`$ and for the ratio of coupling constants satisfying $`g/\sqrt{\lambda }<\sqrt{\lambda }m_{\text{Pl}}/m_\varphi `$ . In particular, this will be the case for supersymmetric models, which motivate the choice $`g^2=2\lambda `$ . In addition, we will show that large values of $`g/\sqrt{\lambda }`$ are in fact inconsistent with the significant amplification of super-Hubble modes. On the other hand, for $`g/\sqrt{\lambda }>\sqrt{\lambda }m_{\text{Pl}}/m_\varphi `$, the theory of “stochastic resonance” for a massive inflaton may need to be applied . It is traditional to separate the inflaton field into a homogeneous, “classical” background $`\varphi (t)`$ and a perturbation $`\delta \varphi (x,t)`$, which begins as sub-Hubble quantum vacuum fluctuations early in inflation. It will be useful to make a similar separation for the field $`\chi `$, although we have no reason to expect a non-zero initial homogeneous $`\chi `$ component in this model. The equations of motion for the homogeneous parts of the inflaton and $`\chi `$ fields are the Klein-Gordon equations, $$\ddot{\varphi }+3H\dot{\varphi }+\lambda \varphi ^3+g^2\chi ^2\varphi =0,$$ (3) $$\ddot{\chi }+3H\dot{\chi }+g^2\varphi ^2\chi =0,$$ (4) with Hubble parameter $`H=\dot{a}/a`$, where $`a`$ is the scale factor. To complete the background dynamics we must specify the evolution of the background spacetime metric. We assume a spatially flat Friedmann-Robertson-Walker universe, for which the 0-0 Einstein equation gives the Friedmann equation $$H^2=\frac{8\pi }{3m_{\text{Pl}}^2}\left[\frac{1}{2}\dot{\varphi }^2+\frac{1}{2}\dot{\chi }^2+V(\varphi ,\chi )\right].$$ (5) For $`\varphi m_{\text{Pl}}`$, the universe undergoes slow-roll inflation, with $`H`$ approximately constant and the scale factor $`a`$ increasing approximately exponentially with time. As slow-roll ends, the “damping” term $`3H\dot{\varphi }`$ becomes less important in Eq. (3) and the field begins to oscillate about $`\varphi =0`$. This marks the start of the preheating period. Averaged over several oscillations, the equation of state (in the absence of backreaction) is very nearly that of a radiation-dominated universe , and the amplitude of the inflaton’s oscillations decays as $`a^1`$. This is a consequence of the (near) conformal invariance of this massless model, which considerably simplifies the treatment of parametric resonance as compared with the massive case . In writing the linearized equations of motion for perturbations about the background, we will use the longitudinal gauge. For this model the metric can be written $$ds^2=(12\mathrm{\Phi })dt^2a^2(t)(1+2\mathrm{\Phi })dx_idx^i,$$ (6) so scalar metric perturbations are described by the single variable $`\mathrm{\Phi }`$. The momentum-space first-order perturbed Einstein and Klein-Gordon equations are $`3H\dot{\mathrm{\Phi }}`$ $`+`$ $`\left({\displaystyle \frac{k^2}{a^2}}+3H^2\right)\mathrm{\Phi }`$ (7) $`=`$ $`{\displaystyle \frac{4\pi }{m_{\text{Pl}}^2}}{\displaystyle \underset{i}{}}\left(\dot{\varphi }_i\dot{\delta \varphi }_i\mathrm{\Phi }\dot{\varphi }_i^2+V_{,i}\delta \varphi _i\right),`$ (8) $$\dot{\mathrm{\Phi }}+H\mathrm{\Phi }=\frac{4\pi }{m_{\text{Pl}}^2}\underset{i}{}\dot{\varphi }_i\delta \varphi _i,$$ (9) $$\ddot{\delta \varphi _i}+3H\dot{\delta \varphi }_i+\frac{k^2}{a^2}\delta \varphi _i+\underset{j}{}V_{,ij}\delta \varphi _j=4\dot{\mathrm{\Phi }}\dot{\varphi }_i2V_{,i}\mathrm{\Phi },$$ (10) where $`\varphi _1\varphi `$, $`\varphi _2\chi `$, $`V_{,i}V/\varphi _i`$, and comoving momentum $`k`$ subscripts have been suppressed for clarity. Equations (8) and (9) can be combined to give $$\mathrm{\Phi }=\frac{_i\left(\dot{\varphi }_i\dot{\delta \varphi }_i+3H\dot{\varphi }_i\delta \varphi _i+V_{,i}\delta \varphi _i\right)}{(m_{\text{Pl}}^2/4\pi )(k/a)^2+_i\dot{\varphi }_i^2},$$ (11) which fixes $`\mathrm{\Phi }`$ once the matter fields are known. An important quantity in the study of the linear evolution of metric perturbations is the Bardeen parameter $$\zeta _k=\mathrm{\Phi }_k\frac{H}{\dot{H}}\left(\dot{\mathrm{\Phi }}_k+H\mathrm{\Phi }_k\right),$$ (12) which for $`k/aH`$ and single field models satisfies the “conservation law” $$(1+w)\dot{\zeta _k}=0,$$ (13) where $`w=P/\rho `$ is the equation of state ($`\rho `$ and $`P`$ denoting energy density and pressure, respectively). When $`\dot{\mathrm{\Phi }}_k`$ can be neglected, Eqs. (12) and (13) can be combined to give the familiar result that the change in $`\mathrm{\Phi }_k`$ on super-Hubble scales over some interval of time is determined solely by the change in equation of state. In the absence of metric perturbations, the linearized dynamics in the model described above is known to exhibit parametric resonance during preheating . To see this, it helps to take advantage of the near conformal invariance of the model and rewrite the equations in terms of conformally scaled fields $`\stackrel{~}{\varphi }_ia\varphi _i`$ and a dimensionless conformal time $`x\sqrt{\lambda }\stackrel{~}{\varphi }_0\eta `$, where $`\eta =𝑑t/a`$ and $`\stackrel{~}{\varphi }_0`$ is the amplitude of inflaton oscillations at the start of preheating. Then, with the $`\chi `$ background and metric perturbations set to zero, the inflaton background and perturbed field equations (3) and (10) become $$\stackrel{~}{\varphi }^{\prime \prime }+\frac{\stackrel{~}{\varphi }^3}{\stackrel{~}{\varphi }_0^2}=0,$$ (14) $$\delta \stackrel{~}{\varphi }_k^{\prime \prime }+\left(\kappa ^2+3\frac{\stackrel{~}{\varphi }^2}{\stackrel{~}{\varphi }_0^2}\right)\delta \stackrel{~}{\varphi }_k=0,$$ (15) $$\delta \stackrel{~}{\chi }_k^{\prime \prime }+\left(\kappa ^2+\frac{g^2}{\lambda }\frac{\stackrel{~}{\varphi }^2}{\stackrel{~}{\varphi }_0^2}\right)\delta \stackrel{~}{\chi }_k=0,$$ (16) where $`\kappa ^2k^2/(\lambda \stackrel{~}{\varphi }_0^2)`$ is a dimensionless comoving momentum and primes denote derivatives with respect to the scaled conformal time $`x`$. Here we have ignored terms proportional to $`a^{\prime \prime }/a`$ since preheating is a nearly-radiation-dominated phase in this model . The conformal field $`\stackrel{~}{\varphi }`$ then undergoes constant amplitude elliptic cosine oscillations, while the perturbation equations are Lamé equations , which are known to exhibit resonance within certain bands in parameter space . In particular, the $`\delta \stackrel{~}{\chi }_k`$ equation exhibits strong resonance for a band that includes $`k=0`$ for the supersymmetric point $`g^2/\lambda =2`$, and weak “narrow” resonance in small-scale bands. On the other hand, $`\delta \stackrel{~}{\varphi }_k`$ exhibits narrow resonance for a sub-Hubble momentum range, independent of the coupling constants. For resonant modes, the growth is a modulated exponential, $`\delta \stackrel{~}{\chi }_ke^{\mu _kx}`$, with Floquet index $`\mu _k`$. If we allow $`g^2/\lambda `$ to vary, we find a sequence of $`\delta \stackrel{~}{\chi }_k`$ resonance bands for $`k=0`$, centred at $`g^2/\lambda =2n^2`$ with width $`2n`$, for positive integral $`n`$ . The Floquet index reaches a maximum value of $`\mu _{\text{max}}0.238`$ at the centre of each $`k=0`$ band. ### B Numerical results For our numerical calculations, we were primarily interested in the behaviour of cosmological-scale matter and metric modes. Thus we evolved a scale which left the Hubble radius (at time $`t_0`$) at about $`N=50`$ $`e`$-folds before the end of inflation. For $`(\lambda /4)\varphi ^4`$ models, the number of $`e`$-folds during slow-roll inflation after initial time $`t_0`$ is $$N\pi \left(\frac{\varphi (t_0)}{m_{\text{Pl}}}\right)^2;$$ (17) thus we used the homogeneous inflaton initial value of $`\varphi (t_0)=4m_{\text{Pl}}`$. We began the calculations with the modes still somewhat inside the Hubble radius, so the initial conditions for the matter field fluctuations were simply given by the conformal vacuum state $$\delta \varphi _{ik}(t_0)=\frac{1}{a^{3/2}(t_0)}\left(\frac{1}{2\omega _i(t_0)}\right)^{1/2},$$ (18) $$\dot{\delta \varphi }_{ik}(t_0)=i\omega _i(t_0)\delta \varphi _i(t_0),$$ (19) with $`\omega _\varphi ^2(t)=(k/a)^2+3\lambda \varphi ^2+g^2\chi ^2`$ and $`\omega _\chi ^2(t)=(k/a)^2+g^2\varphi ^2`$. Physically, the $`a^{3/2}`$ dependence arises because particle number densities $`n_k\|\delta \varphi _{ik}\|^2`$ must decay like $`a^3`$ in the massive, adiabatic regime. The initial metric perturbations were then determined by Eq. (11). To illustrate the dynamics in the absence of backreaction, we numerically integrated the coupled set of background equations (3) and (5) and perturbation equations (9) and (10) using the initial conditions described above, and for $`g^2/\lambda =2`$, $`\lambda =10^{14}`$, and a zero $`\chi `$ background. We used the constraint Eq. (11) as well as the conservation equation (13) to check the accuracy of the calculations. In Fig. 1 we display the evolution of our cosmological modes, together with the Bardeen parameter $`\zeta _k`$, during inflation and preheating. For each of the perturbations $`X_k=\delta \chi _k`$, $`\delta \varphi _k`$, $`\mathrm{\Phi }_k`$, and $`\zeta _k`$ we plot the power spectrum $$𝒫_X(k)=\frac{k^3}{2\pi ^2}\|X_k\|^2,$$ (20) rather than the mode amplitudes, to facilitate comparison with the COBE measured normalization which gives $`𝒫_\mathrm{\Phi }10^{10}`$ . Figure 1 shows how the modes begin early in inflation as sub-Hubble oscillations, and become “frozen in” after they exit the Hubble radius. Note that the $`\delta \chi _k`$ fluctuation experiences some damping late in inflation, when its effective mass squared $`g^2\varphi ^2`$ becomes somewhat greater than $`H^2`$, which decreases like $`\varphi ^4`$ in the slow-roll approximation. The inflaton perturbation $`\delta \varphi _k`$, however, stays roughly constant during inflation even though its effective mass is comparable to that of the $`\delta \chi _k`$ mode. This is because of the coupling between $`\delta \varphi _k`$ and $`\mathrm{\Phi }_k`$ in the linearized perturbation equations. We can also observe a growth of $`\mathrm{\Phi }_k`$ between the time the mode exits the Hubble radius and the beginning of preheating, by a factor of approximately $`20`$, in good agreement with the growth predicted from the “conservation law” Eq. (13). Also note that after this small growth stage the cosmological-scale metric power spectrum ends up close to the $`10^{10}(e^{23})`$ level, as the standard theory predicts for $`\lambda 10^{14}`$ in the absence of parametric resonance . During preheating we observe exponential growth of $`\delta \chi _k`$ while the super-Hubble $`\delta \varphi _k`$ mode does not grow, as expected from the analytical theory. $`\mathrm{\Phi }_k`$ and $`\zeta _k`$ remain constant, since according to Eq. (9) the metric perturbations couple only to $`\delta \varphi _k`$ in the absence of a $`\chi `$ background, at linear level . To observe the effect of including a non-zero homogeneous $`\chi `$ background, we repeated the above calculation with an initial value of $`\chi (t_0)=10^{10}m_{\text{Pl}}`$ (this value illustrates well the various stages of evolution). Figure 2 indicates that $`\delta \chi _k`$ grows as before, but $`\delta \varphi _k`$ and $`\mathrm{\Phi }_k`$ now grow initially with twice the Floquet index of $`\delta \chi _k`$. This is the result of the driving term $`2g^2\varphi \chi \delta \chi _k`$ in the equation of motion for $`\delta \varphi _k`$, Eq. (10), which contains two factors growing like $`e^{\mu _{\text{max}}x}`$ (clearly the evolution of the background $`\chi `$ will be essentially the same as that of $`\delta \chi _k`$ for $`kaH`$). Once the background $`\chi `$ field becomes comparable to the inflaton background, all the perturbations synchronize and grow at the same rate. We will discuss the significance of the homogeneous $`\chi `$ field in relation to the nonlinear evolution of the fields in the next section. ## III Backreaction ### A Equations of motion The linearized equations in the previous section describe unbounded growth of perturbations during resonance. In reality this growth must of course stop at some point, namely when the perturbed field values are on the order of the background values. A full nonlinear simulation will include this effect automatically, but approximation methods can alleviate the computational costs significantly. A common approach to approximate this backreaction on the background and perturbation evolution is to include Hartree terms in the equations of motion . This entails making the replacements $`\varphi _i\varphi _i+\delta \varphi _i`$, $`\delta \varphi _i^2\delta \varphi _i^2`$, and $`\delta \varphi _i^33\delta \varphi _i^2\delta \varphi _i`$. In this approximation, the background equations (3) – (5) become $`H^2`$ $`=`$ $`{\displaystyle \frac{8\pi }{3m_{\text{Pl}}^2}}[V(\varphi ,\chi )`$ (23) $`+{\displaystyle \frac{3}{2}}\lambda \varphi ^2\delta \varphi ^2+{\displaystyle \frac{g^2}{2}}\varphi ^2\delta \chi ^2+{\displaystyle \frac{g^2}{2}}\chi ^2\delta \varphi ^2`$ $`+{\displaystyle \frac{1}{2}}{\displaystyle \underset{i}{}}(\dot{\varphi }_i^2+\dot{\delta \varphi }_i^2+{\displaystyle \frac{1}{a^2}}\delta \varphi _i^2)],`$ $$\ddot{\varphi }+3H\dot{\varphi }+V_{,\varphi }+3\lambda \delta \varphi ^2\varphi +g^2\delta \chi ^2\varphi =0,$$ (24) $$\ddot{\chi }+3H\dot{\chi }+V_{,\chi }+g^2\delta \varphi ^2\chi =0.$$ (25) Similarly, the momentum-space linearized field perturbation equations (10) become $`\ddot{\delta \varphi }_k`$ $`+`$ $`3H\dot{\delta \varphi }_k+\left({\displaystyle \frac{k^2}{a^2}}+3\lambda \delta \varphi ^2+g^2\delta \chi ^2\right)\delta \varphi _k`$ (26) $`+`$ $`{\displaystyle \underset{j}{}}V_{,\varphi j}\delta \varphi _{jk}=4\dot{\varphi }\dot{\mathrm{\Phi }}_k2V_{,\varphi }\mathrm{\Phi }_k,`$ (27) $`\ddot{\delta \chi }_k`$ $`+`$ $`3H\dot{\delta \chi }_k+\left({\displaystyle \frac{k^2}{a^2}}+g^2\delta \varphi ^2\right)\delta \chi _k`$ (28) $`+`$ $`{\displaystyle \underset{j}{}}V_{,\chi j}\delta \varphi _{jk}=4\dot{\chi }\dot{\mathrm{\Phi }}_k2V_{,\chi }\mathrm{\Phi }_k.`$ (29) In this approach, the field fluctuations are calculated self-consistently from the relations $$\delta \varphi _i^2=\frac{1}{(2\pi )^3}d^3k\|\delta \varphi _{ik}\|^2.$$ (30) In practice, the resonance band will provide a natural ultraviolet cutoff. The Hartree terms approximate the full nonlinear dynamics of the fields. To illustrate what this approximation entails, we may consider the exact dynamics of the fields, treating the Klein-Gordon equation as a classical field equation. This should be a good approximation soon after the beginning of the resonance stage, since occupation numbers will grow exponentially with time . As an example, consider the exact evolution equation for $`\delta \varphi `$ in position space, obtained by perturbing the Klein-Gordon equation, setting the background $`\chi `$ to zero, and ignoring metric perturbations, $`\ddot{\delta \varphi }+3H\dot{\delta \varphi }`$ $`+`$ $`{\displaystyle \frac{1}{a^2}}^2\delta \varphi +3\lambda \varphi ^2\delta \varphi +3\lambda \varphi \delta \varphi ^2+\lambda \delta \varphi ^3`$ (31) $`+`$ $`g^2\delta \chi ^2\varphi +g^2\delta \chi ^2\delta \varphi =0.`$ (32) The terms in this equation describing the interaction between the $`\varphi `$ and $`\chi `$ fields become in momentum space $`{\displaystyle \frac{g^2\varphi }{(2\pi )^{3/2}}}{\displaystyle d^3k^{}\delta \chi _𝐤^{}\delta \chi _{𝐤𝐤^{}}}`$ (33) $`+`$ $`{\displaystyle \frac{g^2}{(2\pi )^3}}{\displaystyle d^3k^{}d^3k^{\prime \prime }\delta \chi _𝐤^{}\delta \chi _{𝐤^{\prime \prime }}\delta \varphi _{𝐤𝐤^{}𝐤^{\prime \prime }}}.`$ (34) Thus the Hartree term $`g^2\delta \chi ^2\delta \varphi _k`$ in Eq. (27) corresponds to the second term in expression (34), restricted to $`𝐤^{\prime \prime }=𝐤^{}`$. Physically, this means that only scattering events which do not change the $`\delta \varphi _k`$ momentum are included in the Hartree approximation, and “rescattering” events are ignored. It is important to notice that the first term in (34), which scatters particles from the homogeneous inflaton background into mode $`\delta \varphi _k`$, could be larger than the Hartree term since initially $`\|\varphi \|>\|\delta \varphi \|`$, unless the first term vanishes upon averaging (integrating) over the entire phase space of contributing terms (which is what is assumed in the Hartree approximation). If it does not vanish, the first term in (34) will act as a driving term for the $`\delta \varphi `$ modes in (32). Since some $`\delta \chi `$ modes experience parametric amplification with Floquet exponent $`\mu `$, this term will lead to an important second-order effect, namely the growth of $`\delta \varphi `$ as $`e^{2\mu x}`$. This effect is left out in the Hartree approximation. Because the metric perturbations are coupled to $`\delta \varphi `$ through Eq. (9), we also expect that, with the homogeneous $`\chi `$ set to zero, the Hartree approximation will miss the corresponding growth of $`\mathrm{\Phi }`$. Note, however, that by including the $`\chi `$ background term $`2g^2\chi \varphi \delta \chi `$ in (32), and setting $`\chi ^2\delta \chi ^2`$, we can approximate the effect of the important first term in (34), as we saw in Fig. 2. Just as with the scalar fields, metric fluctuations may grow rapidly in our model. We can account for the backreaction of metric perturbations through the effective energy-momentum tensor formalism of Abramo et al.. This involves expanding the Einstein equations to second order in the perturbations and taking the spatial average to obtain effective background equations. In our case, the metric and inflaton equations (23) and (24) become, with background $`\chi `$ set to zero, $`H^2`$ $`=`$ $`{\displaystyle \frac{8\pi }{3m_{\text{Pl}}^2}}[{\displaystyle \frac{1}{2}}\dot{\varphi }^2+V(\varphi )+{\displaystyle \frac{3}{2}}\lambda \varphi ^2\delta \varphi ^2+{\displaystyle \frac{g^2}{2}}\varphi ^2\delta \chi ^2`$ (37) $`+{\displaystyle \frac{1}{2}}{\displaystyle \underset{i}{}}(\dot{\delta \varphi }_i^2+{\displaystyle \frac{1}{a^2}}(\delta \varphi _i)^2)+2\lambda \varphi ^3\mathrm{\Phi }\delta \varphi ]`$ $`+4H\mathrm{\Phi }\dot{\mathrm{\Phi }}\dot{\mathrm{\Phi }}^2+{\displaystyle \frac{3}{a^2}}(\mathrm{\Phi })^2,`$ $`(\ddot{\varphi }`$ $`+`$ $`3H\dot{\varphi })(1+4\mathrm{\Phi }^2)+\lambda \varphi ^3+3\lambda \delta \varphi ^2\varphi `$ (38) $`+`$ $`g^2\delta \chi ^2\varphi 2\mathrm{\Phi }\ddot{\delta \varphi }4\dot{\mathrm{\Phi }}\dot{\delta \varphi }`$ (39) $``$ $`6H\mathrm{\Phi }\dot{\delta \varphi }+4\dot{\varphi }\dot{\mathrm{\Phi }}\mathrm{\Phi }{\displaystyle \frac{2}{a^2}}\mathrm{\Phi }^2\delta \varphi =0.`$ (40) ### B Analytical estimates #### 1 Evolution of perturbations during inflation Perturbations will grow during parametric resonance until backreaction becomes important. We can analytically estimate the amount of growth by estimating the time at which the Hartree term $`g^2\delta \chi ^2`$ is of the order of the background $`\lambda \varphi ^2`$ (cf Eq. (27)). Note that in the absence of metric fluctuations, such an estimate should be accurate, at least for $`g^2/\lambda 1`$, as nonlinear lattice simulations indicate . In addition, we expect the matter sector to dominate the dynamics. In order to estimate the variance $`\delta \chi ^2`$, we will need to calculate the evolution of $`\delta \chi _k`$ modes, starting from the adiabatic vacuum inside the Hubble radius, continuing through inflation, and finally through preheating. The evolution during inflation is quite complicated, and will have a crucial effect on the final variances, so we will describe the inflationary stage in some detail. We consider general values of $`g^2/\lambda `$, rather than just the supersymmetric point. We will only need to consider the contribution to $`\delta \chi ^2`$ from modes which are super-Hubble at the start of preheating. To see this, first note that for $`g^2/\lambda =2`$, the small-scale boundary of the strongest (and largest-scale) resonance band is at $`k_{\text{max}}/a\sqrt{\lambda }\varphi _0/2`$, where $`\varphi _0(t)`$ is the amplitude of inflaton oscillations during preheating . Next, we can use the Friedmann equation (5) to write the Hubble parameter in terms of $`\varphi _0`$, giving $$H^2=\frac{2\pi }{3m_{\text{Pl}}^2}\lambda \varphi _0^4.$$ (41) (Note that this equation also applies approximately during slow-roll.) Using the value $`\varphi _0=0.2m_{\text{Pl}}`$, we calculate the ratio $`aH/k_{\text{max}}0.6`$ at the start of preheating. Thus the Hubble radius corresponds closely to the smallest resonant scale. This result is not very sensitive to $`g^2/\lambda `$ as long as we are near the centre of a band, i.e. $`g^2/\lambda =2n^2`$, since $`k_{\text{max}}`$ increases only slowly with $`g^2/\lambda `$ in this case . Also, we can ignore the resonance bands at higher $`k`$ values, since they correspond to narrow resonance. To estimate the evolution of $`\delta \chi _k`$ on super-Hubble scales during inflation, we can ignore terms containing the background $`\chi `$ as well as the spatial gradient term in Eq. (10), resulting in a damped harmonic oscillator equation with time-dependent coefficients, $$\ddot{\delta \chi }_k+3H\dot{\delta \chi }_k+g^2\varphi ^2\delta \chi _k=0.$$ (42) During slow-roll, we can use the adiabatic approximation to find solutions to this equation, since $`\|\dot{H}\|H^2`$. Thus for $`g^2\varphi ^2>(3H/2)^2`$ we have underdamped oscillations with damping envelope $$\delta \chi _k\mathrm{exp}\left[(3H/2)𝑑t\right]=a^{3/2}.$$ (43) For $`g^2\varphi ^2<(3H/2)^2`$, we have the overdamped case with two decaying modes. Ignoring the more rapidly decaying mode, we obtain $$\delta \chi _k\mathrm{exp}\left[\left(3H/2\sqrt{9H^2/4g^2\varphi ^2}\right)𝑑t\right].$$ (44) In this case, the fluctuations are very slowly decaying in the massless limit $`g^2\varphi ^2(3H/2)^2`$, while they approach the $`a^{3/2}`$ decay as $`g^2\varphi ^2(3H/2)^2`$. During slow-roll we have $`H^2\varphi ^4`$ (see Eq. (41)), so that $`H^2`$ decreases more rapidly than $`g^2\varphi ^2`$, and there is a transition between the over- and underdamped stages. The two types of behaviour are separated by the critically damped case, $`g^2\varphi ^2=(3H/2)^2`$. Using Eqs. (41) and (17), we can write this critical damping condition in terms of the number of $`e`$-folds after critical damping, $`N_{\text{crit}}`$, as $$N_{\text{crit}}=\mathrm{ln}\left(\frac{a_\text{f}}{a_{\text{crit}}}\right)=\frac{2}{3}\frac{g^2}{\lambda }\mathrm{ln}\left(\frac{k_\text{f}}{k_{\text{crit}}}\right),$$ (45) where subscript “f” refers to the end of inflation and “crit” to the time of critical damping. Wavevectors $`k_{\text{crit}}`$ and $`k_\text{f}`$ leave the Hubble radius at $`t_{\text{crit}}`$ and $`t_\text{f}`$, respectively. We see that as $`g^2/\lambda `$ increases, cosmological scales are damped like $`a^{3/2}`$ during a greater and greater part of inflation. We thus expect that for large enough $`g^2/\lambda `$, the backreaction of the smaller-scale modes will terminate parametric resonance when cosmological-scale $`\delta \chi _k`$ modes are still greatly suppressed. In other words, there will be a maximum value of $`g^2/\lambda `$ for which there is significant amplification of super-Hubble $`\delta \chi _k`$ perturbations, as anticipated in . We first consider the evolution of the modes which leave the Hubble radius after $`t_{\text{crit}}`$, i.e. $`k>k_{\text{crit}}`$ (but which are still super-Hubble at the end of inflation, $`k<k_\text{f}`$). These modes are effectively massive during inflation, and hence we can simply use the adiabatic vacuum state, Eq. (18), which for $`kaH`$ gives $$\|\delta \chi _k(t_\text{f})\|^2=\frac{1}{2a_\text{f}^3g\varphi _\text{f}}.$$ (46) Note that if we define the spectral index $`n`$ through $`𝒫_\chi (k)k^{n1}`$ , then for this part of the spectrum we have $`n=4`$, an extreme blue tilt. Next, we will calculate the evolution of modes which leave the Hubble radius before $`t_{\text{crit}}`$, i.e. modes with $`k<k_{\text{crit}}`$. In this case, the modes are approximately massless when they exit the Hubble radius ($`g^2\varphi ^2<(3H/2)^2`$ for $`t<t_{\text{crit}}`$), so we can use the standard result for a massless inflaton , $$\|\delta \chi _k(t_k)\|^2=\frac{H^2(t_k)}{2k^3},$$ (47) where $`t_k`$ is the time that mode $`\delta \chi _k`$ exits the Hubble radius. We now must use Eq. (44) to evolve the modes during the overdamped period, $`t_k<t<t_{\text{crit}}`$. Writing $`dt=d\varphi /\dot{\varphi }`$, and using the slow-roll approximation $`\dot{\varphi }V_{,\varphi }/3H`$, we can perform the integral to obtain $$\|\delta \chi _k(t_{\text{crit}})\|^2=\frac{H^2(t_k)}{2k^3}e^{3F(N_k)},$$ (48) where $`N_k`$ is the number of $`e`$-folds after time $`t_k`$ and $`F(N_k)`$ $``$ $`N_kN_{\text{crit}}\sqrt{N_k}\sqrt{N_kN_{\text{crit}}}`$ (49) $`+`$ $`N_{\text{crit}}\mathrm{ln}\left({\displaystyle \frac{\sqrt{N_k}+\sqrt{N_kN_{\text{crit}}}}{\sqrt{N_{\text{crit}}}}}\right).`$ (50) Next we can readily propagate the modes through the underdamped period, $`t_{\text{crit}}<t<t_\text{f}`$, using Eqs. (43) and (45), giving $$\|\delta \chi _k(t_\text{f})\|^2=\frac{H^2(t_k)}{2k^3}e^{3F(N_k)2g^2/\lambda }.$$ (51) Since the damping term $`F(N_k)`$ is positive, we see as expected that large-scale modes are strongly damped for large $`g^2/\lambda `$. Finally, we can approximate the conformal time dependence of all super-Hubble modes during parametric resonance as $$\delta \chi _ke^{\mu _{\text{max}}x},$$ (52) if we are near the centre of a resonance band. This is valid since, in this case, the Floquet index $`\mu _k`$ varies only slightly for scales larger than a few times the Hubble radius (i.e. the smallest resonant scale) . #### 2 Variances and total resonant growth Now we can proceed to calculate the field variance, $`\delta \chi ^2`$. We will use Eq. (30), restricting the integral to the resonantly growing modes. We begin with the case $`g^2/\lambda =2`$. Equation (45) tells us that in this case $`N_{\text{crit}}4/3`$, so that essentially all of the evolution during inflation is in the overdamped regime, and we only need to consider modes with $`k<k_{\text{crit}}`$. The variance integral will be dominated by modes with $`N_kN_{\text{crit}}`$, so we may approximate the damping term in Eq. (50) as $$e^{3F(N_k)}\left(\frac{N_{\text{crit}}}{4N_k}\right)^{g^2/\lambda }.$$ (53) For the current case, $`g^2/\lambda =2`$, we can now combine the expression (51) with Eqs. (17), (41), (52), and (53) to obtain for the power spectrum on resonant scales at the end of preheating $$𝒫_\chi (k,t_\text{e})=\frac{\lambda m_{\text{Pl}}^2}{54\pi ^3}e^{2\mu _{\text{max}}x_\text{e}}.$$ (54) Here $`t_\text{e}`$ is the time that the resonance shuts down, and $`x_\text{e}`$ is the corresponding scaled conformal time. As we will see, the important thing about this result is that the power spectrum is essentially Harrison-Zel’dovich (independent of $`k`$), with spectral index $`n=1`$. We can next rewrite the variance integral, Eq. (30), in terms of the power spectrum as $$\delta \chi ^2(t_\text{e})=_0^{N_0}𝑑N_k𝒫_\chi (k,t_\text{e})=N_0𝒫_\chi (t_\text{e}),$$ (55) where $`N_050`$ is the total number of $`e`$-folds during inflation. Finally, the criterion $`g^2\delta \chi ^2(t_\text{e})\lambda \varphi ^2(t_\text{e})`$ gives, using the value $`\varphi (t_\text{e})10^2m_{\text{Pl}}`$, $$𝒫_\chi (t_\text{e})10^6m_{\text{Pl}}^2$$ (56) for the $`\delta \chi _k`$ power spectrum on cosmological scales at the end of preheating. Note that this result used only the $`k`$-independence of the power spectrum (which is a result of the special choice $`g^2/\lambda =2`$), and the values of $`N_0`$ and $`\varphi (t_\text{e})`$. In particular, the result is independent of $`\lambda `$, unless, contrary to our implicit assumption, $`\lambda `$ is so large that $`g^2\delta \chi ^2>\lambda \varphi ^2`$ already at the start of preheating. In this case, Eq. (56) will be an underestimate. According to the results from Section II B, we expect synchronization of the other fields to $`\delta \chi _k`$, so that in particular we expect $`𝒫_\mathrm{\Phi }𝒫_\chi /m_{\text{Pl}}^2`$. Therefore we conclude that, for $`g^2/\lambda =2`$, the metric perturbation amplitude will indeed be considerably larger than the COBE measured value, even including the effect of backreaction. Next we will repeat the preceding analysis for the second super-Hubble resonance band, at $`g^2/\lambda =8`$. In this case we have $`N_{\text{crit}}5`$, so we must consider modes that exit the Hubble radius both before and after $`t_{\text{crit}}`$. For the large-scale modes, $`k<k_{\text{crit}}`$, it will be sufficient to place an upper limit on the variance. Using Eq. (51), but ignoring the damping factor $`e^{3F}`$, we obtain $$𝒫_\chi (k,t_\text{f})<\frac{H^2(t_{\text{crit}})}{(2\pi )^2}e^{2g^2/\lambda }2\times 10^8\lambda m_{\text{Pl}}^2$$ (57) on scales $`k<k_{\text{crit}}`$ at the end of inflation. Thus, using Eq. (55), the contribution to the variance from modes with $`k<k_{\text{crit}}`$ satisfies the (probably very conservative) bound $$\delta \chi ^2(t_\text{f})_{k<k_{\text{crit}}}<9\times 10^7\lambda m_{\text{Pl}}^2.$$ (58) Next we can use Eq. (46) to calculate the contribution to $`\delta \chi ^2`$ from smaller-scale modes with $`k_{\text{crit}}<k<k_\text{f}`$, $`\delta \chi ^2(t_\text{f})_{k_{\text{crit}}<k<k_\text{f}}`$ $`=`$ $`{\displaystyle \frac{1}{16\pi ^3a_\text{f}^3g\varphi _\text{f}}}{\displaystyle _{k_{\text{crit}}}^{k_\text{f}}}d^3k`$ (59) $`=`$ $`{\displaystyle \frac{1}{18}}\sqrt{{\displaystyle \frac{2}{3\pi }}{\displaystyle \frac{\lambda }{g^2}}}{\displaystyle \frac{\lambda \varphi _\text{f}^5}{m_{\text{Pl}}^3}}`$ (60) $``$ $`3\times 10^6\lambda m_{\text{Pl}}^2.`$ (61) Here we have used $`k_{\text{crit}}^3k_\text{f}^3`$ (which follows from Eq. (45) for $`g^2/\lambda =8`$), the relation $`k_\text{f}/a_\text{f}=H(t_\text{f})`$, Eq. (41), and the value $`\varphi _\text{f}=0.2m_{\text{Pl}}`$. This value of the small-scale variance exceeds our upper limit on the large-scale variance in Eq. (58), so we can ignore the contribution from the large-scale modes, $`\delta \chi ^2(t_\text{f})_{k<k_{\text{crit}}}`$. Now we can again apply the condition $`g^2\delta \chi ^2(t_\text{e})\lambda \varphi ^2(t_\text{e})`$, which in this case gives $$e^{2\mu _{\text{max}}x_\text{e}}4\lambda ^1.$$ (62) Finally, we can use Eq. (51) without approximation to calculate the cosmological-scale power spectrum at the end of preheating, for the case $`g^2/\lambda =8`$, $`𝒫_\chi (t_\text{e})`$ $`=`$ $`{\displaystyle \frac{H^2(t_0)}{(2\pi )^2}}\mathrm{exp}\left[3F(N_0)2{\displaystyle \frac{g^2}{\lambda }}+2\mu _{\text{max}}x_\text{e}\right]`$ (63) $``$ $`10^{14}m_{\text{Pl}}^2.`$ (64) In this case the growth stops before the cosmological perturbations exceed the COBE value, and thus parametric resonance does not change the standard predictions for the size of the fluctuations. Therefore, since the damping of super-Hubble $`\delta \chi _k`$ modes increases as $`g^2/\lambda `$ increases, the standard predictions are not modified for all resonance bands beyond the first, i.e. for $`g^2/\lambda 8`$. To close this section, consider the behaviour of the large-scale variance if we suppose that inflation started much earlier than the time that cosmological scales left the Hubble radius, i.e. let $`N_0\mathrm{}`$. In this limit we see from the approximation Eq. (53), the expression $`H(t_k)N_k`$, and Eq. (51) that the variance becomes $$\delta \chi ^2_{k<k_{\text{crit}}}_{N_{\text{crit}}}^{N_0}𝑑N_kN_k^{2g^2/\lambda },$$ (65) which, for $`g^2/\lambda =2`$, is divergent as $`N_0\mathrm{}`$. Such divergences are well-known in inflationary models . Here the divergence suggests that for large enough $`N_0`$, the growth of cosmological-scale modes will stop before they exceed the COBE amplitude, due to the large contribution to the variance from super-cosmological scales. Indeed, for $`N_010^6`$ (a value not out of the question in chaotic inflation ) Eq. (55) gives $`𝒫_\chi (t_\text{e})10^{10}m_{\text{Pl}}^2`$ for $`g^2/\lambda =2`$. For $`g^2/\lambda >3`$ the variance converges, although in this case cosmological-scale modes are already supressed by the mechanism described above. ### C Numerical Results It is straightforward to check our analytical estimates from the previous section by numerically integrating the coupled set of Hartree approximation evolution equations (23) – (30) and metric perturbation equation (9). We now must evolve a set of modes that fill the relevant resonance band. For example, for $`g^2/\lambda =2`$, the first resonance band extends from $`\kappa =0`$ to $`\kappa =0.5`$ . Again we begin each mode’s evolution inside the Hubble radius during inflation, using the initial vacuum state, Eqs. (18) and (19). Each mode is incorporated into the calculation shortly before it leaves the Hubble radius, so that the spatial gradient terms are never too large. The variances are calculated by performing the discretized integrals, Eqs. (30), only over the resonance band; thus they are convergent. Note that the variances are calculated simultaneously with the field backgrounds and perturbations. In Fig. 3 we present the evolution of the $`\delta \chi _k`$, $`\delta \varphi _k`$, and $`\mathrm{\Phi }_k`$ power spectra on the same cosmological scale as was studied in Section II. All parameters are the same as for Fig. 2, except here we use for the initial background value $`\chi (t_0)=10^6m_{\text{Pl}}`$, which means that during preheating $`\chi ^2\delta \chi ^2`$. The evolution is initially similar to that of Fig. 2, only here the growth saturates at $`𝒫_\chi 3\times 10^7m_{\text{Pl}}^2`$, in good agreement with our prediction based on Eq. (56). Also, as expected, the other fields closely follow $`𝒫_\chi `$. Whereas in the linear calculations the Einstein constraint equation (11) was satisfied to extremely good accuracy, with the inclusion of backreaction $`𝒫_\mathrm{\Phi }`$ saturates at a factor of roughly $`10^3`$ higher using Eq. (11) than the illustrated result, which used Eq. (9). Note that a similar result was found in . We suspect that this is a fundamental problem related to our attempt to capture some of the nonlinear dynamics with the Hartree approximation. Regardless of which value is used, the cosmological metric perturbations considerably exceed the COBE normalisation. In addition, we find no significant difference in the results when backreaction of metric perturbations is included using Eqs. (37) and (40), as expected if the matter fields dominate the backreaction. Thus all of our presented results exclude the metric backreaction terms. As discussed above, larger values of $`g^2/\lambda `$ result in increased damping of $`\delta \chi _k`$ on large scales during inflation, and at large enough $`g^2/\lambda `$ we expect insignificant amplification of super-Hubble modes. This is illustrated in Fig. 4. Here we examine the second resonance band at $`g^2/\lambda =8`$, but use otherwise identical parameters to Fig. 3. Resonance stops at $`𝒫_\chi 10^{14}m_{\text{Pl}}^2`$, consistent with our analytical estimate from Eq. (64), and not exceeding the standard predictions for $`\lambda 10^{14}`$ . Note that the small rise in $`𝒫_\mathrm{\Phi }`$ at late times should not be trusted, as our Hartree approximation scheme will not capture the full nonlinear behaviour. For resonance bands at even higher $`g^2/\lambda `$, we find extremely suppressed cosmological $`\delta \chi _k`$ amplitudes, in quantitative agreement with the calculations of the previous section. ## IV Self-interacting $`\chi `$ models ### A Positive coupling We now consider the addition of a quartic self-interaction term for the $`\chi `$ field, so that our potential becomes $$V(\varphi ,\chi )=\frac{\lambda }{4}\varphi ^4+\frac{g^2}{2}\varphi ^2\chi ^2+\frac{\lambda _\chi }{4}\chi ^4,$$ (66) with $`g^2>0`$. The significance of such a term for parametric resonance was studied in lattice simulations and analytically , but in the absence of metric perturbations. Bassett and Viniegra included metric perturbations, but ignored backreaction. Essentially, for $`\lambda _\chi \lambda `$ we expect the $`\chi `$ self-interaction to limit the growth of perturbations as compared with the $`\lambda _\chi =0`$ case studied above, due to the presence of the “potential wall” $`(\lambda _\chi /4)\chi ^4`$. More precisely, the linearized equation of motion for the $`\chi `$ field perturbation becomes, with $`\chi `$ self-interaction but ignoring metric perturbations, $`\ddot{\delta \chi }_k+3H\dot{\delta \chi }_k`$ $`+`$ $`\left({\displaystyle \frac{k^2}{a^2}}+3\lambda _\chi \chi ^2+g^2\varphi ^2\right)\delta \chi _k`$ (67) $`+`$ $`2g^2\varphi \chi \delta \varphi _k=0.`$ (68) Thus for small enough initial $`\chi `$ background, the initial behaviour of the modes will be essentially unchanged from the $`\lambda _\chi =0`$ case. However, when the $`\chi `$ background grows to the point that $`\chi ^2/\varphi ^2g^2/\lambda _\chi `$, the analytical parametric resonance theory of Section II A no longer applies, and we may expect the perturbations to stop growing. Since, as discussed above, for the significant production of super-Hubble modes we require $`g^2\lambda `$, we expect that $`\chi `$ self-interaction will shut down the resonance when $`\chi ^2/\varphi ^2\lambda /\lambda _\chi `$, as long as $`\lambda _\chi \lambda `$. If $`\lambda _\chi <\lambda `$, then the $`\chi ^4`$ interaction term will not lead to a shutdown of the resonance since (based on our numerical simulations) the homogeneous $`\chi `$ field never substantially exceeds the value of the inflaton background. We have confirmed this expectation numerically, and we give an example of our results in Fig. 5. Here we have included Hartree backreaction and metric perturbations, and used coupling constant values $`\lambda =10^{14}`$, $`g^2/\lambda =2`$, and $`\lambda _\chi =10^{10}`$, and initial backgrounds $`\varphi (t_0)=4m_{\text{Pl}}`$ and $`\chi (t_0)=10^6m_{\text{Pl}}`$. We indeed observe the termination of the super-Hubble modes’ growth at approximately the time when $`\chi ^2/\varphi ^2=\lambda /\lambda _\chi `$. Note that, in the absence of backreaction, Bassett and Viniegra observed a continued slow growth of super-Hubble perturbations after the initial termination of the resonance when $`\chi ^2/\varphi ^2\lambda /\lambda _\chi `$ . We confirmed this result; however, we note that when we include the backreaction term $`3\lambda _\chi \delta \chi ^2`$ in the evolution equations, we expect backreaction to become important also at the time that $`\chi ^2/\varphi ^2\lambda /\lambda _\chi `$, with our choice $`\chi ^2\delta \chi ^2`$. Hence, as seen in Fig. 5, the slow growth is completely suppressed. ### B Negative coupling The presence of $`\chi `$ self-coupling means that we no longer require $`g^20`$ for global stability. In fact, for the case $`g^2<0`$, the potential will be bounded from below for $`\lambda \lambda _\chi /g^4>1`$ . This negative coupling case was studied in the absence of metric perturbations using lattice simulations in , and without backreaction in . The behaviour of the fields is qualitatively different in the negative and positive coupling cases. For $`g^2<0`$, potential minima exist with non-zero homogeneous part of the $`\chi `$ field. Thus, assuming the fields fall into these minima, the problem of choice of $`\chi `$ background discussed in previous sections for the positive coupling case is alleviated. For initial homogeneous $`\chi `$ fields large enough ($`\chi (t_0)m_{\text{Pl}}`$), we find numerically that the fields fall into the potential minimum by the end of inflation, and the two fields subsequently evolve in step during preheating. This effectively reduces the system to a single-field system, and hence no resonance is possible on super-Hubble scales. To see this explicitly consider the case $`\lambda _\chi =\lambda `$, for which the symmetry of the potential requires the potential minima to lie along $`\chi ^2=\varphi ^2`$. If we choose the same initial signs for $`\chi `$ and $`\varphi `$, then during preheating the backgrounds lie along the attractor $`\varphi =\chi `$. Similarly, since the behaviour of super-Hubble modes is essentially the same as that of the backgrounds, we have $`\delta \chi =\delta \varphi `$ during preheating. Then the perturbation equation (68) becomes $$\ddot{\delta \chi }_k+3H\dot{\delta \chi }_k+\left[\frac{k^2}{a^2}+3(\lambda +g^2)\varphi ^2\right]\delta \chi _k=0.$$ (69) Thus the effective mass of the $`\delta \chi `$ oscillations is precisely three times the effective mass of the background inflaton oscillations (cf Eq. (3)), so that just as with the case of the inflaton perturbations in Eq. (15), there will be no resonance on super-Hubble scales for all allowed values of $`g^2`$. We have confirmed this numerically; indeed more generally, as long as initially $`\chi (t_0)m_{\text{Pl}}`$ but for any $`\lambda _\chi \lambda `$, the two fields will be proportional during preheating and no super-Hubble resonance will result. This result assumes that during preheating only the “field” $`\delta \chi +\delta \varphi `$ is excited. If orthogonal field excitations $`\delta \chi \delta \varphi `$ are present, they can grow resonantly. The effective squared mass of $`\delta \chi \delta \varphi `$ excitations is $`(3\lambda g^2)\varphi ^2`$, so that according to the analytical parametric resonance theory of Section II A, super-Hubble resonance will occur near $`3\lambda g^2=2n^2(\lambda +g^2)`$, for integral $`n`$ (we require $`n2`$ for negative $`g^2`$). That is, super-Hubble $`\delta \chi \delta \varphi `$ modes will grow for $`g^2\lambda (32n^2)/(1+2n^2)`$. However, numerically we observe only extremely small components $`\delta \chi \delta \varphi `$ by the end of inflation, so their growth is substantially delayed. On the other hand, for small initial homogeneous part $`\chi (t_0)m_{\text{Pl}}`$, we find that the potential minima are not reached by the end of inflation, and the two fields evolve in a very complicated manner during preheating. The analytical theory of parametric resonance cannot be applied, but numerically we do find roughly exponential growth of super-Hubble modes in this case, as found in . The growth rate increases as $`g^2`$ decreases towards the value at which global instability sets in, $`g^2=\sqrt{\lambda \lambda _\chi }`$. We have illustrated this case in Fig. 6, using the parameter values $`\lambda =10^{14}`$, $`g^2=0.5\lambda `$, $`\lambda _\chi =\lambda `$, $`\varphi (t_0)=4m_{\text{Pl}}`$, and $`\chi (t_0)=10^6m_{\text{Pl}}`$. Here the growth rates and final power spectra values are comparable to the $`\lambda _\chi =0`$ case of Fig. 3, though the $`\delta \chi `$ field is not damped during inflation for negative coupling. For $`\lambda _\chi >\lambda `$, the growth is terminated early, just as in the positive coupling case. ## V Summary and Discussion In this paper we have studied backreaction effects on the growth of super-Hubble cosmological fluctuations in a specific class of two field models with a massless inflaton $`\varphi `$ coupled to a scalar field $`\chi `$. Our study was based on the Hartree approximation. For the non-self-coupled $`\chi `$ field case, we found that backreaction has a crucial effect in determining the final amplitude of fluctuations after preheating. For values of the coupling constants satisfying $`g^2/\lambda =2`$ (the ratio predicted in supersymmetric models), the predicted amplitude of the super-Hubble metric perturbations at the end of preheating is too large to be consistent with the COBE normalization, thus apparently ruling out such models. One possible loophole is the backreaction contribution from super-cosmological scale fluctuations. For sufficiently long periods of inflation, the predicted amplitude can be consistent with the COBE normalization. In addition, the final amplitude of the fluctuation spectrum is independent of the coupling constant $`\lambda `$. Note that the growth of inflaton fluctuations $`\delta \varphi _k`$ (and hence metric perturbations $`\mathrm{\Phi }_k`$) occurs in these models either through coupling to $`\delta \chi _k`$ via a homogeneous background $`\chi `$ field or through nonlinear evolution effects. The situation for $`g^2/\lambda 1`$ is very similar to the previously studied case of a massive inflaton in the broad resonance regime . Cosmological-scale $`\delta \chi _k`$ modes are significantly damped during inflation, and the end of resonant growth is determined by the growing small-scale modes. Already for the second resonance band (centred at $`g^2/\lambda =8`$) cosmological metric perturbations are not amplified above the COBE normalization value. This implies that preheating does not alter the standard predictions for the $`\mathrm{\Phi }_k`$ normalization in $`(\lambda /4)\varphi ^4`$ inflation for the second and all higher resonance bands. The important difference between the model we have studied and the massive inflaton case is that, in the massive model, weak super-Hubble suppression at small $`g^2`$ is accompanied by weak resonant growth during preheating , so that no significant super-Hubble amplification is possible. The inclusion of $`\chi `$ field self-interaction alters the evolution in a predictable way: the resonant growth stops when $`\chi ^2/\varphi ^2\lambda /\lambda _\chi `$, as long as $`\lambda _\chi \lambda `$. This means that we are unable to rule out models (on the basis of a too large production of metric perturbations) with $`\lambda _\chi /\lambda 10^4`$. In the negative coupling case, there are two possibilities. For large initial $`\chi `$ backgrounds, $`\chi (t_0)m_{\text{Pl}}`$, the system becomes essentially single-field, and no resonance occurs (at least until late times). For small initial $`\chi `$, exponential growth occurs for large enough allowed $`\|g^2\|`$. The Hartree approximation provides a useful approach for the inclusion of the effects of backreaction. However, as mentioned in Section III, this approximation misses terms which could contribute to the evolution of fluctuations in an important way. We believe, nevertheless, that our results are sufficiently accurate to predict, for the models studied, whether or not the metric perturbation amplitude after preheating is consistent with the COBE measurement. Nonlinear effects, or rescattering, will primarily affect the detailed evolution of matter fields after backreaction is important. Still, it is of great interest to extend our analysis to a full nonlinear treatment, as was done in the absence of gravitational fluctuations in , and including metric fluctuations in for single field models. ###### Acknowledgements. We wish to thank F. Finelli for useful discussions. This research was supported by the Natural Sciences and Engineering Research Council of Canada. The work of R.B. was supported in part by the U.S. Department of Energy under Contract DE-FG02-91ER40688, TASK A.
warning/0007/math0007002.html	ar5iv	text	# Relation between the dimensions of the ring generated by a vector bundle of degree zero on an elliptic curve and a torsor trivializing this bundle ## 1. Introduction and Notations Let $`X`$ be a complete, connected, reduced scheme over a perfect field $`k`$. We define Vect$`(X)`$ to be the set of isomorphism classes $`[V]`$ of vector bundles $`V`$ on $`X`$. We can define an addition and a multiplication on Vect$`(X)`$: $`[V]+[V^{}]=[VV^{}]`$ $`[V][V^{}]=[VV^{}].`$ The (naive) Grothendieck ring $`K(X)`$ (see ) is the ring associated to the additive monoid Vect$`(X)`$, that means $$K(X)=\frac{[\text{Vect}(X)]}{H},$$ where $`H`$ is the subgroup of $`[\text{Vect}(X)]`$ generated by all elements of the form $`[VV^{}][V][V^{}]`$. The indecomposable vector bundles on $`X`$ form a free basis of $`K(X)`$. Since H$`{}_{}{}^{0}(X,\text{End}(V))`$ is finite dimensional, the Krull-Schmidt theorem () holds on $`X`$. This means that a decomposition of a vector bundle in indecomposable components exists and is unique up to isomophism. We want to generalize a theorem of M. Nori on finite vector bundles. A vector bundle $`V`$ on $`X`$ is called finite, if the collection $`S(V)`$ of all indecomposable components of $`V^n`$ for all integers $`n`$ is finite. In the following, we denote by R(V) the $``$-subalgebra of $`K(X)_{}`$ generated by the set $`S(V)`$. Thus a vector bundle $`V`$ is finite if and only if the ring $`R(V)`$ is of Krull dimension zero. In , Nori proves the following theorem: For every finite vector bundle $`V`$ on $`X`$ there exists a finite group scheme $`G`$ and a principal $`G`$-bundle $`\pi :PX`$, such that $`\pi ^{}V`$ is trivial on $`P`$. In particular, the equality $$dimR(V)=dimG(=0)$$ holds. As every vector bundle $`V`$ on $`X`$ of rank $`r`$ trivializes on its associated principal GL($`r`$)-bundle, we can look for a group scheme $`G`$ of smallest dimension and a principal $`G`$-bundle on which the pullback of the vector bundle $`V`$ is trivial. We might also compare the dimension of the group scheme to dim $`R(V)`$. In this article we consider the family of vector bundles of degree zero on an elliptic curve. We will prove in propositions 2 and 3 that they trivialize on a principal $`G`$-bundle with $`G`$ a group scheme of smallest dimension one. As in the situation of Nori’s theorem, this dimension turns out to be equal to the dimension of the ring $`R(V)`$. I am grateful to Hélène Esnault for suggesting the problem treated here and for many useful discussions. ## 2. Dimension relation for vector bundles of degree zero on an elliptic curve Let $`X`$ be an elliptic curve over an algebraically closed field $`k`$ of characteristic zero. We consider vector bundles of degree zero on $`X`$ which can be classified according to Atiyah (see ). By $`(r,0)`$ we denote the set of indecomposable vector bundles of rank $`r`$ and degree zero. ###### Theorem 1. (Atiyah ) 1. There exists a vector bundle $`F_r(r,0)`$, unique up to isomorphism, with $`\mathrm{\Gamma }(X,F_r)0`$. Moreover we have an exact sequence $$\begin{array}{ccccccccc}0& & 𝒪_X& & F_r& & F_{r1}& & 0.\end{array}$$ 2. Let $`E(r,0)`$, then $`ELF_r`$ where $`L`$ is a line bundle of degree zero, unique up to isomorphism (and such that $`L^rdetE`$.) ###### Proposition 2. 1. The $``$-subalgebra $`R(F_r)`$ of $`K(X)_{}`$ generated by $`S(F_r)`$ is $`[x]`$, where $`x=[F_2]`$, if $`r`$ is even, and $`x=[F_3]`$, if $`r`$ is odd. In particular, $`R(F_r)`$ is of Krull dimension zero. 2. There exists a principal $`𝔾_a`$-bundle $`\pi :PX`$ such that $`\pi ^{}(F_r)`$ is trivial for all $`r2`$. Remark: As in Nori’s case we have a correspondence of dimensions $$\text{dim }R(F_r)=\text{dim }𝔾_a=1.$$ Proof: As proved by Atiyah in , the vector bundles $`F_r`$ are self-dual and fulfill the formula $$F_rF_s=F_{rs+1}F_{rs+3}\mathrm{}F_{(rs)+(2s1)}\text{ for }sr.$$ For even $`r`$, it follows by induction that there exist integers $`a_i(n)`$ such that $$F_r^n=a_2(n)F_2a_4(n)F_4\mathrm{}a_{(r1)n1}(n)F_{(r1)n1}F_{(r1)n+1}$$ for odd $`n3`$, and $$F_r^n=a_1(n)𝒪_Xa_3(n)F_3\mathrm{}a_{(r1)n1}(n)F_{(r1)n1}F_{(r1)n+1}$$ for even $`n2`$ . Therefore we obtain $$S(F_r)=\{F_i\|i=1,2,3,\mathrm{}\},\text{ if }r\text{ even },$$ and $`S(F_r)`$ generates the subring $`[F_2]`$ of $`K(X)`$, because inductively we can write every vector bundle $`F_i`$ as $`p(F_2)`$ for some polynomial $`p[x]`$. For odd $`r`$, Atiyah’s multiplication formula gives $$F_r^n=a_1(n)𝒪_Xa_3(n)F_3\mathrm{}a_{(r1)n1}(n)F_{(r1)n1}F_{(r1)n+1}$$ for all $`n2`$. It follows that $$S(F_r)=\{F_i\|i\text{ odd }\},\text{ if }r\text{ odd }.$$ For odd $`r`$, the set $`S(F_r)`$ generates the ring $`R(F_r)=[F_3]`$, as for odd $`i`$ each $`F_i`$ is $`p(F_3)`$ for a polynomial $`p[x]`$. The vector bundle $`F_2`$ is an element of $`H^1(X,\text{GL}(2,𝒪))`$. Because of the exact sequence $$\begin{array}{ccccccccc}0& & 𝒪_X& & F_2& & 𝒪_X& & 0,\end{array}$$ $`F_2`$ is even an element of $`H^1(X,𝔾_a)`$. Here we embed $`𝔾_a`$ into GL($`2,𝒪`$) via $`u\left(\begin{array}{cc}1& u\\ 0& 1\end{array}\right)`$. Hence $`F_2`$ trivializes on a principal $`𝔾_a`$-bundle. As $`F_r=S^{r1}F_2,r3,`$ each $`F_r`$ trivializes on the same principal $`𝔾_a`$-bundle as $`F_2`$. As the classes $`[F_r]`$ are not torsion elements in $`H^1(X,\text{GL}(2,𝒪))`$, none of the bundles $`F_r`$ can trivialize on a principal $`G`$-bundle with $`G`$ a finite group scheme. ∎ Remark: In the given examples of vector bundles $`E`$ there was so far not only a correspondence of the dimensions of the group scheme and the ring $`R(E)`$. The algebra $`R(E)`$ was also the Hopf algebra corresponding to the group scheme. The following proposition shows that this is not true in general. ###### Proposition 3. Let $`ELF_r(r,0)`$ (see theorem 1). 1. If $`L`$ is not torsion, the ring $`R(E)`$ is isomorphic to $`[x,x^1][y]`$ and $`E`$ trivializes on a principal $`𝔾_m\times 𝔾_a`$-bundle. 2. If $`L`$ is torsion, let $`n`$, $`n1`$, be the minimal number such that $`L^n𝒪_X`$. If $`n`$ and $`r`$ are both even, the ring $`R(E)`$ is isomorphic to $$[x]/<x^{n/2}1>[y]$$ and $`E`$ trivializes on a principal $`\mu _n\times 𝔾_a`$-bundle. There is no principal $`\mu _{n/2}\times 𝔾_a`$-bundle where $`E`$ is trivial. If $`n`$ and $`r`$ are not both even, the ring $`R(E)`$ is isomorphic to $$[x]/<x^n1>[y]$$ and $`E`$ trivializes on a principal $`\mu _n\times 𝔾_a`$-bundle. Proof: Let $`E(r,0)`$ with $`\mathrm{\Gamma }(X,E)=0`$. (If $`\mathrm{\Gamma }(X,E)0`$, then $`EF_r`$. This case was already dealt with in proposition 2.) First we consider the case that $`L`$ is not torsion. We must distinguish between odd and even $`r`$. For odd $`r`$, Atiyah’s multiplication formula ( see proof of proposition 4) gives the following result: For $`m`$, $`m2`$, the tensor power $`E^mL^mF_r^m`$ has the indecomposable components $`L^m𝒪_X,L^mF_3,\mathrm{},L^mF_{(r1)m+1}`$, the tensor power $`E^mL^mF_r^m`$ has the indecomposable components $`L^m𝒪_X,L^mF_3,\mathrm{},L^mF_{(r1)m+1}`$. Thus we obtain $$S(E)=\left\{\begin{array}{c}𝒪_X,LF_r,L^1F_r,\hfill \\ L^{\pm i}F_3,L^{\pm i}F_5,\mathrm{},L^{\pm i}F_{(r1)i+1},\text{ i }\hfill \end{array}\right\}.$$ The algebra $`R(E)`$ which is generated by $`S(E)`$ is the subalgebra of $`K(X)_{}`$ generated by $`L`$, $`L^1`$ and$`F_3`$, thus $$R(E)=[L,L^1]_{}[F_3].$$ For even $`r`$, a similar computation gives that $$S(E)=\left\{\begin{array}{c}𝒪_X,LF_r,L^1F_r,\hfill \\ L^{\pm 2i},L^{\pm 2i}F_3,\mathrm{},L^{\pm 2i}F_{(r1)2i+1},\text{ i }\hfill \\ L^{\pm (2i+1)}F_2,L^{\pm (2i+1)}F_4,\mathrm{},\hfill \\ \text{ }L^{\pm (2i+1)}F_{(r1)(2i+1)+1},\text{ i }\hfill \end{array}\right\}.$$ The ring $`R(E)`$, generated by $`S(E)`$, is the subring of $`K(X)_{}`$ which is generated by the elements $`L^2`$, $`L^2`$, $`L^1F_2`$, therefore $$R(E)=[L^2,L^2]_{}[L^1F_2].$$ If $`L`$ is not a torsion bundle, it is clear that $`L`$ trivializes on a principal $`𝔾_m`$-bundle $`P_L`$. The vector bundle $`ELF_2`$ trivializes on the $`𝔾_m\times 𝔾_a`$-bundle $`P_L\times _XP`$, where $`P`$ is the principal $`𝔾_a`$-bundle from proposition 2, where $`F_2`$ and hence all the $`F_r`$ trivialize. Let now L be torsion and $`n`$, $`n2`$, the minimal number with $`L^n𝒪_X`$. As the $`F_r`$ are selfdual and $`L^{n1}=L^1`$, it suffices to consider positive tensor powers. Again we compute the tensor powers using Atiyah’s formula to find the indecomposable components. If $`r`$ is even and $`n`$ is odd, the set $`S(E)`$ contains the following bundles: $$S(E)=\{𝒪_X,L^iF_j\|i=0,1,\mathrm{},n1,j\}.$$ With the help of the multiplication formula for $`F_2`$ it is easy to show that all elements of $`S(E)`$ can be generated by $`L`$ and $`F_2`$. In additon, the relation $`L^n𝒪_X`$ holds. Hence we obtain $$R(E)=\frac{[L]}{<L^n1>}_{}[F_2].$$ If $`r`$ is odd and $`n`$ is even or odd, the result is $$S(E)=\{L^iF_j\|i=0,1,\mathrm{},n1,j\text{ odd}\}.$$ The bundles $`L`$ and $`F_3`$ are in $`S(E)`$ and generate all elements of $`S(E)`$. Because of the relation $`L^n𝒪_X`$, the algebra $`R(E)`$ is $$R(E)=\frac{[L]}{<L^n1>}_{}[F_3].$$ If $`r`$ and $`n`$ are both even $$S(E)=\{L^{2i}F_{2j1},L^{2i+1}F_{2j}\|i=0,1,\mathrm{},n/2,j\}.$$ The algebra R(E) is generated by $`L^2`$ and $`LF_2`$. The generators are subject to the relation $`L^n𝒪_X`$, thus $$R(E)=\frac{[L^2]}{<(L^2)^m1>}[LF_2],$$ where $`m=n/2`$. Recall that $`n2`$ is the minimal number such that $`L^n𝒪_X`$. Thus the bundle $`L`$ trivializes on a $`\mu _n`$-bundle $`P_L`$ and not on a $`\mu _m`$-torsor for $`m<n`$. The bundle $`ELF_r`$ then trivializes on the $`\mu _n\times 𝔾_a`$-bundle $`P_L\times _XP`$, where $`P`$ is again the principal $`𝔾_a`$-bundle from proposition 2. We will now show that the bundle $`E`$ does not trivialize on a $`\mu _{n/2}\times 𝔾_a`$-bundle: If $`ELF_r`$ trivializes on $`Q\times _XP`$, where $`Q`$ is a $`\mu _m`$-torsor and $`P`$ a $`𝔾_a`$-torsor, then det$`(LF_r)=L`$ is the identity element in the group Pic($`Q\times _XP`$). But one has Pic($`Q\times _XP)=`$Pic$`(Q)`$ by homotopy invariance. Thus $`L`$ must trivialize on the $`\mu _m`$-torsor $`Q`$, which is impossible for $`m<n`$. ∎ Remark: The correspondence between the dimension of the “minimal” group scheme and the dimension of the ring $`R(E)`$ also occurs in the case of vector bundles on the projective line, as one easily sees. Let $`X`$ be the complex projective line $`^1`$ and $`E:=𝒪(a)`$ a line bundle. If $`a=0`$ we have $`S(E)=\{𝒪\}`$ and $`R(E)=Q`$. We define the group scheme $`G`$ to be $`G=\text{Spec }`$ and the trivializing torsor is simply $`^1`$. If $`a0`$ we can easily compute that $`S(E)=\{𝒪(\lambda a)\|\lambda \}`$ and $`R(E)=[x,x^1]`$. We define the group scheme to be $`G=𝔾_m=\text{Spec }[x,x^1]`$. The given line bundle $`E`$ trivializes on a principal $`𝔾_m`$-bundle $`P_a`$, which depends on $`a`$. Thus we get the correspondence of $`dimR(E)`$ and $`dimG`$ in the case of a line bundle on $`^1`$. This computation can easily be generalized to the case of vector bundles of higher rank. We illustrate this for bundles of rank two. Let now $`E`$ be a vector bundle of rank 2 on $`^1`$, $`E=𝒪(a)𝒪(b)`$. The case $`(a,b)=(0,0)`$ is trivial. We can see at once that $`S(E)=\{𝒪\}`$ and therefore $`R(E)=`$. The vector bundle $`E`$ trivializes on the principal Spec $``$ \- bundle $`^1`$. If $`(a,b)(0,0)`$ the computation gives that $`S(𝒪(a)𝒪(b))=S(𝒪(c))`$, where $`c=(a,b)`$ (with $`(a,0)=a`$ and $`(0,b)=b`$) and therefore $`R(E)=[x,x^1]`$. $`E`$ trivializes on the principal $`𝔾_m`$-bundle $`P_c`$ that belongs to $`𝒪(c)`$ as $`𝒪(a)=𝒪(c)^\lambda `$ and $`𝒪(b)=𝒪(c)^\mu `$ for appropriate integers $`\lambda `$ and $`\mu `$.
warning/0007/cond-mat0007377.html	ar5iv	text	# 1 Top panel shows differential susceptibility for x=0.125 doping during warming. Bottom panel shows the same during cooling. The first order nature of the CO transition can be seen as a clear difference in the response between warming and cooling. Inset shows a magnified view of the structual transition. rf SUSCEPTIBILITY OF $`La_{1x}Sr_xMnO_3`$ SINGLE CRYSTALS: MAGNETIC SIGNATURES OF STRUCTURAL CHANGES P.V. PARIMI, H. SRIKANTH, M. BAILLEUL, and S. SRIDHAR, Department of Physics, Northeastern University, Boston, MA 02115. R. SURYANARAYANAN, L. PINSARD and A. REVCOLEVSCHI. Laboratoire de Chimie des Solides, UA 446, Universitě Paris-Sud, Orsay 91405, France. ABSTRACT A sensitive tunnel diode oscillator (TDO) operating at $`4MHz`$ is used to probe the dynamic response of $`La_{1x}Sr_xMnO_3`$ single crystals for $`x=0.125,0.175,0.28`$ and $`0.33`$ doping. Systematics of the measured change in reactance $`(\mathrm{\Delta }X)`$ as a function of temperature $`(30K<T<320K)`$ and DC magnetic field $`(0<H<6kOe)`$ reveal distinct temperature and field scales associated with the dynamic response of spin. It is notable that these features are far more striking than the corresponding features in static measurements. The results are discussed in the context of structural changes leading to polaron ordering. INTRODUCTION The perovskite oxides of the form $`\mathrm{R}e_{1x}A_xMnO_3`$ (where $`\mathrm{R}e`$ is a rare earth such as $`La`$ and $`A`$ is a divalent element such as $`Sr`$ or $`Ca`$) have generated considerable interest in recent times because of the discovery of the colossal magnetoresistance (CMR) effect . The CMR is a direct consequence of an unusual paramagnetic insulator (PMI) to ferromagnetic metal (FMM) transition driven mainly by the double exchange mechanism . However, double exchange alone cannot describe the complete phase diagram of the manganites and it has been pointed out that the interplay of strong electron-phonon coupling and double exchange is required to understand the existence of the high temperature insulating phase, the CMR effect and its sensitivity to magnetic field . The deficiency of the double exchange model is the fact that it does not consider spin-lattice or charge-lattice interactions, namely, Jahn-Teller interactions and polarons . Experimental results clearly suggest that lattice contributions are important for a thorough understanding of manganites. Besides MI transition, charge ordering (CO) is one of the characteristic phenomena observed in these materials especially in the low doping regime. CO and stripe correlations of concentrated holes and spins have attracted much attention in recent times, particularly due to their possible role in high T<sub>c</sub> superconductivity. A variety of experiments including structural , transport and thermal measurements have revealed novel features in the $`\mathrm{R}e_{1x}A_xMnO_3`$ directly associated with the interplay between structural, electronic and magnetic properties. Most of the experiments on manganites have been static and there have been relatively few experiments which probe the dynamic response of these systems. Dynamic experiments are likely to provide significant information about the collective response of spin and charge to the oscillating electric and magnetic fields impressed on the materials. In the present work rf dynamic response of La<sub>1-x</sub>Sr<sub>x</sub>MnO<sub>3</sub> for concentrations x=0.125, 0.175,0.28 and 0.33 are reported. We focus on the interplay of between holes and lattice distortions to understand the relation between the magnetic and structural properties. EXPERIMENT Single crystals of La<sub>1-x</sub>Sr<sub>x</sub>MnO<sub>3</sub> were grown using an image furnace technique . Samples used in these measurements had cylindrical disk like shapes with diameter $`5mm`$ and thickness $`2mm`$ with polished surfaces and edges. The rf experiments were performed using a tunnel diode oscillator (TDO) which has very high sensitivity in measuring the electro- and magneto-dynamic properties of materials. The crystal is placed inside a copper coil which forms part of an $`LC`$-tank circuit driven by a stable tunnel diode oscillator. The inductive coil with the sample is mounted at the end of a rigid co-axial cable can be inserted into a continuous flow Helium cryostat. The temperature of this system can be regulated between $`4.2K`$ and $`320K`$ and an electromagnet is used to apply a dc magnetic field up to $`6kOe`$. The resonant frequency ($`f_0)`$ is typically in the range of $`24MHz`$ depending on the geometric characteristics of the inductive coil and sample dimensions. The quantity that is measured, the change in frequency $`\mathrm{\Delta }f=f(T,H)f_0)`$ as a function of $`T`$ and $`H,`$ is proportional to the change in reactance $`\mathrm{\Delta }X`$. For magnetic metals, from elementary considerations and applying Maxwell’s equations, it can be shown that:$`\mathrm{\Delta }X\sqrt{\chi },`$ where $`\chi `$ is the differential susceptibility, $`dM/dH`$ of the material. RESULTS Temperature dependence: La<sub>0.875</sub>Sr<sub>0.125</sub>MnO<sub>3</sub> The high sensitivity of the rf technique enables us to clearly detect a paramagnetic to ferromagnetic transition at T<sub>c</sub>=180K as well as two additional transitions at T<sub>s</sub>=270K and T$`_{co\text{ }}`$=150K, as shown in Fig. 1. Interestingly, this composition is observed to undergo structural transitions which are manifested in the change of lattice parameters at 150K, 180K and 270K . At T<sub>s</sub> the susceptibility shows a dip which is due to a structural phase transition from orthorhombic (pseudo cubic) to a cooperative Jahn-Teller distorted phase at lower temperature. In the presence of a magnetic field two characteristic changes are observed to take place at T<sub>c</sub>. First, the peak disappears and secondly, the transition is broadened. In the absence of dc magnetic field the susceptibility is zero above $`T_c`$ and raises rapidly at T<sub>c</sub>. In the presence of dc field the susceptibility is finite at all temperatures and increases with applied field. Therefore, the sharp transition at T<sub>c</sub> becomes broadened when field is applied. The hump observed at T$`_{co\text{ }}`$is very clear and strong unlike the CO transition observed in resistivity and magnetization measurements. This fact emphasizes the importance of high frequency measurements to detect CO transitions. The hump at T<sub>co</sub> is caused by a magnetic transition accompanied by a change in structure.We also observed hysteretic behavior in the susceptibility around T<sub>co</sub> with decreasing and increasing temperature which indicates the first order nature of this transition. As can be seen from the Fig. 1 the rf reactance shows a dip at T<sub>co</sub> during cooling which is not observed while warming. The hump associated with T<sub>co</sub> appears to be a purely ac response of the charge ordering as the dc response does not show any hump. It is worth mentioning that the CO observed in Nd<sub>0.45</sub>Ca<sub>0.55</sub>MnO<sub>3</sub> at 260K also shows a hump in the ac susceptibility measurement. The reason for the hump only in ac measurement is that in ac measurement the differential susceptibility is measured. The reversible response of the ferromagnetic domains to the rf field just below T<sub>c</sub> gives rise to an increase in the differential susceptibility. With further decrease in T the onset of saturation magnetization locks the individual domain and hence the $`\chi (T)`$ starts decreasing. The key to understanding the contribution of the structural transitions to the electronic and magnetic properties lies in the Mn-O interionic distance of the octahedra. The interionic distances $`m(T),s(T)`$, and $`l(T),`$ which are along $`a,c`$ and $`b`$ axes, respectively, are calculated from the representation $`m^2=0.031(a^2+b^2+c^2)`$, $`s^2=0.125c^2m^2`$ and $`l^2=a^2s^2/(16s^2a^2)`$. In these calculations the rotation of the octahedra with respect to the axes is neglected. Fig 2 shows the temperature dependence of these parameters. As can be seen from the Fig. $`m(T)`$ is constant over the entire temperature range, while $`s(T)`$ and $`l(T)`$ show clear anomalies at 140K and 270K. These results imply that for temperatures below 140K or above 270K there is no contribution of the rhombic J-T Q2 mode to the formation of crystal lattice. The turning on of the Q2 mode as the sample is warmed above 140K results in structural phase transition from low temperature. The response of a ferromagnet in a magnetic field is also important to describe the first order transitions observed at T<sub>str</sub> and T<sub>co</sub>. Below T<sub>str</sub> the system shows a spontaneous cooperative JT distorted phase. The strong dependence of $`\chi (T)`$ on magnetic field suggests magnetoelastic coupling for CO besides coulomb repulsion. An isolated hole in $`LaMnO_3`$ can be considered a small polaron which is given by a localized hole in the $`3d_{x^2y^2}`$ orbital surrounded by inverse Jahn-Teller distortion. The polaron phase is an ordered arrangement of Mn<sup>3+</sup> and Mn<sup>4+</sup> ions for which one of the two alternating atomic layers in the (001) plane contains both Mn<sup>3+</sup> ions, as in pure LaMnO3, while the other layer contains both Mn<sup>3+</sup> and Mn$`^{4+\text{ }}`$ions, i.e. holes. The local distortion is because the hole site Mn<sup>4+</sup> is JT inactive whence the electron-phonon energy is lowered by restoring higher symmetry around the hole. In this picture, at high temperatures the La<sub>1-x</sub>Sr<sub>x</sub>MnO<sub>3</sub> may be viewed as a polaron liquid which will eventually transform into polaron lattice as the temperature is lowered. We, therefore, identify $`T_{co}`$ as the onset point for polaron lattice formation, where holes start to freeze on lattice points. From this point of view La<sub>1-x</sub>Sr<sub>x</sub>MnO<sub>3</sub> is considered to undergo successive transitions from polaron liquid (insulator) to Fermi liquid (metal) to polaron lattice (insulator). Field Dependence: La<sub>1-x</sub>Sr<sub>x</sub>MnO<sub>3</sub> The field dependence of the differential susceptibility, $`\frac{dM}{dh_{ac}}_{H_{dc}},`$ at various temperatures below 300K is shown in Fig. 3. As can be seen from the figure for all the doping levels of Sr the $`\chi (H)`$ response shows an overall decrease with the increase in magnetic field. The magnetization M(H) shows a monotonic increase with field with an eventual saturation at high fields, for all the compositions studied. Therefore, the decrease in the differential susceptibility is not surprising. There are, however, many subtle changes in the $`\chi (H)`$ response at low magnetic fields, H $`<`$ 2000G. In the case of x=0.125 composition $`\chi (H)`$ initially increases, reaches a maximum and starts decreasing forming a peak. For the remaining three compositions the peak is not prominent as can be seen from the figure. The behavior of $`\chi (H)`$ can be understood by a simple picture of domain response to weak and strong magnetic fields. When a weak field is applied the magnetization process is reversible. In the presence of strong fields the domains are locked and tend to form a single domain. Therefore, the response of the domains to the ac field at low dc fields is greater than that at high dc fields, thus contributing to the initial increase. CONCLUSIONS Dynamic rf susceptibility of La<sub>1-x</sub>Sr<sub>x</sub>MnO<sub>3</sub> revealed several magnetic signatures in both temperature and field dependent measurements. These magnetic signatures have direct correlation with structural changes in terms of Mn-O interionic distances of the octahedra, at the corresponding temperatures. Field dependent differential susceptibility is found to decrease monotonically with field with a rich structure at low fields. ACKNOWLEDGMENT This work was supported by a US-NSF- 9711910 and NSF-CNRS grant NSF-INT-9726801.
warning/0007/hep-th0007142.html	ar5iv	text	# Lab/UFR-HEP/0008 Non Trivial Extension of the (1+2)-Poincaré Algebra and Conformal Invariance on the Boundary of AdS₃ ## 1 Introduction Recently a non trivial generalisation of the $`(1+2)`$ dimensional Poincaré algebra going beyond the standard supersymmetric extension has been obtained in . In addition to the usual Poincaré generators, this extension refered herebelow to as the Rausch de Traubenberg-Slupinski algebra (RdTS algebra for short), involves two kinds of conserved charges $`Q_s^\pm `$ transforming as $`so(1,2)`$ Verma modules of spin $`s=\pm \frac{1}{k};k2`$. This construction is interesting first because it goes beyond standard 2d-fractional supersymmetry based on considering k-th roots of the $`so(2)`$ vector and second because it gives a new algebraic structure which a priori is valid for higher rank Lie algebras g where $`so(2)`$ and $`so(1,2)`$ appear just as two special examples. In two dimensions where conformal invariance is infinite we now know, by help of conformal field theory methods and techniques of complex analysis, how to deal with objects type k-th root of $`so(2)`$ vector. For higher space time dimensions however, computations are in general difficult to perform except for some special situations such as the problem we will study herebelow and where RdTS symmetry find applications in low dimensional physical systems. In (1+2) dimensions, representations of the RdTS extension of the $`so(1,2)`$ algebra have quantum states carrying fractional values of the spin and are expected to play a particular role in the exploration of special features of field theoretical models of (1+2)dimensional systems with boundaries. The idea of considering 3d systems with boundaries is crucial. It is motivated by the fact that one can imagine that the RdTS $`so(1,2)`$ extension may naturally be linked to a 2d boundary conformal field theory (BCFT) living on the boundary of the space time. From this view we expect that RdTS construction for $`so(1,2)`$ may be related to known results on integrable deformations of 2d conformal invariance. Recall that representations theory of conformal invariance in two dimensions predict naturally the existence of quantum field operators generating states with exotic spins englobing the $`so(1,2)`$ RdTS ones. It is then an interesting task to check if there exists effectively any relation between the RdTS generalisation of Poincaré invariance in (1+2) dimensions and known results on integrable deformations of $`2d`$ CFT’s. We expect that this relation exists really and its determination may help in understanding the behaviour of physical bulk quantities near the boundary of (1+2) dimensional systems. To study this problem we shall mainly work with $`AdS_3`$ as the $`(1+2)`$ space time with boundary and use recent results on strings propagating on $`AdS_3\times N^d`$, where $`N^d`$ is a d-dimensional compact manifold to be specified later on. The analysis we will develop in this paper might also be adapted to study some features of fractional quantum Hall(FQH) effects ; in particular the understanding of the correspondance between the bulk effective Chern Simons (CS) gauge theory of FQH droplets and the conformal field theory living on its boundary . The aim of this paper is to exhibit explicitly the link between the RdTS analysis and 2d BCFT using recent results on D branes physics on the (1+2) dimensional anti de Sitter space $`AdS_3`$ . We first show that there exists indeed a connection between the RdTS algebra and deformations of 2d space time BCFT. Then we establish the rule of correspondance between the two $`so(1,2)`$ Verma modules, used in constructing the non trivial extension of the (1+2) Poincaré invariance, and primary Virasoro representations of the full conformal algebra on the boundary of $`AdS_3`$. We show moreover that the RdTS supersymmetry, although obtained using an unusual method, has in fact the same origin as standard fractional supersymmetry (FSS) , see also. Both FSS and RdTS algebras are residual subsymmetries of conformal invariance. The presentation of this paper is as follows: In section 2, we review the basic ideas of FSS and RdTS supersymmetry using the conformal field theoretical method for the first and the algebraic approach for the second. We give explicit calculations for the deformation of the $`C=\frac{4}{5}`$ Potts model. In section 3, we review the main lines of RdTS analysis. We also introduce some useful tools for the study of the link between the RdTS modules and highest weight representations(HWR) of the Virasoro algebra. In section 4, we study the relation between RdTS supersymmetry and two dimensional conformal invariance. We show in particular that the two $`so(1,2)`$ modules considered in building supersymmetry are just special HWRs of the conformal invariance on the boundary of $`AdS_3`$. In section 5, we use the spectral flow of $`2d`$ $`N=2`$ and $`N=4`$ superconformal invariances to complete the study of section 2 by giving a new result on FSS. We also take the opportunity of using spectral flow of affine Kac-Moody symmetries to give comments on the k-th roots of the $`su(n)`$ fundamental representations used by RdTS in extending their result for $`so(1,2)`$ for $`su(n)`$. In sections 6 and 7, we give our results and conclusion. ## 2 RdTS supersymmetry. RdTS fractional supersymmetry is a special generalisation of FSS living in two dimensions and considered in many occasions in the past in connection with integrable deformations of conformal invariance and representations of the universal envelopping $`U_qsl(2)`$ quantum ordinary and affine symmetries. Like for FSS, highest weight representations of RdTS algebra carry fractional values of the spin and obey more a less quite similar FSS eqs. We will show throughout this study that, up to some details related to the number of dimensions of space time, RdTS fractional supersymmetry has indeed the same origin as FSS. Both FSS and RdTS invariance describe residual symmetries left after integrable deformations of scale invariance in two dimensions. To better understand the algebraic structure of FSS and RdTS supersymmetry we first propose to describe briefly the main lines of 2d FSS one gets from integrable deformations of conformal invariance. Then we give the RdTS extension of the $`(1+2)`$ dimensional Poincaré invariance as derived in . ### 2.1 2d FSS FSS extends the usual Bose-Fermi symmetry in two dimensions; it exchanges bosons and quasiparticles (parafermions) of fractional spin instead of fermions. In addition to the energy momentum translation operator vector $`P_\pm `$, FSS is generated by conserved charges $`Q_x`$ and $`\overline{Q}_x`$ carrying fractional values of the spin x ( $`x=\frac{l}{k};1<l<kmod[1];k2`$). These charge operators are remanant constants of motion that survive after integrable deformations of conformal invariance. There are various FSS algebras depending on the conformal model one starts with. For the example of the $`Z_k`$ parafermionic invariance of Zamolodchikov and Fateev (ZF)\[18 \],see also , a way to get FSS algebras is as follows. First start from the ZF conformal algebra generated by the energy momentum tensor $`T(z)`$ and the parafermionic currents $`\mathrm{\Psi }{}_{q}{}^{}(z),q=1,\mathrm{}k`$: $`T_\mathrm{\Psi }(z_1)T_\mathrm{\Psi }(z_2)`$ $`=`$ $`c_\mathrm{\Psi }/2z_{12}^4+2z_{12}^2T(z_2)+z_{12}^1T(z_2)+\mathrm{}`$ $`\mathrm{\Psi }_k(z_1)\mathrm{\Psi }_l(z_2)`$ $`=`$ $`C_{k,l}^{k+l}z_{12}^{2kl/N}\{\mathrm{\Psi }_{k+l}(z_2)+\mathrm{}\},(k+l)<N,`$ $`\mathrm{\Psi }_k(z_1)\mathrm{\Psi }_k^+(z_2)`$ $`=`$ $`C_{k,Nl}^{N+kl}z_{12}^{2k(Nl)/N}\{\mathrm{\Psi }_{kl}(z_2)+\mathrm{}\},`$ (2.1) $`\mathrm{\Psi }_k(z_1)\mathrm{\Psi }_k^+(z_2)`$ $`=`$ $`z_{12}^{2k(Nk)/N}[1_{id}+2\mathrm{\Delta }_k/c_kz_{12}^2T_\mathrm{\Psi }(z_2)+\mathrm{}],`$ $`T_\mathrm{\Psi }(z_1)\mathrm{\Psi }_k(z_2)`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Delta }_k}{z_{12}^2}}\mathrm{\Psi }_k(z_2)+{\displaystyle \frac{1}{z_{12}}}_z\mathrm{\Psi }_k(z_2)+\mathrm{},`$ where the parameters $`c_\mathrm{\Psi }`$ and $`C_{k,l}^{k+l}`$ are the central charges and structure constants of the parafermoinic algebra respectively. The $`\mathrm{\Psi }_q(z)`$’s and the $`\overline{\mathrm{\Psi }_q}(\overline{z})`$ have the conformal weights $`\mathrm{\Delta }_q=q\frac{(kq)}{k}`$. Second solve the following operator eqs: $`P_{}={\displaystyle 𝑑zT(z)}`$ $`P_+={\displaystyle 𝑑\overline{z}\overline{T}(\overline{z})},`$ (2.2) where $`T(z)`$ and $`\overline{T}(\overline{z})`$ are replaced by their expressions in terms of the $`\mathrm{\Psi }^\pm (z)`$’s and the $`\overline{\mathrm{\Psi }}^\pm (\overline{z})`$ eqs(1). To solve these eqs, one has to specify the ZF parafermionic primary representations since the mode expansions of the $`\mathrm{\Psi }_q`$’s and the $`\overline{\mathrm{\Psi }}_q`$’s depend on the weight of the ZF primary field operators $`\mathrm{\Phi }_q^p`$. $`\mathrm{\Psi }_k(z_1)\mathrm{\Phi }_p^q(z_2)`$ $`=`$ $`{\displaystyle \underset{nZ}{}}z_{12}^{nkp/Nk}Q_{n+\frac{k(p+k)}{N}}^{k,p}\mathrm{\Phi }_p^q(z_2)`$ $`\mathrm{\Psi }_k^+(z_1)\mathrm{\Phi }_p^q(z_2)`$ $`=`$ $`{\displaystyle \underset{nZ}{}}z_{12}^{n+kp/Nk}Q_{n\frac{k(p+k)}{N}}^{k,p}\mathrm{\Phi }_p^q(z_2),`$ (2.3) where $`Q_{n+\frac{k(p+k)}{N}}^{k,p}`$ and $`Q_{n\frac{k(p+k)}{N}}^{k,p}`$ are the modes of $`\mathrm{\Psi }_k`$ and $`\mathrm{\Psi }_k^+`$ respectively defined by: $`Q_{n+\frac{k(p+k)}{N}}^{k,p}\mathrm{\Phi }_p^q(z_2)`$ $`=`$ $`{\displaystyle 𝑑z_1z_{12}^{n+kp/N+k1}\mathrm{\Psi }(z_1)\mathrm{\Phi }_p^q(z_2)},`$ $`Q_{n\frac{k(p+k)}{N}}^{k,p}\mathrm{\Phi }_p^q(z_2)`$ $`=`$ $`{\displaystyle 𝑑z_1z_{12}^{nkp/N+k1}\mathrm{\Psi }(z_1)\mathrm{\Phi }_p^q(z_2)}.`$ (2.4) To illustrate how things work in practice let us consider an example. The method we will present herebelow applies to all $`Z_k`$ parafermionic models as well as others such as the Tye et al symmetries \[20,21 \]. ### 2.2 Deformation of $`𝐂=\frac{\mathrm{𝟒}}{\mathrm{𝟓}}`$ Potts model To fix the ideas, we consider the $`c=4/5`$ critical Potts model described by the following $`Z_3`$ parafermionic invariance. This is the leading non trivial example having constants of motion carrying fractional values of the spin. The algebra governing the critical behaviour of this model is: $`\mathrm{\Psi }^\pm (z_1)\mathrm{\Psi }^\pm (z_2)`$ $``$ $`z_{12}^{2/3}\mathrm{\Psi }^\pm (z_2)`$ $`\mathrm{\Psi }^+(z_1)\mathrm{\Psi }^{}(z_2)`$ $``$ $`z_{12}^{4/3}[1+5/3z_{12}^2T(z_2)]`$ $`T(z_1)\mathrm{\Psi }^\pm (z_2)`$ $``$ $`{\displaystyle \frac{2/3}{z_{12}^2}}\mathrm{\Psi }^\pm (z_2)+{\displaystyle \frac{1}{z_{12}}}_z\mathrm{\Psi }^\pm (z_2)`$ (2.5) $`T(z_1)T(z_2)`$ $`=`$ $`2/5z_{12}^4+2z_{12}^2T(z_2)+z_{12}^1T(z_2).`$ Similar relations are valid for $`\overline{\mathrm{\Psi }}^\pm (\overline{z})`$’s. The ZF parafermionic currents $`\mathrm{\Psi }^\pm `$ have a spin $`2/3`$ and satisfy $`([\mathrm{\Psi }^\pm (z)]^+=\mathrm{\Psi }^{}(z))`$. The algebra (2.4) has three parafermionic highest weight representations (PHWR)$`[\mathrm{\Phi }_q^q]`$; $`q=0,1,2`$ namely the identity family $`I=[\mathrm{\Phi }_0^0]`$ of highest weight $`h_0=0`$ and two degenerate families $`[\mathrm{\Phi }_1^1]`$ and $`[\mathrm{\Phi }_2^2]`$ of weights $`h_1=h_2=\frac{1}{15}`$. Each one of these PHWRs is reducible into three Virasoro HWRs: $`(\mathrm{\Phi }_q^p);p=q,p=q\pm 2(mod6)`$. These field operators which are rotated omongst others under the action of the parafermionic currents as shown herebelow: $`\mathrm{\Psi }^{}\times \mathrm{\Phi }_q^p`$ $`=`$ $`\mathrm{\Phi }_q^{p\pm 2}`$ $`\mathrm{\Phi }_q^{p\pm 6}`$ $`=`$ $`\mathrm{\Phi }_q^p,`$ (2.6) obey Virasoro and ZF primary conditions: $`L_n\|h`$ $`=`$ $`0,n>o`$ $`Q_{n\pm (p\pm 1)/3}^\pm \|h`$ $`=`$ $`0,n\pm (p\pm 1)/3>0,`$ (2.7) where the $`L_n`$ Virasoro and the $`Q_{n\pm (p\pm 1)/3}^\pm `$ ZF modes are given by: $`L_n\|\mathrm{\Phi }_p^q`$ $`=`$ $`{\displaystyle 𝑑zz^{n+1}T(z)\mathrm{\Phi }_p^q(0)\|0},`$ $`Q_{n\pm (p\pm 1)/3}^\pm \|\mathrm{\Phi }_p^q`$ $`=`$ $`{\displaystyle 𝑑zz^{n\pm p/3}\mathrm{\Psi }^\pm (z)\mathrm{\Phi }_p^q(0)\|0}.`$ (2.8) Note that the the mode expansion of the ZF currents depend on the representation field operator on which they act. This property is manifestly seen on the energies of the creation and annihilation operators $`Q_{n\pm (p\pm 1)/3}^\pm `$ which depend on the quantum number p of the ZF primary field $`\mathrm{\Phi }_q^p(z)`$: $$\mathrm{\Psi }^\pm (z_1)\mathrm{\Phi }_q^p(z_2)=z_{12}^{n1p/3}Q_{n\pm (p\pm 1)/3}^\pm \mathrm{\Phi }_q^p(z_2),$$ (2.9) The ZF primary field operators $`\mathrm{\Phi }_q^p(z)`$ satisfy also braiding properties type: $$z_{12}^\mathrm{\Delta }\mathrm{\Phi }_1(z_1)\mathrm{\Phi }_2(z_2)=z_{21}^\mathrm{\Delta }\mathrm{\Phi }_2(z_2)\mathrm{\Phi }_1(z_1)=\mathrm{\Phi }(z),$$ (2.10) where $`\mathrm{\Delta }=\mathrm{\Delta }_1+\mathrm{\Delta }_2\mathrm{\Delta }_3,\mathrm{\Delta }_i;i=1,2,3,`$ are the conformal weights of the $`\mathrm{\Phi }_i`$ field operators. The second step in the derivation of FSS is to solve the operator eq(2.2) expressing the $`2d`$ energy momentum vector $`P_\pm `$ in terms of the ZF modes $`Q_{n\pm (p\pm 1)/3}^\pm `$: $$\begin{array}{ccc}P_{}=𝑑z\frac{3}{5}z^{2/3}(\mathrm{\Psi }^+(z)\mathrm{\Psi }^{}(0))\hfill & & \\ P_+=𝑑\overline{z}\frac{3}{5}z^{2/3}(\overline{\mathrm{\Psi }}^+(\overline{z})\overline{\mathrm{\Psi }}^{}(0)).\hfill & & \end{array}$$ (2.11) where we replaced $`T(z)`$ and $`\overline{T}(\overline{z})`$ in terms of the $`\mathrm{\Psi }^\pm (z)`$’s and the $`\overline{\mathrm{\Psi }}^\pm (\overline{z})`$ as given by eqs(4 ). The solution of eqs(9) involves three pairs of doublets of the charge operators $`(Q_{\frac{1}{3}}^\pm ,\overline{Q}_{\frac{1}{3}}^\pm )`$, $`(Q_{\frac{2}{3}}^\pm ,\overline{Q}_{\frac{2}{3}}^\pm )`$ and $`(Q_{0}^{}{}_{}{}^{\pm },\overline{Q}_0^\pm )`$. Using the primary highest weight conditions (2.7), one can check by explicit computation that the $`Q,\overline{Q},P_{}`$ and $`P_+`$ charge operators generate the following algebra: $`𝐏`$ $`=`$ $`Q_{1/3}^+Q_0^+Q_{2/3}^+\mathrm{\Pi }_0+Q_{2/3}^+Q_{1/3}^+Q_0^+\mathrm{\Pi }_1+Q_0^+Q_{2/3}^+Q_{1/3}^+\mathrm{\Pi }_1`$ $`[P_\pm ,Q_x]`$ $`=`$ $`0;x=0,1/3,2/3`$ $`\overline{𝐏}`$ $`=`$ $`\overline{Q}_{1/3}^{}{}_{}{}^{+}\overline{Q}_{0}^{}{}_{}{}^{+}\overline{Q}_{2/3}^{}{}_{}{}^{+}\overline{\mathrm{\Pi }}_0+\overline{Q}_{2/3}^{}{}_{}{}^{+}\overline{Q}_{1/3}^{}{}_{}{}^{+}\overline{Q}_{0}^{}{}_{}{}^{+}\overline{\mathrm{\Pi }}_1+\overline{Q}_{0}^{}{}_{}{}^{+}\overline{Q}_{2/3}^{}{}_{}{}^{+}\overline{Q}_{1/3}^{}{}_{}{}^{+}\overline{\mathrm{\Pi }}_1`$ (2.12) $`[P_\pm ,\overline{Q}_{+x}]`$ $`=`$ $`0.`$ In these eqs the $`\mathrm{\Pi }_q`$’s and $`\overline{\mathrm{\Pi }}_q`$’s are projector operators on the q-th ZF primary state $`[\mathrm{\Phi }_{q}^{}{}_{}{}^{q}\times \overline{\mathrm{\Phi }}_{q}^{}{}_{}{}^{q}]`$. The algebra (2.12) may also be obtained by analysing the energy spectrum of the mode operators $`Q_{n\pm (p\pm 1)/3}^\pm `$ and $`\overline{Q}_{n\pm (p\pm 1)/3}^\pm `$, n integer. The $`Q_{n\pm (p\pm 1)/3}^\pm `$’s and $`\overline{Q}_{n\pm (p\pm 1)/3}^\pm `$’s, which depend on the $`p`$ charge, act only on the conformal representation $`\|\mathrm{\Phi }_p^q`$. This property may be interpreted to mean that expect the $`\|\mathrm{\Phi }_p^q`$ family, the action of the $`Q_{n\pm (p\pm 1)/3}^\pm `$’s kills all states $`\|\mathrm{\Phi }_{p}^{}{}_{}{}^{r}`$ with r different from q. For $`q=0`$ for example, the non vanishing actions of $`Q_x^\pm `$ and $`\overline{Q}_x^\pm ,x=0,1/3,2/3`$ on the states $`\|s,p`$ of spin s, $`0s1`$ and charge $`p`$ read as: $`Q_{2/3}^\pm \|0,0`$ $`=`$ $`\|2/3,0`$ $`Q_0^+\|2/3,+2`$ $`=`$ $`\|2/3,2`$ $`Q_0^{}\|2/3,2`$ $`=`$ $`\|2/3,+2`$ (2.13) $`Q_{1/3}^+\|2/3,2`$ $`=`$ $`\|1,0`$ $`Q_{1/3}^{}\|2/3,+2`$ $`=`$ $`\|1,0`$ and similar eqs for the antiholomorphic sector. From these Eqs as well as the expansion(2.3-4) of the ZF currents, one sees that $`Q_{1/3}^\pm `$ and $`Q_0^\pm `$ cannot act directly on the state $`\|0,0`$. Similarly $`Q_{2/3}^\pm `$ cannot operate directly on $`\|2/3,\pm 2`$. This result gives an explicit argument showing that FSS should be generated by more than one $`Q`$ and $`\overline{Q}`$ operators as it was naively used in earlier physical litterature on FSS. It shows moreover that not all the $`Q_x^\pm `$’s are independent since we have: $`Q_{1/3}^{}`$ $`=`$ $`Q_{1/3}^+Q_0^+`$ $`Q_{1/3}^+`$ $`=`$ $`Q_0^{}Q_{1/3}^{}`$ $`Q_{2/3}^{}`$ $`=`$ $`Q_0^+Q_{2/3}^{}`$ (2.15) $`Q_{2/3}^+`$ $`=`$ $`Q_{2/3}^{}Q_0^{}.`$ Similar expressions may be written down for for the antiholomorphic sector. Putting back these relations into eqs(2.12), we find the following linearized algebra. $`2P_1`$ $`=`$ $`\{Q_{2/3}^+,Q_{1/3}^{}\}+\{Q_{1/3}^+,Q_{2/3}^{}\}`$ $`0`$ $`=`$ $`\{Q_{1/3}^\pm ,Q_{1/3}^\pm \}=\{Q_{2/3}^\pm ,Q_{2/3}^\pm \}.`$ (2.16) We shall return to this linearised realisation of FSS in section5 when we discuss the spectral flow of $`N=2`$ and$`N=4`$ superconformal invariance in two dimensions,wher a similar result will be obtained by using special choices of the parameter of the flow. ## 3 More on RdTS supersymmetry In this section we review briefly the derivation of the RdTS extension of the $`(1+2)`$ dimensional Poincaré invariance. We also initiate the study of a field realisation of RdTS supersymmetry which we develop further in the forthcoming section. In this regards we would like to note that as far as $`SO(1,2)`$ group is concerned, we will encounter in our analysis various kinds of $`SO(1,2)`$ symmetries with different physical interpretations. In addition to the $`SO(1,2)`$ Lorentz invariance of the $`(1+2)`$ dimensional space time considered in , we will handle four $`SO(1,2)`$ invariances classified as: (1) Two $`SO(1,2)`$’s given by the zero mode subgroup product $`SO(1,2)\times \overline{SO(1,2)}`$ associated to $`so_k(1,2)\times \overline{so_k(1,2)}`$ affine Kac Moody invariance to be studied in section 4. This subsymmetry will be realised by using the usual $`Sl(2,R)SO(1,2)`$ Wakimoto field theoretical representation. (2) Two other $`SO(1,2)`$ subsymmetries associated to the non anomalous subalgebras of the left and right Virasoro symmetries of some two dimensional BCFT of $`AdS_3`$ to be specified later on. ### 3.1 RdTS extension of $`\mathrm{𝐒𝐎}(\mathrm{𝟏},\mathrm{𝟐})`$ To start consider the Poincaré symmetry in $`(1+2)`$ dimensions generated by the space time translations $`P_\mu `$ and the Lorentz rotations $`J_\alpha `$ satisfying altogether the following closed commutation relations: $`[J_\alpha ,J_\beta ]`$ $`=`$ $`iϵ_{\alpha \beta \gamma }\eta ^{\gamma \delta }P_\delta `$ $`[J_\alpha ,J_\mu ]`$ $`=`$ $`iϵ_{\alpha \beta \gamma }\eta ^{\gamma \delta }P_\delta `$ (3.1) $`[P_\mu ,P_\nu ]`$ $`=`$ $`0.`$ In these eqs, $`\eta _{\alpha \beta }=diag(1,1,1)`$ is the $`(1+2)`$ Minkowski metric and $`ϵ_{\alpha \beta \gamma }`$ is the completly antisymmetric Levi-Civita tensor such that $`ϵ_{012}=1`$. A convenient way to handle eqs(3.1) is to work with an equivalent formulation using the following Cartan basis of generators $`P_{}=P_1\pm iP_2`$ and $`J_{}=J_1\pm iJ_2`$. In this basis eqs(3.1) read as: $`[J_+,J_{}]`$ $`=`$ $`2J_0`$ $`[J_0,J_\pm ]`$ $`=`$ $`\pm J_\pm `$ $`[J_\pm ,P_{}]`$ $`=`$ $`\pm P_0`$ (3.2) $`[J_+,P_+]`$ $`=`$ $`[J_{},P_{}]=0`$ $`[J_0,P_0]`$ $`=`$ $`[P_\pm ,P_{}]=0.`$ The algebra (3.1-2) has two Casimir operators $`P^2=P_{0}^{}{}_{}{}^{2}\frac{1}{2}(P_+P_{}+P_{}P_+)`$ and $`P.J=P_0J_0\frac{1}{2}(P_+J_{}+P_{}J_+)`$. When acting on highest weight states of mass m and spin s, the eigenvalues of these operators are $`m^2`$ and $`ms`$ respectively. For a given s, one distinguishes two classes of irreducible representations: massive and massless representations. To build the $`so(1,2)`$ massive representations, it is convenient to go to the rest frame where the momentum vector $`P_\mu `$ is $`(m,0,0)`$ and the $`SO(1,2)`$ group reduces to its abelian $`SO(2)`$ little subgroup generated by $`J_0`$; $`(J_\pm =0)`$. In this case, massive irreducible representations are one dimensional and are parametrized by a real parameter. For the full $`SO(1,2)`$ group however, the representations are either finite dimensional for $`\|s\|𝐙^+/2`$ or infinite dimensional for the remaning values of $`s`$. Given a primary state $`\|s`$ of spin $`s`$, and using the abovementioned $`SO(1,2)`$ group theoretical properties, one may construct in general two representations HWR(I) and HWR(II) out of this state $`\|s`$. The first representation HWR(I) is a highest weight representation given by: $`J^0\|s`$ $`=`$ $`s\|s`$ $`J_{}\|s`$ $`=`$ $`0`$ $`\|s,n`$ $`=`$ $`\sqrt{{\displaystyle \frac{\mathrm{\Gamma }(2s)}{\mathrm{\Gamma }(2s+n)\mathrm{\Gamma }(n+1)}}}(J_+)^n\|s,n1`$ (3.3) $`J_0\|s,n`$ $`=`$ $`(s+n)\|s,n`$ $`J_+\|s,n`$ $`=`$ $`\sqrt{(2s+n)(n+1)}\|s,n+1`$ $`J_{}\|s,n`$ $`=`$ $`\sqrt{(2s+n1)n}\|s,n1.`$ The second representation is a lowest weight representation which we refer to denote as HWR(II) is defined as: $`\overline{J_0}\|\overline{s}`$ $`=`$ $`s\|\overline{s}`$ $`\overline{J_+}\|\overline{s}`$ $`=`$ $`0`$ $`\|\overline{s,n}`$ $`=`$ $`()^n\sqrt{{\displaystyle \frac{\mathrm{\Gamma }(2s)}{\mathrm{\Gamma }(2s+n)\mathrm{\Gamma }(n+1)}}}(\overline{J_{}})^n\|\overline{s}`$ (3.4) $`\overline{J_0}\|\overline{s,n}`$ $`=`$ $`(s+n)\|\overline{s,n}`$ $`\overline{J_+}\|\overline{s,n}`$ $`=`$ $`\sqrt{(2s+n1)n}\|\overline{s,n+1}.`$ Note in passing that in the second module we have supplemented the generators and the representations states with a bar index. This convention of notation will be justified later on. To fix the ideas, HWR(I) will be identified in section 7 with a left Virasoro Verma module and HWR(II) will be interpred as a right Virasoro one. Note moreover that both HWR(I) and HWR(II) representations have the same $`so(1,2)`$ Casimir $`C_s`$= $`s(s1),s<0`$. For $`s𝐙^{}/\mathrm{𝟐}`$, these representations are finite dimensional and their dimension is $`(2\|s\|+1)`$. For generic real values of $`s`$, the dimension of the representations is however infinite. If one chooses a fractional value of $`s`$ say $`s=\frac{1}{k}`$; each of the two representations (3.3-4) splits a priori into two isomorphic representations respectively denoted as $`D_{\pm 1/k}^+andD_{\pm 1/k}^{}`$. This degeneracy is due to the redundancy in choosing the spin structure of $`\sqrt{2/k}`$ which can be taken either as $`+i\sqrt{2/k}`$ or $`i\sqrt{2/k}`$. These representations are not independent since they are related by conjugations; this why we shall use hereafter the choice of by considering only $`D_{1/k}^+`$ and $`D_{1/k}^{}`$. In this case the two representation generators $`J_{0,\pm }`$ and $`\overline{J}_{0,\pm }`$ are related as: $$\overline{J}_{0,}=(J_{0,\pm })^{}.$$ (3.5) Furthermore taking the tensor product of the primary states $`\|s`$ and $`\|\overline{s}`$ of the two $`so(1,2)`$ modules HWR(I) and HWR(II) and using eqs(3.3-4), it is straightforward to check that it behaves like a scalar under the full charge operator $`J_0\times 1_d+1_d\times \overline{J}_0`$ which we denote simply as $`J_0+\overline{J}_0`$ : $$(J_0+\overline{J}_0)\|s\|\overline{s}=0.$$ (3.6) Eq(3.6) is a familiar relation in the study of primary states of Virasoro algebra. This equation together with the mode operators $`J_{}^n`$ and $`\overline{J}_+^m`$ which act on $`\|s\|\overline{s}`$ as: $`(J_{})^n\|s\|\overline{s}`$ $`=`$ $`0,n1`$ $`(\overline{J}_+)^m\|s\|\overline{s}`$ $`=`$ $`0,m1`$ (3.7) define a highest weight state which looks like a Virasoro primary state of spin 2s and scale dimension $`\mathrm{\Delta }=0`$. We will show later on when we study the primary field representation of the $`2d`$ BCFT of a string propagating on the $`AdS_3`$ background, that eq(3.6-7) correspond indeed to: $`(L_0\overline{L}_0)\mathrm{\Phi }_{h,\overline{h}}(0,0)\|0`$ $`=`$ $`(h\overline{h})\mathrm{\Phi }_{h,\overline{h}}(0,0)\|0`$ $`(L_0+\overline{L}_0)\mathrm{\Phi }_{h,\overline{h}}(0,0)\|0`$ $`=`$ $`(h+\overline{h})\mathrm{\Phi }_{h,\overline{h}}(0,0)\|0`$ (3.8) $`L_n\mathrm{\Phi }_{h,\overline{h}}(0,0)\|0`$ $`=`$ $`0,n1`$ $`\overline{L}_m\mathrm{\Phi }_{h,\overline{h}}(0,0)\|0`$ $`=`$ $`0,m1.`$ where $`L_n`$ and $`\overline{L}_m`$ are respectively the usual left and right Virasoro modes and $`\varphi _{h,\overline{h}}(z,\overline{z})`$ is a primary conformal field representation of conformal scale $`h`$\+ $`\overline{h}`$ and conformal spin $`h`$\- $`\overline{h}`$. This property, which gives the connection between RdTS supersymmetry and conformal invariance, will be explicited in details when we discuss HWRs of the conformal symmetry on the boundary of $`AdS_3`$. The primary $`so(1,2)`$ highest weight states $`\|s`$ and $`\|\overline{s}`$ eqs(3.3-4) are respectively in one to one correspondance with the left Virasoro primary state $`\mathrm{\Phi }_h(0)\|0=\|h`$ and the right Virasoro primary one $`\mathrm{\Phi }_{\overline{h}}(0)\|0=\|\overline{h}`$. On the other hand, if we respectively associate to HWR(I) and HWR(II) the mode operators $`Q_{s+n}^+=Q_{s+n}`$ and $`Q_{sn}^{}=\overline{Q}_{s+n}`$= and using $`SO(1,2)`$ tensor product properties, one may build, under some assumptions, an extension $`𝐒`$ of the $`so(1,2)`$ algebra going beyond the standard supersymmetric one. To do so, note first that the system $`J_0`$, $`J_+`$ ,$`J_{}`$ and $`Q_{s+n}`$ obey the following commutation relations $`(s=1/k)`$. $`[J_0,Q_{s+n}]`$ $`=`$ $`(s+n)Q_{s+n}`$ $`[J_+,Q_{s+n}]`$ $`=`$ $`\sqrt{(2s+n)(n+1)}Q_{s+n+1}`$ (3.9) $`[J_{},Q_{s+n}]`$ $`=`$ $`\sqrt{(2s+n1)n}Q_{s+n1}.`$ Similarly we have for the antiholomorphic sector: $`[\overline{J}_0,\overline{Q}_{s+n}]`$ $`=`$ $`(s+n)\overline{Q}_{s+n}`$ $`[\overline{J}_+,\overline{Q}_{s+n}]`$ $`=`$ $`\sqrt{(2s+n1)n}\overline{Q}_{s+n1}`$ (3.10) $`[\overline{J}_{},\overline{Q}_{s+n}]`$ $`=`$ $`\sqrt{(2s+n)(n+1)}\overline{Q}_{s+n+1}.`$ To close these commutations relations with the $`Q_s`$’s through a k-th order product one should fullfil some constraints. (1) the generalized algebra $`𝐒`$we are looking for should be a generalisation of what is known in two dimensions, i.e a generalisation of FSS. (2) When the charge operator $`Q_{s+n}`$ goes arround an other, say $`Q_{s+m}`$, it picks a phase $`\mathrm{\Phi }=2i\pi /k`$; i.e: $$Q_{s+n}Q_{s+m}=e^{\pm 2i\pi s}Q_{s+m}Q_{s+n}+\mathrm{};s=\frac{1}{k},$$ (3.11) where the dots refer for possible extra charge operators of total $`J_0`$ eigenvalue (2s+n+m). Eq (3.11) shows also that the algebra we are looking for has a $`𝐙_k`$ graduation. Under this discrete symmetry, $`Q_{s+n}`$ carries a $`+1(modk)`$ charge while the $`P_{0,\pm }`$ energy momentum components have a zero charge mod k. (3) the generalised algebra $`𝐒`$ should split into a bosonic B part and an anyonic A and may be written as: $`𝐒=_{r=0}^{k1}A_r=B_{r=1}^{k1}A_r`$ . Since $`A_nA_mA_{(n+m)(modk)}`$ one has: $`\{A_r\mathrm{}A_r\}_k`$ $``$ $`B`$ $`[B,A]`$ $``$ $`B`$ (3.12) $`[B,B]`$ $``$ $`B.`$ In these eqs, $`\{A_r\mathrm{}A_r\}_k`$ means the complete symmetrisation of the k anyonic operators $`A_r`$ and is defined as: $$\{A_{s_r}\mathrm{}A_{s_r}\}_k=\frac{1}{k!}\underset{\sigma \mathrm{\Sigma }}{}(A_{s_{\sigma (1)}}\mathrm{}A_{s_{\sigma (k)}},$$ (3.13) where the sum is carried over the k elements of the permutation group $`\{1,\mathrm{},k\}`$. (4) the algebra $`𝐒`$ should obey generalised Jacobi identities. In particular we should have: $$adB\{A_{s_1}\mathrm{}A_{s_k}\}=0,$$ (3.14) where B stands for the bosonic generators $`J_{0,\pm }`$ or $`P_{0,\pm }`$ of the Poincaré algebra. Using eq(3.12) to write $`\{A_r\mathrm{}A_r\}_k`$ as $`\alpha _\mu P^\mu +\beta _\mu J^\mu `$ where $`\alpha `$ and $`\beta `$ are real constants; then putting back into the above relation we find that $`\{A_r\mathrm{}A_r\}_k`$ is proportional to $`P_\mu `$ only. In other words, $`\beta _\mu `$ should be equal to zero; a property which is easily seen by taking $`B=P_\mu `$ in eq (3.14). Put differently the symmetric product of the $`D_s^\pm `$, denoted hereafter as $`S^k[D_s^\pm ]`$, contains the space time vector representation $`D_1`$ of $`so(1,2)`$ and so the primitive charge operators $`Q_{1/k}`$ and $`\overline{Q}_{1/k}`$ obey: $`[J_0,(Q_{1/k})^k]`$ $`=`$ $`(Q_{1/k})^kP_{}`$ $`[J_{},(Q_{1/k})^k]`$ $`=`$ $`0.`$ (3.15) Similarly we have: $`[\overline{J}_0,(\overline{Q}_{1/k})^k]`$ $`=`$ $`(\overline{Q}_{1/k})^kP_+`$ $`[\overline{J}_+,(\overline{Q}_{1/k})^k]`$ $`=`$ $`0.`$ (3.16) Moreover acting on $`(Q_{1/k})^k`$ by $`adJ_+^n`$ and on $`(\overline{Q}_{1/k})^k`$ by $`ad\overline{J}_+^n`$, one obtains: $`adJ_+(Q_{1/k})^k`$ $``$ $`P_0`$ $`ad\overline{J}_{}(\overline{Q}_{1/k})^k`$ $``$ $`P_0`$ (3.17) $`ad^2J_+(Q_{1/k})^k`$ $``$ $`P_{}`$ $`ad^2\overline{J}_{}(\overline{Q}_{1/k})^k`$ $``$ $`P_+.`$ In summary, starting from $`so(1,2)`$ lorentz algebra (3.1-2) and the two Verma modules HWR(I) and HWR(II)(3.3-4 ), one may build the following new extended symmetry: $`\{Q_{\frac{1}{k}}^\pm ,\mathrm{},Q_{\frac{1}{k}}^\pm \}_k`$ $`=`$ $`P_{}=P_1\pm iP_2`$ $`\{Q_{\frac{1}{k}}^\pm ,\mathrm{},Q_{\frac{1}{k}}^\pm ,Q_{1\frac{1}{k}}^\pm \}_k`$ $`=`$ $`\pm i\sqrt{{\displaystyle \frac{2}{k}}}P_0(k1)\{Q_{\frac{1}{k}}^\pm ,\mathrm{},Q_{\frac{1}{k}}^\pm ,Q_{1\frac{1}{k}}^\pm ,Q_{1\frac{1}{k}}^\pm \}_k`$ $`\pm i\sqrt{k2}\{Q_{\frac{1}{k}}^\pm ,\mathrm{},Q_{\frac{1}{k}}^\pm ,Q_{1\frac{1}{k}}^\pm ,Q_{2\frac{1}{k}}^\pm \}_k`$ $`[J^\pm ,[J^\pm ,[J^\pm ,(Q_{\frac{1}{k}}^\pm )^k]]]`$ $`=`$ $`0.`$ (3.18) Eq(3.18) defines what we have been refering to as RdTS algebra. For more details on this algebraic structure, see . ## 4 Furthermore on RdTS supersymmetry Here we would like to answer the question rised in the introduction concerning the link between RdTS supersymmetry and two dimensional conformal invariance. We have anticipated on the nature of this link by saying that RdST supersymmetry is expected to arise from appropriate deformations of two dimensional CFT’s on the boundary of $`AdS_3`$. The appearence of the $`AdS_3`$ space in this analysis is due to the fact that this geometry has many relevant features for our present study. We give hereafter two useful properties regarding the space time $`SO(1,2)`$ Lorentz group: (1) In its euclidean representation, $`AdS_3`$ has an $`SO(1,3)`$ isometry group containing as subgroup the $`SO(1,2)`$ Lorentz symmetry of the (1+2) space time we are interested in. (2) The two dimensional $`AdS_3`$ boundary space may be realised as a two sphere on which may live boundary conformal field theories, which themselves have $`so(1,2)`$ projective subsymmetries that can be related to the above mentionned $`so(1,2)`$ Lorentz group. Starting from these observations we want to show that the two so(1,2) modules HWR(I) and HWR(II), considered in the building of RdTS supersymmetry, are just special representations of the $`AdS_3`$ BCFT. To proove this relation in a comprehensive manner, let us first review briefly some elements of $`AdS_3`$ geometry. The $`AdS_3`$ space is given by the hyperbolic coset manifold $`Sl(2,C)/SU(2)`$ which may be thought of as the three dimensional hypersurface $`H_{3}^{}{}_{}{}^{+}`$ $$X_{0}^{}{}_{}{}^{2}+X_{1}^{}{}_{}{}^{2}+X_{2}^{}{}_{}{}^{2}+X_{3}^{}{}_{}{}^{2}=l^2,$$ (4.1) embedded in flat $`R^{1,3}`$ with local coordinates $`X^0`$, $`X^1`$, $`X^2`$, $`X^3`$ . This hypersurface describes a space with a constant negative curvature ($`\frac{1}{l^2}`$). The parameter l is choosen to be quantized in terms of the $`l_s`$ fundamental string lenght units; i.e, $`l=l_s\times k`$ where k is an integer to be interpreted later on as the Kac Moody level of the $`so_k(1,2)`$ affine symmetry. To study the field theory on the boundary of $`AdS_3`$, it is convenient to introduce the following set of local coordinates of $`AdS_3`$: $`\varphi `$ $`=`$ $`log(X_0+X_3)/l`$ $`\gamma `$ $`=`$ $`{\displaystyle \frac{X_2+iX_0}{X_0+iX_3}}`$ (4.2) $`\overline{\gamma }`$ $`=`$ $`{\displaystyle \frac{X_2iX_1}{X_0+iX_3}}.`$ An equivalent description of the hypersurface is: $`\gamma `$ $`=`$ $`{\displaystyle \frac{r}{\sqrt{l^2+r^2}}}e^{\tau +i\theta }`$ $`\overline{\gamma }`$ $`=`$ $`{\displaystyle \frac{r}{\sqrt{l^2+r^2}}}e^{\tau i\theta }`$ $`\varphi `$ $`=`$ $`\tau +1/2log(1+r^2/l^2)`$ (4.3) $`r`$ $`=`$ $`le^\varphi \sqrt{\gamma \overline{\gamma }}`$ $`\tau `$ $`=`$ $`\varphi 1/2log(1+e^{2\varphi }\gamma \overline{\gamma })`$ $`\theta `$ $`=`$ $`{\displaystyle \frac{1}{2i}}log(\gamma /\overline{\gamma }),`$ where we have used the change of variables: $`X_0`$ $`=`$ $`X_0(r,\tau )=\sqrt{l^2+r^2}cosh\tau `$ $`X_3`$ $`=`$ $`X_3(r,\tau )=\sqrt{l^2+r^2}sinh\tau `$ (4.4) $`X_1`$ $`=`$ $`X_1(r,\theta )=rsin\theta `$ $`X_2`$ $`=`$ $`X_2(r,\theta )=rcos\theta .`$ In the coordinates $`(\varphi ,\gamma ,\overline{\gamma })`$, the metric of $`H_3^+`$ reads as: $$ds^2=k(d\mathrm{\Phi }^2+e^{2\mathrm{\Phi }}d\gamma d\overline{\gamma }).$$ (4.5) Note that in the $`(\varphi ,\gamma ,\overline{\gamma })`$ frame, the boundary of euclidean $`AdS_3`$ corresponds to take the field $`\mathrm{\Phi }`$ to infinity. As shown on eq(4.3-4),this is a two sphere which is locally isomorphic to the complex plane parametrized by $`(\gamma ,\overline{\gamma })`$ . Quantum field theory on the $`AdS_3`$ space is very special and has very remarkable features governed by the Maldacena correspondence in the zero slope limit of string theory\[24 \]. On this space it has been shown that bulk correlations functions of quantum fields find natural interpretations in the conformal field theory on the boundary of $`AdS_3`$. In algebraic language, this correspondance transforms world sheet symmetries of strings on $`AdS_3`$ into space time infinite dimensional invariances on the boundary of $`AdS_3`$. In what follows we shall review some useful properties of strings on $`AdS_3`$ and $`AdS_3`$. ### 4.1 $`\mathrm{𝐀𝐝𝐒}_\mathrm{𝟑}`$ / CFT correspondence. Strings propagating on the $`AdS_3`$ background are involved in the study of supersymmetric gauge theories with eight supercharges; in particular in the understanding of the Higgs and Coulomb branches near the moduli space singularity\[25 \]. Strings on $`AdS_3`$ have rich symmetries; some of them turn out to be related to the problem we are studing. These symmetries, which may be classified into WS symmetries and space time invariances, carry all relevent informations one needs to know about the string dynamics on $`AdS_3`$. In what follows we want to give some useful relations regarding these two classes of symmetries. To work out explicit field theoretical realisations of these symmetries, we start by recalling that in the presence of the Neveu-Schwarz $`B_{\mu \nu }`$ field with euclidean world sheet parameterized $`(z,\overline{z})`$, the dynamics of the bosonic string on $`AdS_3`$ is described by the following classical lagrangian: $$L=k[\mathrm{\Phi }\overline{}\mathrm{\Phi }+e^{2\mathrm{\Phi }}\gamma \overline{\gamma }].$$ (4.6) In this eq $``$ and $`\overline{}`$ stand for derivatives with respect to z and $`\overline{z}`$ repectively. Introducing two auxiliary variables $`\beta `$ and $`\overline{\beta }`$, the above eq may be put into the following convenient form: $$L^{^{}}=k^2(\mathrm{\Phi }\overline{}\mathrm{\Phi }+\beta \overline{}\gamma +\overline{\beta }\overline{\gamma }e^{2\mathrm{\Phi }}\beta \overline{\beta }.)$$ (4.7) The eqs of motion of the various fields one gets from eq(4.7) read as: $`\overline{}\mathrm{\Phi }2\beta \overline{\beta }e^{2\mathrm{\Phi }}`$ $`=`$ $`0`$ $`\overline{}\gamma \beta e^{2\mathrm{\Phi }}`$ $`=`$ $`0`$ (4.8) $`\overline{\gamma }\overline{\beta }e^{2\mathrm{\Phi }}`$ $`=`$ $`0`$ $`\overline{\beta }`$ $`=`$ $`\overline{}\beta =0.`$ String dynamics on the boundary of $`AdS_3`$ is obtained from the previous bulk eqs by taking the limit $`\mathrm{\Phi }`$ goes to infinity.This gives: $`\overline{}\mathrm{\Phi }`$ $`=`$ $`0`$ $`\overline{}\gamma `$ $`=`$ $`\overline{\gamma }=0`$ (4.9) $`\overline{\beta }`$ $`=`$ $`\overline{}\beta =0.`$ The WS fields $`\mathrm{\Phi },\gamma `$ and $`\overline{\gamma }`$ which had general expressions in the bulk become now holomorphic on the boundary of $`AdS_3`$ and describe a BCFT. Note that consistency of quantum mechanics of the string propagating in space time requires that the target space should be $`AdS_3\times N`$, where N is a (3+n) dimensional compact manifold. To fix the ideas, N may be thought of as $`S^3\times T^n`$ with $`n=20`$ for the bosonic string and $`n=4`$ for superstrings. We shall consider hereafter both of string and superstring cases.Given the big number of relations one may write down, we shall use however a strategy in which we give the strict necessary results. Thus our plan in what follows is: First, we describe some algebraic features of the WS invariance; then we make a pause to give a complement on FSS using spectral flow of $`N=2`$ and $`N=4`$ conformal invariance, after what we return to complete space time symmetries on the boundary of $`AdS_3`$and finally we give our results. ### 4.2 WS Symmetries World sheet invariances include affine Kac-Moody, Virasoro symmetries and their extensions. For a bosonic string propagating on $`AdS_3\times S^3\times T^{20}`$, we have the following: A\- Three kinds of WS affine Kac-Moody invariances: (a)A level $`(k2)`$ $`sl(2)\times \overline{sl(2)}`$ invariance coming from the string propagation on $`AdS_3`$. This invariance is generated by the conserved currents $`J_{sl(2)}^q`$ and $`\overline{J}_{sl(2)}^{}{}_{}{}^{q};q=0,\pm 1`$. In terms of the WS fields $`\mathrm{\Phi },\gamma ,\overline{\gamma },\beta `$ and $`\overline{\beta }`$ of eq(?), the field theoretical realization of these currents is given by the Wakimoto representation: $`J^{}(z)`$ $`=`$ $`\beta (z)`$ $`J^+(z)`$ $`=`$ $`\beta \gamma ^2+\sqrt{2(k2)}\gamma \mathrm{\Phi }+k\gamma `$ $`J^0(z)`$ $`=`$ $`\beta \gamma +1/2\sqrt{2(k2)}\mathrm{\Phi }`$ (4.10) $`\overline{J}^{}(\overline{z})`$ $`=`$ $`\overline{\beta }`$ $`\overline{J}^0(\overline{z})`$ $`=`$ $`\overline{\beta }\overline{\gamma }+1/2\sqrt{2(k2)}\mathrm{\Phi }`$ $`\overline{J}^+(\overline{z})`$ $`=`$ $`\overline{\beta }\overline{\gamma }^2+\sqrt{2(k2)}\overline{\gamma }\mathrm{\Phi }+k\overline{\gamma }.`$ (b)A level $`(k+2)`$ invariance coming from the string propagation on $`S^3`$. The conserved currents are $`J_{su(2)}^q`$ and $`\overline{J}_{su(2)}^q`$. The WS field theoretical realization of these currents is given by the level $`(k+2)`$ WZW $`su(2)`$ model \[26 \]. (c)A $`u(1)^{20}\times \overline{u}(1)^{20}`$ invariance coming from the torus $`T^{20}`$. This symmetry is generated by $`20`$ $`U(1)`$ Kac Moody currents $`J_{u(1)}^i;i=1,\mathrm{},20`$. B-WS Virasoro symmetry This symmetry, which splits into holomorphic and antiholomorphic sectors, is given by the Suggawara construction using quadratic Casimirs of the previous WS affine Kac Moody algebras. For the holomorphic sector, the WS Virasoro currents of a bosonic string on $`AdS_3\times S^3\times T^{20}`$ are: (a) String on $`AdS_3`$: $$T_{sl(2)}^{WS}=\frac{1}{(k2)}[(J_{sl(2)}^0)^2(J_{sl(2)}^1)^2(J_{sl(2)}^2)^2].$$ (4.11) (b) String on $`S^3`$: $$T_{su(2)}^{WS}=\frac{1}{(k+2)}[(J_{su(2)}^0)^2+(J_{su(2)}^1)^2+(J_{su(2)}^2)^2]$$ (4.12) (c) String on $`T^{20}`$: $$T_{u(1)}^{WS}=\underset{i=1}{\overset{20}{}}[J_{u(1)}^i]^2.$$ (4.13) Similar quantities are also valid for the antiholomorphic sector of the conformal invariance. Note that the total WS energy momentum tensor $`T_{tot}^{WS}`$ is given by the sum of $`T_{sl(2)}^{WS},T_{su(2)}^{WS}`$ and $`T_{u(1)}^{WS}`$ eqs(4.11-12-13). In the case of a superstring propagating on $`AdS_3\times S^3\times T^4`$, the above conserved currents are slightly modified by the adjunction of extra terms due to contributions of WS fermions . If we denote by $`\mathrm{\Psi }_{sl(2)}^A,\mathrm{\Psi }_{su(2)}^a`$ and $`\mathrm{\Psi }_{u(1)}^i`$, the $`AdS_3,S^3`$ and $`T^4`$ fermions, the WS theory has a N=1 superconformal theory generated by: $`T(z)`$ $`=`$ $`{\displaystyle \frac{1}{k}}[(J_{sl(2)}^AJ_{sl(2),A}\mathrm{\Psi }_{sl(2)}^A\mathrm{\Psi }_{sl(2),A})+`$ $`(J_{su(2)}^aJ_{su(2),a}\mathrm{\Psi }_{su(2)}^a\mathrm{\Psi }_{su(2),a})]+`$ $`1/2{\displaystyle \underset{i=1}{\overset{4}{}}}((J_{u(1)}^iJ_{u(1)}^i\mathrm{\Psi }_{u(1)}^i\mathrm{\Psi }_{u(1)}^i)`$ $`G(z)`$ $`=`$ $`{\displaystyle \frac{2}{k}}[\mathrm{\Psi }_{sl(2)}^AJ_{sl(2),A}{\displaystyle \frac{i}{3k}}ϵ_{ABC}\mathrm{\Psi }_{sl(2)}^A\mathrm{\Psi }_{sl(2)}^B\mathrm{\Psi }_{sl(2)}^C]+`$ (4.14) $`{\displaystyle \frac{2}{k}}[\mathrm{\Psi }_{su(2)}^aJ_{su(2),a}{\displaystyle \frac{i}{3k}}ϵ_{abc}\mathrm{\Psi }_{su(2)}^a\mathrm{\Psi }_{su(2)}^b\mathrm{\Psi }_{su(2)}^c]`$ $`+{\displaystyle \underset{i=1}{\overset{4}{}}}\mathrm{\Psi }_{u(1)}^iJ_{u(1)}^i.`$ Note that to get a space time supersymetric vaccum, one should enhance the previous N=1 superconformal WS invariance to a N=2 conformal symmetry . This requires the existance of a conserved U(1) current in the world sheet theory under which G splits in two parts $`G^+`$ and $`G^{}`$ with charges +1 and -1 respectively. Skiping the details and denoting by $`G_r^\pm `$ the modes of $`G^\pm (z)`$ N=2 fermions currents; the N=2 U(1) superconformal algebras read as: $`[G_r^{},G_s^+]`$ $`=`$ $`2L_{r+s}(rs)J_{r+s}+(c/3)(r^21/4)\delta _{r+s,0}`$ $`[L_n,L_m]`$ $`=`$ $`(nm)L_{m+n}+{\displaystyle \frac{c}{12}}m(m^21)\delta _{m+n,0}`$ $`[L_n,G_r^\pm ]`$ $`=`$ $`({\displaystyle \frac{n}{2}}r)G_{n+r}^\pm `$ (4.15) $`[L_n,J_m]`$ $`=`$ $`mJ_{m+n}`$ $`[J_m,J_n]`$ $`=`$ $`{\displaystyle \frac{c}{3}}m\delta _{m+n,0}`$ $`[J_n,G_r^\pm ]`$ $`=`$ $`\pm G_{n+r}^\pm ,`$ where the r and s modes take half odd integer values for the Neveu Schwarz(NS) sector and integer ones for the Ramond (R)sector. Before going ahead we would like to make a pause in order to give some relevant features of these algebras. this pause is motivated by the two following: First the $`N=2`$ NS and R conformal algebras have a spectral flow which we want to use in order to complete the study of section 2 on FSS by giving a new result. Second space time symmetry of superstring on $`AdS_3\times S^3\times T^4`$ has a N=4 superconformal invariance which have a spectral flow of the same nature as for N=2 U(1) conformal invariance. Like for the FSS case, the spectral flow of $`N=2`$ and $`N=4`$ conformal invariances may also be used to study RdTS supersymmetry. ## 5 FSS and spectral flow In section 2, we have defined FSS as a hidden finite dimensional invariance which survives after integrable deformations of critical models such as the thermal deformation of $`Z_N`$ models; see eqs(2.11-12). There, we had exposed a method for deriving FSS algebras from parafermionic invariance. In the present section we want to complete the study of section 2 by giving a new way for obtaining FSS using topological field theory ideas . This method is based on using an appropriate choice of the parameter $`\eta `$ of the spectral flow of $`N=2`$ and $`N=4`$ superconformal theories. We will also take the opportunity of analysing the spectral flow of $`N=2`$ and $`N=4`$ conformal symmetries to make a comment on the recent proposal of where a new construction of fractional supersymmetric algebras was derived by using infinite dimensional modules of Lie algebras. To start recall that due to boundary conditions of fermions, the $`2dN=2`$ $`(N=4)`$ superconformal algebra has two sectors: Neveu Schwarz (NS) sector and Ramond (R) sector. These two sectors are not completely independent since they may be related by a continuous spectral flow as shown herebelow: $`U_\theta L_nU_\theta ^1`$ $`=`$ $`L_n+\theta J_n+c/6\theta ^2\delta _{n,0}`$ $`U_\theta J_nU_\theta ^1`$ $`=`$ $`J_n+c/3\theta \delta _{n,0}`$ (5.1) $`U_\theta G_r^+U_\theta ^1`$ $`=`$ $`G_{r+\theta }^+`$ $`U_\theta G_r^{}U_\theta ^1`$ $`=`$ $`G_{r\theta }^{},`$ for N=2 theories $`T_n^3(\eta )`$ $`=`$ $`T_n^3(0){\displaystyle \frac{\eta kp}{2}}\delta _{n,0}`$ $`T_{n_\pm \eta }^\pm (\eta )`$ $`=`$ $`T_n^\pm (0)`$ $`Q_{n+n/2}^1(\eta )`$ $`=`$ $`Q_n^1(0)`$ (5.2) $`Q_{n\eta /2}^2(\eta )`$ $`=`$ $`Q_n^2(0)`$ $`L_n(\eta )`$ $`=`$ $`L_n(\eta )\eta T_n^3(0)+\eta ^2({\displaystyle \frac{kp}{4}})\delta _{n,0}`$ for N=4 superconformal ones. The variable $`\eta `$ is the parameter of the spectral flow. Eqs (5.1-2) mean that $`2dN=2`$ $`(N=4)`$ superconformal algebras have then a continuous one parameter sector interpolating between NS and R algebras. This interpolating sector is generated by mode operators $`G_{r\pm \eta }^\pm `$ and $`\overline{G}_{r\pm \eta }^\pm `$ carrying shifted values of $`L_0`$ and the $`U(1)`$ charge operators. For a generic value of $`\eta `$, the commutation relations of the $`N=2`$ superconformal algebra in two dimensions read as: $`\{G_{r+\eta }^+\overline{G}_{s\eta }^{}\}`$ $`=`$ $`2L_{r+s}(rs+2\eta )J_{r+s}+(c/3)((r+\eta )^21/4)\delta _{r+s,0}`$ $`[L_n,L_m]`$ $`=`$ $`(nm)L_{m+n}+{\displaystyle \frac{c}{12}}m(m^21)\delta _{m+n,0}`$ $`[L_n,G_{r\pm \eta }^\pm ]`$ $`=`$ $`({\displaystyle \frac{n}{2}}r\eta )G_{n+r\pm \eta }^\pm `$ $`[L_n,J_m]`$ $`=`$ $`mJ_{m+n}`$ (5.3) $`[J_m,J_n]`$ $`=`$ $`{\displaystyle \frac{c}{3}}m\delta _{m+n,0}`$ $`[J_n,G_{r+\eta }^\pm ]`$ $`=`$ $`\pm G_{n+r+\eta }^\pm `$ $`\{G_{r+\eta }^+\overline{G}_{s\eta }^+\}`$ $`=`$ $`0`$ $`\{G_{r+\eta }^{}\overline{G}_{s\eta }^{}\}`$ $`=`$ $`0`$ Similar eqs may be written down for the $`N=4`$ case. Eqs(5.3) define a continuous one family parameter superconformal algebra to which we shall refer herebelow to as the $`\eta `$ sector and denote it as $`[(12\eta )NS,2\eta R]`$. For $`\eta =0`$, one discovers the NS algebra and for $`\eta =\frac{1}{2}`$ one gets the R algebra. For $`\eta `$ ranging between zero and $`\frac{1}{2}`$, one has the twisted sector. The $`[(12\eta )NS,2\eta R]`$ twisted conformal algebra plays a crucial role in topological field theories \[28,30 \]and allows to make spectacular transformations such as modifying the spins of the WS field operators by making appropriate choices of $`\eta `$. Taking the spectral parameter $`\eta =\frac{1}{2}`$, a fermion transforms into a boson (scalar or vector) while taking $`\eta =\frac{1}{k}`$, $`k>2`$, it becomes a WS parafermion of spin ($`1\pm \eta `$) depending on the $`U(1)`$ charge of the initial fermion. Putting back $`\eta =\frac{1}{k}`$ into eqs(5.3), one gets amongst others: $$2P_1=\{G_{(k1)/k}^+,G_{1/k}^{}\}+\{G_{1/k}^+,G_{(k1)/k}^{}\},$$ (5.4) together with: $$\begin{array}{ccc}0=\{G_{1/k}^\pm ,G_{1/k}^\pm \}\hfill & & \\ 0=\{G_{(k1)/k}^\pm ,G_{(k1)/k}^\pm \}.\hfill & & \end{array}$$ (5.5) Now comparing these relations with eqs(2.15) we obtained by thermal deformation of the $`Z_k`$ parafermoinic invariance, one discovers that they are quite similar. Eq (5.4) gives just a linearisation form of FSS which coincides with eqs(2.15) by setting $`k=3`$. Moreover eqs(5.5) show that $`G_{1/k}^\pm `$ are anticommuting operators in agreement with the result of . Furthermore starting from eqs(5.4-5) and following the reasoning of section 2 which lead to the derivation of eqs(2.15), one sees that it is possible to reinterpret the minus charge carried by $`G_{\frac{(1k)}{k}}^{}`$ as a $`Z_k`$ charge. So $`G_{\frac{(1k)}{k}}^{}`$ may be viewed as as composite operator given by the product of $`(k1)`$ $`G_{\frac{1}{k}}^+`$. This property is also supported by the fact that the N=2 superconformal currents have mode expansion operators with twisted values. $$G^\pm (z_1)\mathrm{\Phi }_m(z_2)=z_{12}^{n1p/k}G_{n\pm (p\pm 1)/k}^\pm \mathrm{\Phi }_m(z_2).$$ (5.6) Using these modes operators, one may write for $`k=3`$ the following relations $$G_{\frac{2}{3}}^{}=G_{\frac{2}{3}}^+G_0^+.$$ (5.7) Spectral flow of $`N=2`$ superconformal theories gives then an other way to build FSS algebras. In this regards, it is interesting to note that this spectral flow analysis maight also be used to rederive the so called FSUSY algebras considered recently in . We suspect that the fractional quantum numbers considered in when deriving FSUSY from special Verma modules of finite dimensional Lie algebras g could be rederived by taking fractional values of the spectral parameters $`\eta `$ of the corresponding Kac-Moody algebra $`\widehat{g}`$. Recall in passing that under the spectral flow, the step generators $`J_n^\alpha `$ and the Cartan ones $`H_n^i`$ of $`\widehat{g}`$ transform as: $$\begin{array}{ccc}J_n^\alpha J_{n+\eta v.\alpha }^\alpha \hfill & & \\ H_n^iH_n^i+k\eta v^i\delta _{n,0},\hfill & & \end{array}$$ (5.8) where $`\alpha `$ are the roots of $`\widehat{g}`$ and v is a weight vector. This transformation shifts the eigenvalues of the $`H_n^i`$’s Cartan charge operators of $`\widehat{g}`$. By an appropriate choice of the free parameters in the shifted weight $`\frac{2k\eta }{\alpha ^2}\alpha ^iv^i`$ of $`\frac{2}{\alpha ^2}\alpha ^iH_0^i`$, one recovers the fractionality property of the quantum numbers used in the construction of FSUSY algebras . This issue will be exhibited in more details in future occasion. Now we turn to our main topic. ## 6 Space-time invariance To analyse the space-time infinite dimensional symmetries on the boundary of $`AdS_3`$, one may follow the same strategy that we have used for the study of WS invariances. First identify the space time affine Kac-Moody symmetries and then consider the space time conformal invariance and eventually the Casimirs of higher ranks. In this section we shall simplify a little bit the analysis of space-time invariance and focus our attention on the conformal symmetry on $`(AdS_3)`$. Some specific features on space time Kac-Moody symmetries will also be given in due time. We begin by noting that space time infinite invariances on the boundary of $`AdS_3`$ are intimately linked to the WS ones. For the case of a superstring propagating on $`AdS_3\times S^3\times T^4`$, we have already shown that there are various kinds of WS symmetries coming from the propagation on $`AdS_3,S^3`$ and $`T^4`$ respectively. In the $`\varphi `$ infinite limit,we want to show that one may use these WS symmetries to build new space time ones. A. Conformal invariance First of all, note that the global part of the WS $`SO(1,2)\times \overline{SO(1,2)}`$ affine invariance of a bosonic string on $`AdS_3`$, generated by $`J_0^q`$ and $`\overline{J}_0^q;q=0,\pm 1`$ may be realized in different ways. A tricky way, which turns out to be crucial in building space-time conformal invariance, is given by the Wakimoto realization . Classically, this representation reads in terms of the local coordinates $`(\mathrm{\Phi },\gamma ,\overline{\gamma })`$ as follows: $$\begin{array}{ccc}J_0^0=\gamma /\gamma 1/2/\gamma ,\hfill & & \\ J_0^{}=/\gamma ,\hfill & & \\ J_0^+=\gamma ^2/\gamma \gamma /\mathrm{\Phi }e^{2\mathrm{\Phi }}/\gamma .\hfill & & \end{array}$$ (6.1) Similar relations are also valid for $`\overline{J}_0^q`$; they are obtained by substituting $`\gamma `$ by $`\overline{\gamma }`$. Quantum mechanically, the charge operators $`J_0^q`$ and $`\overline{J}_0^q`$ are given in terms of the Laurent mode operators of the quantum fields $`\mathrm{\Phi },\gamma ,\overline{\gamma },\beta `$ and $`\overline{\beta }`$ by using eqs(4.10) and performing the Cauchy integrations: $$\begin{array}{ccc}J_0^q=\frac{dz}{2i\pi }J^q(z)\hfill & & \\ \overline{J}_0^q=\frac{d\overline{z}}{2i\pi }\overline{J}^q(z).\hfill & & \end{array}$$ (6.2) To build the space time conformal invariance on the $`AdS_3`$ boundary, we proceed by steps. First suppose that there exists really a conformal symmetry on the boundary of $`AdS_3`$ and denote the space time Virasoro generators by $`L_n`$ and $`\overline{L}_n,nZ`$. The $`L_n`$ and $`\overline{L}_n`$, which should not be confused with the WS conformal mode generators, satisfy obviously the left and right Virasoro algebras. $$\begin{array}{ccc}[L_n,L_m]=(nm)L_{n+m}+c/12n(n^21)\delta _{n+m}\hfill & & \\ [\overline{L}_n,\overline{L}_m]=(nm)\overline{L}_{n+m}+\overline{c}/12n(n^21)\delta _{n+m}\hfill & & \\ [L_n,\overline{L}_m]=0.\hfill & & \end{array}$$ (6.3) The second step is to solve these eqs by using the string WS fields $`(\mathrm{\Phi },\gamma ,\overline{\gamma })`$ on $`AdS_3`$. To do so, it is convenient to divide the above eqs into two blocks. The first block corresponds to set $`n=0,\pm 1`$ in the generators $`L_n`$ and $`\overline{L}_n`$ of eqs(6.3). It describes the anomaly free projective subsymmetry the Virasoro algebra. The second block concerns the generators associated with the remaining values of n. On the boundary of $`AdS_3`$ obtained by taking the infinite limit of the $`\mathrm{\Phi }`$ field, one solves the projective subsymmetry by natural identification of $`L_q`$ and $`\overline{L}_q`$ ; $`q=0,\pm 1`$ with the zero modes of the WS $`so(1,2)\times \overline{so}(1,2)`$ affine invariance. In other words we have: $$\begin{array}{ccc}L_q\hfill & =& \hfill \frac{dz}{2i\pi }J^q(z)=J_0^q;q=0,\pm 1\\ \overline{L}_q\hfill & =& \hfill \frac{d\overline{z}}{2i\pi }\overline{J}^q(z)=\overline{J}_0^q;q=0,\pm 1.\end{array}$$ (6.4) Note that on the $`AdS_3`$ boundary, viewed as a complex plane parametrized by $`(\gamma ,\overline{\gamma })`$, the charge operators $`J_0^{}`$ ( $`L_1`$ ) and $`\overline{J}_0^{}(\overline{L}_1)`$ taken in the Wakimoto representation coincide respectively with the translation operators $`P_{}`$ and $`\overline{P}_+`$: $$\begin{array}{ccc}P_{}=L_{}=/\gamma \hfill & & \\ P_+=\overline{L}_{}=/\overline{\gamma }.\hfill & & \end{array}$$ (6.5) Eqs(6.4-5)) are interesting; they establish a link between the $`L_{}`$ and $`\overline{L}_{}`$ constants of motion of the boundary conformal field theory on $`AdS_3`$ on one hand and the $`P_{}(=P)`$ and the $`P_+(=\overline{P})`$ translation generators of the ST extension of the $`so(1,2)`$ algebra on the other hand. We will turn to these relations in the discussion of section 7. To get the rigourous solution of the remaining Virasoro charge operators $`L_n`$ and $`\overline{L}_n`$, one has to work hard. This is a lengthy and technical calculation which has been done in in connection with the study of the $`D_1/D_5`$ brane system. Later on we shall give some indications on this method; for the time being we shall use an economic path to work out the solution. This is a less rigourous but tricky way to get the same result. This method is based on trying to extend the $`L_n`$ and $`\overline{L}_n;n=0,\pm 1`$ projective solution to arbitrary integers n using properties of the string WS fields near the boundary, dimensional arguments and similarities with the photon vertex operator in three dimensions. Indeed using the holomorphic property of $`\gamma `$ and $`\overline{\gamma }`$ eqs(4.9 ) as well as the space time dimensional arguments; $$\begin{array}{ccc}[\gamma ]=1;J_{sl(2)}^0=0\hfill & & \\ J_{sl(2)}^{}=1;J_{sl(2)}^+=1,\hfill & & \end{array}$$ (6.6) it is not difficult to check that the following $`L_n(\overline{L}_n)`$ expressions are good condidates: $$𝐋_n=\frac{dz}{2i\pi }[a_0\gamma ^nJ_{sl(2)}^0\frac{a_{}}{2}\gamma ^{n+1}J_{sl(2)}^{}+\frac{a_+}{2}\gamma ^{n1}J_{sl(2)}^+],$$ (6.7) and a similar relation for $`\overline{L}_n`$. To get the $`a_i`$ coefficients, one needs to impose constraints which may be obtained by using results of BRST analysis in QED in three dimensions. Following ,the right constraints one has to impose on the $`a_i`$’ s are: $$\begin{array}{ccc}na_0+(n+1)a_{}+(n1)a_+=0\hfill & & \\ J^0\gamma (1/2)J^{}\gamma ^2(1/2)J^+=0.\hfill & & \end{array}$$ (6.8) The solution of the first constraint of these eqs reproducing the projective generators (6.4) is as follows: $$\begin{array}{ccc}a_0=(n^21)\hfill & & \\ a_{}=n(n1)\hfill & & \\ a_+=n(n+1)\hfill & & \end{array}$$ (6.9) Moreover using the second constraint of eqs(6.8) to express $`J_{sl(2)}^+(z)`$ in terms of $`J_{sl(2)}^0(z)`$ and $`J_{sl(2)}^{}(z)`$; then putting back into eqs(6.7), we find: $$𝐋_n=\frac{dz}{2i\pi }[(n+1)\gamma ^nJ_{sl(2)}^0+n\gamma ^{n+1}J_{sl(2)}^{}].$$ (6.10) Eqs (6.4) and (6.10) define the space time Virasoro algebra on the boundary of $`AdS_3`$. B. comments Having built the $`L_n`$’s space time Virasoro generators, one may be interested in determining the space-time energy momentum tensors $`T(\gamma )`$ and $`\overline{T}(\overline{\gamma })`$ of the BCFT on $`AdS_3`$. It turns out that the right form of the space-time energy momentum tensor depends moreover on auxiliary complex variables $`(y,\overline{y})`$ so that the space time energy momentum tensor has now two arguments; i.e: $`T=T(y,\gamma )`$ and $`\overline{T}=\overline{T}(\overline{y},\overline{\gamma })`$. Following , $`T(y,\gamma )`$ and $`\overline{T}(\overline{y},\overline{\gamma })`$ read as: $$\begin{array}{ccc}T(y,\gamma )=\frac{dz}{2i\pi }[\frac{_yJ(y,\gamma )}{(y\gamma )^2}\frac{_{}^{2}{}_{y}{}^{}J(y,\gamma )}{(y\gamma )}]\hfill & & \\ \overline{T}(\overline{y},\overline{\gamma })=\frac{d\overline{z}}{2i\pi }[\frac{_{\overline{y}J(\overline{y},\overline{\gamma })}}{(\overline{y}\overline{\gamma })^2}\frac{_{}^{2}{}_{\overline{y}}{}^{}J(\overline{y},\overline{\gamma })}{(\overline{y}\overline{\gamma })}],\hfill & & \end{array}$$ (6.11) where the currents $`J(y,\gamma )`$ and $`J(\overline{y},\overline{\gamma })`$ are given by: $$\begin{array}{ccc}J(y,\gamma )=J^+(y,\gamma )=2yJ^0(\gamma )J^+(\gamma )y^2J^{}(\gamma ).\hfill & & \end{array}$$ (6.12) In connection to these eqs, it is interesting to note that the conserved currents $`J^q(y,\gamma )`$ and $`J^q(\overline{y},\overline{\gamma })`$ are related to the WS affine Kac-Moody ones on $`AdS_3`$ as follows: $$\begin{array}{ccc}J^+(y,\gamma )=e^{yJ_0^{}}J^+(\gamma )e^{yJ_0^{}}=J^+(\gamma )2yJ^0(\gamma )+y^2J^{}(\gamma )\hfill & & \\ J^0(y,\gamma )=e^{yJ_0^{}}J^0(\gamma )e^{yJ_0^{}}=J^0(\gamma )yJ^{}(\gamma )=\frac{1}{2}_zJ^+(y,\gamma )\hfill & & \\ J^{}(y,\gamma )=e^{yJ_0^{}}J^{}(\gamma )e^{yJ_0^{}}=J^{}(\gamma )=\frac{1}{2}_z^2J^+(y,\gamma ).\hfill & & \end{array}$$ (6.13) and analogous eqs for $`J^q(\overline{y},\overline{\gamma })`$. Putting eqs(6.12) back into eqs(6.11) and expanding in power series of $`\frac{\gamma }{y}`$, one discovers the $`L_n`$ space time Virasoro generators given by eqs( 6.10). The second comment we want to make concerns the building of space time affine Kac-Moody symmeties out of the WS ones. Staring from WS conserved currents $`E_{ws}^a(z)`$, which may be thought of as $`J_{sl(2)}^q(z)`$, and going to the boundary of $`AdS_3`$, the corresponding space time affine Kac-Moody currents $`E_{spacetime}^a(y,\gamma )`$ read as: $$E_{spacetime}^a(y,\gamma )=\frac{dz}{2i\pi }[\frac{E_{ws}^a(z)}{(y\gamma (z))}].$$ (6.14) Expanding this eq in powers of $`\frac{y}{\gamma }`$ or $`\frac{\gamma }{y}`$, one gets the space time affine Kac-Moody modes: $$E_n^{a,spacetime}=\frac{dz}{2i\pi }[E_{ws}^a(z)\gamma ^n].$$ (6.15) The third comment we want to make concerns superstrings on $`AdS_3\times S^3\times T^4`$. In addition to the bosonic sector,there are moreover contributions coming from the WS fermions $`\psi _{ws}(z)`$. On the $`AdS_3`$ space for which the WS fermions $`\psi _{ws}^q(z)`$, $`q=o,\pm `$ transform in the $`SO(1,2)`$ adjoint, the total level k $`SO(1,2)`$ currents $`J_{sl(2),Total}^q(z)`$ now have two contributions: a level (k+2) bosonic current $`J_{sl(2),Bose}^q(z)`$ and a level (-2) fermionic current $`J_{sl(2),Fermi}^q(z)`$. The same construction may also be done for both $`S^3`$ and $`T^4`$. Note finally that in the limit $`\varphi `$ goes to infinity, the space time conformal symmetry of a superstring propagating on $`AdS_3\times S^3\times T^4`$ form a $`N=4`$ conformal invariance. ## 7 Discussions and Conclusion We have learned hereabove that on $`AdS_3\times N^d`$ may live various boundary conformal field theories depending on the choice of the d-dimensional compact manifold $`N^d`$. In the case of critical models of (super) strings propagating on $`AdS_3\times N^d`$, we have studied two examples: (i) $`N^d`$ is given by $`T^23`$ torus. (ii) $`N^d`$ is given by $`S^3\times T^4`$. The first example describes a bosonic BCFT while the second one describes a $`N=4`$ BCFT. One may also considers other choices of $`N^d`$ and build other BCFT’s. If one forgets about string dynamics as well as the nature of the compact manifold $`N`$ and just retains that on $`(AdS_3)`$ lives a conformal structure, one may consider its highest weight representations which read in general as: $$\begin{array}{ccc}L_0\|h,\overline{h}=h\|h,\overline{h},\hfill & & \\ L_n\|h,\overline{h}=0;n1\hfill & & \\ \overline{L}_0\|h,\overline{h}=\overline{h}\|h,\overline{h},\hfill & & \\ \overline{L}_n\|h,\overline{h}=0;n1,\hfill & & \\ cI\|h,\overline{h}=c\|h,\overline{h},\hfill & & \end{array}$$ (7.1) where $`\|h,\overline{h}`$ are Virasoro primary states. A priori the central charge c and the conformal weights h and $`\overline{h}`$ of these representations are arbitrary. However requiring unitary conditions, the parameters c, h and $`\overline{h}`$ are subject to constraints which become more stronger if one imposes extra symmetries such as supersymmetry or parafermionic invariance. Having these details in mind, one may also build descendant states $`\|h+n,\overline{h}+\overline{n}`$ of $`\|h,\overline{h}`$ from the primary ones as follows: $$\|h+n,\overline{h}+\overline{n}=\underset{\stackrel{n={\scriptscriptstyle \alpha _in_i}}{\overline{n}={\scriptscriptstyle \beta _jn_j}}}{}\lambda _{\{\alpha _i\}\{\beta _j\}}(\mathrm{\Pi }_iL_{n_i}^{\alpha i})(\mathrm{\Pi }_j\overline{L}_{n_j}^{\beta j})\|h,\overline{h}.$$ (7.2) where the $`\alpha _i`$’s and $`\beta _j`$’s are positive integers and $`\lambda _{\alpha \beta }`$ are C-numbers which we use to denote the collective coefficients $`\lambda _{\{\alpha _i\}\{\beta _j\}}`$ of the decomposition eq(7.2). They satisfy the following obvious relations. $$\begin{array}{ccc}L_0\|h+n,\overline{h}+\overline{n}=(h+n)\|h+n,\overline{h}+\overline{n}\hfill & & \\ L_\pm \|h+n,\overline{h}+\overline{n}=a_\pm (h,n)\|h\pm n,\overline{h}\pm \overline{n}\hfill & & \\ \overline{L}_0\|h+n,\overline{h}+\overline{n}=(\overline{h}+\overline{n})\|h+n,\overline{h}+\overline{n}\hfill & & \\ \overline{L}_\pm \|h+n,\overline{h}+\overline{n}=\overline{a}_\pm (\overline{h},\overline{n})\|h\pm n,\overline{h}\pm \overline{n}\hfill & & \end{array}$$ (7.3) where $`a_\pm (h,n)`$ and $`\overline{a}_\pm (\overline{h},\overline{n})`$ are normalization factors. Making an appropriate choice of the $`\lambda _{\alpha \beta }`$ coefficients and taking the $`a_\pm (h,n)`$ and $`\overline{a}_\pm (h,\overline{n})`$ coefficients as given herebelow, $$\begin{array}{ccc}a_{}(h,n)=\sqrt{(2h+n)(n+1)}\hfill & & \\ a_+(h,n)=\sqrt{(2h+n1)n},\hfill & & \end{array}$$ (7.4) one can get the two $`so(1,2)`$ modules used in building RdTS supersymmetry. Note that the descendant states $`\|h+n,\overline{h}+\overline{n}`$ are also eigenstates of the spin $`(L_0\overline{L}_0)`$ and conformal scale $`(L_0+\overline{L}_0)`$ operators of eigenvalues $`s=[(h\overline{h})+(n\overline{n})]`$ and $`\mathrm{\Delta }=[(h+\overline{h})+(n+\overline{n})]`$ respectively. We conclude this study by saying that the RdTS extension of Poincaré invariance in (1+2) dimensions that we have been describing is a special kind of FSS algebra. Like for FSS invariances, the RdTS generalised algebra may be also viewed as a residual symmetry of a boundary conformal invariance living on (1+2) space time manifolds. The RdTS supersymmetry we have described is a special FSS because it is related to a deformation of the space time boundary conformal invariance on $`AdS_3`$. In the end of this study, we should say that the explicit analysis of this paper has been plausible due to the particular properties of the $`AdS_3`$ geometry: (a) the $`AdS_3`$ manifold carries naturally a $`so(1,2)`$ affine invariance which has various realisation ways. (b) the Wakimoto realisation of the $`SO(1,2)`$ affine symmetry which on one hand relates its zero mode to the projective symmetry of a BCFT on $`AdS_3`$ and on the other hand links the $`L_{}`$ and $`\overline{L}_{}`$ to the translation operators on $`AdS_3`$. (c)the correspondance between WS and space time symmetries which plays a crucial role in analysing the various kinds of symmetries living on $`AdS_3`$. At last we would like to note that this study maight find a natural application in FQH systems formulated as an effective Chern-Simon gauge theory. In this model, the physics in the bulk is roughly speaking described by a $`(1+2)`$ dimensional $`U(1)^n`$ gauge while the theory the edge excitations of FQH liquids are described by a boundary conformal field theory. We plan to extend the results of this paper to the case of FQH droplets in a future occasion. ## 8 Aknowledgements The authors would like to thank Dr Rausch de Traubenberg for discussions, suggestions and for reading the manuscript. EHS would like to thank E.M Sahraoui for earlier collaboration on strings on Anti-de Sitter space and A Belhaj for discussions. This research work is supported by the program PARS Physique 27 under contract 372-98 CNR. 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warning/0007/cond-mat0007112.html	ar5iv	text	# On the approximation of Feynman-Kac path integrals for quantum statistical mechanics ## Abstract Discretizations of the Feynman-Kac path integral representation of the quantum mechanical density matrix are investigated. Each infinite-dimensional path integral is approximated by a Riemann integral over a finite-dimensional function space, by restricting the integration to a subspace of all admissible paths. Using this process, a wide class of methods can be derived, with each method corresponding to a different choice for the approximating subspace. The traditional “short-time” approximation and “Fourier discretization” can be recovered from this approach, using linear and spectral basis functions respectively. As an illustration, a novel method is formulated using cubic elements and is shown to have improved convergence properties when applied to a simple model problem. The path integral approach provides a powerful method for studying properties of quantum many-body systems . When applied to statistical mechanics , each element of the quantum density matrix is expressed as an integral over all curves connecting two configurations: $$\rho (𝐛,𝐚)=\text{ }\text{ }_{𝐚:𝐛}𝒟\left[𝐱\left(\tau \right)\right]\mathrm{exp}\left\{\frac{1}{\mathrm{}}\mathrm{\Phi }[𝐱\left(\tau \right);\beta ]\right\}.$$ (1) The symbol $`𝒟\left[𝐱\left(\tau \right)\right]`$ indicates that the integration is performed over the set of all differentiable curves, $`𝐱:[0,\beta \mathrm{}]𝐑^d`$, with $`𝐱\left(0\right)=𝐚`$ and $`𝐱\left(\beta \mathrm{}\right)=𝐛`$. The integer $`d`$ reflects the dimensionality, with $`d=3N`$ for a system of $`N`$-particles in 3-dimensional space. The functional $`\mathrm{\Phi }`$ can be derived from the classical action by introducing a relationship between temperature and imaginary time ($`it=\beta \mathrm{}`$) . In this Letter, we will restrict our attention to the quantum many-body system, for which $`\mathrm{\Phi }`$ takes the following form: $$\mathrm{\Phi }[𝐱\left(\tau \right);\beta ]=_0^\beta \mathrm{}\frac{1}{2}\underset{i=1}{\overset{d}{}}m_i\dot{x}_i\left(\tau \right)^2+V\left[𝐱\left(\tau \right)\right]\mathrm{d}\tau .$$ (2) Calculating the path integral in (1) is a challenging task, which in general cannot be performed analytically. It is only for simple model problems, such as quadratic potentials that an exact solution can be obtained. For more challenging systems, the path integral has traditionally been estimated using either the “short-time” approximation (STA) or “Fourier discretization” (FD) . Many authors have proposed improvements to the standard STA and FD, using techniques such as improved estimators , partial averaging , higher-order exponential splittings , advanced reference potentials , semi-classical expansions , and extrapolation . The fundamental approach is the same in all of these methods: the path integral is reduced to a high (but finite) dimensional Riemann integral, which is approximated using either a Monte Carlo or Molecular Dynamics. The aim of this Letter is to provide a framework for the formulation of a wide class of methods for the discretization of quantum mechanical path integrals. The idea of approximating path integrals using a finite subset of basis functions has been suggested before in the literature. Davison was one of the first to consider the use of orthogonal function expansions in the representation of Feynman path integrals , although he did not explore truncating the expansion. In a related article on Wiener integration, Cameron proposed using a finite set of orthogonal basis functions, and investigated the convergence of Fourier (spectral) elements . In this Letter, we do not require that the basis functions are orthogonal, allowing for the direct comparison of the STA and FD methods. Although other authors have explored fundamental connections between the STA and FD methods, we are unaware of any comparison using the approach investigated here. In addition, our approach allows for the construction of new methods using general classes of orthogonal polynomials or finite elements. To illustrate the flexibility of this approach, we derive a new method, using compactly supported (Hermite) cubic splines (HCS), which is shown to exhibit improved efficiency when applied to model problems. To illustrate how one can use a subspace approximation to discretize the quantum density matrix in (1), we start by introducing a change of variables to simplify the boundary conditions and temperature dependence for each path integral: $`𝐱\left(\tau \right)=𝐚+\left(𝐛𝐚\right)\tau /\beta \mathrm{}+𝐲\left(\tau /\beta \mathrm{}\right)`$. Since the admissible paths, $`𝐱`$, satisfy the boundary conditions $`𝐱\left(0\right)=𝐚`$ and $`𝐱\left(\beta \mathrm{}\right)=𝐛`$, the reduced paths given by $`𝐲`$, will satisfy Dirichlet boundary conditions, $`𝐲\left(0\right)=𝐲\left(1\right)=\mathrm{𝟎}`$, independent of $`𝐚`$, $`𝐛`$, and $`\beta `$. Introducing this change of variables into (1), results in the following: $`\rho (𝐛,𝐚)={\displaystyle }\text{ }\text{ }{\displaystyle _{0:0}}𝒟\left[𝐲\left({\displaystyle \frac{\tau }{\beta \mathrm{}}}\right)\right]\times `$ (4) $`\mathrm{exp}\left\{{\displaystyle \frac{1}{\mathrm{}}}\mathrm{\Phi }[𝐚+\left(𝐛𝐚\right){\displaystyle \frac{\tau }{\beta \mathrm{}}}+𝐲\left({\displaystyle \frac{\tau }{\beta \mathrm{}}}\right);\beta ]\right\}.`$ Note that the $`i`$th component of each reduced path $`𝐲`$, denoted by $`y_i`$, is a real-valued function on the interval $`[0,1]`$, satisfying Dirichlet boundary conditions. For the systems considered in this article, we also require that the derivative of each $`y_i`$ is measurable (i.e., square-integrable). Functions of this form are members of an infinite dimensional Sobolev space , defined by $`𝒮_0^1[0,1]=\left\{w𝒞[0,1]\right\|w\left(0\right)=w\left(1\right)=0\mathrm{a}ndw_𝒮<\mathrm{}\},`$ where $`w_𝒮^2_0^1\dot{w}\left(\xi \right)^2+w\left(\xi \right)^2\mathrm{d}\xi .`$ We proceed in the following manner to discretize (4): Consider a sequence of subspaces of increasing dimension $`𝒱_1,\mathrm{},𝒱_P,\mathrm{}𝒮_0^1[0,1]`$, where each $`𝒱_P`$ is of dimension $`P`$. For convenience, let each subspace be defined as the span of a particular set of basis functions: $`𝒱_P=\mathrm{s}pan\{\psi _1,\mathrm{},\psi _P\}`$. Now, given a component function $`y_i𝒮_0^1[0,1]`$, we can define its projection on $`𝒱_P`$ uniquely by $`y_i^{(P)}\left(\xi \right){\displaystyle \underset{k=1}{\overset{P}{}}}\alpha _{k,i}\psi _k\left(\xi \right),`$ Using the projection $`y_i^{(P)}\left(\xi \right)`$ as an approximation of $`y_i\left(\xi \right)`$ reduces the infinite-dimensional path integral in (4) to a finite-dimensional Riemann integral over the coefficients, $`\alpha _{i,k}`$: $`\stackrel{~}{\rho }(𝐛,𝐚)`$ $`=`$ $`{\displaystyle }\mathrm{d}𝜶J\mathrm{exp}\{{\displaystyle \frac{1}{\mathrm{}}}\mathrm{\Phi }[𝐚+(𝐛𝐚){\displaystyle \frac{\tau }{\beta \mathrm{}}}`$ (6) $`+𝐲^{(P)}\left({\displaystyle \frac{\tau }{\beta \mathrm{}}}\right);\beta ]\}.`$ Here, we have used simplified notation for the multi-dimensional integral, with $`\mathrm{d}𝜶_{k,i}\mathrm{d}\alpha _{k,i}`$. The constant $`J`$ reflects the particular choice of variables, and can be readily calculated (as discussed later). The reader should note that (6) does not depend on the basis functions chosen to represent the approximating subspace. If both $`\{\psi _1,\mathrm{},\psi _P\}`$ and $`\{\stackrel{~}{\psi }_1,\mathrm{},\stackrel{~}{\psi }_P\}`$ span $`𝒱_P`$, then there is an invertible linear transformation (i.e., change of variables) $`𝐔`$ such that $`\stackrel{~}{𝜶}=𝐔𝜶`$. To show in detail how subspace methods can be applied in practice, we consider the case of an $`N`$-body Hamiltonian system: $`\widehat{H}={\displaystyle \frac{1}{2}}{\displaystyle \underset{i=1}{\overset{d}{}}}m_i\widehat{p}_i^2+V[x_1,\mathrm{},x_d].`$ Here, the coordinate and momentum operators are denoted by $`x_i`$ and $`\widehat{p}_i`$ respectively. For this system the functional $`\mathrm{\Phi }`$ is given by (2), which when applied to the projected path, $`𝐱^{(P)}\left(\tau \right)𝐚+(𝐛𝐚)\tau /\beta \mathrm{}+𝐲^{(P)}\left(\tau /\beta \mathrm{}\right)`$, results in $`\mathrm{\Phi }[𝐱^{(P)}\left(\tau \right);\beta ]={\displaystyle _0^\beta \mathrm{}}{\displaystyle \underset{i=1}{\overset{d}{}}}{\displaystyle \frac{m_i}{2}}\left[\dot{x}_i^{(P)}\left(\tau \right)\right]^2+V\left[𝐱^{(P)}\left(\tau \right)\right]\mathrm{d}\tau .`$ After expanding the $`\tau `$-integrals, introducing a change of variables $`\xi =\tau /\beta \mathrm{}`$, and using the boundary conditions of each $`\psi _k`$, we have $`\mathrm{\Phi }[𝐱^{(P)}\left(\beta \mathrm{}\xi \right);\beta ]={\displaystyle \underset{i=1}{\overset{d}{}}}{\displaystyle \frac{m_i}{2\beta \mathrm{}}}\left\{\left(b_ia_i\right)^2+\stackrel{}{\alpha }_i^T𝐊\stackrel{}{\alpha }_i\right\}`$ (7) $`+\beta \mathrm{}{\displaystyle _0^1}V\left[𝐚+(𝐛𝐚)\xi +𝐲^{(P)}\left(\xi \right)\right]d\xi ,`$ (8) where $`\stackrel{}{\alpha }_i\left[\alpha _{1,i}\mathrm{}\alpha _{P,i}\right]^T`$. The “stiffness matrix”, $`𝐊𝐑^{P\times P}`$, has entries given by the inner-product $`K_{j,k}=_0^1\dot{\psi }_j\left(\xi \right)\dot{\psi }_k\left(\xi \right)d\xi `$ . Substituting (8) into (6), we obtain a simplified expression for the approximate density matrix: $`\stackrel{~}{\rho }(𝐛,𝐚)`$ $`=`$ $`\mathrm{exp}\left\{{\displaystyle \underset{i=1}{\overset{d}{}}}{\displaystyle \frac{m_i}{2\beta \mathrm{}^2}}\left(b_ia_i\right)^2\right\}`$ (9) $`\times `$ $`{\displaystyle }\mathrm{d}𝜶J\mathrm{exp}\{{\displaystyle \underset{i=1}{\overset{d}{}}}{\displaystyle \frac{m_i}{2\beta \mathrm{}^2}}\stackrel{}{\alpha }_i^T𝐊\stackrel{}{\alpha }_i`$ (10) $``$ $`\beta {\displaystyle _0^1}V[𝐚+(𝐛𝐚)\xi +𝐲^{(P)}\left(\xi \right)]\mathrm{d}\xi \},`$ (11) For the Fourier case, one typically calculates $`J`$ by requiring that the discretization be exact when applied to an ideal gas (i.e., $`V0`$) . Applying this same technique to a generic subspace method, and assuming that $`𝐊`$ is positive definite, one can solve for $`J`$ in a straightforward manner: $`J={\displaystyle \underset{i=1}{\overset{d}{}}}\sqrt{det𝐊}\left({\displaystyle \frac{m_i}{2\pi \beta \mathrm{}^2}}\right)^{\frac{P+1}{2}}.`$ Before discussing particular choices for basis functions, we should mention that, in general, the one-dimensional $`\xi `$-integral in (11) cannot be performed analytically. This problem has been traditionally circumvented by using a discrete approximation, such as Gaussian quadrature . For example, one can view the primitive STA as using the trapezoidal rule. If the quadrature scheme is of sufficiently high order its use will not reduce the asymptotic rate of convergence of the overall method. An optimal scheme must be efficient, since for nonlinear $`N`$-body systems evaluating $`V`$ may be computationally expensive. As mentioned above, the real benefit of using a general subspace approach is the flexibility afforded through the choice of basis functions. By considering a general class of pseudo-spectral or finite-element basis functions, a diverse group of discretizations can be constructed. Direct comparisons can be made between basis functions of varying smoothness and support. However, for brevity, we restrict our attention in this Letter to three different types of basis functions: linear, spectral, and cubic elements. Representative basis functions from each of these discretizations are shown in Figure 1. The traditional STA method can be constructed by considering polygonal paths, which can be represented by piecewise linear basis functions . For a given number of linear segments, $`P+1`$, we can define an approximating subspace $`𝒱_P`$ as the span of basis functions $`\{\psi _1,\mathrm{},\psi _P\}`$, where each $`\psi _k`$ is defined by the following formula $`\psi _k\left(\xi \right)`$ $`:=`$ $`\varphi ^{\mathrm{l}in}\left(\xi (P+1)k\right),`$ $`\mathrm{w}ith\varphi ^{\mathrm{l}in}\left(u\right)`$ $`:=`$ $`\{\begin{array}{ccc}1\left\|u\right\|& & u[1,1]\hfill \\ 0& & \mathrm{o}therwise\hfill \end{array}.`$ For this discretization, it is routine to show that the elements of the “stiffness” matrix, $`𝐊R^{P\times P}`$, can be determined by $`K_{i,j}=1\delta _{i1,j}+2\delta _{i,j}1\delta _{i+1,j}`$. In a similar manner, the FD method can be derived using the subspace approach by considering spectral basis functions of the form $`\psi _k\left(\xi \right)=1/k\mathrm{sin}\left(k\pi \xi \right)`$. $`𝐊`$ is diagonal for this basis, with entries given by $`K_{i,j}=\pi ^2/2\delta _{i,j}`$. A new method can be constructed by approximating the space of paths using piecewise (Hermite) cubic splines (HCS) . Each spline is defined on an interval of width $`2/P`$, with its shape uniquely determined by its function value and derivative at the ends of the interval. It is assumed here that $`P`$ is an even integer. Each piecewise cubic path has a continuous derivative, and is described by linear combinations of the basis functions $`\psi _k=\{\begin{array}{ccc}\varphi _1^{\mathrm{h}cs}\left(\xi P/2k\right)& & 1k<P/2\hfill \\ \varphi _2^{\mathrm{h}cs}\left(\xi P/2+P/2k\right)& & P/2kP\hfill \end{array},`$ where $`\varphi _1^{\mathrm{h}cs}\left(u\right)`$ $`:=`$ $`\{\begin{array}{ccc}\left(1\left\|u\right\|\right)^2\left(2\left\|u\right\|+1\right)& & u[1,1]\hfill \\ 0& & \mathrm{o}therwise\hfill \end{array}`$ $`\mathrm{a}nd\varphi _2^{\mathrm{h}cs}\left(u\right)`$ $`:=`$ $`\{\begin{array}{ccc}u\left(1\left\|u\right\|\right)^2& & u[1,1]\hfill \\ 0& & \mathrm{o}therwise\hfill \end{array}.`$ One can verify that the reduced path $`y^{(P)}(\xi )=\alpha _k\psi _k(\xi )`$ satisfies Dirichlet boundary conditions, and interpolates the interior grid points $`(2j/P,\alpha _j)`$ for integers $`1j<P/2`$. The derivative of the path at all the grid points is determined by the remaining $`P/2+1`$ coefficients, $`\alpha _k`$. Due to the compact support of the basis functions, the stiffness matrix is banded, with block structure: $`𝐊^{\mathrm{h}cs}={\displaystyle \frac{P}{60}}\left[\begin{array}{ccc}𝐊_1& & 𝐊_3\\ & & \\ 𝐊_3^T& & 𝐊_2\end{array}\right],`$ where the blocks are given by $`𝐊_1`$ $`=`$ $`\left[\begin{array}{cccc}72& 36& 0& .\\ 36& 72& .& 0\\ 0& .& 72& 36\\ .& 0& 36& 72\end{array}\right],`$ $`𝐊_2`$ $`=`$ $`\left[\begin{array}{cccc}4& 1& 0& .\\ 1& 8& .& 0\\ 0& .& 8& 1\\ .& 0& 1& 4\end{array}\right],`$ $`\mathrm{a}nd𝐊_3`$ $`=`$ $`\left[\begin{array}{cccccc}\hfill 3& \hfill 0& \hfill 3& \hfill 0& \hfill .& \hfill .\\ \hfill 0& \hfill 3& \hfill 0& \hfill .& \hfill .& \hfill 0\\ \hfill 0& \hfill .& \hfill .& \hfill 0& \hfill 3& \hfill 0\\ \hfill .& \hfill .& \hfill 0& \hfill 3& \hfill 0& \hfill 3\end{array}\right].`$ Note that the blocks are not all the same size, with $`𝐊_3`$ of dimension $`(P/21)\times (P/2+1)`$. The determinant of $`𝐊^{\mathrm{h}cs}`$ may be calculated exactly, but for most purposes it is enough to know that it is a constant, which will cancel out when (11) is used to calculate averages. As a numerical experiment, we apply each path integral discretization to the problem of calculating the average energy of a particle in a one-dimensional double-well. We have chosen the same double-well potential considered in , which is as follows: $`V\left(x\right)=m\omega ^2x^2/2+A/((x/a)^2+1)`$. The parameter values are all in atomic units, with $`\omega =0.006`$, $`A=0.009`$, $`a=0.09`$, and $`m=1836`$. At low temperatures, the energy is just above $`0.006`$, which is below the barrier height of $`0.009`$. To measure the accuracy of each method, we compute the energy at a fixed temperature of $`T=0.1\mathrm{}\omega /k`$, using Metropolis Monte Carlo to generate the canonically distributed configurations. The one-dimensional line-integrals of the potential are approximated using Simpson’s rule for the FD and HCS methods, and the traditional trapezoidal rule for the STA method. The number of integration nodes is set equal to the number of basis functions, $`P`$, resulting in the same number of potential evaluations for each method. For the STA method this results the potential is evaluated at the end points of each polygonal segment (consistent with its traditional implementation). It has been previously observed that averaged quantities (such as energy) converge at different rates, depending on the system, reference potential, and the form of the estimator . We use a virial estimator of the energy , $`E=V(x)+xV^{}(x)/2`$, which is known to exhibit improved convergence properties in many problems. The accuracy of each average is determined by comparing with the “exact” solution, computed by summing over the 15 lowest energy levels as calculated with Numerov’s method . In Fig. 2, the error in the computed energy is shown as a function of (a) the number of basis functions and (b) normalized CPU time. When the number of basis functions (or potential evaluations) is used as a measure of the work, we find that the FD and HCS methods are comparable, and both are more efficient than the STA method. However, when compared on the basis of CPU time, the HCS method is dramatically more efficient than both other methods. The inefficiency of the FD method for low-dimensional problems can be explained by considering the work required to compute $`P`$ points on the path. This work scales like $`O(P^2)`$ for the FD method, since the spectral basis functions are not compactly supported. On the other hand, for the STA and HCS methods this cost scales linearly with $`P`$. Although for very high dimensional problems, the cost of evaluating the potential should dominate, and we expect that the differences in computational cost would not be as pronounced. In summary, the problem of approximating Feynman-Kac path integrals can be addressed using the finite-dimensional subspace approach. This technique allows for the formulation of new methods through the choice of a suitable set of basis functions. In addition, traditional methods such as the short-time approximation and Fourier discretization methods can be compared using this framework. As an example, by considering (Hermite) cubic splines, a new method can be constructed which exhibits improved efficiency when applied to a one-dimensional double-well problem. The authors were supported by NSF Grant No. DMS-9627330.
warning/0007/math0007035.html	ar5iv	text	# Algebras without noetherian filtrations ## 1. Introduction In this paper all algebras will be defined over a fixed base field $`k`$. Let $`\mathrm{\Gamma }`$ be an $``$-filtration of a $`k`$-algebra $`A`$; thus, $`\mathrm{\Gamma }=\{\mathrm{\Gamma }_i:i\}`$ is an ascending chain of $`k`$-subspaces of $`A`$ satisfying $`1\mathrm{\Gamma }_0`$, $`_j\mathrm{\Gamma }_j=A`$ and $`\mathrm{\Gamma }_i\mathrm{\Gamma }_j\mathrm{\Gamma }_{i+j}`$ for all $`i,j`$. The filtration is finite if $`dim_k\mathrm{\Gamma }_i<\mathrm{}`$ for all $`i`$ and standard if $`\mathrm{\Gamma }_0=k`$ and $`\mathrm{\Gamma }_i=\mathrm{\Gamma }_1^i`$ for all $`i2`$. We say that $`\mathrm{\Gamma }`$ is a (left) noetherian filtration if the associated graded ring $`\mathrm{gr}A=\mathrm{gr}_\mathrm{\Gamma }A=_i\mathrm{\Gamma }_i/\mathrm{\Gamma }_{i1}`$ is (left) noetherian. The algebra $`A`$ is affine if it is finitely generated as a $`k`$-algebra. If an algebra $`A`$ has a left noetherian $``$-filtration, then a standard technique is to pull results back from $`\mathrm{gr}A`$ to $`A`$ since theorems are typically easier to prove in the graded ring $`\mathrm{gr}A`$ than in $`A`$. For related reasons Lorenz asked in \[6, p.436\] and \[7, Question III.4.2\] whether every left noetherian affine algebra $`R`$, satisfying a polynomial identity (PI), admits a left noetherian, standard finite $``$-filtration. This question has been raised again (without the “standard” hypothesis) in \[15, Question 6.16\] because of its importance for dualizing complexes and homological questions: if Lorenz’s question were to have a positive answer then $`R`$ would have a dualizing complex \[15, Corollary 6.9\] and every noetherian affine PI Hopf algebra would have finite injective dimension . The aim of this note is to answer these questions by providing a class of noetherian affine PI algebras which do not admit any noetherian finite $``$-filtration. The basic technique is provided by the following theorem (see Section 2). ###### Theorem 1.1. Suppose that $`I`$ and $`J_1J_2`$ are ideals of an algebra $`R`$ such that $`J_2/J_1`$ is free of rank $`s`$ as a left $`R/I`$-module and free of rank $`t`$ as a right $`R/I`$-module. If $`s<t`$, then there is no finite $``$-filtration $`\mathrm{\Gamma }`$ of $`R`$ such that $`\mathrm{gr}_\mathrm{\Gamma }R`$ is left noetherian. A variant of this theorem also holds if one replaces “free of rank $`x`$” by “of Goldie rank $`x`$.” See Theorem 3.2 for the details. A simple example satisfying the hypotheses (and conclusion) of the theorem is given by the ring (1.2) $$R=\{\left(\begin{array}{cc}f(x)& g(x)\\ 0& f(x^2)\end{array}\right):f,gk[x]\}M_2(k[x]).$$ Here, one takes $`J_2=I`$ to be the ideal of strictly upper triangular matrices and $`J_1=0`$. See Section 4 for more details. The most important case of Theorem 1.1 is when $`\mathrm{\Gamma }`$ is a standard. However, we emphasize that the theorem holds for any filtration satisfying the earlier definitions. The analogue of Theorem 1.1 also holds for $`𝔪`$-adic filtrations (where $`𝔪`$ is the Jacobson radical of a local algebra) and for that reason we prove the result for rings with a Zariskian filtration (see Theorem 2.4). In particular, we provide an example of a local prime noetherian PI ring for which the Jacobson radical does not satisfy the strong AR property. This is given in Section 4 where the reader may find further applications of the main theorem. ## 2. Proof of Theorem 1.1 Since we are also interested in $`𝔪`$-adic filtrations of local rings as well as ascending filtrations, we will prove our main result for $``$-filtrations. First we review some basic facts about filtrations from \[3, Chapter 6\] and . Suppose that $`A`$ is a $`k`$-algebra. For the purposes of this paper a filtration (or more strictly, an exhaustive separated finite $``$-filtration) of $`A`$ is an ascending chain of subspaces $`\mathrm{\Gamma }=\{\mathrm{\Gamma }_i\mathrm{\Gamma }_{i+1}\|i\}`$ of $`A`$, satisfying: 1. $`1\mathrm{\Gamma }_0`$ and $`\mathrm{\Gamma }_i\mathrm{\Gamma }_j\mathrm{\Gamma }_{i+j}`$ for all $`i,j`$; 2. $`\mathrm{\Gamma }`$ is finite in the sense that $`dim\mathrm{\Gamma }_i/\mathrm{\Gamma }_{i1}<\mathrm{}`$ for all $`i`$; 3. $`\mathrm{\Gamma }`$ is exhaustive in the sense that $`A=_i\mathrm{\Gamma }_i`$ and separated in the sense that $`_i\mathrm{\Gamma }_i=0`$. The Rees ring associated to $`\mathrm{\Gamma }`$ is defined to be $`\mathrm{Rees}A=\mathrm{Rees}_\mathrm{\Gamma }A=_i\mathrm{\Gamma }_i`$ and the associated graded ring is $`\mathrm{gr}A=\mathrm{gr}_\mathrm{\Gamma }A=_i\mathrm{\Gamma }_i/\mathrm{\Gamma }_{i1}.`$ Write $`J(A)`$ for the Jacobson radical of $`A`$. Following , a filtered algebra $`A`$ is called (left) Zariskian if (Zar1) $`\mathrm{\Gamma }_1J(\mathrm{\Gamma }_0)`$; (Zar2) $`\mathrm{Rees}_FA`$ is (left) noetherian. Fix a filtration $`\mathrm{\Gamma }`$ of the algebra $`A`$ and a left $`A`$-module $`M`$. The concept of a $``$-filtration (again, finite, exhaustive and separated) $`\mathrm{\Lambda }=\{\mathrm{\Lambda }_i:i\}`$ of $`M`$ is defined analogously; one simply replaces (1) in the above definition by (1) $`\mathrm{\Gamma }_i\mathrm{\Lambda }_j\mathrm{\Lambda }_{i+j}`$ for all $`i,j`$. Corresponding to this filtration one has the Rees module $`\mathrm{Rees}_\mathrm{\Lambda }M=\mathrm{\Lambda }_i`$ and associated graded module $`\mathrm{gr}_\mathrm{\Lambda }M=\mathrm{\Lambda }_i/\mathrm{\Lambda }_{i1}`$. We say that $`\mathrm{\Lambda }`$ is a good filtration if there exist $`\{m_i\mathrm{\Lambda }_{d_i}:1ir<\mathrm{}\}`$ such that $`\mathrm{\Lambda }_n=_{i=1}^r\mathrm{\Gamma }_{nd_i}m_i`$ for all $`n`$. Two filtrations $`\mathrm{\Lambda }`$ and $`\mathrm{\Lambda }^{}`$ of $`M`$ are equivalent, written $`\mathrm{\Lambda }\mathrm{\Lambda }^{}`$, if there is an integer $`q`$ such that $`\mathrm{\Lambda }_i\mathrm{\Lambda }_{i+q}^{}`$ and $`\mathrm{\Lambda }_i^{}\mathrm{\Lambda }_{i+q}`$ for all $`i.`$ The Hilbert function of $`M`$ with respect to $`\mathrm{\Lambda }`$ is defined to be $$H_{M,\mathrm{\Lambda }}(n)=dim\mathrm{\Lambda }_n/\mathrm{\Lambda }_n,\text{for all}n.$$ Let $`H`$ and $`H^{}`$ be two (Hilbert) functions $``$. We say that the growth of $`H`$ is at most the growth of $`H^{}`$, and write $`HH^{}`$, if there is an integer $`q`$ such that $`H(n)H^{}(n+q)`$ for all $`n0`$. We say $`H`$ and $`H^{}`$ are equivalent, written $`HH^{}`$, if both $`HH^{}`$ and $`H^{}H`$ hold. Since we require filtrations to be separated, they need not induce filtrations on factor modules. However, for Zariskian filtrations this is not a problem: ###### Lemma 2.1. Suppose that $`\mathrm{\Gamma }=\{\mathrm{\Gamma }_i\}`$ is a filtration of an algebra $`A`$ and let $`M`$ be a left $`A`$-module with a good filtration $`\mathrm{\Lambda }`$. 1. If $`\mathrm{\Lambda }^{}`$ is another filtration of $`M`$, there exists an integer $`q`$ such that $`\mathrm{\Lambda }_i\mathrm{\Lambda }_{i+q}^{}`$ for all $`i`$. Thus, if $`\mathrm{\Lambda }^{}`$ is good, then $`\mathrm{\Lambda }`$ and $`\mathrm{\Lambda }^{}`$ are equivalent. 2. If $`\mathrm{\Lambda }^{}`$ and $`\mathrm{\Lambda }^{\prime \prime }`$ are equivalent filtrations of $`M`$, then $`H_{M,\mathrm{\Lambda }^{}}`$ and $`H_{M,\mathrm{\Lambda }^{\prime \prime }}`$ are equivalent. 3. Assume that $`\mathrm{\Gamma }`$ is left Zariskian. If $`N`$ is a submodule of $`M`$ then $`\mathrm{\Lambda }`$ induces good filtrations on $`N`$ and $`M/N`$. In particular, $`\mathrm{\Gamma }`$ induces a left Zariskian filtration on every factor ring of $`A`$. ###### Proof. (1) and (2) follow from the definitions while (3) follows from \[4, Theorem 3.3\] or \[5, Theorem II.2.1.2\]. ∎ The key observation in this paper is given by the next proposition. Since our main theorem will come in three slightly different forms, this result will also have three slightly different cases. ###### Proposition 2.2. Let $`A`$ be an algebra with a filtration $`\mathrm{\Gamma }`$ such that the growth of $`H_{A,\mathrm{\Gamma }}`$ is subexponential. Let $`M`$ be an $`A`$-bimodule such that the left module $`{}_{A}{}^{}M`$ is free of rank $`s`$ and the right module $`M_A`$ is free of rank $`t`$. Suppose that there exists a filtration $`\mathrm{\Lambda }`$ on $`M`$ such that $`\mathrm{\Lambda }`$ is a good filtration of $`{}_{A}{}^{}M`$ and a filtration of $`M_A`$. 1. If $`\mathrm{\Lambda }`$ is also a good filtration of $`M_A`$ then $`t=s`$. 2. If $`\mathrm{\Gamma }_n=0`$ for $`n0`$ then $`st`$ 3. If $`\mathrm{\Gamma }_n=A`$ for $`n0`$ then $`st`$. ###### Proof. The beginning of the proof is the same in all three cases. Note that $`{}_{A}{}^{}A`$ has a good filtration, simply because $`\mathrm{\Gamma }_i=\mathrm{\Gamma }_i1`$ for all $`i`$. Thus $`{}_{A}{}^{}A_{}^{(s)}`$ also has a good filtration and so this induces a good filtration $`\mathrm{\Lambda }^{}`$ on $`M`$ such that $`\mathrm{\Lambda }_i^{}\mathrm{\Gamma }_i^{(s)}`$ for all $`i`$. By Lemma 2.1(1), there exists $`q_1`$ such that $`\mathrm{\Lambda }_i\mathrm{\Lambda }_{i+q_1}^{}`$ for all $`i`$. Similarly, the right $`A`$-module structure provides an induced good filtration $`\mathrm{\Lambda }^{\prime \prime }`$ of $`M_A`$ such that $`\mathrm{\Lambda }_i^{\prime \prime }\mathrm{\Gamma }_i^{(t)}`$ for all $`i`$. By Lemma 2.1(1), there exists $`q_2`$ such that $`\mathrm{\Lambda }_i^{\prime \prime }\mathrm{\Lambda }_{i+q_2}`$ for all $`i`$. Hence, for $`p=q_1+q_2`$, (2.3) $$\mathrm{\Lambda }_i^{\prime \prime }\mathrm{\Lambda }_{i+p}^{}\mathrm{for}\mathrm{all}i.$$ We now consider the three cases separately. Under assumption (2), we see that $`\mathrm{\Lambda }_i^{\prime \prime }/\mathrm{\Lambda }_i^{\prime \prime }=\mathrm{\Lambda }_i^{\prime \prime }\mathrm{\Lambda }_{i+p}^{}=\mathrm{\Lambda }_{i+p}^{}/\mathrm{\Lambda }_{(i+p)}^{}`$ for $`i0`$. Thus, $`H_{M,\mathrm{\Lambda }^{\prime \prime }}(n)H_{M,\mathrm{\Lambda }^{}}(n+p)`$. But, by construction, $`H_{M,\mathrm{\Lambda }^{}}=sH_{A,\mathrm{\Gamma }}`$ and $`H_{M,\mathrm{\Lambda }^{\prime \prime }}=tH_{A,\mathrm{\Gamma }}`$ and so $$H_{A,\mathrm{\Gamma }}(n)\frac{s}{t}H_{A,\mathrm{\Gamma }}(n+p)\mathrm{for}\mathrm{all}n.$$ Since $`H_{A,\mathrm{\Gamma }}(n)`$ grows subexponentially this forces $`ts`$. Under assumption (3) we have $`\mathrm{\Lambda }_{i+p}^{\prime \prime }/\mathrm{\Lambda }_{(i+p)}^{\prime \prime }=A/\mathrm{\Lambda }_{ip}^{\prime \prime }A/\mathrm{\Lambda }_i^{}=\mathrm{\Lambda }_i^{}/\mathrm{\Lambda }_i^{}`$ for $`i0`$. Thus, $`H_{M,\mathrm{\Lambda }^{}}(i)H_{M,\mathrm{\Lambda }^{\prime \prime }}(i+p)`$ and repeating the analysis of the last paragraph shows that $`st`$. Finally, assume that (1) holds. In this case, $`\mathrm{\Lambda }^{\prime \prime }\mathrm{\Lambda }`$ and so $`\mathrm{\Lambda }^{\prime \prime }\mathrm{\Lambda }^{}`$. Thus there exists $`p`$ such that $`\mathrm{\Lambda }_i^{}\mathrm{\Lambda }_{i+p}^{\prime \prime }\mathrm{\Lambda }_{i+2p}^{}`$. Now a minor variant of the penultimate paragraph shows that $`ts`$ and hence, by symmetry, that $`s=t`$. ∎ ###### Theorem 2.4. Suppose that $`R`$ is an algebra with ideals $`I`$ and $`J_1J_2`$ such that $`J_2/J_1`$ is free of rank $`s`$ as a left $`R/I`$-module and free of rank $`t`$ as a right $`R/I`$-module. If $`s<t`$, then there is no filtration of $`R`$ such that $`R`$ is both left and right Zariskian. ###### Proof. Let $`\stackrel{~}{\mathrm{\Gamma }}`$ be a left and right Zariskian filtration of $`R`$. Then Lemma 2.1(3) implies that the induced filtration $`\mathrm{\Gamma }`$ on $`A=R/I`$ is left and right Zariskian and so $`\mathrm{gr}_\mathrm{\Gamma }A`$ is left (and right) noetherian. By \[12, Remark after 1.2\] the growth of $`H_{A,\mathrm{\Gamma }}`$ is subexponential. By Lemma 2.1(3), $`\mathrm{\Gamma }`$ induces a good filtration on $`J_2/J_1`$ as a left and a right $`A`$-module, which contradicts to Proposition 2.2(1). ∎ A curious feature of this result is that $`R`$ could be left or right Zariskian; it just cannot be both (see Corollary 4.7). However, with a little more information one can determine which side goes wrong. ###### Theorem 2.5. Suppose that $`R`$ is an algebra with ideals $`I`$ and $`J_1J_2`$ such that $`J_2/J_1`$ is free of rank $`s`$ as a left $`R/I`$-module and free of rank $`t`$ as a right $`R/I`$-module. Assume that $`s<t`$ and that $`\mathrm{\Gamma }`$ is a filtration of $`R`$. 1. If $`\mathrm{\Gamma }_i=0`$ for $`i0`$, then $`\mathrm{\Gamma }`$ is not left Zariskian. 2. If $`\mathrm{\Gamma }_i=A`$ for $`i0`$, then $`\mathrm{\Gamma }`$ is not right Zariskian. ###### Proof. For (1), repeat the proof of Theorem 2.4, but with Proposition 2.2(1) replaced by Proposition 2.2(2). For (2), use the proof of Theorem 2.4, applied to the opposite ring $`R^{\mathrm{op}}`$, with Proposition 2.2(1) replaced by Proposition 2.2(3). ∎ We are now ready to prove Theorem 1.1 from the introduction. Proof of Theorem 1.1. It is easy to see that an $``$-filtration $`\mathrm{\Gamma }`$ of a ring $`R`$ is left Zariskian if and only if $`\mathrm{gr}_\mathrm{\Gamma }`$ is left noetherian (see, for example, \[5, Proposition II.1.2.3\]). Thus, Theorem 1.1 is a special case of Theorem 2.5(1). ∎ The analogue of Theorem 1.1 for complete local rings also holds with much the same proof. To state the result, we need some definitions. Let $`R`$ be a semilocal algebra with Jacobson radical $`J(R)=𝔪`$ and assume that $`dim_kR/𝔪<\mathrm{}`$. A filtration $`\mathrm{\Gamma }`$ of $`R`$ is called a weak adic filtration if it satisfies $`\mathrm{\Gamma }_n=A`$ for all $`n0`$ and $`\mathrm{\Gamma }_1𝔪`$. (We still require the filtration to be finite, separated and exhaustive.) ###### Corollary 2.6. Let $`R`$ be a complete semilocal right noetherian algebra with Jacobson radical $`𝔪`$. Suppose that $`R`$ has ideals $`I`$ and $`J_1J_2`$ such that $`J_2/J_1`$ is free of rank $`s`$ as a left $`R/I`$-module and free of rank $`t`$ as a right $`R/I`$-module. If $`s<t`$, then there is no weak adic filtration $`\mathrm{\Gamma }`$ such that $`\mathrm{gr}_\mathrm{\Gamma }R`$ is right noetherian. ###### Proof. Suppose that such a filtration $`\mathrm{\Gamma }`$ exists. By \[2, Corollary 1.2\] the filtration $`\mathrm{\Gamma }`$ is also complete (in the natural sense that Cauchy sequences modulo the $`\mathrm{\Gamma }_i`$ should converge—see \[5, Definition I.3.3.2\]). Thus \[5, Proposition II.2.2.1\] implies that the filtration is right Zariskian and the result follows Theorem 2.5(2). ∎ An analogue of this corollary also holds for non-complete rings, although the result is less pleasant since one cannot now assume that separated filtrations induce separated filtrations on factor modules. The result becomes the following: Suppose that $`R`$ satisfies the hypotheses of Corollary 2.6, but that $`R`$ is not complete. Let $`\mathrm{\Gamma }`$ be a weak adic filtration of $`R`$ that induces both a filtration on $`R/I`$ and a good right filtration of $`J_2/J_1`$. If $`s<t`$, then $`\mathrm{gr}_\mathrm{\Gamma }R`$ is not right noetherian. ## 3. A partial generalization Although the results of the last section are sufficient for our examples, one can give a version of Theorem 1.1 that works without the assumption that $`J_2/J_1`$ be free, but at the expense of of assuming that $`R/I`$ be a prime Goldie ring. We prove this in this section. Let $`A`$ be a prime left Goldie ring with simple artinian ring of fractions $`Q(A)`$. If $`M`$ is a left $`A`$-module, then the Goldie rank of $`M`$ is defined to be the length of $`Q(A)_AM`$ and written $`\mathrm{Grank}(M)`$. ###### Lemma 3.1. Suppose that $`A`$ is a prime Goldie ring with a left noetherian $``$-filtration $`\mathrm{\Gamma }`$. Let $`M`$ be an $`A`$-bimodule with a filtration $`\mathrm{\Lambda }`$ that is a good filtration of $`{}_{A}{}^{}M`$ and a filtration of $`M_A`$. If $`{}_{A}{}^{}M`$ is torsion, then so is $`M_A`$. ###### Proof. Assume that $`M_A`$ is not torsion. Replacing $`M`$ by $`M^{(n)}`$, for some $`n`$, we may assume that $`A_AM_A`$. Thus, Lemma 2.1(1) implies that there exists $`q0`$ such that $$dim\mathrm{\Gamma }_ndim\mathrm{\Lambda }_{n+q}Adim\mathrm{\Lambda }_{n+q}$$ for all $`n`$. Since $`\mathrm{\Gamma }_n=\mathrm{\Lambda }_n=0`$ for $`n0`$, this implies that $`H_{M,\mathrm{\Lambda }}H_{A,\mathrm{\Gamma }}`$. Now consider $`{}_{A}{}^{}M`$. Since $`\mathrm{\Lambda }`$ is a good filtration, $`{}_{A}{}^{}M`$ is finitely generated. We claim, for all $`p>0`$, that $`H_{M,\mathrm{\Lambda }}\frac{1}{p}H_{A,\mathrm{\Gamma }}`$. Once this has been proved, then the last paragraph implies, for some $`x>0`$, that $`H_{A,\mathrm{\Gamma }}(n)\frac{1}{p}H_{A,\mathrm{\Gamma }}(n+x)`$. This contradicts the fact that $`H_{A,\mathrm{\Gamma }}`$ grows subexponentially \[12, Remark after 1.2\] and proves the lemma. In order to prove the claim we ignore the right-hand structure of $`M`$ and so, by induction, it suffices to prove it for a cyclic module $`M=A/I`$. Since $`{}_{A}{}^{}M`$ is torsion and $`A`$ is Goldie, $`I`$ contains a regular element, $`a`$ say, of $`A`$. We will still write $`\mathrm{\Gamma }`$ for the good filtration on any subfactor of $`{}_{A}{}^{}A`$ induced from $`\mathrm{\Gamma }`$. Now, for any $`n`$, $`M`$ is a homomorphic image of $`L(i)=Aa^{i1}/Aa^iA/Aa`$ and so $`H_{M,\mathrm{\Lambda }}H_{L(i),\mathrm{\Gamma }}`$. Since Hilbert series are additive on short exact sequences, this implies that $`pH_{M,\mathrm{\Lambda }}H_{A/Aa^p,\mathrm{\Gamma }}H_{A,\mathrm{\Gamma }}`$, for any $`p`$. ∎ ###### Theorem 3.2. Suppose that $`I`$ and $`J_1J_2`$ are ideals of an algebra $`R`$ such that $`A=R/I`$ is a prime Goldie ring. Assume that $`J_2/J_1`$ has Goldie rank $`s`$ as a left $`A`$-module and Goldie rank $`t`$ as a right $`A`$-module, for some $`s<t`$. Then there is no $``$-filtration $`\mathrm{\Gamma }^{}`$ of $`R`$ such that $`\mathrm{gr}_\mathrm{\Gamma }^{}R`$ is left noetherian. ###### Proof. Suppose that such a filtration $`\mathrm{\Gamma }^{}`$ exists. Let $`\mathrm{\Gamma }`$ be the induced filtration on $`A`$ and $`\mathrm{\Lambda }`$ the induced filtration on $`M=J_2/J_1`$. Then $`\mathrm{\Lambda }`$ is a good left filtration and so $`{}_{A}{}^{}M`$ is finitely generated. The torsion submodule $`T`$ of $`M_A`$ is an $`A`$-bimodule and so we may pass to $`M/T`$ without affecting the hypotheses (although $`s`$ may decrease). By Lemma 3.1 $`{}_{A}{}^{}M`$ is also torsion-free. If $`\mathrm{Grank}(A)=u`$, replace $`M`$ by $`M^{(u)}`$; thus $`A_A^{(t)}XM_A`$ and $`{}_{A}{}^{}MY{}_{A}{}^{}A_{}^{(s)}`$. Let $`\mathrm{\Theta }`$, respectively $`\mathrm{\Phi }`$, be the filtrations of $`X`$ and $`Y`$ induced from $`\mathrm{\Gamma }`$. Since $`\mathrm{\Theta }`$ equals the direct sum $`\mathrm{\Gamma }^{(t)}`$ of $`t`$ copies of $`\mathrm{\Gamma }`$, certainly $`H_{X,\mathrm{\Theta }}=tH_{A,\mathrm{\Gamma }}`$. Similarly, $`H_{Y,\mathrm{\Phi }}=sH_{A,\mathrm{\Gamma }}`$. Since $`\mathrm{\Lambda }`$ is a good left filtration and $`\mathrm{\Theta }`$ is a good right filtration, by Lemma 2.1(1) there exists $`q0`$ such that $$dim\mathrm{\Theta }_ndim(X\mathrm{\Lambda }_{n+q})dim\mathrm{\Lambda }_{n+q}$$ and $$dim\mathrm{\Lambda }_ndim(M\mathrm{\Phi }_{n+q})dim\mathrm{\Phi }_{n+q},$$ for all $`n0`$. Hence, $$tH_{A,\mathrm{\Gamma }}(n)=H_{X,\mathrm{\Theta }}(n)H_{Y,\mathrm{\Phi }}(n+2q)=sH_{A,\mathrm{\Gamma }}(n+2q),$$ for all $`n0`$. Since $`H_{A,\mathrm{\Gamma }}`$ grows subexponentially, this forces $`ts`$. ∎ With a rather more complicated argument one can prove an analogous version of Theorem 2.4 using Goldie ranks. However, we do not know how to prove the analogous version of Theorem 2.5(2) or Corollary 2.6 without extra hypotheses. ## 4. Examples In this section we use the results of the Section 2 to provide examples of rings without noetherian associated graded rings. ###### Example 4.1. Let $$S=\{\left(\begin{array}{cc}f(x)& g(x)\\ 0& f(x^2)\end{array}\right):f,gk[x]\}+yM_2(k[x,y])M_2(k[x,y]).$$ Then $`S`$ is an affine noetherian prime PI algebra without any left noetherian finite $``$-filtrations. ###### Proof. The diagonal matrices of the form $$\left(\begin{array}{cc}f(x,y)& 0\\ 0& f(x^2,y)\end{array}\right):f(x,y)k[x,y]$$ clearly form a subring $`C`$ of $`S`$ isomorphic to $`k[x,y]`$. Since $`S`$ is finitely generated as a left or right $`C`$-module, it follows that $`S`$ is a finitely generated noetherian PI algebra. It is prime since it contains a nonzero ideal of the prime ring $`M_2(k[x,y])`$. Suppose that $`AB`$ is a surjective homomorphism of algebras and that $`B`$ does not admit a left noetherian $``$-filtration. Then an immediate consequence of Lemma 2.1(3) is that $`A`$ also has no such filtration. Thus, it suffices to prove the final assertion for a factor ring, and we use (4.2) $$R=S/yM_2(k[x,y])\{\left(\begin{array}{cc}f(x)& g(x)\\ 0& f(x^2)\end{array}\right):f,gk[x]\}M_2(k[x]).$$ Notice that this is the ring from (1.2). The nilradical $`N=N(R)`$ is just the set of strictly upper triangular matrices. It is routine to check that $`N`$ is a free left $`R/N`$-module of rank one but a free right $`R/N`$-module of rank two. Hence, by Theorem 1.1, neither $`R`$ nor $`S`$ can have a left noetherian finite $``$-filtrations. ∎ ###### Remark 4.3. The ring $`R`$ from (4.2) clearly has Gelfand-Kirillov dimension one. By inspecting the proof of Theorem 2.4, this shows that there is no finite filtration of $`R`$ that induces a good filtration of $`N(R)`$ as a left $`R`$-module. Other examples of noetherian rings with no noetherian $``$-filtration are given in . However, those examples require that the ring in question have infinite Gelfand-Kirillov dimension and this is impossible for affine PI algebras (see \[8, Proposition 13.10.6\]). It is readily checked that this ring $`R`$ does have a right noetherian, finite filtration (see Corollary 4.7). By slightly modifying the example one can get an example that “works” on both sides. ###### Example 4.4. Let $$T=\{\left(\begin{array}{ccc}f(x)& g(x)& h(x)\\ 0& f(x^2)& l(x)\\ 0& 0& f(x)\end{array}\right):f(x),\mathrm{},l(x)k[x]\}+yM_3(k[x,y]).$$ Then $`T`$ is a affine noetherian prime PI algebra such that, for every $``$-filtration, the associated graded ring is neither left nor right noetherian. ###### Proof. Notice that $`T/(e_{13}T+e_{23}T)R`$, the ring from (4.2), and so $`T`$ has no left noetherian finite filtration. Since $`TT^{op}`$, the same is true on the right. ∎ These examples can be easily modified into complete local rings and we give one example that is analogous to the ring $`R`$ from Example 4.1. Before giving the example, we need a definition. An ideal $`I`$ of a ring $`R`$ is said to satisfy the strong AR property if the associated Rees ring $`\mathrm{Rees}_IR=_{j0}I^j`$ is noetherian. The significance of this condition is that, if $`I`$ satisfies the strong AR property, then it also satisfies the usual Artin-Rees (AR) property; indeed this is the standard way of proving the latter condition in commutative algebra. The same idea has been useful in noncommutative algebra (see, for example, \[10, §2\] and ). Given that the Jacobson radical of a semi-local noetherian PI algebra is automatically AR (\[1, Theorems 3.1.13 and 7.2.5\]), it is natural to ask if it is also strongly AR. The next example shows that it does not. Notice that, although the ring in question is just a completion of the ring $`R`$ from Example 4.1, the Zariskian property fails on the opposite side. ###### Example 4.5. Let $$R=\{\left(\begin{array}{cc}f(x)& g(x)\\ 0& f(x^2)\end{array}\right):f,gk[[x]]\}+yM_2(k[[x,y]])M_2(k[[x,y]]).$$ Then $`R`$ a prime noetherian complete local PI algebra over $`k`$. Moreover, $`R`$ has no weak adic filtration $`\mathrm{\Gamma }`$ such that $`\mathrm{gr}_\mathrm{\Gamma }R`$ is right noetherian. In particular, $`J(R)`$ does not satisfy the strong AR property. ###### Proof. The proof that $`R`$ is prime noetherian PI algebra over $`k`$ is analogous to that for Example 4.1 and is left to the reader. Clearly $`R`$ is a complete local ring. In order to complete the proof it suffices, by Corollary 2.6, to work in a factor ring and we chose: (4.6) $$R^{}=R/yM_2(k[[x,y]])=\{\left(\begin{array}{cc}f(x)& g(x)\\ 0& f(x^2)\end{array}\right):f(x),g(x)k[[x]]\}.$$ The nilradical $`N`$ of $`R^{}`$ is the set of strictly upper triangular matrices and is free of rank $`1`$ as a left $`R^{}/N`$-module and free of rank $`2`$ as a right $`R^{}/N`$-module. Now apply Corollary 2.6. ∎ Finally we justify a remark made after Theorem 2.4: in that theorem it is possible to have a filtration that is Zariskian on either the left or the right. ###### Corollary 4.7. Consider noetherian, PI algebras $`R`$ with ideals $`I`$ and $`J_1J_2`$ such that $`J_2/J_1`$ is a free left $`R/I`$-module of rank one and a free right $`R/I`$-module of rank $`2`$. Then, there exists an example $`R_1`$ of such a ring with a left Zariskian filtration $`\mathrm{\Gamma }_1`$ and an example $`R_2`$ of such a ring with a right Zariskian filtration $`\mathrm{\Gamma }_2`$. ###### Proof. The ring $`R_2`$ is just the ring $`R`$ from (1.2), and so does satisfy the hypotheses of the first paragraph. Set (4.8) $$\alpha =\left(\begin{array}{cc}x& 0\\ 0& x^2\end{array}\right)\mathrm{and}\beta =\left(\begin{array}{cc}0& 1\\ 0& 0\end{array}\right).$$ It is easy to see that $`R`$ is generated by $`\alpha `$ and $`\beta `$ and hence one has the standard $``$-filtration $`\mathrm{\Gamma }_0=k`$, $`\mathrm{\Gamma }_1=k+k\alpha +k\beta `$ and $`\mathrm{\Gamma }_n=\mathrm{\Gamma }_1^n`$ for $`n2`$. Let $`a`$ and $`b`$ be the images of $`\alpha `$ and $`\beta `$ in $`\mathrm{gr}_\mathrm{\Gamma }R`$. Then, $`\mathrm{gr}_\mathrm{\Gamma }R`$ is generated by $`a`$ and $`b`$ and satisfies the relations $`a^2b=0=b\mathrm{gr}_\mathrm{\Gamma }Rb`$. Thus $`\mathrm{gr}_\mathrm{\Gamma }R`$ is spanned by the elements $`\{a^n,ba^n,aba^n:n0\}`$. As such, $`\mathrm{gr}_\mathrm{\Gamma }R`$ is a finitely generated right $`k[a]`$-module, and so is right noetherian. As in the proof of Theorem 1.1, this suffices to prove that $`\mathrm{\Gamma }`$ is right Zariskian. For $`R_1`$ we take the ring $`R^{}`$ from (4.6). We define $`\alpha `$ and $`\beta `$ by (4.8) and observe that $`𝔪=J(R^{})=R^{}\alpha +R^{}\beta `$. Now use the $`𝔪`$-adic filtration $`\mathrm{\Gamma }`$ defined by $`\mathrm{\Gamma }_i=R^{}`$ if $`i0`$ but $`\mathrm{\Gamma }_i=𝔪^i`$ if $`i0`$. Then $`\beta \alpha =\alpha ^2\beta 𝔪^3`$. Hence in the associated graded ring one finds that the images of these elements satisfy $`ba=0=b^2`$. The argument of the last paragraph shows that $`\mathrm{gr}_\mathrm{\Gamma }R^{}`$ is left noetherian and hence, by \[5, Proposition II.1.2.3\], that $`\mathrm{\Gamma }`$ is left Zariskian. ∎ ## 5. A Dualizing Module As was remarked in the introduction, if an affine noetherian PI algebra $`R`$ has a finite noetherian $``$-filtration, then $`R`$ has a dualizing complex. Thus, one can ask whether the ideas of the last section can be used to provide examples of PI rings which do not have such a complex. This appears not to be the case; in this section we check that the ring $`R`$ from (1.2) does indeed have such a complex. One advantage of this example is that we can work with modules rather than complexes, and so we define an $`(R,R)`$-bimodule $`D`$ to be a dualizing module if (i) $`D`$ is finitely generated and of finite injective dimension on both sides and (ii) the natural maps $`R\mathrm{End}(D_R)`$ and $`R^{\mathrm{op}}\mathrm{End}(_RD)`$ are isomorphisms. A dualizing module viewed as a complex is a dualizing complex in the sense of Yekutieli . The way we find a dualizing module for the ring $`R`$ from (1.2) is through the following observation: Identify $`C=k[x]`$ with the diagonal matrices in $`R`$ and set (5.1) $$D_1=\mathrm{Hom}_C({}_{C}{}^{}R,C)\mathrm{and}D_2=\mathrm{Hom}_C(R_C,C).$$ Thus, $`D_1`$ is an $`(R,C)`$-bimodule and $`D_2`$ is a $`(C,R)`$-bimodule. The key to our construction is the following easy lemma. ###### Lemma 5.2. Let $`RC`$ be rings such that $`R`$ is a finitely generated projective $`C`$-module on both sides and $`C`$ is a commutative noetherian algebra of finite injective dimension. Define modules $`D_i`$ by (5.1). Suppose that one has ring isomorphisms $`\mathrm{End}_R(D_1)R^{\mathrm{op}}`$ and $`\mathrm{End}_R(D_2)R`$ through which $`D_1D_2`$ as $`R`$-bimodules. Then, $`D=D_1`$ is a dualizing module for $`R`$. ###### Proof. It is clear that $`C`$ is a dualizing module for itself. By \[9, Theorem 11.66\] one has natural isomorphisms: (5.3) $$\mathrm{Ext}_R^i({}_{R}{}^{}N,{}_{R}{}^{}D_{1}^{})=\mathrm{Ext}_R^i({}_{R}{}^{}N,\mathrm{Hom}_C({}_{C}{}^{}R,C))\mathrm{Ext}_C^i({}_{C}{}^{}N,C),$$ for any finitely generated left $`R`$-module $`N`$. This implies that the injective dimension of $`{}_{R}{}^{}D_{1}^{}`$ is bounded by the injective dimension of $`{}_{C}{}^{}C`$. A similar assertion holds for $`(D_2)_R`$. It follows from \[11, Theorem 3.5\] that $`{}_{R}{}^{}D_{1}^{}`$ and $`(D_2)_R`$ are finitely generated. Therefore $`D=D_1`$ is a dualizing module. ∎ ###### Proposition 5.4. Let $`R`$ be the ring defined by (1.2) and define modules $`D_i`$ by (5.1). Then, $`D=D_1`$ is a dualizing module for $`R`$. ###### Proof. We check that the hypotheses of the lemma are satisfied for $`C=k[x]`$. Since $`R`$ is a free left $`k[x]`$-module of rank two and a free right $`k[x]`$-module of rank three, the proof is a routine computation which we will only outline. One first checks that, under the natural module structures, $$D_1\frac{RR}{R(xe_{12},e_{12})}\mathrm{and}D_2R/e_{12}R,$$ where $`e_{12}=\left(\begin{array}{cc}0& 1\\ 0& 0\end{array}\right)`$ is the matrix unit. It follows that $$\mathrm{End}_R(D_2)II(e_{12}R)/e_{12}R\mathrm{where}II(aR)=\{\theta R:\theta aRaR\}.$$ Computing this out one finds that $`\mathrm{End}_R(D_2)`$ $`\{\left(\begin{array}{cc}f(x^2)& g(x)\\ 0& f(x^4)\end{array}\right):f,gk[x]\}/e_{12}k[x^2]`$ $``$ $`\{\left(\begin{array}{cc}f(x^2)& g(x^2)\\ 0& f(x^4)\end{array}\right):f,gk[x]\}R.`$ Finally, under this isomorphism $`E=\mathrm{End}_R(D_2)R`$ one finds that the modules $`{}_{E}{}^{}D_{2}^{}`$ and $`{}_{R}{}^{}D_{1}^{}`$ become isomorphic, as is required to prove the proposition. ∎ It would be interesting to know whether this proof can be extended to work for any ring $`S`$ that is finitely generated as a module over a commutative subring. The fact that the present proof depends upon the “lucky” isomorphism $`\mathrm{End}_R(D_2)R`$ makes this seem unlikely. By using ideas from , it can at least be extended to Hopf algebras finitely generated as modules over commutative subalgebras. Acknowledgment The authors thank Quanshui Wu for several conversations on the subject.
warning/0007/hep-ph0007142.html	ar5iv	text	# 1 Introduction ## 1 Introduction Higher-order computations in perturbative QCD can be performed by using three main tools: exact calculations at a fixed order in the QCD coupling $`\alpha _\mathrm{S}=g_\mathrm{S}^2/4\pi `$, analytic resummed calculations and parton shower event generators. At present, the accuracy of these tools is respectively limited to the next-to-leading order (NLO), to next-to-leading logarithmic (NLL) accuracy and to the dominant soft and collinear enhanced contributions. An extensive and updated list of references can be found in Refs. . Apart from important conceptual and technical differences, these tools are based on a common ingredient: the universal factorization properties of QCD amplitudes in the infrared (soft and collinear) region. The lowest-order version of the soft and collinear factorization formulae was indeed exploited to develop these tools to their present theoretical accuracy. Higher-order versions of the factorization formulae are required to progress towards next-to-next-to-leading order (NNLO) calculations, resummation of next-to-next-to-leading logarithmic (NNLL) terms and inclusion of subdominant contributions in parton showers. In recent years several groups have contributed to extending infrared factorization to higher perturbative orders . The general factorization properties of tree-level and loop amplitudes in the limit where two or more partons become collinear were studied in Refs. . At tree level, the singular factors for the collinear splitting of one parton into three were explicitly computed in Refs. . The splitting function corresponding to the clustering of four collinear gluons is also known . The one-loop kernels for the collinear splitting of one parton into two were obtained in Refs. . The mixed soft-collinear limit can be studied by exploiting the coherence properties of QCD radiation . Using QCD coherence, the singular behaviour in the soft-collinear limit can be treated by combining the singular factors that separately control the collinear and soft limits. The soft limit is physically more involved than the collinear limit. Long-wavelength (soft) gluons can spread the colour flow over large distances, thus leading to (non-local) colour correlations. Kinematics and colour factors turn out to be deeply entangled in the soft-factorization formulae. The tree-level factorization formulae for the emission of two soft gluons were independently derived in Refs. and . The soft (and collinear) singular behaviour of two-loop amplitudes was studied in Ref. . In this paper we consider the limit in which a soft gluon is radiated from one-loop amplitudes. The limit was first investigated in Refs. . The formalism used by this group is based on the decomposition of the one-loop matrix elements in colour subamplitudes . The colour-subamplitude decomposition depends on the type of external partons, and these authors derived the explicit expressions of the one-loop soft-gluon contribution to colour subamplitudes with $`m`$ external gluons and with $`m`$ external gluons plus a $`q\overline{q}`$ pair . In the present paper the soft limit is studied by means of a completely independent and general method. We apply the eikonal approximation and soft-gluon insertion rules to perform infrared factorization directly in colour space. Within this formalism, the soft limit of tree-level amplitudes is described by a factorization formula written in terms of a soft-gluon (or eikonal) current (see Ref. and Sect. 2) that describes colour radiation in the classical approximation. We show that the factorization formula can be extended to loop amplitudes by introducing a generalized soft-gluon current that embodies (non-abelian) quantum corrections. As at tree level, the soft-gluon current only depends on the colour charges and momenta of the external partons in the loop amplitude. We compute the explicit expression of the soft-gluon current at one-loop order. Our results for the soft limit of one-loop amplitudes agree with those in Refs. for the particular cases considered therein. The one-loop gluon current that we obtain is derived and presented in a general and process-independent way; it can be applied to any one-loop amplitude. In particular, we can easily show that colour and kinematic factors can be completely disentangled in the computation of the soft limit of the square of one-loop matrix elements with two or three external QCD partons. This simplified factorization structure is particularly useful for the NNLO calculation of 2-jet and 3-jet cross sections in $`e^+e^{}`$ annihilation. Infrared-factorization properties at the lowest perturbative order have been known for a long time . Combining the results obtained here with those in Refs. , the general structure of the infrared singularities at the next perturbative order is also completely and explicitly known. It can be used to improve the accuracy of perturbative QCD calculations. The paper is organized as follows. In Sect. 2 we recall the known results for the soft limit of tree-level amplitudes. In Sect. 3 we present the process-independent factorization formula at higher perturbative orders and we discuss in detail its features at one-loop order. In Sect. 4 we prove the factorization formula at one-loop order and we derive the explicit expression of the one-loop soft current. In Sect. 5 we apply our results to the squared amplitudes of processes with two or three hard partons. ## 2 Soft-gluon factorization at tree level We consider a generic scattering process that involves $`m`$ external QCD partons (massless quarks and gluons) with momenta $`p_1,\mathrm{},p_m`$ and an arbitrary number and type of particles with no colour (photons, leptons, vector bosons, …). Note that, by definition, we always consider incoming and outgoing parton momenta in the physical region, i.e. any $`p_i`$ is massless, with positive-definite energy (in particular, $`p_ip_j>0`$). The corresponding matrix element is denoted by $`(p_1,\mathrm{},p_m)`$ and the dependence on the momenta and quantum numbers of non-QCD particles is always understood. The matrix element has the following loop expansion: $$(p_1,\mathrm{},p_m)=^{(0)}(p_1,\mathrm{},p_m)+^{(1)}(p_1,\mathrm{},p_m)+\mathrm{},$$ (1) where $`^{(0)}`$ denotes the tree-level contribution, $`^{(1)}`$ denotes the one-loop contribution, and the dots stand for higher-loop corrections. Note that we always consider unrenormalized matrix elements. Thus Eq. (1) is a power series expansion in the bare QCD coupling $`g_\mathrm{S}`$ and, in particular, $`^{(1)}`$ is the unrenormalized one-loop amplitude. We simultaneously regularize ultraviolet and infrared singularities by using dimensional regularization. Apart from introducing the dimensional-regularization scale $`\mu `$ through the replacement $`g_\mathrm{S}g_\mathrm{S}\mu ^ϵ`$, the key ingredient of dimensional regularization is the analytic continuation of loop momenta to $`d=42ϵ`$ space-time dimensions. Having done this, we are left with some freedom regarding the dimensionality of the momenta of the external particles as well as the number of polarizations of both external and internal particles. This leads to different regularization schemes within the dimensional-regularization prescription. The regularization-scheme dependence of one-loop amplitudes was studied in detail in Refs. . All the results on the soft behaviour presented in this paper do not explicitly depend on the dimensional-regularization scheme (see Sect. 4 for a brief discussion of different regularization schemes): the scheme dependence is implicitly embodied in the expressions of the tree-level and one-loop matrix elements $`^{(0)}`$ and $`^{(1)}`$. The emission of a soft gluon does not affect the momenta and spins of the radiating hard partons. However, it does affect their colour because the gluon always carries away some colour charge, no matter how soft it is. Unlike the case of soft-photon emission in QED, soft-gluon emission thus does not factorize exactly and leads to colour correlations. To take into account the colour structure without referring to any particular choice of basis colour vectors (such as, for instance, the decomposition in colour subamplitudes ), we use a general notation (see e.g. Ref. ). The dependence of the matrix element on the colour indices $`c_1,\mathrm{},c_m`$ of the QCD partons is written as $$_{c_1,\mathrm{},c_m}(p_1,\mathrm{},p_m)c_1,\mathrm{},c_m\|(p_1,\mathrm{},p_m).$$ (2) Thus $`\{\|c_1,\mathrm{},c_m\}`$ is an abstract basis in colour space and the ket $`\|(p_1,\mathrm{},p_m)`$ is a vector in this space. According to this notation, the matrix element squared $`\|\|^2`$ (summed over the colours and spins of the partons) can be written as $$\|(p_1,\mathrm{},p_m)\|^2=(p_1,\mathrm{},p_m)\|(p_1,\mathrm{},p_m).$$ (3) To describe the colour correlations produced by soft-gluon emission, we associate a colour charge $`𝑻_i`$ with the emission of a gluon from each parton $`i`$. If the emitted gluon has colour index $`a`$ ($`a=1,\mathrm{},`$ $`N_c^21`$), the colour-charge operator is: $$𝑻_ia\|T_i^a$$ (4) and its action onto the colour space is defined by $$a,c_1,\mathrm{},c_i,\mathrm{},c_m\|𝑻_i\|b_1,\mathrm{},b_i,\mathrm{},b_m=\delta _{c_1b_1}\mathrm{}T_{c_ib_i}^a\mathrm{}\delta _{c_mb_m},$$ (5) where $`T_{cb}^aif_{cab}`$ (colour-charge matrix in the adjoint representation) if the emitting parton $`i`$ is a gluon and $`T_{\alpha \beta }^at_{\alpha \beta }^a`$ (colour-charge matrix in the fundamental representation with $`\alpha ,\beta =1,\mathrm{},N_c`$) if the emitting particle $`i`$ is a final-state quark or an initial-state antiquark ($`T_{\alpha \beta }^a\overline{t}_{\alpha \beta }^a=t_{\beta \alpha }^a`$, in the case of a final-state antiquark or an initial-state quark). The colour-charge algebra is<sup>§</sup><sup>§</sup>§More details on the colour algebra and useful colour-matrix relations can be found in Appendix A of Ref. .: $$T_i^aT_j^a𝑻_i𝑻_j=𝑻_j𝑻_i\mathrm{if}ij;𝑻_i^2=C_i,$$ (6) where $`C_i`$ is the Casimir operator, i.e. $`C_i=C_A=N_c`$ if $`i`$ is a gluon and $`C_i=C_F=(N_c^21)/2N_c`$ if $`i`$ is a quark or antiquark. Note that, by definition, each vector $`\|(p_1,\mathrm{},p_m)`$ is a colour-singlet state. Therefore colour conservation is simply $$\underset{i=1}{\overset{m}{}}𝑻_i\|(p_1,\mathrm{},p_m)=0.$$ (7) We can now recall the behaviour of the tree-level matrix element $`^{(0)}(q,p_1,\mathrm{},p_m)`$ in the limit where the momentum $`q`$ of the gluon becomes soft. Denoting by $`a`$ and $`\epsilon ^\mu (q)`$ the colour and the polarization vector of the soft gluon, the matrix element fulfils the following factorization formula $$a\|^{(0)}(q,p_1,\mathrm{},p_m)g_\mathrm{S}\mu ^ϵ\epsilon ^\mu (q)J_\mu ^{a(0)}(q)\|^{(0)}(p_1,\mathrm{},p_m),$$ (8) where $`\|^{(0)}(p_1,\mathrm{},p_m)`$ is obtained from the original matrix element by simply removing the soft gluon $`q`$. The factor $`𝑱_\mu ^{(0)}(q)`$ is the tree-level soft-gluon current $$𝑱^{\mu (0)}(q)=\underset{i=1}{\overset{m}{}}𝑻_i\frac{p_i^\mu }{p_iq},$$ (9) which depends on the momenta and colour charges of the hard partons in the matrix element on the right-hand side of Eq. (8). The symbol ‘ $``$ ’ means that on the right-hand side we have neglected contributions that are less singular than $`1/q`$ in the soft limit $`q0`$. Note that Eq. (8) is valid in any number $`d=42ϵ`$ of space-time dimensions, and the sole dependence on $`d`$ is in the overall factor $`\mu ^ϵ`$. The factorization formula (8) can be derived in a simple way by working in a physical gauge and using the following soft-gluon insertion rules. The coupling of the gluon to any internal (i.e. highly off-shell) parton in the amplitude $`^{(0)}(q,p_1,\mathrm{},p_m)`$ is not singular in the soft limit; it can thus be neglected. The soft-gluon coupling to any external or, in general, nearly on-shell parton with colour charge $`𝑻`$ and momentum $`p`$ can be factorized by using the eikonal approximation, that is by extracting the contribution $`g_S\mu ^ϵ2p^\mu 𝑻`$ for the vertex and the contribution $`1/(p+q)^21/(p^2+2pq)`$ for the propagator. Note that the eikonal vertex only depends on the momentum and colour charge of the radiating parton: it does not depend on either the soft momentum or the spin of the parton. Using the eikonal propagator simply amounts to neglecting the terms that are quadratic in the soft momentum. An important property of the soft-gluon current is current conservation. Multiplying Eq. (9) by $`q^\mu `$, we obtain $$q^\mu 𝑱_\mu ^{(0)}(q)=\underset{i=1}{\overset{m}{}}𝑻_i,$$ (10) and thus $$q^\mu 𝑱_\mu ^{(0)}(q)\|^{(0)}(p_1,\mathrm{},p_m)=\underset{i=1}{\overset{m}{}}𝑻_i\|^{(0)}(p_1,\mathrm{},p_m)=0,$$ (11) where the last equality follows from colour conservation as in Eq. (7). Although the factorization formula (8) is most easily derived by working in a physical gauge, the conservation of the soft-gluon current implies that Eq. (8) is actually gauge invariant. Any gauge transformation is equivalent to an addition of a longitudinal component to the polarization vector of the soft gluon through the replacement $`\epsilon ^\mu (q)\epsilon ^\mu (q)+\lambda q^\mu `$. Nonetheless the factorization formula (8) is invariant under this replacement, because of Eq. (11). Squaring Eq. (8) and summing over the gluon polarizations leads to the well-known soft-gluon factorization formula at $`𝒪(g_\mathrm{S}^2)`$ for the squared amplitude : $`\|^{(0)}(q,p_1,\mathrm{},p_m)\|^2g_\mathrm{S}^2\mu ^{2ϵ}\mathrm{\hspace{0.17em}2}{\displaystyle \underset{i,j=1}{\overset{m}{}}}𝒮_{ij}(q)\|_{(i,j)}^{(0)}(p_1,\mathrm{},p_m)\|^2,`$ (12) where the eikonal function $`𝒮_{ij}(q)`$ can be written in terms of two-particle sub-energies $`s_{ij}=(p_i+p_j)^2`$ as follows $$𝒮_{ij}(q)=\frac{p_ip_j}{2(p_iq)(p_jq)}=\frac{s_{ij}}{s_{iq}s_{jq}}.$$ (13) The colour correlations produced at tree level by the emission of a soft gluon are taken into account by the square of the colour-correlated tree-amplitude $`\|_{(i,j)}^{(0)}\|^2`$ on the right-hand side. This is given by $`\|_{(i,j)}^{(0)}(p_1,\mathrm{},p_m)\|^2`$ $``$ $`^{(0)}(p_1,\mathrm{},p_m)\|𝑻_i𝑻_j\|^{(0)}(p_1,\mathrm{},p_m)`$ $`=`$ $`\left[_{c_1..b_i\mathrm{}b_j\mathrm{}c_m}^{(0)}(p_1,\mathrm{},p_m)\right]^{}T_{b_id_i}^aT_{b_jd_j}^a_{c_1..d_i\mathrm{}d_j\mathrm{}c_m}^{(0)}(p_1,\mathrm{},p_m).`$ ## 3 Soft-gluon factorization at one loop The factorized structure of the soft limit of the QCD amplitudes at tree level can be generalized to higher loops. The soft limit of the one-loop amplitudes is studied in detail in Sect. 4. In this section we anticipate and discuss the final results. Our analysis is consistent with the following factorization formula $$a\|(q,p_1,\mathrm{},p_m)\epsilon ^\mu (q)J_\mu ^a(q,ϵ)\|(p_1,\mathrm{},p_m)\left[1+𝒪(g_\mathrm{S}^4)\right],$$ (15) where the symbol ‘ $``$ ’ has the same meaning as in Eq. (8). The matrix element on the right-hand side is the all-loop amplitude in Eq. (1) and the singular dependence on $`q`$ is embodied in the (unrenormalized) soft-gluon current $`J_\mu ^a(q,ϵ)`$, which can be expanded in loop contributions, i.e. in powers of $`g_\mathrm{S}^2`$: $$J_\mu ^a(q,ϵ)=g_\mathrm{S}\mu ^ϵ\left[J_\mu ^{a(0)}(q)+g_\mathrm{S}^2\mu ^{2ϵ}J_\mu ^{a(1)}(q,ϵ)+\mathrm{}\right].$$ (16) The term $`J_\mu ^{a(0)}(q)`$ is the tree-level current in Eq. (9), the term $`J_\mu ^{a(1)}(q,ϵ)`$ is its one-loop correction, and so forth. The soft current contains the entire singular dependence in the soft limit and no further approximation is performed in Eq. (15). By this we mean that $`J_\mu ^a(q,ϵ)`$ behaves as $`1/q`$ (such as at tree level) modulo any possible enhancement proportional to powers of $`\mathrm{ln}q`$ coming from the loop contributions. In particular, the $`ϵ`$-dependence of the current can be evaluated exactly without performing any $`ϵ`$-expansion and, thus, by keeping all the powers of $`\mathrm{ln}q`$ coming from higher-loop contributions of the type $`(q)^ϵ=1+ϵ\mathrm{ln}q+\mathrm{}`$. The discussion in Sect. 4 suggests that the factorization formula (15) is valid to any loop order. Nonetheless, we have included the term $`𝒪(g_\mathrm{S}^4)`$ on the right-hand side of Eq. (15) to indicate that our explicit proof and calculation do not extend beyond the one-loop order. Expanding both sides of Eq. (15) to one-loop accuracy and using Eq. (8), we can obtain the factorization formula for the soft limit of the one-loop amplitudes: $`a\|^{(1)}(q,p_1,\mathrm{},p_m)g_S\mu ^ϵ\epsilon ^\mu (q)`$ $`[J_\mu ^{a(0)}(q)\|^{(1)}(p_1,\mathrm{},p_n)`$ (17) $`+g_\mathrm{S}^2\mu ^{2ϵ}J_\mu ^{a(1)}(q,ϵ)\|^{(0)}(p_1,\mathrm{},p_n)].`$ The explicit expression for the one-loop current is $`J_a^{\mu (1)}(q,ϵ)`$ $`={\displaystyle \frac{1}{16\pi ^2}}{\displaystyle \frac{1}{ϵ^2}}{\displaystyle \frac{\mathrm{\Gamma }^3(1ϵ)\mathrm{\Gamma }^2(1+ϵ)}{\mathrm{\Gamma }(12ϵ)}}`$ $`if_{abc}{\displaystyle \underset{ij}{}}T_i^bT_j^c\left({\displaystyle \frac{p_i^\mu }{p_iq}}{\displaystyle \frac{p_j^\mu }{p_jq}}\right)\left[{\displaystyle \frac{4\pi p_ip_je^{i\lambda _{ij}\pi }}{2(p_iq)(p_jq)e^{i\lambda _{iq}\pi }e^{i\lambda _{jq}\pi }}}\right]^ϵ.`$ (18) We remind the reader that in our notation all the incoming and outgoing momenta are in the physical region (any $`p_i`$ has positive-definite energy and $`p_ip_j>0`$). Thus the complex factors $`e^{i\pi \lambda _{AB}}`$ ($`\lambda _{AB}=+1`$ if $`A`$ and $`B`$ are both incoming or outgoing, and $`\lambda _{AB}=0`$ otherwise) in Eq. (18) are the unitarity phases related to the analytic continuation from unphysical to physical momenta. The result in Eq. (18) explicitly shows that the soft limit of the one-loop amplitudes is process-independent, meaning that it does not depend on the momentum and colour flows of the internal partons (including the parton circulating in the loop) in the matrix element. This simple structure has several interesting features that we comment below. The one-loop soft current is proportional to the structure constants $`f_{abc}`$ of the gauge group, and thus it is purely non-abelian. This is in agreement with the absence of higher-loop corrections to the soft current in massless QED . Since Eq. (18) is proportional to the factor $`(p_i^\mu /p_iqp_j^\mu /p_jq)`$, the one-loop contribution to the soft current is conserved: $$q_\mu J_a^{\mu (1)}(q,ϵ)=0.$$ (19) Combined with the analogous property at tree level, this guarantees that the soft-gluon factorization formula (15) is manifestly gauge-invariant. The double pole $`1/ϵ^2`$ in Eq. (18) is the infrared singularity produced by a soft and collinear virtual gluon. To double-pole accuracy, we can use colour conservation (see Eq. (7)) to show that the one-loop current is simply proportional to the tree-level current: $`J_\mu ^{a(1)}(q,ϵ)\|(q,p_1,\mathrm{},p_m)={\displaystyle \frac{1}{16\pi ^2}}\left[{\displaystyle \frac{C_A}{ϵ^2}}J_\mu ^{a(0)}(q)+𝒪\left({\displaystyle \frac{1}{ϵ}}\right)\right]\|(q,p_1,\mathrm{},p_m).`$ (20) This behaviour is consistent with the known singularity structure of the one-loop amplitudes . Beyond the double-pole approximation, the one-loop current contains two-particle colour correlations. The correlations are induced by the last factor on the right-hand side of Eq. (18). This factor fully embodies the logarithmic dependence on the soft-gluon momentum $`q`$ and has a simple kinematic interpretation, being related to the transverse component $`q_{,ij}`$ of the gluon momentum with respect to the longitudinal direction singled out by the momenta $`p_i`$ and $`p_j`$ of the colour-correlated hard partons: $$q_{,ij}^2=\frac{2(p_iq)(p_jq)}{p_ip_j}.$$ (21) The derivation of the result in Eq. (18) (see Sect. 4) suggests that multiparticle colour correlations will appear in higher-loop contributions to the soft-gluon current (16). For instance, at two-loop order the soft current contains colour correlations of the type $$f_{abe}f_{ecd}T_i^bT_j^cT_k^d$$ (22) between three different hard partons $`i,j`$ and $`k`$. The soft behaviour of the one-loop amplitudes was first investigated in Refs. by using the colour-subamplitude formalism. The explicit expressions of the one-loop soft-gluon contribution to colour subamplitudes, with $`m`$ external gluons and with $`m`$ external gluons plus a $`q\overline{q}`$ pair, were derived in Refs. and , respectively. The reader can straightforwardly check that the result in Eq. (18) agrees with those in Refs. for the particular cases considered therein. The general result in Eq. (18) shows that, although colour and kinematics are deeply entangled in the soft region, the soft limit of the one-loop amplitudes can be factorized in colour-space in a way that is both (relatively) simple and process-independent (in particular, independent of the flavour of the external partons). In particular, this general structure is quite useful to show (see below and Sect. 5) that colour and kinematics can be completely disentangled in the computation of the soft limit of the square of one-loop matrix elements with two or three external QCD partons. The factorization formula (17) can be used to compute the soft limit of the one-loop contribution to the square of the matrix elementIn the following the dependence of the matrix element on the momenta $`p_1,\mathrm{},p_m`$ of the hard partons is denoted by $`\{p\}`$. $`(q,\{p\})`$. Summing over the polarizations of the soft gluon and using Eqs. (8) and (17), we have $``$ $`^{(0)}(q,\{p\})\|^{(1)}(q,\{p\})+\mathrm{c}.\mathrm{c}.(g_S\mu ^ϵ)^2`$ $``$ $`\{[^{(0)}(\{p\})\|𝑱_\mu ^{(0)}(q)𝑱^{\mu (0)}(q)\|^{(1)}(\{p\})+\mathrm{c}.\mathrm{c}.]`$ $`+`$ $`\left(g_S\mu ^ϵ\right)^2[^{(0)}(\{p\})\|𝑱_\mu ^{(0)}(q)𝑱^{\mu (1)}(q,ϵ)\|^{(0)}(\{p\})+\mathrm{c}.\mathrm{c}.]\},`$ where c.c. denotes the complex conjugate. The first term on the right-hand side, which can be evaluated by using the expression (9) of the tree-level gluon current, has the same structure as Eq. (12): $$^{(0)}(\{p\})\|𝑱_\mu ^{(0)}(q)𝑱^{\mu (0)}(q)\|^{(1)}(\{p\})+\mathrm{c}.\mathrm{c}.=2\underset{i,j=1}{\overset{m}{}}𝒮_{ij}(q)\|_{(i,j)}^{(1)}(\{p\})\|^2,$$ (24) where $`𝒮_{ij}(q)`$ is the eikonal function in Eq. (13). In particular, the two-particle colour correlations on the right-hand side are completely analogous to those at tree level (see Eqs. (12) and (2)) and have been taken into account by defining the colour-correlated one-loop amplitude $$\|_{(i,j)}^{(1)}(\{p\})\|^2^{(0)}(\{p\})\|𝑻_i𝑻_j\|^{(1)}(\{p\})+\mathrm{c}.\mathrm{c}..$$ (25) Using expression (18) for the one-loop contribution to the soft-gluon current, the second term on the right-hand side of Eq. (3) can be written as $`^{(0)}(\{p\})\|𝑱_\mu ^{(0)}(q)`$ $`𝑱^{\mu (1)}(q,ϵ)\|^{(0)}(\{p\})+\mathrm{c}.\mathrm{c}.={\displaystyle \frac{1}{4\pi ^2}}{\displaystyle \frac{(4\pi )^ϵ}{ϵ^2}}{\displaystyle \frac{\mathrm{\Gamma }^3(1ϵ)\mathrm{\Gamma }^2(1+ϵ)}{\mathrm{\Gamma }(12ϵ)}}`$ $`\{C_A\mathrm{cos}(\pi ϵ){\displaystyle \underset{i,j}{}^{}}\left[𝒮_{ij}(q)\right]^{1+ϵ}\|_{(i,j)}^{(0)}(\{p\})\|^2`$ (26) $`+2\mathrm{sin}(\pi ϵ){\displaystyle \underset{i,j,k}{}^{}}𝒮_{ki}(q)\left[𝒮_{ij}(q)\right]^ϵ(\lambda _{ij}\lambda _{iq}\lambda _{jq})\|_{(k,i,j)}^{(0)}(\{p\})\|^2\},`$ where we have used $`\left(\lambda _{ij}\lambda _{iq}\lambda _{jq}\right)^2=1`$ for any possible configuration of incoming and outgoing momenta, and the notation $`^{}`$ stands for the sum over the different values of the indices $`(ij,jk,ki)`$. The one-loop contribution on the right-hand side of Eq. (26) contains two terms. The first term only involves colour correlations between two hard partons, which are taken into account by the colour-correlated tree-amplitude $`_{(i,j)}^{(0)}(\{p\})`$ defined in Eq. (2). Thus, apart from an overall $`ϵ`$-dependent factor, its effect simply amounts to the rescaling $`𝒮_{ij}(q)\left[𝒮_{ij}(q)\right]^{1+ϵ}`$ in the tree-level factorization formula (12). The second term instead has a different structure, because it leads to colour correlations between three different hard partons. These are included in the three-parton correlated tree-amplitude: $$\|_{(k,i,j)}^{(0)}(\{p\})\|^2f_{abc}^{(0)}(\{p\})\|T_k^aT_i^bT_j^c\|^{(0)}(\{p\}).$$ (27) Note that the second term contributes only when there are four or more hard partons, because in the case of three partons colour conservation (see Eq. (7)) can be used to show that the three-parton correlations vanish: $$f_{abc}T_1^aT_2^bT_3^c\|^{(0)}(p_1,p_2,p_3)=f_{abc}T_1^aT_2^b(T_1^c+T_2^c)\|^{(0)}(p_1,p_2,p_3)=0.$$ (28) The absence of these correlations extremely simplifies the structure of the soft limit of the squared matrix elements with two or three hard partons. As a matter of fact, in these cases any product $`𝑻_i𝑻_j`$ can be expressed as a linear combination of Casimir operators (see Sect. 5). Therefore, the colour algebra can explicitly be carried out and the soft limit of $`\|(q,\{p\})\|^2`$ is directly proportional to $`\|(\{p\})\|^2`$ at one-loop accuracy. The corresponding explicit expressions are presented in Sect. 5. Note also that in the limit $`ϵ0`$ the three-particle correlation contribution to Eq. (26) is at most as singular as $`1/ϵ`$. More precisely, it gives rise to single poles $`1/ϵ`$ only when there are two or more incoming partons. In fact, in the case of an outgoing soft gluon $`q`$, we have $`\lambda _{ij}\lambda _{iq}\lambda _{jq}=+1`$ when $`i`$ and $`j`$ are both incoming momenta and $`\lambda _{ij}\lambda _{iq}\lambda _{jq}=1`$ otherwise. Therefore, we can rewrite the second term in the curly bracket of Eq. (26) as $`+4\mathrm{sin}(\pi ϵ){\displaystyle \underset{i,j(\mathrm{in}),k}{}^{}}𝒮_{ki}(q)\left[𝒮_{ij}(q)\right]^ϵf_{abc}^{(0)}(\{p\})\|T_k^aT_i^bT_j^c\|^{(0)}(\{p\})`$ $`2\mathrm{sin}(\pi ϵ){\displaystyle \underset{i,j,k}{}^{}}𝒮_{ki}(q)\left[𝒮_{ij}(q)\right]^ϵf_{abc}^{(0)}(\{p\})\|T_k^aT_i^bT_j^c\|^{(0)}(\{p\}),`$ (29) where the sum $`_{i,j(\mathrm{in}),k}^{}`$ denotes the restriction of $`_{i,j,k}^{}`$ to the indices $`i,j`$ of the incoming partons. Then, it is easy to show that the second term in Eq. (29) is of $`𝒪(ϵ^2)`$ in the limit $`ϵ0`$: $`2\mathrm{sin}(\pi ϵ){\displaystyle \underset{i,j,k}{}^{}}𝒮_{ki}(q)\left[𝒮_{ij}(q)\right]^ϵf_{abc}^{(0)}(\{p\})\|T_k^aT_i^bT_j^c\|^{(0)}(\{p\})=`$ (30) $`2\pi ϵ{\displaystyle \underset{i,j,k}{}^{}}𝒮_{ki}(q)f_{abc}^{(0)}(\{p\})\|T_k^aT_i^bT_j^c\|^{(0)}(\{p\})+𝒪(ϵ^2)=`$ $`+2\pi ϵ{\displaystyle \underset{i,k}{}^{}}𝒮_{ki}(q)f_{abc}^{(0)}(\{p\})\|T_k^aT_i^b(T_i^c+T_k^c)\|^{(0)}(\{p\})+𝒪(ϵ^2)=𝒪(ϵ^2).`$ (31) Here we have set $`\left[𝒮_{ij}(q)\right]^ϵ1`$. Then we have used colour conservation, $`_{ji,k}T_j^c=(T_i^c+T_k^c)`$, and the identity $`f_{abc}T_k^aT_i^b(T_i^c+T_k^c)=0`$ for $`ik`$. ## 4 Proof of factorization and <br>calculation of the one-loop current In this section we derive the factorization formula (15) at one-loop order and we explicitly compute the one-loop contribution to the soft-gluon current. To simplify the analysis, it is convenient to work in a gauge with only physical gluon polarizations. We use the axial gauge $`nA=0`$ with a light-like ($`n^2=0`$) gauge vector $`n^\mu `$. The polarization vectors $`\epsilon ^\mu (k)`$ of a gluon with momentum $`k`$ thus fulfil the relations $`n\epsilon (k)=0`$ and $`k\epsilon (k)k^2`$. The sum over the gluon polarizations leads to the polarization tensor $`d^{\mu \nu }`$: $$d^{\mu \nu }(k)=\underset{\mathrm{pol}.}{}\epsilon ^\mu (k)\epsilon ^\nu (k)=g^{\mu \nu }+\frac{k^\mu n^\nu +n^\mu k^\nu }{nk}.$$ (32) The expression on the right-hand side corresponds to a dimensional-regularization scheme with $`d2=22ϵ`$ gluon polarizations. For the sake of definiteness we use this explicit expression in all the intermediate steps of the calculation. However, we shall show that the final results are regularization-scheme independent by pointing out in which steps the scheme dependence might arise. ### 4.1 Proof of factorization We study the soft behaviour of one-loop matrix elements by using the eikonal approximation and the soft-gluon insertion rulesWe use the eikonal approximation both for real and for virtual soft gluons. In the case of virtual gluons, the eikonal approximation is not always justified for each single Feynman diagram. Nonetheless, it is valid for any gauge-invariant set of Feynman diagrams, such as those computed in this section. This can be shown by simply using time-ordered perturbation theory, where the eikonal approximation is valid on a graph-by-graph basis. recalled above Eq. (10). For this purpose, it is useful to decompose the one-loop matrix element $`^{(1)}`$ in three contributions, $$^{(1)}=_{hard}^{(1)}+_{coll}^{(1)}+_{soft}^{(1)},$$ (33) which respectively represent the kinematic regions where the momentum $`k`$ circulating in the loop is * hard: its momentum components are of the same order as those of the hard external momenta $`p`$, * collinear: $`k`$ is parallel to one of the hard external momenta, * soft: its momentum components are much smaller than those of the hard external momenta $`p`$. For the purpose of the following discussion, we recall that the soft virtual behaviour of the one-loop matrix element $`^{(1)}(\{p\})`$ can be computed by using the soft-gluon insertion rules. The soft virtual gluon of momentum $`k`$ is inserted (emitted and reabsorbed) on all the external legs (Fig. 1) of the tree-level amplitude $`^{(0)}(\{p\})`$, leading to the expression : $$\|_{soft}^{(1)}(\{p\})=\frac{1}{2}g_\mathrm{S}^2\mu ^{2ϵ}\frac{d^dk}{(2\pi )^d}\frac{i}{k^2+i0}\left[J_\mu ^{a(0)}(k)\right]^{}J^{\mu a(0)}(k)\|^{(0)}(\{p\}),$$ (34) where $`J_\mu ^{a(0)}(k)`$ is the (tree-level) soft-gluon current in Eq. (9). We now consider the one-loop matrix element $`^{(1)}(q,\{p\})`$ when the momentum $`q`$ of the external gluon becomes soft. To apply the soft-gluon insertion rules, we perform the decomposition in Eq. (33) and discuss the three different kinematic regions of the virtual momentum $`k`$ in turn. When $`k`$ is in the hard region, all the internal lines in $`^{(1)}(q,\{p\})`$ are highly off-shell. Thus the real soft gluon $`q`$ can couple only to the external legs and we can factorize its contribution as in the case of tree-level amplitudes. Neglecting terms that are not singular in the soft limit $`q0`$, we get $$a\|_{hard}^{(1)}(q,\{p\})g_S\mu ^ϵ\epsilon ^\mu (q)J_\mu ^{a(0)}(q)\|_{hard}^{(1)}(\{p\}),$$ (35) where $`J_\mu ^{a(0)}(q)`$ is the soft-gluon current in Eq. (9). We now consider the region where the loop momentum $`k`$ is collinear to the momentum of one, say $`p_i`$, of the hard external legs. Since we work in a physical gauge, the only diagram that is not dynamically suppressed is that in which the loop leads to a self-energy contribution on the external leg $`p_i`$ (e.g. the diagram in Fig. 1 $`(b)`$). The soft gluon can then be inserted in this diagram in all possible ways: on the line $`p_i`$ before and after the self-energy contribution and on the self-energy lines themselves. However, we can exploit the colour coherence properties of QCD radiation. Since the loop-momentum $`k`$ is parallel to $`p_i`$, the soft gluon $`q`$ cannot distinguish the self-energy lines from a single line with total colour charge $`T_i^a`$. The sum of the insertions of $`q`$ on this diagram is thus insensitive to the presence of the collinear loop momentum and leads (see, for instance, Sect. 3.4 in Ref. for a detailed similar discussion) to the same factor, $`T_i^ap_i\epsilon (q)/p_iq`$, as in the soft-gluon factorization at tree level. Considering the insertions of $`q`$ on all the other external legs $`ji`$, we obtain a factorization formula of tree-level type also in the collinear region: $$a\|_{coll}^{(1)}(q,\{p\})g_S\mu ^ϵ\epsilon ^\mu (q)J_\mu ^{a(0)}(q)\|_{coll}^{(1)}(\{p\}).$$ (36) We finally have to deal with the region in which the loop momentum $`k`$ is carried by a soft gluon<sup>\</sup><sup>\</sup>\**Quarks loops are dynamically suppressed when their momentum become soft. Without loss of infrared accuracy, we thus consider the quark loops as included in the hard or collinear regions.. Unlike the case of the hard and collinear regions, where the one-loop effects can be factorized with respect to the tree-level current $`J_\mu ^{a(0)}(q)`$ (see Eqs. (35) and (36)), new ‘non-factorizable’ contributions appear when the loop momentum is soft. To single out these new contributions, we write the following identity: $`\|_{soft}^{(1)}(q,\{p\})`$ $`=g_S\mu ^ϵ\epsilon ^\mu (q)𝑱_\mu ^{(0)}(q)\|_{soft}^{(1)}(\{p\})`$ $`+\left(\|_{soft}^{(1)}(q,\{p\})g_S\mu ^ϵ\epsilon ^\mu (q)𝑱_\mu ^{(0)}(q)\|_{soft}^{(1)}(\{p\})\right),`$ (37) where we have added and subtracted the ‘factorized’ contribution. Then we combine the contributions from the hard, collinear and soft regions by adding Eqs. (35), (36) and (4.1), and we obtain $`\|^{(1)}(q,\{p\})`$ $`=g_S\mu ^ϵ\epsilon ^\mu (q)𝑱_\mu ^{(0)}(q)\|^{(1)}(\{p\})`$ $`+\left(\|_{soft}^{(1)}(q,\{p\})g_S\mu ^ϵ\epsilon ^\mu (q)𝑱_\mu ^{(0)}(q)\|_{soft}^{(1)}(\{p\})\right).`$ (38) The first term on the right-hand side of Eq. (4.1) together with the contributions from Eqs. (35) and (36) have reconstructed the first term on the right-hand side of Eq. (4.1), which is exactly the first term on the right-hand side of the factorization formula (17). What remains to be done to prove the factorization formula is to relate the second term on the right-hand side of Eq. (17) with the contribution in the round bracket of Eq. (4.1). For this purpose, we first note that when the real gluon $`q`$ and the virtual gluon $`k`$ are both soft, they can couple only to the external hard lines. In the corresponding Feynman diagrams, which are schematically represented by the first graph in Fig. 2, the tree-level amplitude $`^{(0)}(\{p\})`$ is factorized in the soft limit. We can write: $$\|_{soft}^{(1)}(q,\{p\})\left(g_S\mu ^ϵ\right)^3\epsilon ^\mu (q)𝑲_\mu ^{(1)}(q,ϵ)\|^{(0)}(\{p\}),$$ (39) where the kernel $`𝑲^{(1)}`$ (represented by the box in Fig. 2) denotes all the soft-gluon insertions of $`q`$ and $`k`$ on the hard-momentum lines. Then, we note that $`^{(0)}(\{p\})`$ is factorized also in the expression (34) for $`_{soft}^{(1)}(\{p\})`$. Therefore, the contribution in the round bracket of Eq. (4.1) can be recast in the form of the second term on the right-hand side of the factorization formula (17). Moreover, using Eqs. (39) and (34), we obtain the following explicit representation of the one-loop contribution $`𝑱^{(1)}`$ to the soft-gluon current (Fig. 2): $$\epsilon ^\mu (q)𝑱_\mu ^{(1)}(q,ϵ)=\epsilon ^\mu (q)\left\{𝑲_\mu ^{(1)}(q,ϵ)𝑱_\mu ^{(0)}(q)\frac{1}{2}\frac{d^dk}{(2\pi )^d}\frac{i}{k^2+i0}\left[𝑱_\nu ^{(0)}(k)\right]^{}𝑱^{\nu (0)}(k)\right\}.$$ (40) ### 4.2 Calculation of the one-loop current We now proceed to the explicit calculation of the one-loop soft current $`𝑱^{(1)}`$. Using the eikonal approximation, we have to evaluate the Feynman diagrams of the kernel $`𝑲^{(1)}`$ and to subtract those corresponding to the second term on the right-hand side of Eq. (40). We divide the diagrams in two classes: $`(A)`$ the diagrams that involve interactions with a single hard line (Fig. 3), and $`(B)`$ all the remaining diagrams (Fig. 4). We first consider class $`(B)`$. The diagrams $`(a),(b),(c)`$ and $`(d)`$ in Fig. 4 come from the kernel $`𝑲^{(1)}`$, while Fig. 4 $`(e)`$ represents the corresponding subtractions. We see that within this class there are diagrams that involve interactions between three different external lines $`i,j`$ and $`l`$. These diagrams cancel. Indeed, the diagram in Fig. 4 $`(d)`$ is exactly cancelled by that $`(e)`$-contribution in which $`q`$ is emitted by the line $`l`$ (the sum over the emissions of $`q`$ is included in the factor $`𝑱^{(0)}(q)`$). Thus the only non-vanishing terms coming from class $`(B)`$ are those that involve interactions between two external lines, namely the diagrams $`(a),(b),(c)`$ and the corresponding subtractions in $`(e)`$. The diagrams $`(a),(c)`$ and the subtractions in $`(e)`$ are very similar: they can be combined in a simple way, because they only differ by the momentum and colour flow along the hard line $`i`$. The different factors coming from the line $`i`$ in the diagrams $`(c)`$ and $`(e)`$ give $$+T_i^aT_i^b\frac{1}{p_iq+i0}\frac{1}{p_i(k+q)+i0}T_i^aT_i^b\frac{1}{p_iq+i0}\frac{1}{p_ik+i0},$$ (41) while the corresponding factor from the diagram $`(a)`$ is $$+T_i^bT_i^a\frac{1}{p_ik+i0}\frac{1}{p_i(k+q)+i0}.$$ (42) Decomposing the colour factor in Eq. (42) in its non-abelian and abelian components, $`T_i^bT_i^a=if_{bac}T_i^c+T_i^aT_i^b`$, and adding Eq. (41), we see that the abelian component $`T_i^aT_i^b`$ of the diagram $`(a)`$ exactly cancels the diagrams $`(c)`$ and $`(e)`$. In conclusion, the total contribution to the one-loop current from class $`(B)`$ simply amounts to computing the diagram $`(b)`$, which is non-abelian, and the non-abelian part of the diagram $`(a)`$, which is obtained by the replacement $`T_j^bT_i^bT_i^aif_{bac}T_j^bT_i^c`$ ($`b`$ is the colour index of the virtual gluon) in its overall colour factor. The calculation of the diagrams in Fig. 4 $`(a)`$ and $`(b)`$ is quite simple (see below). Then we have to add the diagrams of class $`(A)`$ (Fig. 3). An argument similar to that in Eqs. (41) and (42) can be used to show that the abelian parts of these diagrams cancel<sup>††</sup><sup>††</sup>††More precisely, the abelian diagrams in Figs. 3 $`(a)`$ and $`(b)`$ are cancelled by the abelian part of the diagram in Fig. 3 $`(c)`$.. Thus, only their non-abelian part has to be evaluated. Although straightforward, this calculation is quite cumbersome, in particular because it has to be carried out in the axial gauge $`nA=0`$. This cumbersome calculation can be short-circuited by exploiting gauge invariance. We first recall that the non-vanishing contributions from class $`(B)`$ depend on the momenta and charges of (at most) two hard partons. Since the diagrams of class $`(A)`$ involve interactions with a single hard line, we can split the one-loop soft current in two terms, $`J_{1P}`$ and $`J_{2P}`$, that respectively denote the contributions that depend on one and two hard momenta: $$J_\mu ^{a(1)}(q,ϵ)=J_{\mu ;1P}^{a(1)}(q,ϵ)+J_{\mu ;2P}^{a(1)}(q,ϵ).$$ (43) All the diagrams of class $`(A)`$ are included in $`J_{1P}`$. Those of class $`(B)`$ contribute to $`J_{2P}`$ and (because of possible cancellations in the dependence of one hard momentum) to $`J_{1P}`$. We shall explicitly compute $`J_{2P}`$ and show that it is gauge-invariant. Thus $`J_{1P}`$ has to be gauge invariant as well, that is, it cannot depend on the gauge vector $`n^\nu `$. Power counting and the fact that the one-loop current is non-abelian are then sufficient to determine $`J_{1P}`$, apart from an overall factor $`f^{(1)}(ϵ)`$ that can only depends on $`ϵ`$: $$J_{\mu ;1P}^{a(1)}(q,ϵ)=C_Af^{(1)}(ϵ)\underset{i=1}{\overset{m}{}}(p_iq)^ϵ\frac{p_{i\mu }}{p_iq}T_i^a.$$ (44) This overall factor can then be obtained by imposing that the soft-gluon factorization formula be gauge-invariant (with respect to the polarizations of the real gluon $`q`$) or, equivalently, that the complete one-loop current in Eq. (43) be conserved. This leads to the constraint $$q^\mu J_\mu ^{a(1)}(q,ϵ)=q^\mu J_{\mu ;2P}^{a(1)}(q,ϵ)+C_Af^{(1)}(ϵ)\underset{i=1}{\overset{m}{}}(p_iq)^ϵT_i^a=0,$$ (45) which can be used to get $`f^{(1)}(ϵ)`$ (and, hence, $`J_{\mu ;1P}^{a(1)}`$) from the two-parton contribution $`J_{\mu ;2P}^{a(1)}`$. In particular, the calculation of the diagrams of class $`(A)`$ can be avoided. We anticipate that our explicit calculation of the two-parton contribution $`J_{\mu ;2P}^{a(1)}`$ gives $`q^\mu J_{\mu ;2P}^{a(1)}(q,ϵ)=0`$. Thus the single-parton contribution $`J_{\mu ;1P}^{a(1)}`$ vanishes. Note that this does not mean that the diagrams of class $`(A)`$ give a vanishing contribution, but rather that their contribution is cancelled by the terms of class $`(B)`$ that depend on a single hard momentum. We can now complete our calculation of the one-loop current by explicitly computing the non-abelian part of the diagrams in Figs. 4 $`(a)`$ and $`(b)`$, and, more precisely, their contribution to $`J_{\mu ;2P}^{a(1)}`$. For the sake of definiteness, we write down the diagrams for the case in which all external momenta are outgoing. The diagram in Fig. 4 $`(a)`$ depends on the polarization tensor $`d(k)`$ of the virtual gluon. Using the expression in Eq. (32) and performing the Lorentz algebra of the spin numerator, the non-abelian part of the diagram in Fig. 4 $`(a)`$ gives $`f_{abc}{\displaystyle \underset{ij}{}}T_i^cT_j^b{\displaystyle }`$ $`{\displaystyle \frac{d^dk}{(2\pi )^d}}{\displaystyle \frac{1}{(k^2+i0)(p_jki0)}}{\displaystyle \frac{p_i\epsilon (q)}{p_i(k+q)+i0}}`$ $`{\displaystyle \frac{1}{p_ik+i0}}\left(p_ip_jp_jk{\displaystyle \frac{p_in}{kn}}p_ik{\displaystyle \frac{p_jn}{kn}}\right).`$ (46) The diagram<sup>‡‡</sup><sup>‡‡</sup>‡‡Note that we use the eikonal approximation for the vertices and propagators of the hard-momentum lines $`i`$ and $`j`$, while gluon propagators and the three-gluon vertex must be treated exactly. in Fig. 4 $`(b)`$ depends on the polarization tensors $`d(k)`$ and $`d(k+q)`$ of the virtual-gluon lines. Performing the Lorentz algebra of the spin numerator and exploiting the symmetry of the integrand under the transformation $`ij,k(k+q)`$, the diagram in Fig. 4 $`(b)`$ gives $`f_{abc}{\displaystyle \underset{ij}{}}T_i^cT_j^b{\displaystyle \frac{d^dk}{(2\pi )^d}\frac{1}{(k^2+i0)(p_jki0)}\frac{1}{p_i(k+q)+i0}\frac{1}{(k+q)^2+i0}\epsilon _\mu (q)}`$ $`\left[k^\mu p_ip_jp_i^\mu p_j(k+2q)p_jk\left(2k^\mu {\displaystyle \frac{p_in}{kn}}p_i^\mu {\displaystyle \frac{(k+2q)n}{kn}}\right)+p_i^\mu (k+q)^2{\displaystyle \frac{p_jn}{kn}}\right].`$ (47) By inspection of Eqs. (4.2) and (4.2), we see that the last term in the round bracket of Eq. (4.2) cancels the last term in the square bracket of Eq. (4.2). The remaining $`n`$-dependence of the integrands is proportional to $`p_jk`$. This factor cancels the fermion propagator $`1/(p_jki0)`$, thus leaving a contribution that, although $`n`$-dependent, depends on a single hard momentum $`p_i`$. This term can then be included, together with the class $`(A)`$ diagrams, in the single-particle contribution $`J_{\mu ;1P}^{(1)}(q,ϵ)`$ to the one-loop current. We conclude that (the contribution of Eqs. (4.2) and (4.2) to) the two-particle current $`J_{\mu ;2P}^{(1)}(q,ϵ)`$ is explicitly independent of the gauge vector $`n^\mu `$. Note also that $`n`$-dependence cancels at the integrand level, and thus, the cancellation is completely insensitive to the actual prescription used to regularize the gauge pole $`1/(nk)`$ of the gluon polarization tensor (32). Although the expression (32), which we have used for the polarization tensor, corresponds to $`d2=22ϵ`$ helicity states for the gluon, our calculation also allows us to discuss the (in)dependence on the dimensional-regularization scheme. Other regularization schemes use 2 helicity states for the gluon. At one-loop order, the difference eventually amounts to set $`ϵ=0`$ in the integrands after having performed the Lorentz algebra of the spin numerators. Since the integrands of Eqs. (4.2) and (4.2) have no explicit $`ϵ`$-dependence, our result for $`\epsilon (q)J_{2P}^{(1)}`$ (and hence for the complete current $`\epsilon (q)J^{(1)}`$) does not depend on the scheme used to implement dimensional regularization. We proceed to the computation of $`\epsilon (q)J_{2P}^{(1)}`$ by adding the two-particle contributions from Eqs. (4.2) and (4.2). As for Eq. (4.2), the term proportional to $`k^\mu `$ can be reduced to scalar integrals. Then using $`q\epsilon (q)=0`$ and exploiting the symmetry with respect to the transformation $`ij,k(k+q)`$, we obtain $`\epsilon (q)J_{2P}^{a(1)}(q,ϵ)`$ $`=f_{abc}{\displaystyle \underset{ij}{}}T_i^cT_j^b{\displaystyle \frac{d^dk}{(2\pi )^d}\frac{p_i\epsilon (q)}{p_i(k+q)+i0}\frac{1}{(k^2+i0)(p_jki0)}}`$ $`\left[{\displaystyle \frac{p_ip_j}{p_ik+i0}}+{\displaystyle \frac{1}{(k+q)^2+i0}}\left({\displaystyle \frac{p_ip_j}{p_iq}}kq2p_jq{\displaystyle \frac{p_jq}{p_iq}}p_ik\right)\right].`$ (48) The ultraviolet-finite integral can be easily performed in the collinear frame in which $`p_i`$ and $`p_j`$ are directed along the ‘$`+`$’ and ‘$``$’ light-cone directions, respectively. In this frame, the integration over the $`k_+`$-complex plane receives contributions from the pole $`1/(p_jki0)`$ in the fermion propagators and from the poles in the gluon propagators. Closing the integration contour on the lower half-plane, and using the residue theorem, we can select only the gluon poles, which amounts to the following replacements in Eq. (4.2): $$\frac{1}{k^2+i0}2\pi i\delta _+(k^2),\frac{1}{(k+q)^2+i0}2\pi i\delta _+((k+q)^2).$$ (49) Therefore Eq. (4.2) can be expressed as follows $`\epsilon (q)J_{2P}^{a(1)}(q,ϵ)`$ $`={\displaystyle \frac{i}{2}}f_{abc}{\displaystyle \underset{ij}{}}T_i^cT_j^b{\displaystyle \frac{p_i\epsilon (q)}{p_iq}}[I(p_i,p_j;p_i;2p_iq)+I(p_i,p_j;p_j;2p_jqi0)`$ $`I(q,p_j;p_j;2p_jqi0)I(p_j,q;p_i;2p_iq)I(p_i,q;p_j;2p_jqi0)],`$ (50) in terms of the master integral $$I(p,\overline{p};r;s)=\frac{d^dk}{(2\pi )^{d1}}\delta _+(k^2)\frac{p\overline{p}}{(pk)(\overline{p}k)}\frac{s}{2rk+s},$$ (51) which depends on the light-like momenta $`p^\mu ,\overline{p}^\mu `$ and $`r^\mu `$ and on the complex scalar $`s`$. The evaluation of the $`d`$-dimensional integration in Eq. (51) over the on-shell gluon momentum $`k`$ is straightforward and gives $$I(p,\overline{p};r;s)=\frac{1}{8\pi ^2}\left[\frac{8\pi (pr)(\overline{p}r)}{s^2(p\overline{p})}\right]^ϵ\frac{1}{ϵ^2}\mathrm{\Gamma }^2(1ϵ)\mathrm{\Gamma }(1+ϵ)\mathrm{\Gamma }(1+2ϵ).$$ (52) Since $`I(p,\overline{p};r;s)`$ vanishes when $`r=p`$ or $`r=\overline{p}`$, inserting Eq. (52) in Eq. (50), we finally obtain $`\epsilon (q)J_{2P}^{a(1)}(q,ϵ)`$ $`={\displaystyle \frac{1}{16\pi ^2}}if_{abc}{\displaystyle \underset{ij}{}}T_i^cT_j^b{\displaystyle \frac{p_i\epsilon (q)}{p_iq}}\left({\displaystyle \frac{4\pi p_ip_j}{2p_iqp_jq}}\right)^ϵ`$ $`{\displaystyle \frac{1}{ϵ^2}}\mathrm{\Gamma }^2(1ϵ)\mathrm{\Gamma }(1+ϵ)\mathrm{\Gamma }(1+2ϵ)\left[1+\left({\displaystyle \frac{1}{e^{2\pi i}}}\right)^ϵ\right]`$ $`={\displaystyle \frac{1}{16\pi ^2}}if_{abc}{\displaystyle \underset{ij}{}}T_i^bT_j^c\epsilon _\mu (q)\left({\displaystyle \frac{p_i^\mu }{p_iq}}{\displaystyle \frac{p_j^\mu }{p_jq}}\right)\left({\displaystyle \frac{4\pi p_ip_j}{2p_iqp_jqe^{i\pi }}}\right)^ϵ`$ $`{\displaystyle \frac{1}{ϵ^2}}{\displaystyle \frac{\mathrm{\Gamma }^3(1ϵ)\mathrm{\Gamma }^2(1+ϵ)}{\mathrm{\Gamma }(12ϵ)}},`$ (53) where we have used the identity $$\mathrm{cos}(\pi ϵ)=\frac{\mathrm{\Gamma }(1+ϵ)\mathrm{\Gamma }(1ϵ)}{\mathrm{\Gamma }(1+2ϵ)\mathrm{\Gamma }(12ϵ)},$$ (54) and the symmetry under the replacement $`ij`$. Note that the current in Eq. (53) is conserved, $`q^\mu J_{\mu ;2P}^{a(1)}(q,ϵ)=0`$. Therefore, as discussed below Eq. (45), the single-particle term $`J_{\mu ;1P}^{a(1)}`$ vanishes and Eq. (53) gives the complete one-loop contribution to the soft-gluon current. We conclude that we have explicitly derived the result in Eq. (18) in the case of outgoing parton momenta $`p_i,p_j`$ (when $`\lambda _{ij}=\lambda _{iq}=\lambda _{jq}=+1).`$ The completely general result in Eq. (18) is straightforwardly obtained by crossing ($`p_ip_i`$ and/or $`p_jp_j`$) the parton momenta in Eq. (4.2) and repeating the steps that lead to Eq. (53). ## 5 Processes with two and three partons <br>and $`e^+e^{}`$$`3`$ jets at NNLO In general soft factorization formulae involve colour correlations (see Eqs. (12) and (26)). As discussed at the end of Sect. 3, in the case of processes with two or three hard external partons, the correlations are completely given in terms of the products of the colour-charge factors $`𝑻_i𝑻_j`$. In this case, moreover, these colour-charge factors can be expressed in terms of the Casimir operators of the hard partons. Therefore the colour algebra can explicitly be carried out and we can obtain QED-like factorization formulae. We first recall the relations between colour-charge factors and Casimir operators and then give the factorization formulae. When there are only two hard partons in the amplitude, we can use colour conservation $`(𝑻_1=𝑻_2)`$ to obtain $$𝑻_1𝑻_2\|(p_1,p_2)=𝑻_1^2\|(p_1,p_2)=C_1\|(p_1,p_2)=C_2\|(p_1,p_2),$$ (55) where $`C_1=C_2`$ is the Casimir of the hard partons. When the hard partons are three, the colour algebra can still be performed in closed form. Let us consider, for instance, the correlation term $`𝑻_1𝑻_2`$. Using colour conservation $`(𝑻_3=𝑻_1𝑻_2)`$, its action on the ket $`\|(p_1,p_2,p_3)`$ is $$2𝑻_1𝑻_2\|(p_1,p_2,p_3)=(𝑻_3^2𝑻_1^2𝑻_2^2)\|(p_1,p_2,p_3)=(C_3C_1C_2)\|(p_1,p_2,p_3)$$ (56) and thus $`𝑻_1𝑻_2`$ is again given in terms of the Casimir invariants $`C_i`$. The same can be done for the other products $`𝑻_1𝑻_3`$ and $`𝑻_2𝑻_3`$. When the number of hard partons is four or more, colour correlations cannot be avoided. In fact the number of independent equations following from colour conservation is not sufficient to express the products $`𝑻_i𝑻_j`$ in terms of Casimir invariants. The colour-algebra results in Eqs. (55) and (56) can be inserted in Eqs. (12), (24) and (26) to straightforwardly obtain the soft limit of the squared amplitudes with a soft gluon and two or three hard partons (plus any number of colourless partons). In both cases, the contributions from Eqs. (12) and (3) reconstruct the one-loop expansion of the all-order squared amplitude $`\|(\{p\})\|^2`$ in Eq. (1), and we obtain the following factorization formulae $$\|(q,p_1,p_2)\|^2g_\mathrm{S}^2\mu ^{2ϵ}\mathrm{\hspace{0.17em}4}C_1_{12}(q)\|(p_1,p_2)\|^2,$$ (57) $$\|(q,p_1,p_2,p_3)\|^2g_\mathrm{S}^2\mu ^{2ϵ}\mathrm{\hspace{0.17em}2}\left[2C_1_{12}(q)+C_A\left(_{23}(q)+_{13}(q)_{12}(q)\right)\right]\|(p_1,p_2,p_3)\|^2,$$ (58) where the kinematic factor $`_{ij}(q)`$ includes the loop corrections to the eikonal function $`𝒮_{ij}(q)`$ in Eq. (13). To one-loop accuracy, we find $$_{ij}(q)=𝒮_{ij}(q)\left\{1C_A\frac{g_\mathrm{S}^2}{8\pi ^2}\frac{1}{ϵ^2}\frac{\mathrm{\Gamma }^4(1ϵ)\mathrm{\Gamma }^3(1+ϵ)}{\mathrm{\Gamma }^2(12ϵ)\mathrm{\Gamma }(1+2ϵ)}\left[4\pi \mu ^2𝒮_{ij}(q)\right]^ϵ+𝒪(g_\mathrm{S}^4)\right\},$$ (59) where we have used Eq. (54). Note that in the amplitude with three hard partons, two of them (say, the ones with momenta $`p_1`$ and $`p_2`$) have to form a particle–antiparticle pair (either a $`q\overline{q}`$ pair or a gluon pair) and the third one (with momentum $`p_3`$) has to be a gluon. In Eq. (58) we have therefore set $`C_1=C_2`$ and $`C_3=C_A`$. The results in Eqs. (57) and (58), when combined with the analogous QED-like factorization formulae for the emission of two soft gluons and of two and three collinear partons, can be used to perform cross section calculations at NNLO for several processes. For instance, Eq. (57) is relevant to 2-jet production in $`e^+e^{}`$ annihilation and for the production of Drell–Yan, photon and vector-boson pairs in hadron collisions, while Eq. (58) can be used for the production of 3 jets in $`e^+e^{}`$ annihilation, (2+1) jets in deep-inelastic lepton–hadron collisions and (vector boson + jet) in hadron collisions. 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warning/0007/quant-ph0007111.html	ar5iv	text	# NONLINEAR QUANTUM EVOLUTION WITH MAXIMAL ENTROPY PRODUCTION ## I INTRODUCTION A number of recent, independent experiments have provided impressive bounds on the possible deviations from a linear and unitary propagation of pure quantum states, at least on a laboratory accessible space-time scale. The limits imposed in this way on potential generalizations of the standard unitary quantum equations of motions, as sought in relation to Hawking’s blackhole evaporation process , are likewise severe. Certainly, there always remains the possibility of modified dynamical laws on the (inaccessible) Plank scale , as well as under the extreme physical environment characteristic of singular cosmological phenomena. Related models of open system dynamics due to alleged statistical perturbations, e.g. from the space-time foam, have enjoyed considerable attention lately . But in case the unitarity of pure state propagation holds under universal conditions, one is necessarily lead to a quest for genuine nonlinear extensions for isolated systems, possibly involving an explicit arrow of time. Indeed, it has been pointed out in a fairly general ansatz that if the pure states happen to be attractors of a nonlinear evolution, then testing the unitary propagation of pure states alone cannot rule out a nonlinear propagation of mixtures. This situation has been noted recently in the context of certain nonlinear Lie-Poisson dynamics , wherein pure states still propagate in the usual hamiltonian way, while density matrices evolve nonlinearly, but preserving a time-independent spectrum. Unfortunately, the underlying physics remains rather obscure in these theories and the selection of particular realizations relevant to various experimental setups is, in general, a matter of guesswork. In the following we show that a physically meaningful nonlinear extension emerges when the fundamental postulates of quantum mechanics are supplemented by the first and second principles of thermodynamics, at the sole expense of ignoring the constraint of a linear, unitary evolution in time. The result is a largely irreversible, highly nonlinear generalization of the non-relativistic quantum Liouville equation, of a form closely related to the ansatz of Ref.(but not in the Lie-Poisson class), which features a number of rather intriguing properties. In particular, pure states are still propagated unitarily into pure states according to the usual (time-reversible) hamiltonian dynamics. The same is true of mixed states characterized by an initial equiprobable distribution on a (finite) set of uncorrelated (orthogonal) states. Non-pure states evolve so as to maximize the entropy production at each moment in time and to reach stationary states of maximum entropy (or minimum entropy production, according to Prigogine’s nonequilibrium principle ) on the shortest path in the appropriate state space. Precisely, mixed states arbitrarily distributed on a finite set of uncorrelated states evolve into mixed states distributed on an equal number of uncorrelated states, have a time-dependent eigenspectrum and eventually attain stationarity on a subset of energy eigenstates. A similar statement can be inferred, by extension, for mixtures of an infinite set of uncorrelated pure states. It follows as well that the probability distribution at equilibrium, on (a subset of) energy eigenstates, has a canonical-like dependence on the energy eigenvalues. For mixtures with an infinite energy range, the corresponding temperature is, of course, strictly positive, whereas for mixtures of a finite set of pure states the stationary state may display a “negative-temperature” distribution, in analogy to systems with a finite-dimensional state space. The above mentioned properties are endorsed by the positivity of the underlying evolution equation, which ensues by construction despite the high degree of nonlinearity involved. The nature of this essentially irreversible propagation becomes evident in the close-to-equilibrium limit, when the matrix elements of the density operator between energy eigenstates are found to undergo simple exponential decays to the canonical equilibrium values. Finally, proper (non-relativistic) invariance and conservation properties under the symmetry group of the hamiltonian are also accounted for. However, in the absence of an explicit general law of entropy increase, the time scale for thermal relaxation is set by one multiplication factor, a scalar functional, which is yet to be given a specific expression. Unlike the nonlinear Lie-Poisson dynamics , our framework apparently challenges the notion of separability of isolated, non-interacting systems, lack of which has long been thought to be unacceptable . We argue, nevertheless, that in a nonlinear theory it is necessary to refine the operational definition of isolation and to acknowledge that the mutual isolation of two non-interacting systems prohibits entanglement, if individual time-translation invariance is to be preserved. When this restriction is properly taken into account in the formulation of the corresponding equation of motion, separability can be easily recovered. On the other hand, the case where non-interacting subsystems are allowed to develop correlations spontaneously and eventually exchange energy (heat) is shown to correspond in our ansatz to the phenomenon of ideal thermal contact. From a precise technical perspective, the effect has its origin in that the second principle applies, as usual, to the total entropy of a compound system and not to the entropies of individual subsystems. This necessarily results in such redistribution of probabilities and energy as to maximize the overall entropy. In physical terms, an ideal gas is allowed to relax spontaneously to thermal equilibrium. The formalism can be adapted straightforwardly to cover nonstandard forms for the entropy and energy functionals. As immediate examples, we construct a generalization of the Lie-Poisson dynamics with maximal entropy production and a nonlinear extension of the standard von Neumann evolution with maximal increase of the nonextensive Tsallis q-entropy . ## II THE MODIFIED EQUATION OF MOTION Following an earlier suggestion , the state of a quantum system will be represented by a generalized “square-root” $`\gamma `$ of the density matrix $`\rho `$, defined by $$\rho =\gamma \gamma ^+.$$ (1) In analogy to the common terminology, the operator $`\gamma `$ (not necessarily hermitian) will be called here a state operator. Note that the above decomposition is always well-defined, although not unique, for any hermitian and positive definite $`\rho `$. On the other hand, to any given $`\gamma `$ there corresponds a unique hermitian and positively defined $`\rho `$. We also adopt the standard inner product on the associated Hilbert space of operators, $$(\beta \|\gamma )=Tr(\beta ^+\gamma ),$$ (2) such that for $`\gamma `$ normalized, $`(\gamma \|\gamma )=Tr(\gamma ^+\gamma )=1`$, the average of an observable O becomes the bilinear form $$(\gamma \|\text{O}\|\gamma )=(\gamma ^+O\gamma )=Tr(O\rho ),$$ (3) with O the super-operator defined by O, $$\text{O}\|\gamma )=\|O\gamma ).$$ (4) It is further convenient to define the tilde-conjugate $`\stackrel{~}{\text{A}}`$ of an arbitrary, and not necessarily linear, super-operator A , by $$(\text{A}\|\alpha ))^+=\stackrel{~}{\text{A}}\|\alpha ^+).$$ (5) It can be immediately verified that the super-operator A maps hermitian operators $`\alpha =\alpha ^+`$ into hermitian operators $`\beta =\beta ^+=\text{A}\|\alpha )`$ if and only if it is tilde-symmetric, $`\text{A}=\stackrel{~}{\text{A}}`$. For a super-operator generated by a linear operator, such as in Eq (4) above, the tilde-conjugate is given by $$\stackrel{~}{\text{A}}\|\alpha )=\|\alpha A^+).$$ (6) In particular, for the hermitian observable O it reads $$\stackrel{~}{\text{O}}\|\alpha )=\|\alpha O).$$ (7) The tilde operation is distributive against the addition and multiplication of super-operators, $`\stackrel{~}{(\text{A+B})}=\stackrel{~}{\text{A}}+\stackrel{~}{\text{B}}`$, $`\stackrel{~}{\text{AB}}=\stackrel{~}{\text{A}}\stackrel{~}{\text{B}}`$, and is anti-linear against multiplication by scalars, $`\stackrel{~}{(a\text{A})}=a^{}\stackrel{~}{\text{A}}`$. Let us consider now a massive isolated system characterized by an energy operator (hamiltonian) H and a state operator $`\gamma `$ (density matrix $`\rho =\gamma \gamma ^+`$), in an inertial reference system where its center-of-mass is at rest. We wish to find an equation of motion for this system which is first-order differential in time and such that : 1) Probability is conserved: $$\frac{d}{dt}(\gamma \|\gamma )=\frac{d}{dt}Tr(\rho )=0$$ (8) or $$(\dot{\gamma }\|\gamma )+(\gamma \|\dot{\gamma })=0,\dot{\gamma }=\frac{d}{dt}\gamma .$$ (9) 2) Energy is conserved (first principle of thermodynamics): $$\frac{d}{dt}(\gamma \|\text{H}\|\gamma )=\frac{d}{dt}Tr(H\rho )=0$$ (10) or $$(\dot{\gamma }\|\text{H}\|\gamma )+(\gamma \|\text{H}\|\dot{\gamma })=0.$$ (11) 3) The entropy production is always positive (second principle of thermodynamics in non-equilibrium form), $$\frac{d}{dt}\text{S}(\text{t})0.$$ (12) or $`\dot{\text{S}}(\text{t})=[Tr(\dot{\rho }\mathrm{ln}\rho )+Tr(\dot{\rho })]=[(\dot{\gamma }\|\mathrm{ln}(\gamma \gamma ^+)\|\gamma )+(\gamma \|\mathrm{ln}(\gamma \gamma ^+)\|\dot{\gamma })+`$ (13) $`+(\dot{\gamma }\|\gamma )+(\gamma \|\dot{\gamma })]0.`$ () where we adopt the standard entropy expression for a normalized state ($`Tr(\rho )=1`$) $$\text{S}(\text{t})=\text{k}_BTr\left[\rho (\text{t})\mathrm{ln}\rho (\text{t})\right]=\text{k}_B(\gamma (\text{t})\|\mathrm{ln}\left[\gamma (\text{t})\gamma ^+(\text{t})\right]\|\gamma (\text{t})),$$ (14) with k<sub>B</sub> the Boltzmann constant. In order to construct the desired equation of motion, we find it convenient to consider a stronger form of the second principle, by requiring that the entropy, as a functional of $`\gamma `$, increase in time along a path of maximum ascent. In other words, let the entropy production (10a) be maximized, for any given state $`\gamma `$, against variations of the time derivative $`\dot{\gamma }`$, under the constraints (8a) of conservation of probability and (9a) of conservation of energy. Note that the variation of $`\dot{\gamma }`$ must avoid the simple multiplication by a positive scalar, i.e. a trivial norm increase, since $`\dot{\text{S}}(\text{t})`$ increases then unconditionally. Hence the entropy production must be maximized against the “direction” of $`\dot{\gamma }`$, that is, against derivatives $`\dot{\gamma }`$ of equal, but otherwise arbitrary norm. This amounts to deriving the equation of motion from the following variational principle with constraints $`\delta \{(\dot{\gamma }\|\mathrm{ln}(\gamma \gamma ^+)\|\gamma )+(\gamma \|\mathrm{ln}(\gamma \gamma ^+)\|\dot{\gamma })+2\zeta (\dot{\gamma }\|\text{H}\|\gamma )+2\zeta ^{}(\gamma \|\text{H}\|\dot{\gamma })+`$ (15) $`+\xi [(\dot{\gamma }\|\gamma )+(\gamma \|\dot{\gamma })]+{\displaystyle \frac{2}{\sigma }}(\dot{\gamma }\|\dot{\gamma })\}=0.`$ (16) The variation refers to $`\dot{\gamma }`$ and $`\dot{\gamma }^+`$ only and the form of the Lagrange multipliers $`\zeta `$, $`\xi `$, $`\sigma `$ has been chosen for later convenience. $`\sigma `$ and $`\xi `$ are real scalars on account of their corresponding real functionals, while $`\zeta `$ is allowed to span complex values. Upon taking the variation of $`\dot{\gamma }`$, $`\dot{\gamma }^+`$, one is left with $$\|\dot{\gamma })=\sigma [\frac{1}{2}[\mathrm{ln}(\gamma \gamma ^+)]\|\gamma )+\zeta \text{H}\|\gamma )+\frac{\xi }{2}\|\gamma )]$$ (17) and the hermitian conjugate. Using Eq.(17) into conditions (9) and (11) immediately gives $$Re\zeta =\frac{1}{2}\frac{(\gamma \|\text{H}\mathrm{ln}(\gamma \gamma ^+)\|\gamma )+\text{E}\frac{\text{S}}{\text{k}_B}}{\overline{\mathrm{\Delta }H^2}},$$ (19) $$\xi =\frac{\text{S}}{\text{k}_B(\gamma \|\gamma )}2Re\zeta \text{E},$$ (20) with $`\left(\gamma \|\gamma \right)=1,\text{E}=(\gamma \|\text{H}\|\gamma )/(\gamma \|\gamma )`$ the average energy of the system, $`\text{S}0`$ the entropy and $`\overline{\mathrm{\Delta }H^2}=(\gamma \|H^2\|\gamma )\text{E}^2`$ the squared energy deviation. One can also check condition (()II) and find that $$\frac{\dot{\text{S}}}{k_B}=\sigma \left(\theta \|\theta \right),$$ (21) $$\|\theta )=\mathrm{ln}(\gamma \gamma ^+)\|\gamma )+2\zeta \text{H}\|\gamma )+\xi \|\gamma ),$$ (22) hence inequality (()II) is satisfied provided $$\sigma 0.$$ (23) In deriving expression (21) we used the fact that for $`\|\theta )`$ as in Eq.(22), and $`Re\zeta `$, $`\xi `$ given by eqs.(II), it is also true that $$(\gamma \|\text{H}\|\theta )=\mathrm{\hspace{0.33em}0},(\gamma \|\theta )=0.$$ (24) Let us stress at once that, unlike the usual stationary action principle, our variational principle Eq.(15) does not involve variations of functionals over an extended interval of time, but only variations against $`\dot{\gamma }`$ which are local in time, at each given instant t. As a result, the Lagrange parameters $`\zeta `$, $`\xi `$, $`\sigma `$ need only be constants against these same variations of $`\dot{\gamma }`$ and not constants of time or $`\gamma `$ itself. Likewise, condition (23) for $`\sigma `$ only guarantees the positivity of the entropy production, but does not make $`\dot{S}`$ independent of time. Hence all parameters in the equation of motion (17) for $`\gamma `$, as well as the entropy production and the entropy itself, are time dependent through their dependence on $`\gamma `$. Furthermore, note that $`Re\zeta `$, $`\xi `$ are really functionals of $`\rho `$ and H only and therefore are invariant under transformations of the kind $$\gamma \gamma U,UU^+=U^+U=I,$$ (25) which leave the density matrix unchanged, $$\rho \rho =\gamma UU^+\gamma ^+=\gamma \gamma ^+.$$ (26) Eq.(17) will be invariant in its entirety under transformations (26) provided $`\sigma `$ and $`Im\zeta `$ are likewise invariant as functionals of $`\rho `$ and H. In that case the entropy production Eq.(21) will also be invariant under transformations (26), as should be expected on physical grounds. Now let us introduce the equivalent equation of motion for the density matrix, starting from $$\dot{\rho }=\dot{\gamma }\gamma ^++\gamma \dot{\gamma }^+.$$ (27) It follows at once that $$\dot{\rho }=\sigma \left[\rho \mathrm{ln}\rho +Re\zeta \{H\text{E},\rho \}\rho Tr(\rho \mathrm{ln}\rho )\right]+i\sigma (Im\zeta )[\rho ,H],$$ (28) where {,} denotes the anticommutator, as usual. The commutator on the right hand side of Eq.(28) obviously provides the unitary hamiltonian limit, and the standard Liouville equation suggests $$\sigma (Im\zeta )=\frac{1}{\mathrm{}}.$$ (29) Setting now, for simplicity, $`Re\zeta \zeta `$, the final form of our equation of motion for the density matrix is found to be, in common notation, $$\dot{\rho }=\sigma \left[\rho \mathrm{ln}\rho +\zeta (\rho ,H\text{E})\{H\text{E},\rho \}\rho \frac{Tr(\rho \mathrm{ln}\rho )}{Tr(\rho )}\right]+\frac{i}{\mathrm{}}[\rho ,H],$$ (30) where $`\zeta (\rho ,H\text{E})={\displaystyle \frac{1}{2}}{\displaystyle \frac{Tr[(H\mathrm{E})\rho \mathrm{ln}\rho ]}{Tr[(H\mathrm{E})^2\rho ]}},`$ $`\sigma (\rho ,H\text{E})0,`$ $`Tr(\rho )=const.(=1),`$ $`\text{E}={\displaystyle \frac{Tr(H\rho )}{Tr(\rho )}}=const.,`$ $`\dot{\text{S}}=\text{k}_B{\displaystyle \frac{d}{dt}}Tr(\rho \mathrm{ln}\rho )0.`$ The scale setting parameter $`\sigma `$ remains unspecified so far, and will be regarded in the following as a functional of $`\rho `$ and H. In order to secure that Eq.(30) is invariant under a scaling $`\rho a\rho `$, it must be assumed that $`\sigma (a\rho ,H)=\sigma (\rho ,H)`$, in which case scaling invariance is verified straightforwardly. Moreover, since eq.(30) should not show a dependence on the zero-point of the energy, it may be assumed also, as above, that $`\sigma =\sigma (\rho ,H\text{E})`$. For simplicity, it will be understood throughout the following that $`Tr(\rho )=1`$. It is interesting to note that Eq.(30) can be recovered from a modified form of the nonlinear ansatz proposed in Ref., $`\dot{\rho }={\displaystyle \frac{i}{\mathrm{}}}[\rho ,H]{\displaystyle \frac{a}{T}}\left[f(\rho )\rho {\displaystyle \frac{Tr(f(\rho ))}{Tr(\rho )}}\right],`$ with the obvious substitutions $`{\displaystyle \frac{a}{T}}\sigma `$, $`f(\rho )\rho \mathrm{ln}\rho +\zeta \{H,\rho \}.`$ ## III FUNDAMENTAL PROPERTIES OF THE NONLINEAR EVOLUTION Eq.(30) secures the hermiticity and positivity of the density matrix by construction, since it has been generated from an equation for the state operator $`\gamma `$. Conversely, Eq.(30) can be easily decomposed into the corresponding equations for $`\gamma `$ and $`\gamma ^+`$ by using the substitution $`\rho =\gamma \gamma ^+`$, hence the equations of motion for $`\rho `$ and $`\gamma `$ are indeed equivalent. Assuming again a well-behaved $`\sigma `$, Eq.(30) is seen to be covariant under time-independent unitary transformations, $`\rho \stackrel{~}{\rho }=U^+\rho U,H\stackrel{~}{H}=U^+HU,`$ and, in particular, invariant under the (time-independent) symmetry group of the hamiltonian, $`[U,H]=0`$. But an observable O which commutes with H, $`[H,O]=0`$, is not, in general, an integral of motion. More details on the problem follow in Sec. 4. It is convenient to absorb the hamiltonian commutator term by setting, in analogy to the usual Heisenberg representation, $$\rho (\text{t})=\mathrm{exp}\left[\frac{i}{\mathrm{}}H\text{t}\right]\overline{\rho }(\text{t})\mathrm{exp}\left[\frac{i}{\mathrm{}}H\text{t}\right].$$ (31) Upon substituting expression (31), eq.(30) becomes $$\dot{\overline{\rho }}=\sigma \left[\overline{\rho }\mathrm{ln}\overline{\rho }+\zeta \{H\text{E},\overline{\rho }\}\overline{\rho }Tr(\overline{\rho }\mathrm{ln}\overline{\rho })\right].$$ (32) Now note that for $`\overline{\rho }`$ corresponding to a pure state, $`\overline{\rho }=\overline{\rho }^2=\|\mathrm{\Psi }\mathrm{\Psi }\|`$, the entropy operator vanishes together with the coefficient $`\zeta `$, i.e. $`\overline{\rho }\mathrm{ln}\overline{\rho }0`$, $`\zeta (H\text{E})0`$, such that $`\dot{\overline{\rho }}(\text{t})=0`$ and $`\overline{\rho }(\text{t})=\overline{\rho }(0)=\|\mathrm{\Psi }\mathrm{\Psi }\|`$, if $`\sigma `$ is also finite in this limit. From Eq.(31) it follows then that a pure state evolves into a pure state according to the usual hamiltonian law: $$\rho (\text{t})=\rho ^2(\text{t})=\mathrm{exp}\left[\frac{i}{\mathrm{}}H\text{t}\right]\|\mathrm{\Psi }\mathrm{\Psi }\|\mathrm{exp}\left[\frac{i}{\mathrm{}}H\text{t}\right].$$ (33) Another situation where the nonlinear evolution reduces to the hamiltonian law is found for uniform (equiprobable) distributions $`\rho _{unif}`$, when the eigenvalues of the density matrix are all identical. In this case one has the identity $`\overline{\rho }_{unif}\mathrm{ln}\overline{\rho }_{unif}=\overline{\rho }_{unif}Tr(\overline{\rho }_{unif}\mathrm{ln}\overline{\rho }_{unif})`$ and $`\zeta (\rho _{unif},H\text{E})0`$, wherefrom $`\dot{\overline{\rho }}_{unif}(\text{t})=0,\overline{\rho }_{unif}(\text{t})=\rho _{unif}(0)`$ and $$\rho _{unif}(\text{t})=\mathrm{exp}\left[\frac{i}{\mathrm{}}H\text{t}\right]\rho _{unif}(0)\mathrm{exp}\left[\frac{i}{\mathrm{}}H\text{t}\right].$$ (34) Recall that under unitary propagation the cardinality of the set of nonzero eigenvalues of the density matrix is preserved in time. The same holds true if the density matrix evolves according to Eq.(30). In order to see this, let $`P_\nu =\|\varphi _\nu \varphi _\nu \|`$ be the projector on some eigenstate of $`\overline{\rho }(\text{t}),\overline{\rho }P_\nu =\rho _\nu P_\nu `$, where $`\rho _\nu =Tr(P_\nu \overline{\rho })`$ denotes the corresponding eigenvalue. Since $`Tr(\dot{\overline{\rho }}P_\nu )=\dot{\rho }_\nu `$, multiplying Eq.(32) by $`P_\nu `$ and taking the trace yields $$\dot{\rho }_\nu =\sigma [\rho _\nu \mathrm{ln}\rho _\nu +\alpha _\nu (\overline{\rho },H)\rho _\nu ],$$ (36) $$\alpha _\nu (\overline{\rho },H)=2\zeta (\overline{\rho },H)Tr[P_\nu (\text{t})(H\text{E})]+\frac{\text{S}(\text{t})}{\text{k}_B}.$$ (37) Taking $`\rho _\nu \mathrm{ln}\rho _\nu 0`$ for $`\rho _\nu =0`$ gives $`\dot{\rho }_\nu =0`$ and $`\rho _\nu (\text{t})=0`$, i.e. a zero eigenvalue evolves into a zero eigenvalue. As an immediate corollary, density matrices with a finite number of “occupied” state vectors (i.e. a finite number of nonzero eigenvalues) are necessarily driven towards a stationary state with a thermal-like distribution on a finite set of energy eigenstates. Indeed, in this case the entropy, as a functional of the eigenvalues $`\rho _\nu `$ and under the constraint of conserved energy and probability, has a finite absolute maximum. For this reason, and because $`\dot{\text{S}}(\text{t})0`$ at all times, it can only evolve towards a stationary value less or equal to that maximum. But, as will be shown, $`\dot{S}(\text{t})=0`$ implies in fact $`\dot{\overline{\rho }}=0`$ and $`[\overline{\rho },H]=0`$, and the stationary version of Eq.(IIIa) gives then the thermal-like distribution. Let us prove now that $`\dot{S}(\text{t})=0`$ implies stationarity. We begin by making a change of variables, $`\rho _\nu =e^{\eta _\nu },\eta _\nu 0`$, such as to write $$\frac{\text{S}}{\text{k}_B}=\underset{\nu }{}\eta _\nu e^{\eta _\nu }$$ (38) and $$\frac{\dot{\text{S}}}{\text{k}_B}=\underset{\nu }{}(\dot{\eta _\nu }\eta _\nu \dot{\eta _\nu })e^{\eta _\nu }=\underset{\nu }{}\eta _\nu \dot{\eta _\nu }e^{\eta _\nu },$$ (39) since $`\underset{\nu }{}\dot{\eta }_\nu e^{\eta _\nu }=\underset{\nu }{}\dot{\rho }_\nu =0`$. Also, Eqs.(III) give $$\dot{\eta }_\nu =\sigma [\eta _\nu \alpha _\nu ],$$ (40) which taken into Eq.(39) produces $$\frac{\dot{\text{S}}}{\text{k}_B}=\sigma \underset{\nu }{}\left[\eta _\nu ^2\alpha _\nu \eta _\nu \right]e^{\eta _\nu }=\sigma \underset{\nu }{}\left[\alpha _\nu \eta _\nu \alpha _\nu ^2\right]e^{\eta _\nu }+\frac{1}{\sigma }\underset{\nu }{}(\dot{\eta }_\nu ^2)e^{\eta _\nu }.$$ (41) Further, use of the explicit expression for $`\alpha _\nu `$, Eq.(IIIb), will show that $`{\displaystyle \underset{\nu }{}}\left[\alpha _\nu \eta _\nu \alpha _\nu ^2\right]e^{\eta _\nu }={\displaystyle \underset{\nu }{}}\left[2\zeta \eta _\nu Tr(P_\nu (H\mathrm{E}))\mathrm{e}^{\eta _\nu }+{\displaystyle \frac{\text{S}}{\text{k}_\mathrm{B}}}\eta _\nu \mathrm{e}^{\eta _\nu }\right]`$ $`{\displaystyle \underset{\nu }{}}\left[4\zeta ^2[Tr(P_\nu (H\text{E}))]^2e^{\eta _\nu }+4\zeta {\displaystyle \frac{\text{S}}{\text{k}_B}}Tr(P_\nu (H\text{E}))e^{\eta _\nu }+\left({\displaystyle \frac{\text{S}}{\text{k}_B}}\right)^2e^{\eta _\nu }\right]=`$ $`=2\zeta Tr((H\text{E})\rho \mathrm{ln}\rho )+\left({\displaystyle \frac{\text{S}}{\text{k}_B}}\right)^24\zeta ^2Tr((H\text{E})^2\rho )`$ $$4\zeta \frac{\text{S}}{\text{k}_B}Tr((H\text{E})\rho )\left(\frac{\text{S}}{\text{k}_B}\right)^2=0,$$ (42) where we have used the explicit expression of $`\zeta `$, Eq.(IIa). Accounting for Eq.(42) in Eq.(41) shows that $$\frac{\dot{\text{S}}}{\text{k}_B}=\frac{1}{\sigma }\underset{\nu }{}\dot{\eta }_\nu ^2e^{\eta _\nu },$$ (43) wherefrom it follows that $`\dot{\text{S}}=0`$ if and only if $`\dot{\eta }_\nu =0`$ or, equivalently, $`\dot{\rho }_\nu =0`$. Consider now that the system is evolving in an asymptotic region where $`\dot{\text{S}}(\mathrm{t})0`$ for all t $`>`$ 0 . Since necessarily $`\dot{\rho }_\nu 0`$, $`\dot{\overline{\rho }}`$ must be driven by a unitary evolution, $`\overline{\rho }(\mathrm{t}\mathrm{t}_0)=\mathrm{U}(\mathrm{t})\overline{\rho }(\mathrm{t}_0)\mathrm{U}^+(\mathrm{t})`$. But for $`\dot{\rho }_\nu 0`$, Eq.(39) gives $`\mathrm{ln}\rho _\nu =\alpha _\nu `$, which in turn shows that $`\overline{\rho }\mathrm{ln}\overline{\rho }={\displaystyle \underset{\nu }{}}\rho _\nu \alpha _\nu P_\nu =2\zeta {\displaystyle \underset{\nu }{}}\rho _\nu P_\nu Tr(P_\nu (H\text{E})){\displaystyle \underset{\nu }{}}\rho _\nu P_\nu {\displaystyle \frac{\text{S}}{\text{k}_B}}=`$ (44) $`=2\zeta \{H_D\text{E},\overline{\rho }\}{\displaystyle \frac{\text{S}}{\text{k}_B}}\overline{\rho },`$ (45) where $`H_D=\underset{\nu }{}P_\nu Tr(P_\nu H)`$ is the diagonal part of H in the eigenbasis of $`\overline{\rho }`$, $`[H_D,\overline{\rho }]=0`$. Introducing the above result into Eq.(32), one is lead to $$\dot{\overline{\rho }}=\sigma \zeta \{H_{ND},\overline{\rho }\},$$ (47) with $`H_{ND}=HH_D`$ the non-diagonal part of H relative to $`\overline{\rho }`$. But Eq.(47) cannot generate a unitary evolution unless $`H_{ND}=0`$, which implies that stationary entropy over an extended period of time is equivalent to $$H_D=H,$$ (48) hence $`[\overline{\rho },H]=0`$ and $`\dot{\overline{\rho }}=0`$. In other words, the density matrix of the system (see also Eq.(31)) is stationary and also diagonal over energy eigenstates. The explicit form of the occupation probability corresponding to an (occupied) energy state of energy $`\mathrm{E}_\nu `$ follows from Eqs.(III), $$\rho _\nu ^{eq}=\mathrm{exp}\left[2\zeta ^{eq}(\text{E}_\nu \text{E})\frac{\text{S}^{eq}}{\text{k}_B}\right],$$ (49) and can be brought to the recognizable thermal form $$\rho _\nu ^{eq}=\frac{1}{\mathrm{Z}}e^{\beta \mathrm{E}_\nu },$$ (50) with $`\beta =2\zeta ^{eq}`$ and $`\text{Z}=\beta \text{E}+(\text{S}^{eq}/\text{k}_B)`$. Surprisingly, the parameter $`\zeta `$ is seen to become at equilibrium, up to a factor of 2, the reciprocal temperature $`\beta =1/\text{k}_B\text{T}`$. It should be noted, nevertheless, that according to our initial assumptions Eq.(50) applies only to a finite number of energy eigenstates and therefore does not refer to a canonical equilibrium distribution. More precisely, the sign of $`\zeta ^{eq}`$, and of the generalized temperature T, is not necessarily positive. For instance, let the occupied energy eigenstates be labeled by $`\nu `$ in order of their increasing energy $`\text{E}_\nu `$ and let their total number be N. If the conserved average energy E is such that $$\text{E}\frac{1}{N}\underset{\nu =1}{\overset{N}{}}\text{E}_\nu ,$$ (51) a simple calculation will verify that the entropy will have an (absolute) maximum, corresponding to the equilibrium state, on a distribution characterized by a negative $`\zeta ^{eq}`$, hence a “negative temperature”. At this point, let us examine more closely the restrictive assumption of a finite number of non-vanishing eigenvalues for the density matrix. It can be noted that it has entered the argument developed above solely by way of the related assumption of a finite absolute maximum for the entropy, at the given value E for the average energy. However, there is good reason to assume that such an absolute maximum exists at least for a large class of distributions over infinite sets of (orthogonal) state vectors. If we can extend this “finite absolute maximum” conjecture to all distributions with a finite average energy, it becomes possible to generalize the results in Eqs.(48-50) and state that the nonlinear dynamics described by Eqs.(30,32) drives the system towards an equilibrium state on energy eigenstates, with thermal-like occupation probabilities. Of course, when the range of occupied energy eigenvalues extends to infinity, relation (51) can no longer be satisfied for any finite E, and the corresponding temperature can only be positive. Finally, we wish to clarify the consistency of the present nonlinear dynamics, which follows a path of maximal entropy production, with Prigogine’s celebrated principle of minimum entropy production. Let us recall that, according to the latter, physical systems evolve towards stationary states which have minimum entropy production compared to slightly displaced neighboring states. Given that the entropy is a convex functional on the state (configuration) space, bounded from above for any finite average energy, this implies that the physical evolution will take the system towards a local maximum of the entropy or at least towards a ridge. Indeed, in a small enough vicinity of a maximum of the entropy, or of a ridge, any evolution with positive entropy production will eventually enter a regime where $`\dot{S}`$ decreases in time until it vanishes in the equilibrium state or is minimized for the stationary states corresponding to a ridge. The variational principle Eq.(15) only complements this picture by stating that the evolution should follow the shortest route to a state of maximum entropy, i.e. the direction of the physical path is selected from among all directions satisfying $`\dot{S}0`$ by the requirement that the increase in entropy be maximized at each point in time. In this case it can be said that the entropy production evolves towards a minimum of the maximum, to be attained on a local maximum or a ridge of the entropy (hyper)surface in state space. ## IV THE LINEAR NEAR-EQUILIBRIUM LIMIT It is natural to anticipate a linear limit for any nonlinear dynamics evolving sufficiently close to a canonical thermal equilibrium state, at least in the high-temperature limit. For the modified equation of motion proposed here, the linearization process entails essentially the approximation of the entropy operator $`\overline{\rho }\mathrm{ln}\overline{\rho }`$ to first order in $`\mathrm{\Delta }(\overline{\rho }\rho ^{eq})`$ around the target equilibrium state $$\rho ^{eq}=\frac{1}{\text{Z}}e^{\beta H},\mathrm{ln}\text{Z}=\beta \text{E}+\frac{\text{S}^{eq}}{\text{k}_B},$$ (52) with given average energy E and reciprocal temperature $`\beta `$. We proceed from the exact expansion $$\mathrm{ln}\overline{\rho }=\underset{n=1}{\overset{\mathrm{}}{}}\frac{1}{n}(I\overline{\rho })^n,$$ (53) which gives for $`\overline{\rho }=\rho ^{eq}+\mathrm{\Delta }\overline{\rho }`$, and in symmetrized form, $$(\rho ^{eq}+\mathrm{\Delta }\overline{\rho })\mathrm{ln}(\rho ^{eq}+\mathrm{\Delta }\overline{\rho })=\frac{1}{2}\{(\rho ^{eq}+\mathrm{\Delta }\overline{\rho }),\underset{n=1}{\overset{\mathrm{}}{}}\frac{1}{n}(I\rho ^{eq}\mathrm{\Delta }\overline{\rho })^n\}.$$ (54) Separation of the zero- and first-order terms in $`\mathrm{\Delta }\overline{\rho }`$ yields $$(\rho ^{eq}+\mathrm{\Delta }\overline{\rho })\mathrm{ln}(\rho ^{eq}+\mathrm{\Delta }\overline{\rho })=\rho ^{eq}\mathrm{ln}\rho ^{eq}\frac{1}{2}\{\mathrm{\Delta }\overline{\rho },\mathrm{ln}\rho ^{eq}\}+\frac{1}{2}\{\rho ^{eq},\mathrm{\Lambda }(\mathrm{\Delta }\overline{\rho })\},$$ (55) where $`\mathrm{\Lambda }(\mathrm{\Delta }\overline{\rho })`$ represents the collection of all terms first-order in $`\mathrm{\Delta }\overline{\rho }`$ from the infinite sum on the right hand side of Eq.(54). In order to calculate $`\mathrm{\Lambda }(\mathrm{\Delta }\overline{\rho })`$, it is convenient to define the superoperator R and its tilde-conjugate $`\stackrel{~}{\text{R}}`$ by $$\text{R}\mathrm{\Delta }\overline{\rho }=(I\rho ^{eq})\mathrm{\Delta }\overline{\rho },$$ (57) $$\stackrel{~}{\text{R}}\mathrm{\Delta }\overline{\rho }=\mathrm{\Delta }\overline{\rho }(I\rho ^{eq}),$$ (58) $`[\text{R},\stackrel{~}{\text{R}}]=0.`$ The expression of $`\mathrm{\Lambda }(\mathrm{\Delta }\overline{\rho })`$ can be obtained now in the compact form $$\mathrm{\Lambda }(\mathrm{\Delta }\overline{\rho })=\underset{n=1}{\overset{\mathrm{}}{}}\frac{1}{n}\underset{m=0}{\overset{n1}{}}\text{R}^m\stackrel{~}{\text{R}}^{nm1}\mathrm{\Delta }\overline{\rho }.$$ (59) But R has a well-defined inverse and the super-operator sum in the above expression can be rewritten as $`{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\stackrel{~}{\text{R}}^{n1}}{\text{n}}}{\displaystyle \underset{m=0}{\overset{n1}{}}}\left(\text{R}\stackrel{~}{\text{R}}^1\right)^m=\left(\text{I}\text{R}\stackrel{~}{\text{R}}^1\right)^1{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\stackrel{~}{\text{R}}^{n1}}{\text{n}}}\left(\text{I}\left(\text{R}\stackrel{~}{\text{R}}^1\right)^n\right)=`$ (60) $`=\left(\stackrel{~}{\text{R}}\text{R}\right)^1{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}\left({\displaystyle \frac{\stackrel{~}{\text{R}}^n}{\text{n}}}{\displaystyle \frac{\text{R}^n}{\text{n}}}\right)=\left(\stackrel{~}{\text{R}}\text{R}\right)^1\left[\mathrm{ln}\left(\text{I}\text{R}\right)\mathrm{ln}\left(\text{I}\stackrel{~}{\text{R}}\right)\right].`$ (61) Taking also into account that $$\rho ^{eq}\mathrm{\Delta }\overline{\rho }=(\text{I}\text{R})\mathrm{\Delta }\overline{\rho },\mathrm{\Delta }\overline{\rho }\rho ^{eq}=(\text{I}\stackrel{~}{\text{R}})\mathrm{\Delta }\overline{\rho }$$ (63) and $$\mathrm{ln}(\text{I}\text{R})=\beta \text{H}(\mathrm{ln}\text{Z})\text{I},\mathrm{ln}(\text{I}\stackrel{~}{\text{R}})=\beta \stackrel{~}{\text{H}}(\mathrm{ln}\text{Z})\text{I},$$ (64) where $$\text{H}\mathrm{\Delta }\overline{\rho }=H\mathrm{\Delta }\overline{\rho },\stackrel{~}{\text{H}}\mathrm{\Delta }\overline{\rho }=\mathrm{\Delta }\overline{\rho }H,$$ (65) we are lead to: $$\frac{1}{2}\{\rho ^{eq},\mathrm{\Lambda }(\mathrm{\Delta }\overline{\rho })\}=\frac{\beta }{2}\left(\text{H}\stackrel{~}{\text{H}}\right)\mathrm{coth}\left[\frac{\beta }{2}\left(\text{H}\stackrel{~}{\text{H}}\right)\right]\mathrm{\Delta }\overline{\rho }.$$ (66) Returning to Eq.(55), the first-order in $`\mathrm{\Delta }\overline{\rho }`$ approximation to the entropy operator reads now $`(\rho ^{eq}+\mathrm{\Delta }\overline{\rho })\mathrm{ln}(\rho ^{eq}+\mathrm{\Delta }\overline{\rho })=\rho ^{eq}\mathrm{ln}\rho ^{eq}{\displaystyle \frac{\beta }{2}}\{\mathrm{\Delta }\overline{\rho },H\text{E}\}{\displaystyle \frac{\text{S}^{eq}}{\text{k}_B}}\mathrm{\Delta }\overline{\rho }`$ (67) $`{\displaystyle \frac{\beta }{2}}\left(\text{H}\stackrel{~}{\text{H}}\right)\mathrm{coth}\left[{\displaystyle \frac{\beta }{2}}\left(\text{H}\stackrel{~}{\text{H}}\right)\right]\mathrm{\Delta }\overline{\rho }.`$ (68) Note that taking the trace in Eq.(53) gives $`\text{S}(\text{t})\text{S}^{eq}`$ in this regime. Similarly, a simple calculation shows that $`\zeta \beta /2`$. Assuming also that $`\sigma \sigma ^{eq}=const.(\text{E},\beta )`$ and inserting everything into Eq.(32) yields the linearized equation of motion $$\dot{\overline{\rho }}=\sigma ^{eq}\frac{\beta }{2}\left(\text{H}\stackrel{~}{\text{H}}\right)\mathrm{coth}\left[\frac{\beta }{2}\left(\text{H}\stackrel{~}{\text{H}}\right)\right]\mathrm{\Delta }\overline{\rho }$$ (69) or , as well, $$\mathrm{\Delta }\dot{\overline{\rho }}=\sigma ^{eq}\frac{\beta }{2}\left(\text{H}\stackrel{~}{\text{H}}\right)\mathrm{coth}\left[\frac{\beta }{2}\left(\text{H}\stackrel{~}{\text{H}}\right)\right]\mathrm{\Delta }\overline{\rho }.$$ (70) The general solution of Eq.(69) is given by $$\overline{\rho }\left(\text{t}\right)=\text{e}^{\sigma ^{eq}\left(\beta \right)\mathrm{t}}\text{e}^{𝐆\mathrm{t}}\overline{\rho }\left(0\right)+\left(1\text{e}^{\sigma ^{eq}\left(\beta \right)\mathrm{t}}\right)\rho ^{eq},$$ (71) where $$\text{G}=\sigma ^{eq}(\beta )\frac{\beta }{2}\left(\text{H}\stackrel{~}{\text{H}}\right)\mathrm{coth}\left[\frac{\beta }{2}\left(\text{H}\stackrel{~}{\text{H}}\right)\right]\text{I}.$$ (72) We observe immediately that G is tilde-symmetric, hence it maps any hermitian operator into a hermitian operator, and that it preserves probability, since $`Tr[\text{G}\overline{\rho }]=0`$, $`Tr[\text{e}^{𝐆\mathrm{t}}\overline{\rho }(0)]=Tr[\overline{\rho }(0)]=1`$. This is entirely sufficient to secure the hermiticity of $`\overline{\rho }`$ and the overall conservation of probability. Unfortunately, the action of G does not always preserve positivity and G cannot be identified as a generator of Lindblad type . But the positive domain of G does include the small neighborhood of $`\rho ^{eq}`$ identified as the near-equilibrium domain. Indeed, note first that in the diagonal representation of the hamiltonian, the matrix elements of $`\overline{\rho }^0(\text{t})=\text{e}^{𝐆\mathrm{t}}\overline{\rho }(0)`$ obey the simple damping law $$\overline{\rho }_{\mu \nu }^0(\mathrm{t})=\text{e}^{\gamma _{\mu \nu }(\beta )\text{t}}\overline{\rho }_{\mu \nu }(0),$$ (73) where the (temperature-dependent) relaxation coefficient $`\gamma _{\mu \nu }`$ is given by $$\gamma _{\mu \nu }(\beta )=\sigma ^{eq}(\beta )\left[\frac{\beta }{2}\left(\text{E}_\mu \text{E}_\nu \right)\mathrm{coth}\left[\frac{\beta }{2}(\text{E}_\mu \text{E}_\nu )\right]1\right],$$ (74) $`\gamma _{\nu \nu }=0,\gamma _{\nu \mu }=\gamma _{\mu \nu }.`$ If we consider now an arbitrary state vector $`\|\mathrm{\Psi }=\underset{\nu =0}{\overset{\mathrm{}}{}}\text{E}_\nu \|\mathrm{\Psi }\|\text{E}_\nu `$ and the matrix element $$\mathrm{\Psi }\|\overline{\rho }^0(\text{t})\|\mathrm{\Psi }=\underset{\nu =0}{\overset{\mathrm{}}{}}\mathrm{\Psi }\|\text{E}_\nu \overline{\rho }_{\nu \nu }(0)\text{E}_\nu \|\mathrm{\Psi }+\underset{\genfrac{}{}{0pt}{}{\mu ,\nu =0}{\mu >\nu }}{\overset{\mathrm{}}{}}Re\left[\mathrm{\Psi }\|\text{E}_\mu \overline{\rho }_{\mu \nu }(0)\text{E}_\nu \|\mathrm{\Psi }\right]\text{e}^{\gamma _{\mu \nu }\mathrm{t}}.$$ (75) it is easily seen that $`\overline{\rho }(\text{t})`$ remains positive for $`\text{t}>0`$ if the initial off-diagonal matrix elements $`\overline{\rho }_{\mu \nu }(0)`$, $`\mu \nu `$, are sufficiently small, as expected for the near-equilibrium regime. On the other hand, one can resort to the equation of motion for the state operator $`\gamma `$, Eq(17), and derive a linear approximation in $`\mathrm{\Delta }\gamma =\gamma \gamma ^{eq}`$, $`\mathrm{\Delta }\gamma ^+=\gamma ^+(\gamma ^{eq})^+`$ by the same procedure as above. The result reads $$\dot{\mathrm{\Delta }\gamma }=\left[\left(\sigma \beta +\frac{i}{\mathrm{}}\right)\left(\text{H}\stackrel{~}{\text{H}}\right)\left(\stackrel{~}{\text{R}}\text{R}\right)^1\left(\gamma ^{eq}\mathrm{\Delta }\gamma ^++\mathrm{\Delta }\gamma \left(\gamma ^{eq}\right)^+\right)\right]\gamma ^{eq}$$ (77) $$\dot{\mathrm{\Delta }\gamma }^+=\left(\sigma \beta \frac{i}{\mathrm{}}\right)\left(\gamma ^{eq}\right)^+\left[\left(\text{H}\stackrel{~}{\text{H}}\right)\left(\stackrel{~}{\text{R}}\text{R}\right)^1\left(\gamma ^{eq}\mathrm{\Delta }\gamma ^++\mathrm{\Delta }\gamma \left(\gamma ^{eq}\right)^+\right)\right]$$ (78) and shows that, up to first-order terms in $`\mathrm{\Delta }\gamma `$, $$\rho (\text{t})=\rho ^{eq}+\mathrm{\Delta }\rho \left[\gamma ^{eq}+\mathrm{\Delta }\gamma \right]\left[\left(\gamma ^{eq}\right)^++\mathrm{\Delta }\gamma ^+\right]$$ (79) such that $`\mathrm{\Delta }\rho =\mathrm{\Delta }\gamma \left(\gamma ^{eq}\right)^++\gamma ^{eq}\mathrm{\Delta }\gamma ^+`$ evolves according to Eq.(70) derived above. Furthermore, the conservation of energy follows from $$Tr\left(H\mathrm{\Delta }\dot{\overline{\rho }}\right)=\sigma ^{eq}Tr\left(H\mathrm{\Delta }\overline{\rho }\right)$$ (80) upon recalling that, according to the original equation of motion, the initial state necessarily has the same average energy E as the asymptotic equilibrium state. The initial conditions for Eq.(71) are so restricted to $`Tr\left(H\mathrm{\Delta }\overline{\rho }(0)\right)=0`$, which implies of course $`Tr\left(H\overline{\rho }(0)\right)=\text{E}`$. As a general feature of the underlying physics, it follows from Eqs.(71), (73) and (74) that the greater the energy gap between two energy eigenstates, the faster the quantum correlation between them is destroyed as the system evolves towards equilibrium. On the other hand, the relaxation of the occupation probabilities for each of the energy states proceeds at a common rate, independent of the corresponding energy level, since $`\overline{\rho }_{\nu \nu }(\mathrm{t})=\frac{\text{e}^{\beta \mathrm{E}_\nu }}{\text{Z}}\left(1\text{e}^{\sigma ^{\mathrm{eq}}\left(\beta \right)\mathrm{t}}\right)+\text{e}^{\sigma ^{\mathrm{eq}}\left(\beta \right)\mathrm{t}}\overline{\rho }_{\nu \nu }(0)`$. As a corrolary, the same holds true for the average of any observable O which commutes with the hamiltonian, $`[H,O]=0`$, since $`O(\text{t})=Tr[O\rho (\text{t})]=Tr\left[O\text{e}^{\frac{i}{\mathrm{}}H\mathrm{t}}\overline{\rho }(\text{t})e^{\frac{i}{\mathrm{}}H\mathrm{t}}\right]`$ (see Eq.(32)) will involve only $`\overline{\rho }_{\nu \nu }`$-s. The same result can be obtained in a formal manner from a generalized Heisenberg representation for Eq.(69), in which the observables evolve in time according to $$\dot{O}_\mathrm{\Delta }(\text{t})=\left[\sigma ^{eq}(\beta )\left(\text{G}\left(\frac{\beta }{2}\left(\text{H}\stackrel{~}{\text{H}}\right)\right)+\text{I}\right)\frac{i}{\mathrm{}}\left(\text{H}\stackrel{~}{\text{H}}\right)\right]O_\mathrm{\Delta }(\text{t})$$ (82) $$O_\mathrm{\Delta }(\text{t})=\mathrm{exp}\left\{\left[\sigma ^{eq}(\beta )\left(\text{G}\left(\frac{\beta }{2}\left(\text{H}\stackrel{~}{\text{H}}\right)\right)+\text{I}\right)\frac{i}{\mathrm{}}\left(\text{H}\stackrel{~}{\text{H}}\right)\right]\text{t}\right\}O_\mathrm{\Delta }(0).$$ (83) Here $`\left(\text{H}\stackrel{~}{\text{H}}\right)O=[H,O]`$ and the lower label $`\mathrm{\Delta }`$ reminds that all averages are to be calculated with $`\mathrm{\Delta }\rho (0)=\rho (0)\rho ^{eq}`$. From Eq.(IV) above it is immediate that $`[H,O]=\left(\text{H}\stackrel{~}{\text{H}}\right)O=0`$ yields $$\dot{O}_\mathrm{\Delta }(\text{t})=\sigma ^{eq}O_\mathrm{\Delta }(\text{t}),$$ (85) $$O_\mathrm{\Delta }(\text{t})=\mathrm{exp}\left[\sigma ^{eq}\text{t}\right]O_\mathrm{\Delta }(0),$$ (86) in agreement with the observation above. An unexpected outcome of this result is that the average of an observable which commutes with the hamiltonian is conserved throughout the evolution, provided the initial average value is identical to the equilibrium average. In other words, $`O_{Delta}(0)`$ implies $`O(\text{t})=O^{eq}`$. The conservation of energy, Eq.(80), is seen to be in fact just a particular realization of this feature. Furthermore, for operators satisfying commutation relations of the form $`[H,A]=\epsilon A`$, Eqs.(IV) lead to $$\dot{A}_\mathrm{\Delta }\left(t\right)=\left[\sigma ^{eq}\left(\beta \right)\left(G\left(\frac{\beta \epsilon }{2}\right)+1\right)\frac{i}{\mathrm{}}\epsilon \right]A_\mathrm{\Delta }\left(t\right),$$ (88) $$A_\mathrm{\Delta }\left(t\right)=\mathrm{exp}\left[\left[\sigma ^{eq}\left(\beta \right)\left(G\left(\frac{\beta \epsilon }{2}\right)+1\right)\frac{i}{\mathrm{}}\epsilon \right]t\right]A_\mathrm{\Delta }\left(0\right),$$ (89) where $`\text{G}(x)+1=x\mathrm{coth}(x)`$. Eqs.(IV) and (IV) allow us to provide a handful of instant examples. 1) A two-level atom, with the hamiltonian $`H=\text{E}_1\|11\|+\text{E}_2\|22\|`$ and the occupation numbers $`n_1=1\|\rho \|1,n_2=2\|\rho \|2,n_1+n_2=1,`$ obeys a simple relaxation law which follows from Eq.(IVa): $$\dot{n}_1=\sigma ^{eq}(\beta )\left(n_1n_1^{eq}(\beta )\right)$$ (91) $$\dot{n}_2=\sigma ^{eq}(\beta )\left(n_2n_2^{eq}(\beta )\right)$$ (92) If Eqs.(IV) are rearranged into the kinetic-like form $$\dot{\text{n}}_1=\text{k}_{12}\text{n}_1+\text{k}_{21}\text{n}_2,$$ (94) $$\dot{\text{n}}_2=\text{k}_{12}\text{n}_1\text{k}_{21}\text{n}_2,$$ (95) the corresponding (thermal) transition rates $`\text{k}_{12}=\sigma ^{eq}(\beta )\text{n}_2^{eq}(\beta ),\text{k}_{21}=\sigma ^{eq}(\beta )\text{n}_1^{eq}(\beta )`$, are seen to have, up to the factor of $`\sigma ^{eq}`$, an Arrhenius-like dependence on the temperature. 2) For a harmonic oscillator of unit mass and frequency $`\omega `$, described by $`H={\displaystyle \frac{p^2}{2}}+{\displaystyle \frac{\omega ^2q^2}{2}},`$ $`p_{eq}=0,q_{eq}=0,`$ one can apply Eq.(IVa) to the annihilation and creation operators, $`a=\sqrt{{\displaystyle \frac{\omega }{2\mathrm{}}}}\left(q+i{\displaystyle \frac{p}{\omega }}\right),a^+=\sqrt{{\displaystyle \frac{\omega }{2\mathrm{}}}}\left(qi{\displaystyle \frac{p}{\omega }}\right),`$ to recover a coupled system of equations for the average momentum and the average coordinate, $$\dot{p}=\gamma (\omega ,\beta )p\omega ^2q,$$ (97) $$\dot{q}=p\gamma (\omega ,\beta )q,$$ (98) where $`\gamma (\omega ,\beta )=\sigma ^{eq}\left(\beta \right)\left[1+G\left(\beta \mathrm{}\omega /2\right)\right]`$. We recognize a typical damped motion, driven by the classical Langevin equation $$\ddot{q}+2\gamma (\omega ,\beta )\dot{q}+\left[\omega ^2+\gamma ^2(\omega ,\beta )\right]q=0,$$ (99) which is obtained by elimination of the momentum variables from Eqs.(IV). 3) For the non-relativistic free particle hamiltonian $`H={\displaystyle \frac{p^2}{2m}}`$ Eq.(IVa) gives the relaxation law $$\dot{p}=\sigma ^{eq}\left[pp^{eq}\right],$$ (100) which shows a (thermal) friction force linear in momentum. In case the initial momentum average coincides with the final thermal average, one obviously obtains conservation of the average momentum. More details can be extracted from the Wigner function $$f_\mathrm{\Delta }(\stackrel{}{\mathrm{p}},\stackrel{}{\mathrm{r}},\mathrm{t})=\underset{\stackrel{}{\mathrm{q}}}{}e^{\frac{i}{\mathrm{}}\stackrel{}{\mathrm{q}}\stackrel{}{\mathrm{r}}}\stackrel{}{\mathrm{p}}\frac{\stackrel{}{\mathrm{q}}}{2}\left\|\rho _\mathrm{\Delta }(\mathrm{t})\right\|\stackrel{}{\mathrm{p}}+\frac{\stackrel{}{\mathrm{q}}}{2}.$$ (101) Differentiation of Eq.(101) on time and use of Eqs.(71), (73), (74) yields $$\dot{f}_\mathrm{\Delta }(\stackrel{}{p},\stackrel{}{r},t)=\underset{\stackrel{}{q}}{}e^{\frac{i}{\mathrm{}}\stackrel{}{q}\stackrel{}{r}}\left[\sigma ^{eq}\left(\beta \right)\left(G\left(\beta \frac{\stackrel{}{p}\stackrel{}{q}}{m}\right)+1\right)+\frac{i}{\mathrm{}}\frac{\stackrel{}{p}\stackrel{}{q}}{m}\right]\stackrel{}{p}\frac{\stackrel{}{q}}{2}\left\|\rho _\mathrm{\Delta }\left(t\right)\right\|\stackrel{}{p}+\frac{\stackrel{}{q}}{2}.$$ (102) But note that $$\frac{\stackrel{}{p}}{m}_\stackrel{}{r}f_\mathrm{\Delta }(\stackrel{}{p},\stackrel{}{r},t)=\underset{\stackrel{}{q}}{}e^{\frac{i}{\mathrm{}}\stackrel{}{q}\stackrel{}{r}}\left(\frac{i}{\mathrm{}}\frac{\stackrel{}{p}\stackrel{}{q}}{m}\right)\stackrel{}{p}\frac{\stackrel{}{q}}{2}\left\|\rho _\mathrm{\Delta }\left(t\right)\right\|\stackrel{}{p}+\frac{\stackrel{}{q}}{2}$$ (103) and rewrite the right hand side of Eq.(102) in differential form to obtain $$\dot{f}_\mathrm{\Delta }+\frac{\stackrel{}{p}}{m}_\stackrel{}{r}f_\mathrm{\Delta }=\sigma ^{eq}\left(\beta \right)\left(G\left(i\frac{\mathrm{}\beta }{m}\stackrel{}{p}_\stackrel{}{r}\right)+1\right)f_\mathrm{\Delta }.$$ (104) The operatorial expression on the right hand side is to be understood in terms of the power expansion $`G\left(x\right)+1=x\mathrm{coth}\left(x\right)=1+2\underset{n=1}{\overset{\mathrm{}}{}}\left(1\right)^{n1}\zeta _R\left(2n\right)(x/\pi )^{2n}`$ , where $`\zeta _R(\mathrm{s})=\underset{\mathrm{k}=1}{\overset{\mathrm{}}{}}\mathrm{k}^\mathrm{s}`$ is the Riemann zeta function. Hence Eq.(104) reads, in explicit form, $$\dot{f}_\mathrm{\Delta }+\frac{\stackrel{}{p}}{m}_\stackrel{}{r}f_\mathrm{\Delta }=\sigma ^{eq}\left(\beta \right)\left[1+2\underset{n=1}{\overset{\mathrm{}}{}}\zeta _R\left(2n\right)\left(\frac{\mathrm{}\beta }{\pi m}\stackrel{}{p}_\stackrel{}{r}\right)^{2n}\right]f_\mathrm{\Delta }$$ (105) and proves to be a Burnett-type (or generalized Fokker-Plank) equation. Note further that the dependence of $`f_\mathrm{\Delta }`$ on momentum and coordinate variables can be separated in Eq.(105), and one can integrate over momentum to obtain an exact equation in coordinate space. It can be safely assumed also that the momentum distribution does not deviate significantly from equilibrium, such that one can write $`f_\mathrm{\Delta }(\stackrel{}{p},\stackrel{}{r},t)n_\mathrm{\Delta }(\stackrel{}{r},t)f^{eq}\left(\left\|\stackrel{}{p}\right\|\right)`$, where $`n_\mathrm{\Delta }(\stackrel{}{r},t)=n(\stackrel{}{r},t)n^{eq}(\stackrel{}{r},t)`$ is the deviation from the equilibrium value of the localization probability and $`f^{eq}\left(\left\|\stackrel{}{p}\right\|\right)`$ is the equilibrium momentum distribution. In that case, in the high-temperature limit when only contributions to leading order in $`\beta `$ survive, integration over momentum leads apparently to a diffusion-like equation, $$\dot{\mathrm{n}}_\mathrm{\Delta }=\mathrm{D}\left(\beta \right)\mathrm{\Delta }\mathrm{n}_\mathrm{\Delta }\sigma ^{eq}(\beta )\mathrm{n}_\mathrm{\Delta },$$ (106) with the diffusion coefficient $$D\left(\beta \right)=\sigma ^{eq}\left(\beta \right)\frac{\mathrm{}^2\beta }{3m},$$ (107) where it has been taken into account that $`\zeta _R(2)=\pi ^2/6`$. But let us recall that $`\lambda _T=\sqrt{(\mathrm{}^2\beta )/(3m)}`$ is just the de Broglie wavelength corresponding to the root-mean-square momentum $`\sqrt{\stackrel{}{p}^2^{eq}}`$, such that in fact $`\text{D}(\beta )=\sigma ^{eq}(\beta )\left(\lambda _T\right)^2`$. It follows necessarily that the diffusion term in Eq.(106) can give significant contributions only if the localization probability varies substantially on the scale of the thermal de Broglie wavelength $`\lambda `$<sub>T</sub>, regardless of the specific value of $`\sigma ^{eq}`$($`\beta `$). But since states with such variations do not belong to the high-temperature, near-equilibrium regime, we are forced to recognize that Eq.(106) actually reduces to $$\dot{n}_\mathrm{\Delta }=\sigma ^{eq}\left(\beta \right)n_\mathrm{\Delta }.$$ (108) The linearization procedure developed in this Section can be extended without significant modifications to equilibrium states other than the thermal canonical distribution. It can be shown that the relaxation laws for the elements of the density matrix in the diagonal representation of the hamiltonian are similar to those found here for the canonical case. A detailed account of this issue will be given elsewhere. ## V SYMMETRY INVARIANCE, CONSERVATION LAWS AND SEPARABILITY It has already been pointed out in Sec.3 that Eq.(30) is invariant under any time-independent unitary transformations that leave the hamiltonian unchanged. It is also obviously invariant against time translations, albeit this operation can no longer be associated with a unitary transformation. The same is not true, in this form, of time-dependent transformations relating different observers in relative motion. But at least in the non-relativistic case, this deficiency can be easily corrected so that invariance under the complete dynamical group of the system is recovered. Indeed, let us rewrite Eq.(30) in the form $$\dot{\rho }=\sigma \left[\rho ln\rho +\{D\left(\rho \right),\rho \}\rho Tr\left(\rho ln\rho \right)\right]+\frac{i}{\mathrm{}}[\rho ,H],$$ (109) where $`D(\rho )`$ replaces $`\zeta (\rho ,H\text{E})(H\text{E})`$, and let us consider the invariance conditions for Eq.(109) under a time-dependent unitary transformation $`U(\text{t})`$, $`U\left(t\right)U^+\left(t\right)=U^+\left(t\right)U\left(t\right)=I`$. As usual, the density matrix becomes $`\rho ^{}\left(t\right)=U\left(t\right)\rho \left(t\right)U^+\left(t\right)`$, hence $$\dot{\rho }^{}(\text{t})=U(\text{t})\dot{\rho }(\text{t})U^+(\text{t})[\rho ^{}(\text{t}),\dot{U}(\text{t})U^+(\text{t})],$$ (110) while multiplication of Eq.(109) by $`U(\text{t})`$ on the left and $`U^+(\text{t})`$ on the right, followed by use of Eq.(110) gives $`\dot{\rho }^{}=\sigma \left[\rho ^{}\mathrm{ln}\rho ^{}+\{U(\text{t})D(\rho )U^+(\text{t}),\rho ^{}\}\rho ^{}Tr\left(\rho ^{}\mathrm{ln}\rho ^{}\right)\right]+`$ (111) $`+{\displaystyle \frac{i}{\mathrm{}}}[\rho ^{},U(\text{t})HU^+(\text{t})+i\mathrm{}\dot{U}(\text{t})U^+(\text{t})].`$ (112) It is easily seen that Eq.(111) will regain the form of Eq.(109) provided H is invariant under $`U(\text{t})`$ in the customary sense, $$U(\text{t})HU^+(\text{t})+i\mathrm{}\dot{U}(\text{t})U^+(\text{t})=H$$ (113) and if, in addition, $$U(\text{t})D(\rho )U^+(\text{t})=D(\rho ^{}),$$ (115) $$\sigma (\rho ,H\text{E})=\sigma (\rho ^{},H\text{E}^{}),$$ (116) In the absence of any evidence to the contrary, the functional $`\sigma `$ will be assumed in the following to have all necessary invariance properties. From Eq.(113) it follows in the customary way that if $`U(\text{t})`$ spans a Lie group of order n, such that $`U(\text{t})=\mathrm{exp}\left[(i/\mathrm{})\lambda ^jK_j(\text{t})\right]`$, with $`\lambda ^j,j=1,\mathrm{\hspace{0.33em}2},\mathrm{}n`$ the group parameters and summation over repeated indices understood, then the corresponding infinitesimal, hermitian generators $`K_j(\text{t}),j=1,\mathrm{\hspace{0.33em}2},\mathrm{}n`$ satisfy the familiar commutation relations $$[K_j\left(t\right),H]\frac{}{t}K_j\left(t\right)=0.$$ (117) Note that a conservation law is not yet implied. But let us assume further that the transformations $`U(\text{t})`$ are such that $$\dot{U}(\text{t})U^+(\text{t})=\frac{i}{\mathrm{}}\left(a^jC_j(\text{t})+b\right),$$ (119) $$U(\text{t})C_j(\text{t})U^+(\text{t})=c_j^lC_l(\text{t})+f_j,$$ (120) where all parameters $`a^j,b,c_j^l,f_j`$ are real functions of the group parameters $`\lambda ^j`$ and time, and the $`C_j`$-s are hermitian operators (observables). In that case, if the conservation of energy is to be invariant under all transformations $`U(\text{t})`$, it follows from the expression of the transformed average energy $$\text{E}^{}=Tr\left(H\rho ^{}\left(t\right)\right)=\text{E}+i\mathrm{}Tr\left(\dot{U}\left(t\right)U^+\left(t\right)\rho ^{}\left(t\right)\right)=\text{E}a^jC_j^{}b,$$ (121) that a conservation law is required for each $`C_j`$. Unfortunately, Eq.(30) does not account for such supplementary constants of motion and simple algebra reveals that $`D(\rho )=\zeta (\rho ,H\text{E})(H\text{E})`$ does not satisfy the first of Eqs.(V), despite an invariant hamiltonian, since $`U\left(t\right)\left[\zeta (\rho ,H\text{E})\left(H\text{E}\right)\right]U^+\left(t\right)=`$ $`=\zeta (\rho ^{},Hi\mathrm{}\dot{U}\left(t\right)U^+\left(t\right)\text{E})\left(Hi\mathrm{}\dot{U}\left(t\right)U^+\left(t\right)\text{E}\right)`$ $$\zeta (\rho ^{},H\text{E}^{})\left(H\text{E}^{}\right).$$ (122) Let us examine now whether modifying Eq.(109) to include conservation of the quantities C<sub>j</sub> brings about the desired invariance under the transformations of the given Lie group. Let the conservation of each C<sub>j</sub> be added to the set of constraints accounted for in the original variational principle, such that Eq.(15) is brought to the form $`\delta [(\dot{\gamma }\|\mathrm{ln}(\gamma \gamma ^+)\|\gamma )+(\gamma \|\mathrm{ln}(\gamma \gamma ^+)\|\dot{\gamma })+2\zeta (\dot{\gamma }\|𝐇\|\gamma )+2\zeta ^{}(\gamma \|𝐇\|\dot{\gamma })+`$ (123) (124) $`+\xi ((\dot{\gamma }\|\gamma )+(\gamma \|\dot{\gamma }))+2\eta ^j((\dot{\gamma }\|𝐂_j\|\gamma )+(\gamma \|𝐂_j\|\dot{\gamma }))+{\displaystyle \frac{2}{\sigma }}(\dot{\gamma }\|\dot{\gamma })]=0,`$ (125) with the new parameters $`\eta ^j`$ assumed real, since the corresponding terms will not contribute to the hamiltonian part of the equation of motion. Taking again the variation with respect to $`\dot{\gamma },\dot{\gamma }^+`$ yields $$\|\dot{\gamma })=\sigma [\frac{1}{2}\left[\mathrm{ln}\left(\gamma \gamma ^+\right)\right]\|\gamma )+\zeta \text{H}\|\gamma )+\eta ^j\text{C}_j\|\gamma )+\frac{\xi }{2}\|\gamma )]$$ (126) and the corresponding equation of motion for the density matrix, $$\dot{\rho }=\sigma \left[\rho \mathrm{ln}\rho +\{\zeta (H\text{E})+\eta ^j\left(C_jC_j\right),\rho \}\rho Tr(\rho \mathrm{ln}\rho )\right]+\frac{i}{\mathrm{}}[\rho ,H].$$ (127) Here $`C_j=Tr\left[C_j\rho \right]`$ is the conserved average of $`C_j`$, $`\zeta `$ and $`\eta ^j`$ are solutions of $$Tr\left[(H\text{E})\rho \mathrm{ln}\rho \right]+2\zeta Tr\left[\left(H\text{E}\right)^2\rho \right]+\eta ^jTr\left[\{H\text{E},C_jC_j\}\rho \right]=0,$$ (129) $`Tr\left[\left(C_jC_j\right)\rho \mathrm{ln}\rho \right](i/\mathrm{}\sigma )Tr\left[[C_jC_j,H\text{E}]\rho \right]+`$ $$+\zeta Tr\left[\{C_jC_j,H\text{E}\}\rho \right]+\eta ^lTr\left[\{C_jC_j,C_lC_l\}\rho \right]=0,$$ (130) $`j=1,\mathrm{\hspace{0.33em}2},\mathrm{}n`$ and we have identified $`\sigma (Im\zeta )=(1/\mathrm{})`$, $`Re\zeta \zeta `$, $`\xi =\left[\mathrm{Tr}\left(\rho \mathrm{ln}\rho \right)+2\left(Re\zeta \right)\mathrm{E}+2\eta ^jC_j\right]`$. Eqs.(V) always have solution, as the matrix of coefficients for the unknowns $`\zeta `$, $`\eta ^j`$ is recognized to be the positively defined covariance matrix for the hamiltonian and the operators $`C_j`$. If we presume the invariance of the hamiltonian as defined by Eq.(113), in accordance with the discussion above, it is now straightforward to verify that $`\overline{D}=\zeta (H\text{E})+\eta ^j\left(C_jC_j\right)`$ is invariant as well in the sense of Eq.(Va), provided $`\zeta `$, $`\eta ^j`$ change as $$\zeta ^{}=\zeta ,$$ (132) $$\eta ^j=\eta ^lc_l^j+\zeta a^j.$$ (133) In deriving Eqs.(V) use has been made of Eqs.(V) and the following transformation of $`C_j`$ under the action of $`U(\text{t})`$: $$C_jTr\left[C_j\rho \right]=Tr\left[U\left(t\right)C_jU^+\left(t\right)\rho ^{}\right]=c_j^lC_j^{}+f_j.$$ (134) The complete invariance of Eq.(127) requires, of course, that $`\zeta ^{},\stackrel{}{\eta }^{}`$ defined in Eqs.(V) be solutions of the transformed Eqs.(V), obtained upon substituting $`\rho ^{},E^{}and\stackrel{}{P}^{}`$ for $`\rho ,Eand\stackrel{}{P}`$, respectively. But substitution of $`\rho \left(\mathrm{t}\right)=U\left(\mathrm{t}\right)\rho ^{}\left(\mathrm{t}\right)U^+\left(\mathrm{t}\right)`$, followed by rearrangement of U, $`U^+`$ over observables and use of the relations $$U\left(t\right)HU^+\left(t\right)E=HE^{}+a^j\left(C_jC_j^{}\right),$$ (136) $$U\left(t\right)C_jU^+\left(t\right)C_j=c_j^l\left(C_lC_l\right),$$ (137) $`j=1,\mathrm{\hspace{0.33em}2},\mathrm{}n`$ obtained from Eqs.(113), (V), (121) and (134), leads to $`Tr\left[\left(HE^{}\right)\rho ^{}ln\rho ^{}\right]+2\zeta Tr\left[\left(HE^{}\right)^2\rho ^{}\right]+\left(\eta ^lc_l^j+\zeta a^j\right)Tr\left[\{HE^{},C_jC_j^{}\}\rho ^{}\right]`$ $$+(i/\mathrm{}\sigma )a^jTr\left[[C_jC_j^{},HE^{}]\rho ^{}\right]=0,$$ (139) $`Tr\left[\left(C_jC_j^{}\right)\rho ^{}ln\rho ^{}\right](i/\mathrm{}\sigma )Tr\left[[C_jC_j^{},HE^{}]\rho ^{}\right]+`$ $`+\zeta Tr\left[\{C_jC_j^{},HE^{}\}\rho ^{}\right]+\left(\eta ^mc_m^l+\zeta a^l\right)Tr\left[\{C_jC_j^{},C_lC_l^{}\}\rho ^{}\right]`$ $$(i/\mathrm{}\sigma )a^lTr\left[[C_jC_j^{},C_lC_l^{}]\rho ^{}\right]=0,$$ (140) $`j=1,\mathrm{\hspace{0.33em}2},\mathrm{}n.`$ The first of these equations displays the required invariance only if the last term vanishes identically, which demands $$[C_j\left(t\right),H]=0,$$ (141) for $`j=1,\mathrm{\hspace{0.33em}2},\mathrm{}n`$, while the second equation is seen to be invariant provided $$[C_j(\mathrm{t}),C_l(\mathrm{t})]=0,$$ (142) for $`j=1,\mathrm{\hspace{0.33em}2},\mathrm{}n`$, $`l=1,\mathrm{\hspace{0.33em}2},\mathrm{}n`$. We conclude that Eq.(91) is invariant if and only if Eqs.(113), (141) and (142) are simultaneously verified, in which case the parameters $`\zeta `$, $`\eta ^j`$ transform according to Eqs.(V). The generating variational principle, Eq.(123), is invariant, of course, under the same conditions. Let us substitute now for $`U(\text{t})`$ the special Galilei boost of velocity $`\stackrel{}{v}_0`$, $$U(t;\stackrel{}{v}_0)=\mathrm{exp}\left[\frac{i}{\mathrm{}}\left(\stackrel{}{P}tm\stackrel{}{X}\right)\stackrel{}{v}_0\right],$$ (143) where m is the total mass of the system, $`\stackrel{}{X}`$ is the position of the center of mass and $`\stackrel{}{P}`$ denotes the total momentum. Expression (143) obviously prompts the identifications $`C_j=P_j`$, $`a^j=1`$, $`b=(m\stackrel{}{v}_0)/2`$, $`c_j^l=\delta _{jl}`$, $`f_j=m\stackrel{}{v}_0`$, which introduced in Eqs.(141) and (142) lead to the recognizable commutation relations $$[H,P_j]=[P_j,P_l]=0.$$ (144) Subsequent substitution in Eq.(91) gives the corresponding equation of motion in the form $$\dot{\rho }=\sigma \left[\rho \mathrm{ln}\rho +\{\zeta \left(HE\right)+\eta ^j\left(P_jP_j\right),\rho \}\rho Tr\left(\rho \mathrm{ln}\rho \right)\right]+\frac{i}{\mathrm{}}[\rho ,H].$$ (145) We recover thus in an unexpected manner the celebrated result that the Galilei-invariance of the appropriate non-relativistic equation of motion is equivalent to the corresponding invariance of the hamiltonian, the conservation of total momentum and the commutation of the hamiltonian and the total momentum operators. Eq.(145) reduces to the original Eq.(30) in the center-of-mass referential, where only states corresponding to an eigenstate of zero total momentum for the center-of-mass coordinates need be considered and the dissipative momentum terms vanish. It also retains the fundamental features previously outlined for Eq.(30). In particular, it can be checked that pure states evolve according to the usual hamiltonian dynamics, the entropy of mixed states increases and the nature of the asymptotic equilibrium states is preserved, up to a slight change of form which accounts for the conservation of momentum. It is also evident that Eq.(145) is invariant under time-independent symmetry transformations which leave the hamiltonian and the dissipator $`\overline{D}`$ invariant, provided the time-scale parameter $`\sigma `$ has the same property. In particular, if the hamiltonian commutes with the total angular momentum, Eq.(145) is invariant under finite rotations. However, as for the linear momentum, rotational invariance alone does not imply, in general, a conservation law for the angular momentum. The latter can be brought into view by requiring that the equation of motion for the density matrix be covariant with respect to all reference frames where the conservation of energy is a valid physical law. In particular we should consider translations to observers in uniform rotational motion around an axis at rest in some inertial frame. The rather cumbersome details of adding this supplementary constraint will be left aside, since nothing new will be gained for the formalism. A more interesting lack of symmetry for Eq.(145), or better, the simpler Eq.(30), lies concealed in the apparent absence of separability. Indeed, let the system described by the hamiltonian H be composed of two noninteracting subsystems, such that $`H=H_1+H_2`$, $`[H_1,H_2]=0`$, and consider the situation of a separable initial state $`\rho (0)=\rho _1(0)\rho _2(0)`$, of energy $`\text{E}=\text{E}_1+\text{E}_2`$. Direct inspection of Eq.(30) shows that the energies of the two subsystems cannot be separately conserved and a completely separable solution is thus prohibited. But we also observe that lifting the constraint of separate conservation of energy allows a pseudo-separable solution $`\rho (\text{t})=\rho _1(\text{t})\rho _2(\text{t})`$ given by the coupled system $$\dot{\rho }_1=\sigma \left[\rho _1\mathrm{ln}\rho _1+\zeta \{H_1\frac{Tr\left(H_1\rho _1\right)}{Tr\left(\rho _1\right)},\rho _1\}\rho _1\frac{Tr\left(\rho _1\mathrm{ln}\rho _1\right)}{Tr\left(\rho _1\right)}\right]+\frac{i}{\mathrm{}}[\rho _1,H_1],$$ (147) $$\dot{\rho }_2=\sigma \left[\rho _2\mathrm{ln}\rho _2+\zeta \{H_2\frac{Tr\left(H_2\rho _2\right)}{Tr\left(\rho _2\right)},\rho _2\}\rho _2\frac{Tr\left(\rho _2\mathrm{ln}\rho _2\right)}{Tr\left(\rho _2\right)}\right]+\frac{i}{\mathrm{}}[\rho _2,H_2].$$ (148) In this case probability is independently conserved for each subsystem, since $`Tr(\dot{\rho }_i)=0`$, while energy is only conserved globally, $`{\displaystyle \frac{Tr\left(H_1\rho _1\right)}{Tr\left(\rho _1\right)}}+{\displaystyle \frac{Tr\left(H_2\rho _2\right)}{Tr\left(\rho _2\right)}}=\text{E}.`$ The coupling between the (noninteracting) subsystems appears to be as instantaneous and nonlocal as usual quantum entanglement, but unlike the latter, it involves an unorthodox exchange of energy. The significance of this unusual outcome follows from the observation that, according to Eqs.(V), the equilibrium of the compound system is attained for values of $`\sigma `$ and $`\zeta `$ common to both subsystems, hence for a common generalized temperature. Imagine now that the initial states for the two subsystems are chosen as individual equilibrium states with different corresponding temperatures. It follows that the dynamics given by Eq.(30) will drive the total system towards a new state of equilibrium, with a temperature common to both components. We cannot but concede the obvious similarity of this unconventional effect with the classical process of equilibration by thermal contact. Its origin lies in the very assumption of maximal entropy increase on which Eq.(30) has been derived. Indeed, even when the entropy of each subsystem is already maximal under individual isolation, if states of larger total entropy are available, probabilities and energy (heat) will be necessarily redistributed so as to enforce a further increase of the overall entropy. Whether this entropic entanglement, or ideal thermal contact, is or not an element of reality appears equivalent to accepting or rejecting the conjecture that an isolated, perfectly ideal gas can undergo relaxation towards equilibrium. We can provide formal support towards the positive by pointing out that the effect of entropic entanglement does not necessarily interfere with the concept of separable evolution for mutually isolated systems. First let us note that explicitly specifying an adiabatic separation (in the thermodynamic sense) of the noninteracting systems, and hence allowing for separate conservation of energy, removes most of the entropic entanglement. In this case the resulting equation of motion will display distinct $`\zeta `$-s for each of the systems, but a common time-scale parameter, i.e. $$\dot{\rho }=\sigma \left[\rho \mathrm{ln}\rho +\zeta _1\{H_1E_1,\rho \}+\zeta _2\{H_2E_2,\rho \}\rho Tr\left(\rho \mathrm{ln}\rho \right)\right]+\frac{i}{\mathrm{}}[\rho ,H_1+H_2],$$ (149) with $`E_i={\displaystyle \frac{Tr(H_i\rho )}{Tr(\rho )}}=const.,`$ for $`i=1,2`$. As before, it proves possible to extract a pseudo-separable solution $`\rho (t)=\rho _1(t)\rho _2(t)`$, but Eqs.(V) are replaced by $$\dot{\rho }_1=\sigma \left[\rho _1\mathrm{ln}\rho _1+\zeta _1\{H_1E_1,\rho _1\}\rho _1\frac{Tr\left(\rho _1\mathrm{ln}\rho _1\right)}{Tr\left(\rho _1\right)}\right]+\frac{i}{\mathrm{}}[\rho _1,H_1],$$ (151) $$\dot{\rho }_2=\sigma \left[\rho _2\mathrm{ln}\rho _2+\zeta _2\{H_2E_2,\rho _2\}\rho _2\frac{Tr\left(\rho _2\mathrm{ln}\rho _2\right)}{Tr\left(\rho _2\right)}\right]+\frac{i}{\mathrm{}}[\rho _2,H_2].$$ (152) where this time the $`\zeta _i`$ parameters, $`i=1,2`$, will be found to depend only on the corresponding $`\rho _i`$ and $`H_i`$, in exactly the manner obtained for a single isolated system. Yet the two evolutions remain tethered by the time-scale parameter $`\sigma `$, thus retaining a weaker form of entropic entanglement. The simple presence of other noninteracting, adiabatically separated systems appears to alter the time-scale of dissipative relaxation for any given system. If $`\sigma `$ is assumed variable in time, e.g. through a dependence on $`\rho `$, this influence will be time-dependent unless all other systems have reached equilibrium. But since $`\sigma `$ does not affect the nature of the asymptotic equilibrium state, the equilibrium of any one system will not be disturbed by other systems and will display an individual temperature determined solely by the corresponding energy content. A careful examination will trace the above type of nonseparability to the fact that the corresponding variational principle selects the direction of maximum entropy increase by referring to the time derivative of the total (entangled) state operator, and not to disentangled, individual state operators separately. However, this pitfall can be avoided if it is recognized that true mutual isolation precludes entanglement on invariance grounds. Indeed, regardless of the nature of the underlying dynamics, the evolution of two mutually isolated systems should remain invariant under every transformation pertaining to the individual symmetry groups. In particular, it should be invariant under individual time translations. Since entangled states certainly do not possess this invariance, they do not describe truly isolated systems. In other words, the restricted subspace of the state space that can be spanned by the dynamics of mutually isolated systems should contain only non-entangled states and the evolution of each of the factor states should be driven independently. In our nonlinear setting, where this subspace is selected by means of the generating variational principle, this restriction has to be correctly built in the variational functional itself. Hence one has to account both for individual conservation laws, excluding thus any energy exchange, as well as for vanishing entanglement. The latter imposes a separable state operator $`\gamma (t)=\gamma _1(t)\gamma _2(t)`$ and also requires that the entropy production be maximized separately with respect to variations of $`\dot{\gamma }_1`$ and $`\dot{\gamma }_2`$, i.e. the $`\sigma `$ term in the variational principle should be replaced according to $`{\displaystyle \frac{2}{\sigma }}(\dot{\gamma }\|\dot{\gamma }){\displaystyle \frac{2}{\sigma _1}}(\dot{\gamma }_1\gamma _2\|\dot{\gamma }_1\gamma _2)+{\displaystyle \frac{2}{\sigma _2}}(\gamma _1\dot{\gamma }_2\|\gamma _1\dot{\gamma }_2)`$ with each $`\sigma _i`$ a functional only of $`\gamma _i`$ and $`H_i`$. But then the variational principle takes the form $$\delta \{\left(\gamma _2\|\gamma _2\right)F_1+\left(\gamma _1\|\gamma _1\right)F_2\}=\mathrm{\hspace{0.33em}0},$$ (154) $`F_i=(\dot{\gamma }_i\|\mathrm{ln}(\gamma _i\gamma _i^+)\|\gamma _i)+(\gamma _i\|\mathrm{ln}(\gamma _i\gamma _i^+)\|\dot{\gamma }_i)+2\zeta _i(\dot{\gamma }_i\|\text{H}_i\|\gamma _i)+2\zeta _i^{}(\gamma _i\|\text{H}_i\|\dot{\gamma }_i)+`$ $$+\left[\overline{\xi }_i(\dot{\gamma }_i\|\gamma _i)+\overline{\xi }_i^{}(\gamma _i\|\dot{\gamma }_i)\right]+\frac{2}{\sigma _i}(\dot{\gamma }_i\|\dot{\gamma }_i),$$ (155) $`i=1,\mathrm{\hspace{0.33em}2},`$ where $`\overline{\xi }_1=\xi _1+\left(\xi _2+\zeta _2\mathrm{E}_2{\displaystyle \frac{\mathrm{S}_2}{\mathrm{k}_\mathrm{B}(\gamma _2\|\gamma _2)}}\right)`$ $`\overline{\xi }_2=\xi _2+\left(\xi _1+\zeta _1\mathrm{E}_1{\displaystyle \frac{\mathrm{S}_1}{\mathrm{k}_\mathrm{B}(\gamma _1\|\gamma _1)}}\right)`$ Independent variation on $`\dot{\gamma }_1`$ and $`\dot{\gamma }_2`$, followed by extraction of the Lagrange parameters from the corresponding conservation conditions leads now to the desired separate equations of motion for $`\rho _i=\gamma _i\gamma _i^+`$, $$\dot{\rho }_i=\sigma _i\left[\rho _i\mathrm{ln}\rho _i+\zeta _i\{H_i\text{E}_i,\rho _i\}\rho _i\frac{Tr(\rho _i\mathrm{ln}\rho _i)}{Tr(\rho _i)}\right]+\frac{i}{\mathrm{}}[\rho _i,H_i],$$ (156) $`i=1,\mathrm{\hspace{0.33em}2},`$ with $`\sigma _i=\sigma _i(\rho _i,H_i)0,`$ $`\zeta _i={\displaystyle \frac{1}{2}}{\displaystyle \frac{Tr[(H_i\mathrm{E}_\mathrm{i})\rho _\mathrm{i}\mathrm{ln}\rho _\mathrm{i}]}{Tr[(H_i\mathrm{E}_\mathrm{i})^2\rho _\mathrm{i}]}},`$ $`\text{E}_i={\displaystyle \frac{Tr(H_i\rho _i)}{Tr(\rho _i)}}=const..`$ Obviously, the invariance of the nonlinear dynamics under the symmetry group of each component subsystem is so restored, provided the $`\sigma _i`$-s are invariant also. ## VI GENERALIZATION TO ARBITRARY ENTROPY AND ENERGY FUNCTIONAL FORMS The framework developed in the previous Secs. can be easily expanded to accommodate non-standard entropy functionals and/or energy forms with a nonlinear dependence on the density matrix $`\rho `$. This generalized formalism can then provide nonlinear extensions for, e.g., the Lie-Poisson dynamics or a standard hamiltonian evolution supplemented by a nonextensive Tsallis entropy , appropriate for systems with fractal properties. We sketch here only the derivation of the generalized equation of motion, since a detailed analysis exceeds the purpose of the present work. To this end, let us start with a Lie-Poisson equation of motion in the form $$\dot{\rho }=\frac{i}{\mathrm{}}[\rho ,\widehat{H}(\rho )]$$ (157) where $`\widehat{H}(\rho )`$ is in general a hermitian, nonlinear functional of $`\rho `$. The energy conservation law is now replaced by $`Tr\left(\widehat{H}(\rho )\dot{\rho }\right)=0`$ or in terms of the state operator $`\gamma `$, $$(\dot{\gamma }\|\widehat{𝐇}(\rho )\|\gamma )+(\gamma \|\widehat{𝐇}(\rho )\|\dot{\gamma })=0$$ (158) The law of probability conservation, on the other hand, remains unchanged since $`Tr(\dot{\rho })=0`$ or $$(\dot{\gamma }\|\gamma )+(\gamma \|\dot{\gamma })=0$$ (159) Let us search now for a nonlinear evolution that observes the above conservation constraints, Eqs.(158) and (159), and is also subject to a second principle based on some unspecified, positive definite entropy functional $`\frac{\text{S}}{\text{k}_\mathrm{B}}=Tr(\widehat{S}(\rho ))`$, such that $`\dot{\text{S}}=Tr\left((\delta \widehat{S}/\delta \rho )\dot{\rho }\right)0`$ or $$(\dot{\gamma }\|\frac{\delta \widehat{𝐒}}{\delta \rho }\|\gamma )+(\gamma \|\frac{\delta \widehat{𝐒}}{\delta \rho }\|\dot{\gamma })0$$ (160) Here the operator $`\widehat{S}(\rho )`$ is assumed hermitian and $`(\delta \widehat{S}/\delta \rho )`$ denotes its hermitian functional derivative with respect to $`\rho `$. The corresponding variational principle is now written $$\delta \left\{(\dot{\gamma }\|\frac{\delta \widehat{𝐒}}{\delta \rho }\|\gamma )(\gamma \|\frac{\delta \widehat{𝐒}}{\delta \rho }\|\dot{\gamma })+2\zeta (\dot{\gamma }\|\widehat{𝐇}(\rho )\|\gamma )+2\zeta ^{}(\gamma \|\widehat{𝐇}(\rho )\|\dot{\gamma })+\xi [(\dot{\gamma }\|\gamma )+(\gamma \|\dot{\gamma })]+\frac{2}{\sigma }(\dot{\gamma }\|\dot{\gamma })\right\}=0.$$ (161) and can be verified to generate the following equation of motion: $$\dot{\rho }=\sigma \left[\frac{\delta \widehat{S}}{\delta \rho }\rho +\zeta \{\widehat{H}(\rho )\widehat{H}(\rho ),\rho \}+\frac{\delta \widehat{S}}{\delta \rho }\rho \right]+\frac{i}{\mathrm{}}[\rho ,\widehat{H}(\rho )],$$ (162) where $`A={\displaystyle \frac{Tr(A\rho )}{Tr(\rho )}},`$ and $`\zeta ={\displaystyle \frac{1}{2}}{\displaystyle \frac{\left(\widehat{H}(\rho )\widehat{H}(\rho )\right)(\delta \widehat{S}/\delta \rho )}{\left(\widehat{H}(\rho )\widehat{H}(\rho )\right)^2}},`$ $`\sigma =\sigma (\rho ,\widehat{H}(\rho )\widehat{H}(\rho ))0.`$ We note that if $`(\delta \widehat{S}/\delta \rho )\rho =0`$ for pure states, $`\rho =\rho ^2`$, then $`\delta \widehat{S}/\delta \rho =0`$, $`\zeta =0`$, and the pure state dynamics reduces to that prescribed by Eq.(157). When the energy functional reduces to the hamiltonian, $`\widehat{H}(\rho )=H`$, and the entropy is given the standard von Neumann expression, such that $`\widehat{S}(\rho )=\rho \mathrm{ln}\rho `$, $`(\delta \widehat{S}/\delta \rho )\rho =\rho \mathrm{ln}\rho \rho `$, we recover the basic Eq.(30). A $`\rho `$-dependent $`\widehat{H}(\rho )`$, complemented by the standard entropy, leads to a nonlinear extension of the Lie-Poisson dynamics, $$\dot{\rho }=\sigma \left[\rho \mathrm{ln}\rho +\zeta \{\widehat{H}(\rho )\widehat{H}(\rho ),\rho \}\frac{Tr(\rho \mathrm{ln}\rho )}{Tr(\rho )}\rho \right]+\frac{i}{\mathrm{}}[\rho ,\widehat{H}(\rho )],$$ (163) with $`\zeta ={\displaystyle \frac{1}{2}}{\displaystyle \frac{Tr\left[\left(\widehat{H}\widehat{H}(\rho )\right)\rho \mathrm{ln}\rho \right]}{Tr\left[\left(\widehat{H}(\rho )\widehat{H}(\rho )\right)^2\rho \right]}}.`$ If $`\widehat{H}(\rho )`$ is reduced to the standard hamiltonian $`H`$, but the entropy is given a Tsallis form, with $`\widehat{S}(\rho )={\displaystyle \frac{\rho \rho ^q}{q1}},`$ $`{\displaystyle \frac{\delta \widehat{S}}{\delta \rho }}\rho ={\displaystyle \frac{\rho q\rho ^q}{q1}},`$ for given real q, the result will be a nonlinear extension of the von Neumann dynamics under Tsallis q-thermostatistics, which reads, after a few elementary manipulations, $$\dot{\rho }=\sigma \left[\frac{q}{q1}\rho ^q+\zeta \{HE,\rho \}\frac{q}{q1}\frac{Tr\left(\rho ^q\right)}{Tr(\rho )}\rho \right]+\frac{i}{\mathrm{}}[\rho ,H],$$ (164) where $`\zeta ={\displaystyle \frac{1}{2}}{\displaystyle \frac{q}{q1}}{\displaystyle \frac{Tr\left[(H\mathrm{E})\rho ^\mathrm{q}\right]}{Tr\left[(H\mathrm{E})^2\rho \right]}}`$ Situations where the standard averages have to be replaced by q-averages can be approached in the same fashion, by appropriately redefining the conserved functionals. ## VII CONCLUSION We have constructed and analyzed a non-relativistic nonlinear extension of the quantum law of evolution, which accounts for the second principle of thermodynamics and is not at odds with the factual linearity of pure state propagation. The theoretical existence of such an extension confirms that the linear and unitary evolution of pure states is not in itself sufficient proof for the general linearity of quantum mechanics . One must conclude that the linear propagation of mixed states also has to be corroborated experimentally, to comparable precision, before a definitive conclusion can be drawn. It is hoped that the formal study developed here provides a meaningful benchmark in this sense. Our main result is Eq.(30), which defines the modified time evolution of the density matrix. The equation of motion was extracted from a variational principle on the space of state operators, rather than the space of density matrices, as a trajectory of maximal entropy production under the constraint of energy and probability conservation, augmented eventually by the requirement of Galilei invariance (see Eq.(145)). Should we drop the requirement of entropy increase, the parameters $`Re\zeta ,\xi `$ vanish and the equation of motion reduces automatically to the common hamiltonian form. The outlined procedure may not be unique, but is encouraging in its consistency. In addition, it applies as well to alternate theories which use nonstandard energy or entropy forms. It is notable that the variational principle has sense only in terms of state operators, whereas the equation of motion can be stated simply in terms of the conventional density matrix. A peculiar and unexpected idea brought forth in our ansatz is that a maximal increase of entropy does not necessarily result in maximal decoherence, to the effect that a pure state of a perfectly isolated system is not allowed to evolve into a mixed state. On the contrary, the proposed quantum equivalent of the second principle of thermodynamics is seen to introduce only a limited degree of decoherence, in the sense that the cardinality of the set of nonzero eigenvalues of the density matrix is preserved. As already mentioned, for the particular case of a pure initial state this leads to the usual unitary evolution. The same property also supports, aside from canonical equilibrium states, a rich class of “negative-temperature” equilibrium states, which bring to mind the notion of thermal coherence. Furthermore, the ideal thermal contact phenomenon discussed in Sec.5 abides by the same rule and, according to Eqs.(V), a system in an initially pure state will remain in a pure state even if it is in contact with, but not necessarily interacting with, other systems. However, in that case the pure state undergoes relaxation according to a dynamics of Gisin type , as seen by taking, e.g., $`\rho _1=\rho _1^2`$, $`\rho _1\mathrm{ln}\rho _1=0`$, in Eq.(Va), $$\dot{\rho _1}=\sigma \zeta \{H_1\frac{Tr\left(H_1\rho _1\right)}{Tr\left(\rho _1\right)},\rho _1\}+\frac{i}{\mathrm{}}[\rho _1,H_1].$$ (165) Depending on the sign of $`\zeta `$, the asymptotic stationary state is an energy eigenstate for the lowest (if $`\zeta >0`$) or for the highest (if $`\zeta <0`$) energy level contributing to the initial state $`\rho _1(0)`$. As detailed in Sec.5, the state of thermal contact is not to be mistaken for a state of mutual isolation, despite the absence of explicit interactions. We find it promising that bending quantum dynamics to account for classical phenomenological irreversibility suggests a rather unified picture of both reversibility and irreversibility, as well as coherence and decoherence, while preserving such fundamental features as symmetry invariance. However, the self-consistency of the theory is limited at this point by the need for an explicit expression for the entropy production, which means that the equation of motion remains determined up to the scale setting functional $`\sigma `$. We leave the resolution of this problem for future consideration, although a definite expression for $`\sigma `$ certainly conditions the consistency of our results. For instance, Eq.(74) for the near-equilibrium damping constants of the density matrix elements between energy eigenstates shows an acceptable dependence on the energy gap between the states, but the wrong temperature dependence ($`\gamma _{\mu \nu }0as\beta 0`$, $`\gamma _{\mu \nu }\mathrm{}as\beta \mathrm{}`$) if $`\sigma `$ is assumed temperature independent. In the least, this observation serves to hint that $`\sigma `$ should behave like $`\beta ^{(2+\delta )},\delta >0`$, in the vicinity of canonical equilibrium, which in turn can be used, of course, as a theoretical benchmark.
warning/0007/quant-ph0007054.html	ar5iv	text	# Quantum Message Disruption: A Two-State Model ## 1 References Meyer, David A. Quantum Strategies. Physical Review Letters 82(5)1052-1055. von Neumann, John. Mathematical foundations of quantum mechanics. translated from the German edition by Robert T. Beyer. Princeton University Press, 1996. Princeton, N.J. 445 p. von Neumann, John, Morgenstern, Oskar . Theory of games and economic behavior. \[3d ed.\], Princeton University Press, 1953 \[c1944\]. Princeton
warning/0007/cond-mat0007349.html	ar5iv	text	# Sliding Phases via Magnetic Fields ## I introduction In 1979, Efetov suggested that it would be possible to extend the low temperature Kosterlitz-Thouless (KT) phase of a two-dimensional superconductor to three dimensions by stacking two-dimensional systems in the presence of a parallel magnetic field. The underlying idea, most simply understood in a particular gauge for the field which we specify below, is that the interlayer Josephson coupling which would ordinarily be relevant even when weak, is now spatially modulated and no longer gives rise to divergences. It turns out that this does not work. As pointed out by Korshunov and Larkin, the modulated Josephson coupling gives rise to a coupling that is not modulated, and although it is of higher scaling dimension than the zero field Josephson coupling, it is still relevant everywhere within the KT phase. However, recent work on cationic lipid-DNA complexes by O’Hern and Lubensky and by Golubovic and Golubovic, and then on XY systems themselves by O’Hern, Lubensky and Toner (OLT) has found a different way of obtaining analogs of KT phases in three dimensions. In this approach, additional derivative couplings leave the phases in the different planes free to rotate globally with respect to each other (hence “sliding phases”) while extending the region of irrelevant vortex fugacity to a range where the interlayer Josephson coupling is now irrelevant. Emery et al., and Vishwanath and Carpentier have applied this insight to quantum problems and obtained an analog of the Luttinger liquid in two dimensions. Our purpose in this note is to point out that one can combine Efetov’s insight with the more recent work and considerably extend the domain of these sliding phases by reducing the dimension of a large class of relevant operators via the action of a parallel magnetic field. This is of considerable interest for the full class of perturbations in such problems can be quite constraining, even though it is reasonable that most of them are not realized with substantial amplitude . We will be especially interested in “sliding Luttinger liquids” or “smectic metals” which have been argued to arise in the cuprate superconductors on account of the stripe instability of a doped antiferromagnet . We should note that there is a close connection between our work and that on the striped phases in high Landau levels even though our point of departure (Efetov’s conjecture) is very different. In the Landau level problem, the field is built in at the first step and is central in giving rise to the striped phase in the first instance, while for us it can be variable in magnitude and give rise to both gapped quantum Hall and gapless smectic behavior and an interesting phase transition between them. Nevertheless, in both cases the field serves to constrain the available set of relevant operators in very similar fashion. We will begin in Section II with a quick account of the “dimensional reduction” of the Josephson coupling produced by a parallel field, the genesis of the sliding phase and its enlargement by the field. Next (Section III) we discuss the application of these ideas to coupled Luttinger liquids and present contrasting phase diagrams for a model studied by Emery et al. In this discussion we also show how the integer quantum Hall states are rediscovered by perturbation theory about a smectic metal if the interactions are not too strong. We close with a brief summary and a discussion of possible experimental implications for the cuprates. ## II Sliding XY Phases in Parallel Fields We begin with a brief summary of the genesis of the sliding phase in a three dimensional stack of layers characterized by and XY order parameter. We largely follow OLT and their notation for ease of comparison. The Hamiltonians of the sliding phase fixed points (the plural is warranted) belong to the family, $$H_S=\frac{1}{2}\underset{nn^{}}{}d^2rK_{nn^{}}\mathbf{}_{}\theta _n(𝐫)\mathbf{}_{}\theta _n^{}(𝐫),$$ (1) where $`K_{nn^{}}=Kf_{nn^{}}`$ with $`f_n=(1+_m\gamma _m)\delta _{n,0}\frac{1}{2}_m\gamma _m(\delta _{n,m}+\delta _{n,m})`$ and and $`_{}\theta _n(𝐫)`$ denotes the in-layer gradient of the XY variable in layer $`n`$. We take $`𝐫(x,y)`$ and set the separation of successive layers along the $`z`$-axis to 1. One can check that $`H_S`$ is invariant under shifts $`\theta _n(𝐫)\theta _n(𝐫)+\psi _n`$ for any choice of $`\psi _n`$. This freedom to globally rotate the angle in one layer relative to another, even in the presence of interlayer couplings in $`H_S`$, is the hallmark of the sliding phase. Note that $`H_S`$ returns to itself under a renormalization group (RG) transformation that “lives” in two dimensions and treats the layer index $`n`$ as an internal or flavor index on the fields $`\theta _n`$. In order to identify functions $`K_{nn^{}}`$ that would govern stable fixed points under this RG, we need to examine the behavior of vortex fugacities and Josephson couplings. The former yield, for a vortex configuration $`\{\sigma _n\}`$ in which a net vorticity $`\sigma _n`$ occurs in layer $`n`$, the scaling dimension $$\mathrm{\Delta }_v[\sigma _n]=\frac{\pi K}{T}\underset{n,n^{}}{}f_{nn^{}}\sigma _n\sigma _n^{},$$ (2) which signals a KT unbinding transition at a temperature $`T_{KT}[\sigma _n]`$, upon exceeding the value 2 appropriate to a two-dimensional RG. The generalized Josephson couplings, $$H_J[s_n]=V_J[s_n]d^2r\mathrm{cos}\left[\underset{p}{}s_p\theta _{n+p}(𝐫)\right],$$ (3) where the $`s_n`$ are integers that satisfy $`_ns_n=0`$, are readily shown to have the scaling dimension, $$\mathrm{\Delta }_J[s_n]=\frac{T}{4\pi K}\underset{n,n^{}}{}s_ns_n^{}f_{nn^{}}^1;$$ (4) where the inverse couplings $$f_p^1=\frac{1}{\pi }_0^\pi 𝑑k\frac{\mathrm{cos}kp}{f(k)}.$$ (5) are defined via the Fourier transform $$f(k)=1+\underset{m}{}\gamma _m(1\mathrm{cos}km)$$ (6) of the scaled couplings $`f_n`$. The Josephson couplings are irrelevant above a decoupling temperature $`T_d[s_n]`$ at which $`\mathrm{\Delta }_J[s_n]=2`$. If $`\mathrm{min}_{\sigma _n}T_{KT}[\sigma _n]>\mathrm{max}_{s_n}T_d[s_n]`$ for some choice of $`K_{nn^{}}`$ then we obtain a sliding phase. In the sliding phase the spin correlations are algebraically long ranged in a given layer and vanish between layers, $$\mathrm{cos}[\theta _n(𝐫)\theta _m(0)]\frac{\delta _{nm}}{r^\eta },$$ (7) where $`\eta =\frac{T}{2\pi K}f_0^1`$. Including a parallel magnetic field: We now consider the inclusion of a magnetic field parallel to the layers, appropriate to instances where the $`\theta _n`$ are phases of a superconducting order parameter; without loss of generality, we take $`𝐁=B\widehat{y}`$. It is convenient to work in the gauge $`A_z(𝐫)=Bx`$. In this gauge, the sliding phase Hamiltonians are of the same form, and the computation of the scaling of the vortex fugacity is unchanged. However, the Josephson couplings are modified by the replacements $$\theta _n\theta _n+2nq_Bx$$ (8) where $`q_B=\frac{eB}{\mathrm{}c}`$ is a characteristic wavevector introduced by the field. The key observation regarding the effect of the field is this: for those Josephson couplings for which $`_pps_p0`$, there is an explicit oscillating term in the argument of the cosine that will render them less relevant. Most straightforwardly, consider treating such a term in perturbation theory. In zero field, we would discover that the term was relevant upon finding divergences in perturbation theory. The inclusion of the field will attenuate these divergences due to the oscillation of the correlation functions of the perturbation. However, the net result is not always to render the perturbation theory convergent. Higher order graphs can involve regions where products of the oscillating couplings nevertheless give rise to operators that do not oscillate. For example, the Josephson coupling for layers at distance $`p`$, $`\mathrm{cos}[\theta _{n+p}(𝐫)\theta _n(𝐫)`$ $`+`$ $`2pq_Bx]`$ (9) $``$ $`e^{i[\theta _{n+p}(𝐫)\theta _n(𝐫)]}e^{i2pq_Bx}+\mathrm{c}.\mathrm{c}.`$ (10) will give rise to $$\mathrm{cos}[\theta _{n+p}(𝐫)2\theta _n(𝐫)+\theta _{np}(𝐫)]$$ (11) which can then produce divergences of its own. Indeed, this particular generation is exactly what invalidates Efetov’s original conjecture, for the operators (11) are relevant everywhere in the KT phase of decoupled XY layers. Nevertheless, the application of the field does effect a “dimensional reduction” in that “charged” operators that have a net $`_pps_p`$ (microscopically these arise from hopping processes that move a net charge up or down the stack) can only affect the result through the generation of “neutral” operators for which $`_pps_p=0`$. At the (unstable) decoupled 2D XY fixed line, the latter have higher dimension and we might expect that this will be true at sliding fixed points as well. While that is not always the case, as will be clear by the following example, it is still the case that knocking out the charged operators improves the stability of the sliding phase—after all, the neutral operators were present anyway! To illustrate this effect, we consider the example used by OLT with first and second neighbor couplings. The coupling function $`f_n`$ has a Fourier transform $$f(k)=1+\gamma _1(1\mathrm{cos}k)+\gamma _2(1\mathrm{cos}2k)$$ (12) that is required to take its minimum value at $`k=k_o`$: $$f(k_o)=\delta ,f^{}(k_o)=0\mathrm{and}f^{\prime \prime }(k_o)=2C.$$ (13) Sliding phases arise when $`\delta `$ and $`k_o`$ are chosen so that the system is close to an incommensurate transverse ordering instability, as has been discussed nicely by Vishwanath and Carpentier . At small $`\delta `$, the asymptotic form, $$f_p^1\frac{\mathrm{cos}(pk_0)e^{p\sqrt{\delta /C}}}{\sqrt{C\delta }}$$ (14) enables easy numerical calculation of the scaling dimensions of the Josephson couplings and thence of the temperatures $`T_d`$. In Fig. 1, we plot the ratio, $`\beta =\mathrm{min}_{\sigma _n}T_{KT}[\sigma _n]/\mathrm{max}_{s_n}T_d`$ where the $`s_n`$ are restricted to the two layer Josephson couplings (10) and the three layer terms that they generate (11). The value of $`k_o`$ where both are greater than 1 support a sliding phase in zero field , while the latter alone determines the sliding phase in a magnetic field. The expansion of the phase is clear. (We have not attempted to include all operators that might be allowed by symmetry. As noted in Ref. in the Luttinger liquid context, higher order operators allow increasingly finer instabilities. We do not know of a proof that all such operators allow or exclude a connected sliding phase, but assume that in a given system a finite set will be important over some reasonable range of length scales. Regardless, the magnetic field will improve matters by knocking out all the charged operators.) ## III Smectic Metals in Transverse Fields In this section we discuss coupled one-dimensional (1D) Luttinger liquids (LL) in the presence of a magnetic field, with the field $`𝐁`$ transverse to the plane in which the 1D chains are placed. It was known that at the decoupled LL fixed points, the transverse inter-chain coupling is always relevant, in one of three channels: single electron hopping, (Cooper) pair hopping, and inter chain $`2k_F`$ back scattering. As a consequence, the decoupled LL phase is always unstable and driven toward the Fermi liquid , superconducting, or charge/spin density wave (CDW/SDW) phases. It was recently pointed out that adding strong interchain forward scattering terms (which are exactly marginal) to the decoupled LL fixed point can drive all these interchain couplings irrelevant. The resulting stable, non-Fermi liquid, smectic metal phase is the quantum analog of the classical sliding phase. Since the single electron and Cooper pair hopping processes involve charge transfer between neighboring chains, the presence of a magnetic field has a similar effect, as before, of increasing the scaling dimensions of the operators corresponding to these processes that can perturb the smectic metal fixed points, and hence increasing the range of stability of the smectic metal phase (for simplicity we will neglect the Zeeman effect of the field in this paper). In the following we present an explicit analysis of this effect. Following Emery et al., three different types of smectic metal fixed points need to be distinguished and analyzed in turn: (i) a spinful smectic metal with a spin gap; (ii) a spinful smectic metal without a spin gap; and (iii) a spinless smectic metal. For simplicity we will only include nearest neighbor interchain couplings and their immediate descendants. Spin-gapped smectic metal. In this case the fixed point action in Euclidean space takes the form $`S`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{Q}{}}[W_0(k_{})\omega ^2+W_1(k_{})k^2]\|\varphi (Q)\|^2`$ (15) $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{Q}{}}[{\displaystyle \frac{\omega ^2}{W_0(k_{})}}+{\displaystyle \frac{k^2}{W_1(k_{})}}]\|\theta (Q)\|^2,`$ (16) where $`Q=(\omega ,k,k_{})`$, for each chain the 2-current is $`j_\mu =\frac{1}{\sqrt{\pi }}ϵ_{\mu \nu }^\nu \varphi `$, and $`\theta `$ is the dual field of $`\varphi `$. The scaling dimensions of various local operators are determined by the dimensionless Luttinger coupling function $$w(k_{})=\sqrt{W_0(k_{})W_1(k_{})},$$ (17) which is periodic in $`(k_{})`$ with period $`2\pi `$ as we have set the interchain distance to 1. As in Ref., we consider the simplified model in which $`w(k_{})`$ takes the form $$w(k_{})=K_0+K_1\mathrm{cos}(k_{})=K_0[1+\lambda \mathrm{cos}(k_{})].$$ (18) Stability requires $`\|\lambda \|<1`$. In the presence of a spin gap single electron hopping is irrelevant, and the magnetic field has no effect on $`2k_F`$ back scattering which does not involve charge transfer between chains. We thus focus on the singlet pair hopping process, which in the presence of a magnetic field is described by the following perturbing Hamiltonian (near neighbor hopping only): $`H_{sc}`$ $`=`$ $`t_J{\displaystyle 𝑑xh_{sc}(x)},`$ (19) $`h_{sc}(x)`$ $`=`$ $`{\displaystyle \underset{j}{}}\mathrm{cos}[\sqrt{2\pi }(\theta _j(x)\theta _{j+1}(x))+2q_Bx],`$ (20) where $`t_J`$ is the Josephson coupling strength and $`q_B=eB/\mathrm{}c`$ as before. As in the previous section, the field adds an oscillatory phase to the pair hopping term, which renders $`h_{sc}`$ irrelevant by itself. However, as it flows, it again generates terms in which the oscillatory phases cancel. The most relevant of these is $`\stackrel{~}{H}_{sc}`$ $``$ $`t_J^2{\displaystyle 𝑑x\stackrel{~}{h}_{sc}(x)},`$ (21) $`\stackrel{~}{h}_{sc}(x)`$ $`=`$ $`{\displaystyle \underset{j}{}}\mathrm{cos}[\sqrt{2\pi }(2\theta _j(x)\theta _{j+1}(x)\theta _{j1}(x))],`$ (22) which is generated at second order in $`t_J`$. The scaling dimension of this term is $`\stackrel{~}{\mathrm{\Delta }}_{sc}`$ $`=`$ $`{\displaystyle \frac{2K_0}{2\pi }}{\displaystyle _0^{2\pi }}𝑑k_{}(1+\lambda \mathrm{cos}(k_{}))(1\mathrm{cos}(k_{}))^2`$ (23) $`=`$ $`(32\lambda )K_0.`$ (24) Combining the knowledge of $`\stackrel{~}{\mathrm{\Delta }}_{sc}`$ with the scaling dimension of the $`2k_F`$ back scattering operator $$\mathrm{\Delta }_{CDW}=\frac{2}{K_0(1\lambda +\sqrt{1\lambda ^2})},$$ (25) we can determine the phase diagram of the model (16) in the presence of a magnetic field and near neighbor interchain couplings, and subject to weak, generic perturbations, using the criteria that the smectic phase is stable when $`\stackrel{~}{\mathrm{\Delta }}_{sc}>2`$ and $`\mathrm{\Delta }_{CDW}>2`$; otherwise the system is in the stripe crystal/superconducting phase for $`\mathrm{\Delta }_{CDW}`$ smaller/bigger than $`\stackrel{~}{\mathrm{\Delta }}_{sc}`$. (In this identification we have made the natural assumption that the coupling (22) will govern the properties of the phase when it grows most rapidly. By itself, it will produce a vortex lattice .) The phase diagram is plotted in Fig. 2. For comparison we have also included the phase boundaries separating the superconducting phase from the smectic metal and stripe crystal phases in the absence of a magnetic field as dotted lines. It is quite obvious that both the smectic metal and stripe crystal phases get expanded by the magnetic field, which suppresses interchain Josephson coupling and increases the scaling dimension of operators involving pair hopping. Spin-ungapped smectic metal. In this case the fixed point action has contributions from both the charge and spin sectors: $`S=S_\rho +S_\sigma `$, where we take $`S_\rho `$ to have the same form as Eq. (16), and $$S_\sigma =\frac{K_\sigma }{2}\underset{j}{}[\frac{1}{v}(\tau \varphi _{j\sigma })^2+v(_x\varphi _{j\sigma })^2],$$ (26) in which we assume that there is no interchain coupling among spin fields, as in Ref. . Spin rotation invariance (also assumed here) requires $`K_\sigma =1`$. The analysis of pair hopping is similar to the previous case and it is easy to show that the most relevant operator generated by the pair hopping has scaling dimension $`\stackrel{~}{\mathrm{\Delta }}_{sc}^{nogap}=\stackrel{~}{\mathrm{\Delta }}_{sc}^{gap}+3>2`$; i.e., operators generated by pair hopping are always irrelevant here. Low-energy single electron hopping, which is allowed, is on the other hand more complicated and interesting. In terms of the original electron operators it takes the form $`H_e`$ $`=`$ $`t_e{\displaystyle 𝑑xh_e(x)},`$ (27) $`h_e(x)`$ $`=`$ $`{\displaystyle \underset{j\sigma }{}}(\psi _{j\sigma }^{}(x)\psi _{j+1\sigma }(x)e^{iq_Bx}+h.c.).`$ (28) We need to distinguish two different cases here. (i) $`k_F`$ and $`q_B`$ are incommensurate. In this case single-electron processes all involve an oscillating phase, and the most relevant process without an oscillating phase generated by $`H_e`$ is $`\stackrel{~}{H}_e`$ $``$ $`t_e^2{\displaystyle 𝑑x\stackrel{~}{h}_e(x)},`$ (29) $`\stackrel{~}{h}_e(x)`$ $`=`$ $`{\displaystyle \underset{i}{}}[\psi _j^L(x)\psi _j^R(x)\psi _{j+1}^L\psi _{j1}^R+h.c.+\mathrm{}],`$ (30) where $`L/R`$ stands for left/right mover, and $`\mathrm{}`$ stands for terms of similar structure. In bosonized form, $`\stackrel{~}{h}_e(x)`$ (31) $`\mathrm{cos}\{\sqrt{2\pi }[2(\theta _{\rho i}+\varphi _{\sigma i})\theta _{\rho i+1}\theta _{\rho i1}\varphi _{\sigma i+1}\varphi _{\sigma i1}]\},`$ (32) which has the scaling dimension $$\stackrel{~}{\mathrm{\Delta }}_e=1+\frac{K_0}{2}(\frac{3}{2}\lambda )+\frac{1\sqrt{1\lambda ^2}}{2K_0\lambda ^2\sqrt{1\lambda ^2}}.$$ (33) We assume that the system behaves as a Fermi liquid in a magnetic field when this term dominates. This identification is suggested if we note that at the non-interacting point, $`\lambda =0`$ and $`K_0=1`$, this term is marginal. This leads to the phase diagram Fig. 3, which is qualitatively different from the phase diagram in the absence of the field, Fig. 2 of Ref. . There are two particularly interesting differences: i) The superconducting phase gets completely squeezed out by the field; ii) The smectic metal phase now extends all the way to $`\lambda =0`$, which corresponds to the decoupled LL fixed point, a situation impossible without the field. (ii) $`2k_F=nq_B`$ where $`n`$ is an integer. In this case $`H_e`$, or its higher order descendents in the low energy theory, can turn a left mover on the Fermi point of the $`j`$th chain to a right mover on the Fermi point of the $`j+n`$th chain; this is a low energy single electron hopping process that does not involve an oscillatory phase, which takes the form $`H_e^{}`$ $`=`$ $`t_e{\displaystyle 𝑑xh^{}(x)},`$ (34) $`h^{}(x)`$ $`=`$ $`{\displaystyle \underset{j\sigma }{}}(\psi _{jL}^{}\psi _{j+nR}+h.c.)`$ (35) $``$ $`\mathrm{cos}\sqrt{{\displaystyle \frac{\pi }{2}}}(\theta _{\rho i}\theta _{\rho i+n}+\varphi _{\rho i}+\varphi _{\rho i+n})`$ (36) $`\times `$ $`\mathrm{cos}\sqrt{{\displaystyle \frac{\pi }{2}}}(\theta _{\sigma i}\theta _{\sigma i+n}+\varphi _{\sigma i}+\varphi _{\sigma i+n}).`$ (37) The scaling dimension of this operator $`\mathrm{\Delta }_{e,n}^{}`$ for $`n=1`$ is $$\mathrm{\Delta }_{e,1}^{}=\frac{K_0}{4}(1\frac{\lambda }{2})+\frac{1}{2K_0(1+\lambda +\sqrt{1\lambda ^2})}+\frac{1}{2}.$$ (38) In regions of parameter space where this is the most relevant operator, we expect that the system develops a gap that is largely single particle in character. The identification of the resulting state is easy once we recognize that the condition $`2k_F=nq_B`$ is precisely that the Landau level filling of the system is $`\nu =2n`$—i.e. the electrons (inclusive of their spin degeneracy) occupy $`n`$ Landau bands and form an integer quantum Hall state! In Fig. 4 we show the phase diagram for the case of $`n=1`$ ($`\nu =2`$). As the transition between the quantum Hall state and the smectic metal happens via the hopping going irrelevant, it is a continuous transition. To our knowledge, this is the first instance of a continuous transition between a quantum Hall state and a metallic state. We should note that the persistence of the quantum Hall phase upto the upper boundary at $`\lambda =1`$ is non-generic; it arises in the particular model studied upon a cancellation between numerator and denominator that will not typically take place. Finally, higher order commensurations between $`k_F`$ and $`q_B`$ are possible when lattice effects are strong on the chains and the electron operator has pieces oscillating at higher multiples of $`k_F`$. We have not investigated these. Spinless smectic metal. In this case we only have charged fields as in the spin gapped case, but single electron processes need to be considered as in the spin ungapped case. The analysis of perturbing operators is very similar to the spin ungapped case, which leads to the phase diagram Fig. 5 when $`k_F`$ and $`q_B`$ are incommensurate. Integer quantum Hall cases are, of course, allowed here as well, when $`2k_F=nq_B`$. Disorder: Following Giamarchi and Schulz , one can also analyze the scaling of weak single-particle randomness. We have not done this systematically, but will content ourselves with a couple of remarks. First, in all cases it is possible to find subsets of the smectic metal where both intrachain random backscattering and interchain random hopping are irrelevant—hence the system is a perfect, albeit completely anisotropic, metal in the long-wavelength limit. Second, it is possible to find sections of the phase boundary between the quantum Hall states and the smectic metal where disorder is still irrelevant, e.g. in the spin ungapped problem this happens both near $`\lambda =0`$ and near $`\lambda =1`$. In these cases we find an analytically tractable fixed point governing a transition out of a quantum Hall state in the presence of interactions and disorder that warrants further analysis . ## IV Summary Achieving a “dimensional continuation” of strong correlation physics from low dimensions by weakly coupling an infinite set of systems is an appealing strategy in the study of higher dimensional systems . The application of a magnetic field has been conjectured previously to be useful in this task. In addition to the work of Efetov, we should also mention the suggestion of Strong, Clarke and Anderson that a two-dimensional non-Fermi liquid phase could be induced in this fashion in a layered system. Striking experiments in the organic superconductors that are evidence for this point of view have been discussed at some length . In this paper, we have shown that this decoupling effect of the magnetic field can be given precise meaning in the context of two-dimensional sliding phases, via its reduction of the dimension of the most relevant charged operators that perturb them. This significantly expands the size of the sliding phases. As a bonus we find, in the quantum version of the problem, quantum Hall phases at commensurate fields that undergo a novel continuous transition to a smectic metal. In the underdoped region of the cuprates, it has been argued that the stripe instability leads to a smectic metal state and that it may already have been observed. In this setting, the spin gapped phase discussed here is the one at issue, whence we anticipate that the Zeeman coupling (ignored in our analysis) will not be important. We suggest that the field sensitivity of the phase diagram in this region would be an interesting test of the smectic hypothesis—essentially, one should look for the expansion of the metal or the onset of a CDW. The parameters needed to see this effect should ensure that the interchain hopping is weaker than the field, $`t_e<v_Fq_B`$ ($`v_F`$ is the on-chain Fermi velocity) and that the temperature does not wash out the phases induced by the field. The latter condition can be translated, via the on-chain smectic correlation length $`\xi ϵ_F/nT`$ ($`ϵ_F`$ is the on chain Fermi energy and $`n`$ is the linear density of electrons), to the statement $`Ba\xi \varphi _o`$ where $`a`$ is the interchain spacing and $`\varphi _o`$ is the flux quantum. ###### Acknowledgements. We are grateful to Ashvin Vishwanath, Tom Lubensky and Steve Kivelson for useful discussions. We are especially grateful to David Huse for many illuminating discussions and especially for producing a lucid, intuitive explanation of the content of Ref. . This work was supported by NSF DMR-9971541 and the Sloan Foundation (KY), and NSF DMR-9978074 and the Sloan and Packard Foundations (SLS).
warning/0007/cond-mat0007136.html	ar5iv	text	# Dzyaloshinskii-Moriya interaction in the paramagnetic state and the polarized neutron scattering. ## Abstract Dzyaloshinskii-Moriya (DM) interaction in the paramagnetic state leads to the incommensurate spin fluctuations with incommensurate vector proportional to the relative strength of the DM interaction. We show that the DM interaction leads to helical spin fluctuations which may be observed by the polarized neutron scattering. , thanks: Corresponding author. E-mail: aristov@thd.pnpi.spb.ru In the case of inelastic magnetic scattering of polarized neutrons the cross section consists of two terms. The first one is independent on the initial neutron polarization $`𝐏_0`$ and determined by the symmetric part of the generalized magnetic susceptibility $`\chi _{\alpha \beta }(Q,\omega )`$. The second one is proportional to $`𝐏_0`$ and connected with the antisymmetric part of $`\chi _{\alpha \beta }`$. This antisymmetric part of the susceptibility appears if the system is characterized by an axial vector. There are two possibilities. i) External magnetic field or the sample magnetization (see and references therein). ii) Some intrinsic axial- vector interaction which is connected to the noncentrosymmetry of the system . In this paper we consider the case of the Dzyaloshinskii–Moriya interaction (DMI) and demonstrate that the dependence of the magnetic scattering on $`𝐏_0`$ may appear in the paramagnetic phase along with the incommensurate peaks in both parts of the scattering cross section. We demonstrate it using the DMI as perturbation in the three-dimensional (3D) case. Then we confirm these results by exact solution of the $`1D`$ problem. It should be noted here that the incommensurate $`𝐏_0`$-dependent paramagnetic scattering was observed in MnSi . To the best of our knowledge it is the only experimental study of this problem. The DMI has the following form $$V_{DM}=\frac{1}{2}\underset{l,m}{}𝐃_{lm}[𝐒_l\times 𝐒_m],$$ (1) where $`𝐃_{lm}=𝐃_{ml}`$ is the DM axial vector. We begin with the case when $`𝐃_{lm}`$ is invariant under translations on the lattice, $`𝐃_{lm}=𝐃_{𝐥+𝐚_1,𝐦+𝐚_2}`$, and assume that $`𝐃`$ is directed along the $`z`$ axis. After Fourier transform Eq.(1) may be represented as follows $$V_{DM}=i\underset{q}{}d_𝐪S_𝐪^xS_𝐪^y,$$ (2) where $$d_𝐪=d_𝐪=i\underset{l}{}D_{lm}e^{i\mathrm{𝐪𝐑}_{lm}}.$$ (3) We assume now that the paramagnetic spin fluctuations are isotropic, if one neglects the DMI. In this case the spin Green function has the form $`G_{\alpha \beta }^0(q,\omega )=\delta _{\alpha \beta }G^0(q,\omega )`$. Using interaction (2) as small perturbation we obtain $`G_{xx}`$ $`=`$ $`G^0+iG^0d_qG_{yx}`$ (4) $`G_{yx}`$ $`=`$ $`iG^0d_qG_{xx}.`$ As a result we get $`G_{xx}`$ $`=`$ $`G_{yy}=G^0\left(1d_q^2(G^0)^2\right)^1`$ (5) $`G_{xy}`$ $`=`$ $`G_{yx}=iG^0d_q\left(1d_q^2(G^0)^2\right)^1.`$ We see that the DMI leads to the nondiagonal antisymmetric components of the spin Green function. To clarify these expressions let us consider the static approximation $`(\omega =0)`$ and choose $`G(q,0)`$ in the conventional Ornstein–Zernike form $$G(q,0)=G(q)=A(q^2+\kappa ^2)^1,$$ (6) where $`\kappa `$ is the inverse correlation length and $`A(T_ca^2)^1`$ where $`T_c`$ is the transition temperature to the ordered state and $`a`$ is the interatomic spacing. Having in mind that $`d_0=0`$ and at small $`q`$ one has $`Ad_q=2\alpha (𝐪\widehat{n})`$ where $`\widehat{n}`$ is the direction of the bonds, along which the DM interaction is present, and $`\alpha ADaa^1`$ , we obtain from Eqs. (5) and (6) $`G_{xx}`$ $`=`$ $`G_{yy}={\displaystyle \frac{A}{2}}([\kappa _1^2+(𝐪\alpha \widehat{n})^2]^1`$ (7) $`+[\kappa _1^2+(𝐪+\alpha \widehat{n})^2]^1),`$ $`G_{xy}`$ $`=`$ $`G_{yx}=i{\displaystyle \frac{A}{2}}([\kappa _1^2+(𝐪+\alpha \widehat{n})^2]^1`$ $`[\kappa _1^2+(𝐪+\alpha \widehat{n})^2]^1).`$ with $`\kappa _1^2=\kappa ^2\alpha ^2`$ ; these expressions describe incommensurate spin fluctuations at $`𝐪=\pm \alpha \widehat{n}`$. These expressions are the result of the first order perturbation theory in the DMI value and there should be additional terms of order $`\alpha ^2`$ in the denominators, which we did not evaluate. According to Ref. due to the DMI the phase transition to the ordered state should be the first order one. Experimental study of this problem would be very interesting. The possible candidates for such study could be the systems MnSi, FeG, Fe<sub>2</sub>O<sub>3</sub> and quasi-1D antiferromagnet CsCuCl<sub>3</sub>. As was stated above the antisymmetric part of the spin Green function gives rise to the $`𝐏_0`$-dependent part of the cross section. In our case it may be represented as $`G_{\alpha \beta }=iϵ_{\alpha \beta \gamma }\widehat{z}_\gamma G^A`$, where $`\widehat{𝐳}`$ is the unit vector along the $`z`$axis. In this case the $`𝐏_0`$-dependent part of the cross section has the form (cf.) $`\left({\displaystyle \frac{d\sigma }{d\mathrm{\Omega }d\omega }}\right)_{P_0}`$ $`=`$ $`{\displaystyle \frac{2}{\pi }}r^2f^2(q){\displaystyle \frac{k_f}{k_i}}(1e^{\omega /T})^1`$ (8) $`\times (𝐏_0\widehat{𝐪})(\widehat{𝐪}\widehat{𝐳})\text{ Im }G^A(𝐪,\omega ),`$ where $`r^2=0.292`$ barns, $`f(q)`$ is the magnetic form-factor and $`\widehat{𝐪}=𝐪/q`$. If the asymmetry of $`G_{\alpha \beta }`$ is determined by the magnetic field, Im$`G^A`$ is an even function of $`\omega `$ and, provided $`\omega T`$, we have $`𝑑\omega (d\sigma /d\mathrm{\Omega }d\omega )_{𝐏_0}=0`$. It is a consequence of the $`t`$-oddness of the magnetic field . In our case the vector $`D`$ is $`t`$-even and for the static contribution one has $`\left({\displaystyle \frac{d\sigma }{d\mathrm{\Omega }}}\right)_{P_0}`$ $`=`$ $`{\displaystyle \frac{T}{\pi }}Ar^2f^2(q)(𝐏_0\widehat{𝐪})(\widehat{𝐪}\widehat{𝐳})[{\displaystyle \frac{1}{\kappa _1^2+\left(𝐪+\alpha \widehat{n}\right)^2}}`$ (9) $`{\displaystyle \frac{1}{\kappa _1^2+\left(𝐪\alpha \widehat{n}\right)^2}}].`$ Up to now we discussed the translationally invariant DMI. In the case of the staggered DMI, one has $`D_{l+b_x,m+b_y}=D_{lm}`$, where $`𝐛`$ is the minimal vector along the bond where the DMI is present. In this case, instead of Eq.(2), we have $$V_{DM}=i\underset{q}{}d_𝐪S_𝐪^xS_{𝐪𝐤_0}^y,$$ where $`𝐤_0`$ is the AF reciprocal wave vector along $`𝐛`$. As a result $`G_{xy}`$ depends on $`𝐪`$ and $`𝐪+𝐤_0`$ and cannot be determined in the neutron scattering experiments. In this case the cross section is commensurate and independent of $`𝐏_0`$. The above results were obtained in the perturbation theory. We present now an exact solution of the problem in the 1D case. We consider the spin chain Hamiltonian of the form $`H`$ $`=`$ $`{\displaystyle \underset{l=1}{\overset{L}{}}}(J𝐒_l𝐒_{l+1}+𝐃[𝐒_l\times 𝐒_{l+1}])`$ (10) with AF Heisenberg coupling $`J>0`$ and the Dzyaloshinskii-Moriya term $`𝐃`$. It is convenient to introduce here the quantity $`\delta =\mathrm{tan}^1(D/J)`$. We observe that $`H`$ is simplified upon a canonical transformation $`He^{iU}He^{iU}`$ with $`U=\delta _{l=1}^LlS_l^z.`$ One can easily see that for the combinations $`S_j^\pm =S_j^x\pm iS_j^y`$ we get $$\stackrel{~}{S}_l^\pm e^{iU}S_l^\pm e^{iU}=S_l^\pm e^{il\delta },\stackrel{~}{S}_l^z=S_l^z,$$ (11) and the Hamiltonian is reduced to the $`XXZ`$ model: $$H=\underset{l=1}{\overset{L}{}}(J_x(\stackrel{~}{S}_l^x\stackrel{~}{S}_{l+1}^x+\stackrel{~}{S}_l^y\stackrel{~}{S}_{l+1}^y)+J\stackrel{~}{S}_l^z\stackrel{~}{S}_{l+1}^z)$$ (12) with $`J_x=J/\mathrm{cos}\delta `$. It follows then, that the spectrum of the initial problem (10) coincides with the one of (12). The observables in the initial system are recalculated with the use of (11) from the observables in the XXZ model (12). In the latter model one distinguishes the longitudinal ($`G^{}(k,\omega )`$) and the transverse ($`G^{}(k,\omega )`$) spin correlations, for the $`z`$ and $`x`$ components of spin, respectively. The difference between these Green functions is small in the considered limit, $`\delta 0`$. First we note that $`U`$ does not affect the $`z`$component of spins. Therefore the “longitudinal” Green function $`G_{zz}(k,\omega )=G^{}(k,\omega )`$ has a commensurate antiferromagnetic modulation. The transverse spin susceptibilities look a bit more complicated. Some calculation shows that $`G_{xx}(l,m,\omega )`$ $`=`$ $`G_{yy}(l,m,\omega )`$ (13) $`=`$ $`G^{}(l,m,\omega )\mathrm{cos}\delta (lm)`$ $`G_{yx}(l,m,\omega )`$ $`=`$ $`G_{xy}(l,m,\omega )`$ (14) $`=`$ $`G^{}(l,m,\omega )\mathrm{sin}\delta (lm)`$ In terms of the Fourier transform this reads as $`G_{xx}(q,\omega )`$ $`=`$ $`G_{yy}(q,\omega )`$ $`=`$ $`{\displaystyle \frac{1}{2}}[G^{}(q+\delta ,\omega )+G^{}(q\delta ,\omega )],`$ $`G_{xy}(q,\omega )`$ $`=`$ $`G_{yx}(q,\omega )`$ $`=`$ $`{\displaystyle \frac{i}{2}}[G^{}(q+\delta ,\omega )G^{}(q\delta ,\omega )].`$ From these expressions we see that the transverse and chiral fluctuations are incommensurate along the chain and in the limit $`D/J1`$ the incommensurate vector coincides with that determined by Eq.(7). However the complete solution of the problem (12) can be found in literature (see, e.g., Ref. ). In the 1D case we know the exact $`\omega `$ and $`q`$dependence of all types of the spin fluctuations. Note that in the quasi-1D compounds the value of $`\delta D/J`$, determining the incommensurate wavevector of the fluctuations, may be sufficiently large. For instance, one has $`\delta 0.18`$ in the CsCuCl<sub>3</sub> and $`\delta 0.05`$ in copper benzoate. In the latter compound, however, the presumably staggered variant of DMI should not lead to consequences, observable by the polarized neutron scattering. In conclusion, we demonstrate that the DM interaction in the paramagnetic state leads to the incommensurate spin fluctuations with incommensurate vector proportional to the strength of the DM interaction relative to the exchange one. It is shown also that DMI leads to the helical spin fluctuations which may be observed by the polarized neutron scattering. This work was supported by Russian State Program for Statistical Physics (Grant VIII-2), RFBR Grants No. 00-02-16873 and 00-15-96814, and the Russian Program ”Neutron Studies of Condensed Matter”.
warning/0007/hep-lat0007034.html	ar5iv	text	# A study of the ’t Hooft loop in 𝑆⁢𝑈⁢(2) Yang–Mills theory ## Abstract We study the behaviour of the spatial and temporal ’t Hooft loop at zero and finite temperature in the $`4D`$ $`SU(2)`$ Yang-Mills theory, using a new numerical method. In the deconfined phase $`T>T_c`$, the spatial ’t Hooft loop exhibits a dual string tension, which vanishes at $`T_c`$ with $`3D`$ Ising-like critical exponent. The $`4D`$ $`SU(2)`$ Yang-Mills theory undergoes a transition between a cold confined phase and a hot deconfined phase at a critical temperature $`T_c`$. An order parameter widely used to characterize this transition is the Polyakov loop. It develops a non-vanishing expectation value in the deconfined phase. However, the corresponding operator creates a single fundamental static color source, which does not belong to the physical Hilbert space of the theory; it cannot be defined at zero temperature; and it is afflicted by ultra-violet divergences in the continuum limit . A long time ago it was proposed to consider instead the ’t Hooft loop operator as an order parameter to characterize this transition. It is the purpose of this paper to study this order parameter. The ’t Hooft loop, $`\stackrel{~}{W}(C)`$, is an operator associated with a given closed contour $`C`$, and is defined in the continuum $`SU(N)`$ theory by the following equal–time commutation relations $`[W(C),W(C^{})]=[\stackrel{~}{W}(C),\stackrel{~}{W}(C^{})]=0`$ (1) $`\stackrel{~}{W}^{}(C)W(C^{})\stackrel{~}{W}(C)=e^{i\frac{2\pi }{N}n_{CC^{}}}W(C^{})`$ (2) where $`W(C^{})`$ is the Wilson loop associated with the closed contour $`C^{}`$ and $`n_{CC^{}}`$ is the linking number of $`C`$ and $`C^{}`$. Just like the Wilson loop creates an elementary electric flux along $`C^{}`$, the ’t Hooft loop creates an elementary magnetic flux along the closed path $`C`$ affecting any Wilson loop “pierced” by $`C`$. In that sense, the two types of loop are dual to each other. At zero temperature, it has been shown that also the ’t Hooft loop behaviour is dual to that of the Wilson loop: in the absence of massless excitations, an area law behaviour for one implies a perimeter law for the other, and vice versa. Hence, at $`T=0`$ the ’t Hooft loop obeys a perimeter law. Several analytical and numerical studies have been carried out in order to investigate this issue of duality at finite temperature. At $`T>0`$, the Lorentz symmetry is broken, so spatial and temporal loops can have different behaviours. Because the spatial string tension persists also above $`T_c`$ for the Wilson loop, temporal ’t Hooft loops are expected to show a perimeter law in both phases; spatial ’t Hooft loops are expected to obey a perimeter law in the confined phase and an area law – defining a dual string tension (strictly speaking it is an action density) – in the deconfined phase. On the lattice the ’t Hooft loop is defined as follows. Let us consider the $`SU(2)`$ lattice gauge theory with the usual Wilson plaquette action, $$S(\beta )=\beta \underset{P}{}(1\frac{1}{2}\text{Tr}(U_P)),$$ (3) where the sum extends over all the plaquettes $`P`$, and $`U_P`$ is the path–ordered product of the links around $`P`$. Starting from $`S(\beta )`$, one defines the partition function $$Z(\beta )=[dU]\mathrm{exp}(S(\beta )).$$ (4) Let us now switch on “by hand” an elementary magnetic flux along a closed contour $`C`$ defined on the dual lattice. To create this magnetic flux, we have to multiply $`U_P`$ by a non-trivial element of the center group for the plaquettes $`P`$ dual to a given surface $`𝒮`$ supported by $`C`$. For the $`SU(2)`$ gauge group, this means flipping the sign of the coupling, since the only non-trivial element of the center $`Z_2`$ is $`1`$. Call $`𝒫(𝒮)`$ the set of plaquettes whose coupling is flipped $`\beta \beta `$. Then the action of the system where an elementary flux along a closed contour $`C`$ has been switched on is given, up to an additive constant, by $$S_𝒮(\beta )=\frac{1}{2}\beta \left(\underset{P𝒫(𝒮)}{}\text{Tr}(U_P)\underset{P𝒫(𝒮)}{}\text{Tr}(U_P)\right),$$ (5) and the partition function is $$Z_C(\beta )=[dU]\mathrm{exp}(S_𝒮(\beta )).$$ (6) $`Z_C(\beta )`$ does not depend on the particular chosen surface $`𝒮`$, since different choices are related by a change of integration variables. The simplest choice for $`𝒮`$ is the minimal surface spanning $`C`$. Thus, if $`C`$ is an $`R_x\times R_y`$ rectangle in the $`(x,y)`$ plane (spatial ’t Hooft loop), one flips the coupling of the $`(z,t)`$ plaquettes dual to the plaquettes belonging to the rectangular area. The ’t Hooft loop expectation value is related to the free energy cost needed to create the magnetic flux along the contour $`C`$ and is given by $$\stackrel{~}{W}(C)=Z_C(\beta )/Z(\beta ).$$ (7) This expression can be rewritten in the form $$\stackrel{~}{W}(C)=\mathrm{exp}\left(\beta \underset{P𝒫(𝒮)}{}\text{Tr}(U_P)\right),$$ (8) with the average taken with the standard Wilson action. In this form, the difficulty of measuring the ’t Hooft loop becomes clear: the observable is exponentially suppressed on typical configurations of the statistical ensemble, and gets a significant contribution only from configurations having an extremely small statistical weight. Therefore the numerical evaluation of $`\stackrel{~}{W}(C)`$ represents a difficult sampling problem, increasingly so with the loop size. Recently, the ’t Hooft loop, or special cases of it, have been studied numerically on the lattice. In the sampling problem was overcome by using a multihistogram method, where one performs several different simulations in which the coupling associated with the plaquettes in $`𝒫(𝒮)`$ is gradually changed from $`\beta `$ to $`\beta `$. In , instead, the derivative $`\mathrm{d}/\mathrm{d}\beta \mathrm{ln}\stackrel{~}{W}(C)`$ has been determined, which is a much simpler numerical task. In this letter we report on a similar numerical study where, by measuring spatial and temporal ’t Hooft loops at zero and finite temperature, we confirm the role of the ’t Hooft loop as a dual order parameter for confinement. We adopt a new numerical method, rewriting the ratio $`Z_C(\beta )/Z(\beta )`$ as a product of intermediate ratios, each easily measurable. We establish the perimeter law behaviour of the ’t Hooft loop at zero temperature. We measure the free energy of two center monopoles as a function of their separation: as expected from duality, temporal ’t Hooft loops show screening behaviour at all temperatures, while spatial ’t Hooft loops exhibit a dual string tension above $`T_c`$. Moreover, this dual string tension vanishes at $`T_c`$ with a critical exponent $`\nu `$ very close to that of the $`3D`$ Ising model, consistent with universality. The numerical technique – We now describe our numerical method to measure the expectation value of the ’t Hooft operator. Because the relevant contributions to $`Z_C(\beta )`$ and to $`Z(\beta )`$ come from regions of the phase space with very poor overlap, a direct evaluation of (7), (8) by a single Monte Carlo simulation is not a reliable way to compute $`\stackrel{~}{W}(C)`$. It is necessary to consider a sequence of intermediate partition functions which interpolate between $`Z(\beta )`$ and $`Z_C(\beta )`$. The approach used in consists of interpolating in the coupling of the flipped plaquettes in $`𝒫(𝒮)`$, from $`\beta `$ to $`\beta `$. We interpolate instead in the number of flipped plaquettes, and rewrite $`Z_C(\beta )/Z(\beta )`$ as the product of ratios of partition functions where the number of plaquettes with flipped coupling is progressively reduced. $`N`$ being the number of plaquettes in $`𝒫(𝒮)`$, we use the identity $$\frac{Z_C(\beta )}{Z(\beta )}=\frac{Z_N(\beta )}{Z_{N1}(\beta )}\frac{Z_{N1}(\beta )}{Z_{N2}(\beta )}\mathrm{}\frac{Z_1(\beta )}{Z_0(\beta )}$$ (9) where $`Z_k(\beta )`$, $`k=0,\mathrm{},N`$ ($`Z_NZ_C`$ and $`Z_0Z`$) is the partition function of the system where only the first $`k`$ plaquettes in $`𝒫(𝒮)`$ have flipped coupling. Every ratio $`Z_k/Z_{k1}`$ can be now computed efficiently by a single Monte Carlo simulation, due to the good overlap of the relevant phase space in the two partition functions. Furthermore, in the practical implementation, it is useful to reexpress $`Z_k(\beta )/Z_{k1}(\beta )`$ as a ratio of expectation values $$\frac{Z_k(\beta )}{Z_{k1}(\beta )}=\frac{\mathrm{exp}(\frac{1}{2}\beta \text{Tr}(U_{P_k}))_k}{\mathrm{exp}(+\frac{1}{2}\beta \text{Tr}(U_{P_k}))_k}$$ (10) The averages $`_k`$ are computed with respect to an action where the $`k`$th plaquette $`U_{P_k}`$ of $`𝒫(𝒮)`$ has zero coupling, the first $`(k1)`$ plaquettes have coupling $`\beta `$ and the remaining ones coupling $`+\beta `$. The benefits over simply computing $`\mathrm{exp}(\beta \text{Tr}(U_{P_k}))`$ using the distribution corresponding to $`Z_{k1}(\beta )`$ are two-fold. $`(i)`$ Using the intermediate weight, $`_k`$, allows a much better sampling for both $`\mathrm{exp}(\frac{1}{2}\beta \text{Tr}(U_{P_k}))_k`$ and $`\mathrm{exp}(+\frac{1}{2}\beta \text{Tr}(U_{P_k}))_k`$, with reduced errors. $`(ii)`$ Since the quantity $`Z_k(\beta )/Z_{k1}(\beta )`$ refers to a single plaquette of the lattice, it is very useful to perform a partial integration after each updating sweep of the whole lattice, and measure $`\mathrm{exp}(+\frac{\beta }{2}\text{Tr}(U_{P_k}))`$ and $`\mathrm{exp}(\frac{\beta }{2}\text{Tr}(U_{P_k}))`$ several times while updating only the four links belonging to plaquette $`P_k`$. Since this plaquette has zero coupling, these four links are decoupled from each other. Therefore each link can be updated independently $`N_{\mathrm{hit}}`$ times and the different copies of each link can be combined to obtain $`N_{\mathrm{hit}}^4`$ measurements of $`\mathrm{exp}(\pm \frac{\beta }{2}\text{Tr}(U_{P_k}))`$. Although these measurements are correlated, the variance reduction is important. The statistical error on ratio (10) must be evaluated with care, since the two averages are computed from the same sample of configurations. We use a jackknife analysis. Then, after each $`Z_k(\beta )/Z_{k1}(\beta )`$ has been computed, the ’t Hooft loop expectation value is evaluated according to Eq.(9). The final statistical error is obtained by standard error propagation since each $`Z_k(\beta )/Z_{k1}(\beta )`$ comes from an independent Monte Carlo simulation. An advantage of our method over the multihistogram technique is that the products in (9) give us information on smaller ’t Hooft loops for free; moreover, the error analysis is simpler and less delicate than in a multihistogram analysis. Results – We focus on the free energy $`F(R)`$ of a pair of static center monopoles as a function of their separation $`R`$. It can be obtained as $`\mathrm{lim}_{R_t\mathrm{}}\mathrm{Ln}[\stackrel{~}{W}(R,R_t)]/R_t`$ by taking elongated $`R\times R_t`$ rectangular loops, in the same way as one extracts the static potential between two chromo-electric charges. We take $`R_t`$ as large as possible, i.e. equal to the lattice size $`L`$. This is analogous to measuring the correlation of two Polyakov loops, and is the correct approach at finite temperature. Therefore we must flip the coupling of $`R\times L`$ plaquettes. We do this by scanning first the $`R_t`$, then the $`R`$direction. With this ordering, the intermediate partition functions $`Z_L,Z_{2L},..,Z_{R\times L}`$ in (9) provide us with the free energy at separations $`1,2,..,R`$ respectively. The final ratio $`Z_{L\times L}/Z_0`$ gives the free energy of a center vortex as computed in . Our lattice sizes range from $`10^3\times 2`$ to $`20^3\times 10`$, at couplings $`\beta =\frac{4}{g^2}`$ from $`2.3`$ to $`2.8`$, for which the lattice spacing varies by a factor $`5`$. We perform 5-10000 (multi-hit) measurements of each ratio (10). At zero temperature the ’t Hooft loop is expected to obey a perimeter law: $`Z_k/Z_0e^{c\stackrel{~}{P}_k}`$, where $`\stackrel{~}{P}_k`$ is the length of the contour $`C_k`$. For the sequence of flipped plaquettes defined above, $`\stackrel{~}{P}_k=2L`$ if $`\mathrm{mod}(k,L)=0`$, $`(2L+2)`$ otherwise, unless $`k<L`$ or $`k>L(L1)`$. $`Z_k`$ should therefore center around two values only. This is exactly what appears in Fig.1, on a $`10^4`$ lattice. The free energy is clearly insensitive to changes in the ’t Hooft loop area. This is also confirmed by a direct measurement of Creutz ratios $`\chi (R,R)`$, which estimate the force at distance $`R`$ between the two magnetic charges and quickly drop to zero as $`R`$ increases. As with the Wilson loop, the coefficient $`c`$ of the perimeter term is affected by UV divergences. We verified that it increases as the lattice spacing is reduced. The free energy $`F(r)`$ can be fitted by a Yukawa form $`\frac{e^{mr}}{r}`$, up to an irrelevant additive constant. However the screening mass $`m`$ is rather large, so that the signal quickly dies out. Moreover, consistency of data at different values of the lattice spacing is only fair. This is presumably caused by short-distance lattice distortions of the Yukawa potential, which could be explicitly taken into account. In any case, Fig.2 shows the free energy as a function of the magnetic charge separation. The curve corresponds to a screening mass of $`2`$ GeV. It makes sense to match this mass with the lightest gluonic excitation, the scalar glueball ($`1.65`$ GeV), as attempted in ; but this issue awaits a more extensive numerical study. At finite temperature we must distinguish between spatial and temporal ’t Hooft loops. In the case of electric charges, the spatial string tension persists above $`T_c`$, while the correlation of time-like Polyakov loops develops a disconnected part showing saturation of the static potential. Here the temporal ’t Hooft loop should show screening, while the spatial ’t Hooft loop should obey an area law. It is clear that the latter will cost more in free energy, since the center vortex created by each flipped $`(z,t)`$ plaquette can only spread over a limited time extent $`T^1`$. Indeed, this is precisely what we observe. Fig.3 shows the same Yukawa form for a temporal ’t Hooft loop at all temperatures. The only effect of temperature is to increase the screening mass. The fitting coefficient of the Yukawa potential is smaller than the short-distance perturbative prediction $`1/4\pi (2\pi /g)^2`$ , but grows towards this value at higher $`\beta `$. An attempt to fit the data with the ansatz $`F_0+c\frac{e^{mr}}{r}+\sigma r`$, including a linear term, gives a dual string tension $`\sigma `$ consistent with zero. In contrast, this linear term is required to obtain an acceptable fit above $`T_c`$ for the spatial ’t Hooft loop: a dual string tension appears. As a practical consequence, the screening mass becomes yet harder to determine. It seems little affected by temperature, as Fig.4 (top) shows, unlike the temporal screening mass (bottom) which rises more or less linearly with $`T`$, much like the glueball excitation which it presumably represents . Precise quantitative statements about these dependences on $`T`$ require a more accurate numerical investigation. The dual string tension $`\sigma `$ depends on temperature and must vanish at $`T_c`$. Fig.5 shows that it does so as $`\sigma \left(\frac{TT_c}{T_c}\right)^{2\nu }`$. The critical exponent $`\nu `$, associated with the correlation length $`\xi =\sigma ^{1/2}`$, comes out very close to that of the $`3D`$ Ising model: $`0.66(3)`$ vs. $`0.63`$. Indeed this should be expected since both models are in the same universality class, although the Ising exponent is typically extracted from the divergence of $`\xi `$ in the symmetric (confined) phase. This dual string tension can then be taken as order parameter for the restoration of the (magnetic) $`Z_N`$ symmetry, corresponding to deconfinement . Fig.5 also shows high $`T`$, perturbative results. Notice that a spatial $`(x,y)`$ ’t Hooft loop of maximal size $`L\times L`$ introduces a flipped plaquette in every $`(z,t)`$ plane, which is equivalent to enforcing twisted boundary conditions in the $`(z,t)`$ directions, thus creating a $`Z_N`$ interface with tension $`(\sigma T)`$. This interface tension has been calculated to two loops . Taking the running coupling $`g(T)`$ from , one obtains the two curves in Fig.5 for the leading and next order. The numerical data lie in between. In summary, by using a dual observable, we have measured a dual string (or interface) tension in the deconfined phase. The corresponding correlation length diverges at $`T_c`$ with the $`3D`$ Ising critical exponent $`\nu `$ as expected from universality. An extension to $`SU(3)`$ is straightforward. In that case, since the transition is first-order, the dual string tension will persist all the way to $`T_c`$. Its value at $`T_c`$ has been the object of many studies . Finally, it would be very desirable to measure more precisely the screening mass and to clarify its physical origin. Acknowledgments: We thank C. Alexandrou, P. Butera, L. Del Debbio, A. Di Giacomo, C. Korthals–Altes, M. Laine and O. Philipsen for many helpful discussions, and acknowledge communication with T. Kovács, T. Tomboulis and C. Rebbi.
warning/0007/hep-ph0007034.html	ar5iv	text	# 1 Leading contributions to 𝜋̃⁰→𝛾̃⁢𝛾̃ in the CFL phase. ## Abstract QCD superconductors in the color-flavor-locked (CFL) phase support light excitations (generalized pions) in the form of particle-particle or hole-hole excitations. We analyze the generalized process $`\stackrel{~}{\pi }^0\stackrel{~}{\gamma }\stackrel{~}{\gamma }`$ in the weak coupling limit and show that it is related to the recently suggested Wess-Zumino-Witten (WZW) term. In dense QCD, the radiative decay of the generalized pion is constrained by geometry and vanishes at large density. KIAS-P00043 FZJ-IKP(TH)-2000-16 $`\stackrel{~}{\pi }^0\stackrel{~}{\gamma }\stackrel{~}{\gamma }`$ in Dense QCD Maciej A. Nowak<sup>a,b</sup><sup>1</sup><sup>1</sup>1E-mail: nowak@kiwi.if.uj.edu.pl , Mannque Rho<sup>a,c</sup><sup>2</sup><sup>2</sup>2E-mail: rho@spht.saclay.cea.fr, Andreas Wirzba<sup>d</sup><sup>3</sup><sup>3</sup>3E-mail: a.wirzba@fz-juelich.de and Ismail Zahed<sup>a,e</sup><sup>4</sup><sup>4</sup>4E-mail: zahed@zahed.physics.sunysb.edu <sup>a</sup> School of Physics, Korea Institute for Advanced Study, Seoul 130-012, Korea <sup>b</sup> M. Smoluchowski Institute of Physics, Jagellonian University, Cracow, Poland <sup>c</sup> Service de Physique Théorique, CE Saclay, 91191 Gif-sur-Yvette, France <sup>d</sup>FZ Jülich, Institut für Kernphysik, D-52425 Jülich, Germany <sup>e</sup> Department of Physics and Astronomy, SUNY-Stony-Brook, NY 11794, U. S. A. 1. At high density, QCD exhibits a superconducting phase with novel and nonperturbative phenomena . For degenerate quark masses the QCD superconductor breaks color and flavor spontaneously with the occurrence of Goldstone modes. These modes were analyzed recently using effective Lagrangians (zero size) and bound state equations (finite size). In the color-flavor-locked (CFL) phase the Goldstone modes are characterized by a Wess-Zumino-Witten (WZW) term with a generalized axial-anomaly . In the present letter we analyze the decay process $`\stackrel{~}{\pi }^0\stackrel{~}{\gamma }\stackrel{~}{\gamma }`$ of the generalized pion $`\stackrel{~}{\pi }^0`$ into a pair of modified photons $`\stackrel{~}{\gamma }`$ using weak-coupling arguments at large quark chemical potential. After recalling some features of the QCD superconductor, we proceed to an analysis of the triangle anomaly in the CFL phase. We derive an explicit result for the amplitude of $`\stackrel{~}{\pi }^0\stackrel{~}{\gamma }\stackrel{~}{\gamma }`$ in the limit of large quark density. We show that our result agrees with the recently suggested WZW term . The decay amplitude and the color-flavor anomalies in the CFL phase are related at asymptotic densities. 2. In the QCD superconductor, the quarks are gapped. Their propagation in the chiral limit is given in the Nambu-Gorkov formalism by $`S_{11}(q)`$ $`=`$ $`{\displaystyle \frac{1}{i}}\psi (q)\overline{\psi }(q)=\left[\mathrm{\Lambda }^+(𝐪)\gamma ^0{\displaystyle \frac{q_0+q_{\|\|}}{q_0^2ϵ_q^2}}\mathrm{\Lambda }^{}(𝐪)+\mathrm{\Lambda }^{}(𝐪)\gamma ^0{\displaystyle \frac{q_0q_{\|\|}2\mu }{q_0^2\overline{ϵ}_q^2}}\mathrm{\Lambda }^+(𝐪)\right]`$ $`S_{12}(q)`$ $`=`$ $`{\displaystyle \frac{1}{i}}\psi (q)\overline{\psi }_C(q)=[\mathrm{\Lambda }^+𝐪)𝐌^{}{\displaystyle \frac{G^{}(q)}{q_0^2ϵ_q^2}}\mathrm{\Lambda }^+(𝐪)+\mathrm{\Lambda }^{}(𝐪)𝐌^{}{\displaystyle \frac{\overline{G^{}}(q)}{q_0^2\overline{ϵ}_q^2}}\mathrm{\Lambda }^{}(𝐪)]`$ $`S_{21}(q)`$ $`=`$ $`{\displaystyle \frac{1}{i}}\psi _C(q)\overline{\psi }(q)=\left[\mathrm{\Lambda }^{}(𝐪)𝐌{\displaystyle \frac{G(q)}{q_0^2ϵ_q^2}}\mathrm{\Lambda }^{}(𝐪)+\mathrm{\Lambda }^+(𝐪)𝐌{\displaystyle \frac{\overline{G}(q)}{q_0^2\overline{ϵ}_q^2}}\mathrm{\Lambda }^+(𝐪)\right]`$ $`S_{22}(q)`$ $`=`$ $`{\displaystyle \frac{1}{i}}\psi _C(q)\overline{\psi }_C(q)=\left[\mathrm{\Lambda }^{}(𝐪)\gamma ^0{\displaystyle \frac{q_0q_{\|\|}}{q_0^2ϵ_q^2}}\mathrm{\Lambda }^+(𝐪)+\mathrm{\Lambda }^+(𝐪)\gamma ^0{\displaystyle \frac{q_0+q_{\|\|}+2\mu }{q_0^2\overline{ϵ}_q^2}}\mathrm{\Lambda }^{}(𝐪)\right]`$ with the momentum $`q_{\|\|}=(\|𝐪\|\mu )`$ measured parallel and relative to the Fermi momentum $`\|𝐩_F\|=\mu `$, the squared particle/hole energies $`ϵ_q^2=q_{\|\|}^2+𝐌^{}𝐌\|G(q)\|^2`$ and the squared anti-particle/anti-hole energies $`\overline{ϵ}_q^2=(q_{\|\|}+2\mu )^2+𝐌^{}𝐌\|\overline{G}(q)\|^2`$. The operators $`\mathrm{\Lambda }^\pm (𝐪)=\frac{1}{2}(1\pm 𝜶\widehat{𝐪})`$ are the positive and negative energy projectors. In the CFL phase, $`𝐌=ϵ_f^aϵ_c^a\gamma _5=𝐌^{}`$ is the locking matrix in the color-flavor ($`c`$-$`f`$) sector with $`(ϵ^a)^{bc}=ϵ^{abc}`$ the totally antisymmetric tensor. To leading logarithm accuracy, the gap function $`G(q)`$ is solely a real-valued function of $`q_{\|\|}`$ given by $$G(x)=G_0\mathrm{sin}\left(\frac{\pi x}{2x_0}\right)=G_0\mathrm{sin}\left(h_{}x/\sqrt{3}\right).$$ (1) Here the logarithmic scales are defined as $`x=\mathrm{ln}(\mathrm{\Lambda }_{}/q_{\|\|})`$ and $`x_0=\mathrm{ln}(\mathrm{\Lambda }_{}/G_0)`$, with $`\mathrm{\Lambda }_{}=(4\mathrm{\Lambda }_{}^6/\pi m_E^5)`$ in terms of the transversal cutoff $`\mathrm{\Lambda }_{}=2\mu `$ and the electric screening mass $`m_E=\sqrt{\frac{N_F}{2\pi ^2}}g\mu `$, where $`g`$ is the strong coupling constant and $`N_F`$ the number of flavors. Also, we have defined $`h_{}x_0/\sqrt{3}=\pi /2`$, where $`h_{}=g/(\sqrt{6}\pi )`$, and $`G_0`$ as $`G_0\left({\displaystyle \frac{4\mathrm{\Lambda }_{}^6}{\pi m_E^5}}\right)e^{\frac{\sqrt{3}\pi }{2h_{}}}.`$ (2) The wave-functions of the pseudoscalar excitations in the CFL phase have been discussed in . In particular, in leading logarithm approximation, the pseudoscalar vertex operator of ‘isospin’ $`A`$ for a pair of particles or holes with momenta $`P/2\pm p`$ is given by $`𝚪_\pi ^A(p,P)={\displaystyle \frac{1}{F_T}}\left(\begin{array}{cc}0& i\gamma _5\mathrm{\Gamma }_{PS}^{}(p,P)\left(𝐌^A\right)^{}\\ i\gamma _5\mathrm{\Gamma }_{PS}(p,P)𝐌^A& 0\end{array}\right)`$ (5) with $`𝐌^A=𝐌^{a\alpha }(\tau ^A)^{a\alpha }`$ and $`𝐌^{a\alpha }=ϵ_f^aϵ_c^\alpha \gamma _5`$. Note that $`\left(𝐌^A\right)^{}=ϵ_f^aϵ_c^\alpha \gamma _5(\tau _{}^{A}{}_{}{}^{})^{a\alpha }=ϵ_f^aϵ_c^\alpha \gamma _5(\tau ^A)^{\alpha a}`$. In the chiral limit, the pseudoscalar vertex reduces to the gap, i.e. $`\mathrm{\Gamma }(p,0)=G(p)`$. Throughout, $`F_T=\mu /\pi `$ refers to the temporal pion decay constant . The CFL phase supports other collective excitations as well . They are not needed for the rest of our discussion. 3. In the CFL phase the ordinary photon is screened, and the gluons are either screened or higgsed. However, it was pointed out in that the CFL phase is transparent to a modified or tilde photon, $$\stackrel{~}{A}_\mu =A_\mu \mathrm{cos}\theta +H_\mu \mathrm{sin}\theta ,$$ (6) where $`\mathrm{cos}\theta =g/\sqrt{e^2+g^2}`$, $`A_\mu `$ is the photon field coupling to the charge matrix $`e𝐐_{em}`$ of the quarks and $`H_\mu `$ is the gluon field for $`U(1)_Y`$ where $`g𝐘`$ is the color-hypercharge matrix <sup>#1</sup><sup>#1</sup>#1Without loss of generality, the representation of $`𝐘`$ can be chosen to be identical to the one of $`𝐐_{em}`$. and $`\mathrm{sin}\theta =e/\sqrt{e^2+g^2}`$. $`\stackrel{~}{A}_\mu `$ carries color-flavor and tags to the charges of the Goldstone modes. The quark coupling to the tilde photon is in units of $`\stackrel{~}{e}=e\mathrm{cos}\theta `$. As a result, the CFL phase is characterized by generalized flavor-color anomalies . In this section, we show how the triangle ‘anomaly’ emerges from a direct calculation of the $`\stackrel{~}{\pi }^0\stackrel{~}{\gamma }\stackrel{~}{\gamma }`$ in the leading logarithm approximation. The contribution of (6) to the triangle graph is $`𝒯_{\mu \nu }^A(K_1,K_2)`$ $`=`$ $`{\displaystyle \frac{d^4q}{(2\pi )^4}\mathrm{Tr}\left(i𝚪_\pi ^A(q,P)i𝐒(q\frac{P}{2})i\stackrel{~}{𝐐}_\mu i𝐒(q+\frac{Q}{2})i\stackrel{~}{𝐐}_\nu i𝐒(q+\frac{P}{2})\right)}`$ (7) $`+\left[(K_1K_2)\mathrm{and}(\mu \nu )\right]`$ with $`P=K_1+K_2`$ and $`Q=K_2K_1`$. The spin-color-flavor vertex due to (6) reads $`\stackrel{~}{𝐐}_\mu =\stackrel{~}{e}𝝆_3\gamma _\mu \mathrm{diag}(\stackrel{~}{𝐐},\stackrel{~}{𝐐}^T)=\stackrel{~}{e}\gamma _\mu 𝝆_0\stackrel{~}{𝐐}=\stackrel{~}{e}\gamma _\mu 𝝆_0\left(𝐐_{em}\mathrm{𝟏}_c\mathrm{𝟏}_f𝐘\right).`$ (8) Here $`𝝆_3`$ and $`𝝆_0`$ are Pauli matrices acting on the Nambu-Gorkov indices. We assume that the indices $`\mu `$ and $`\nu `$ are eventually contracted with the polarizations $`_T^\mu `$ and $`_T^\nu `$ (space-like). The result for (7) after inserting (S0.Ex4) and (5) and lengthy algebra is $`𝒯_{\mu \nu }^A(K_1,K_2)`$ $`=`$ $`2i{\displaystyle \frac{\stackrel{~}{e}^2}{F_T}}{\displaystyle \frac{d^4q}{(2\pi )^4}\mathrm{\Gamma }_{PS}(p,P)\left[\mathrm{\Sigma }_{\mu \nu }^{A(0)}(P_i,𝐌)+\mathrm{\Sigma }_{\mu \nu }^{A(1)}(P_i,𝐌)+𝒪(1/\mu ^3)\right]}`$ $`+(P_{i}^{}{}_{+}{}^{}P_{i}^{}{}_{}{}^{}\text{or}P_{i}^{}{}_{}{}^{}2\mu P_{i}^{}{}_{+}{}^{}+2\mu \text{and}\mathrm{\Lambda }^+\mathrm{\Lambda }^{},𝐌^{(A)}𝐌_{}^{(A)}{}_{}{}^{})`$ with $`P_1qP/2`$, $`P_2q+Q/2`$, $`P_3q+P/2`$, $`P_{i}^{}{}_{\pm }{}^{}P_{i}^{}{}_{0}{}^{}\pm P_{i}^{}{}_{\|\|}{}^{}`$. The explicit form of $`\mathrm{\Sigma }_{\mu \nu }^{A(0)}`$ is $`\mathrm{\Sigma }_{\mu \nu }^{A(0)}`$ $`=`$ $`\{\mathrm{Tr}_{cf}[𝐌^A\stackrel{~}{𝐐}\stackrel{~}{𝐐}𝐌^{}]P_{1}^{}{}_{+}{}^{}P_{2}^{}{}_{+}{}^{}G(P_3)`$ (10) $`\mathrm{Tr}_{cf}[𝐌^A\stackrel{~}{𝐐}𝐌^{}\stackrel{~}{𝐐}]P_{1}^{}{}_{+}{}^{}G(P_2)P_{3}^{}{}_{}{}^{}`$ $`+\mathrm{Tr}_{cf}[𝐌^A𝐌^{}\stackrel{~}{𝐐}\stackrel{~}{𝐐}]G(P_1)P_{2}^{}{}_{}{}^{}P_{3}^{}{}_{}{}^{}`$ $`\mathrm{Tr}_{cf}[𝐌^A𝐌^{}\stackrel{~}{𝐐}𝐌\stackrel{~}{𝐐}𝐌^{}]G(P_1)G(P_2)G(P_3)\}`$ $`\times {\displaystyle \frac{1}{\mathrm{\Delta }(P_1)\mathrm{\Delta }(P_2)\mathrm{\Delta }(P_3)}}\mathrm{Tr}[\gamma _5\mathrm{\Lambda }^+(𝐏_1)\gamma _\mu \mathrm{\Lambda }^{}(𝐏_2)\gamma _\nu \mathrm{\Lambda }^+(𝐏_3)].`$ The explicit form of $`\mathrm{\Sigma }_{\mu \nu }^{A(1)}`$ is $`\mathrm{\Sigma }_{\mu \nu }^{A(1)}`$ $`=`$ $`\{{\displaystyle \frac{1}{\mathrm{\Delta }(P_1)\overline{\mathrm{\Delta }}(P_2)\mathrm{\Delta }(P_3)}}\mathrm{Tr}_{cf}[𝐌^A\stackrel{~}{𝐐}\stackrel{~}{𝐐}𝐌^{}]P_{1}^{}{}_{+}{}^{}(P_{2}^{}{}_{}{}^{}2\mu )G(P_3)`$ (11) $`{\displaystyle \frac{1}{\overline{\mathrm{\Delta }}(P_1)\mathrm{\Delta }(P_2)\overline{\mathrm{\Delta }}(P_3)}}\mathrm{Tr}_{cf}[𝐌^A\stackrel{~}{𝐐}𝐌^{}\stackrel{~}{𝐐}](P_{1}^{}{}_{}{}^{}2\mu )G(P_2)(P_{3}^{}{}_{+}{}^{}+2\mu )`$ $`{\displaystyle \frac{1}{\mathrm{\Delta }(P_1)\overline{\mathrm{\Delta }}(P_2)\mathrm{\Delta }(P_3)}}\mathrm{Tr}_{cf}[𝐌^A\stackrel{~}{𝐐}𝐌^{}\stackrel{~}{𝐐}]P_{1}^{}{}_{+}{}^{}\overline{G}(P_2)P_{3}^{}{}_{}{}^{}`$ $`+{\displaystyle \frac{1}{\mathrm{\Delta }(P_1)\overline{\mathrm{\Delta }}(P_2)\mathrm{\Delta }(P_3)}}\mathrm{Tr}_{cf}[𝐌^A𝐌^{}\stackrel{~}{𝐐}\stackrel{~}{𝐐}]G(P_1)(P_{2}^{}{}_{+}{}^{}+2\mu )P_{3}^{}{}_{}{}^{}`$ $`{\displaystyle \frac{1}{\mathrm{\Delta }(P_1)\overline{\mathrm{\Delta }}(P_2)\mathrm{\Delta }(P_3)}}\mathrm{Tr}_{cf}[𝐌^A𝐌^{}\stackrel{~}{𝐐}𝐌\stackrel{~}{𝐐}𝐌^{}]G(P_1)\overline{G}(P_2)G(P_3)\}`$ $`\times \mathrm{Tr}[\gamma _5\mathrm{\Lambda }^+(𝐏_1)\gamma _\mu \mathrm{\Lambda }^+(𝐏_2)\gamma _\nu \mathrm{\Lambda }^+(𝐏_3)].`$ We have defined $`\mathrm{\Delta }(P)=P^2ϵ_P^2`$, $`\overline{\mathrm{\Delta }}(P)=P^2\overline{ϵ}_P^2`$, and made use of the fact that in the generalized pion CM frame $`𝐏_1=𝐏_3`$, thereby reducing the spin contributions <sup>#2</sup><sup>#2</sup>#2Because of $`\gamma _5\mathrm{\Lambda }^\pm (𝐏_1)=\gamma _5\mathrm{\Lambda }^\pm (𝐏_3)=\mathrm{\Lambda }^\pm (𝐏_3)\gamma _5`$, the left-out spin contributions project to zero after a cyclic permutation under the Dirac trace.. The various contributions in $`\mathrm{\Sigma }^{(0)}`$ arise from particles and holes. Naively, $`\mathrm{\Sigma }^{(0)}`$ generates a term of order $`\mu ^2`$ as it receives contribution from the full Fermi surface. We show below that this contribution is identically zero after integration. The various contributions to $`\mathrm{\Sigma }^{(1)}`$ involve at least an antiparticle. Since $`\overline{\mathrm{\Delta }}4\mu ^2`$, the overall contribution is of order $`\mu ^0`$ after integrating over the Fermi surface <sup>#3</sup><sup>#3</sup>#3The terms linear and quadratic in $`\mu `$ cancel when (11) is inserted into (LABEL:MAIN).. It can be shown that the non-vanishing contributions arising from the various terms in $`\mathrm{\Sigma }^{(1)}`$ are those displayed in Fig. 1, where the cross refers to the pair-condensate insertion. They correspond to the first and fourth line of (11) (and first and third line of (10)) when inserted into (LABEL:MAIN), and they saturate the photo-decay of the pion in dense QCD. Finally, we note the different structure in the Lorentz projectors in (10) and (11). To the order considered, we use $`𝐌^{}𝐌=1`$ in $`\mathrm{\Delta }`$ and $`\overline{\mathrm{\Delta }}`$ <sup>#4</sup><sup>#4</sup>#4Under the assumption that $`U_A(1)`$ is broken by the anomaly.. The overall factor of 2 in front of (LABEL:MAIN) comes from the triangle graph with crossed tilde-photon legs. The various color-flavor traces yield equal contributions modulo overall signs. Specifically, $`\mathrm{Tr}_{cf}\left(𝐌^𝐀\stackrel{~}{𝐐}\stackrel{~}{𝐐}𝐌^{}\right)=2N_f\mathrm{Tr}\left(\tau ^A\stackrel{~}{𝐐}_{}^2\right)`$ $`\mathrm{Tr}_{cf}\left(𝐌^𝐀\stackrel{~}{𝐐}𝐌^{}\stackrel{~}{𝐐}\right)=+2N_f\mathrm{Tr}\left(\tau ^A\stackrel{~}{𝐐}_{}^2\right)`$ $`\mathrm{Tr}_{cf}\left(𝐌^𝐀𝐌^{}\stackrel{~}{𝐐}\stackrel{~}{𝐐}\right)=2N_f\mathrm{Tr}\left(\tau ^A\stackrel{~}{𝐐}_{}^2\right)`$ $`\mathrm{Tr}_{cf}\left(𝐌^𝐀𝐌^{}\stackrel{~}{𝐐}𝐌\stackrel{~}{𝐐}𝐌^{}\right)=+2N_f\mathrm{Tr}\left(\tau ^A\stackrel{~}{𝐐}_{}^2\right),`$ (12) where we have defined $`\stackrel{~}{𝐐}_{}={\displaystyle \frac{1}{3}}\left(\begin{array}{ccc}2& 0& 0\\ 0& 1& 0\\ 0& 0& 1\end{array}\right).`$ (16) The overall factor of 2 in (12) stems from an equal contribution due to either flavor or color, since in the CFL phase $`N_f=N_c`$. Inserting (12) into (10) and performing the energy integration over the Fermi surface yield $`𝒯_{\mu \nu }^{A(0)}(K_1,K_2)=_0_0{\displaystyle \frac{\stackrel{~}{e}^2}{2\pi ^2F_T}}\mathrm{Tr}\left(\tau ^A\stackrel{~}{𝐐}_{}^2\right)ϵ_{\mu \nu \alpha \beta }K_1^\alpha K_2^\beta ,`$ (17) where $`_0`$ is an overall constant and $`_0`$ is an integral of the form $`{\displaystyle d^4q\frac{G^2(q)}{\mathrm{\Delta }^3(q)}}{\displaystyle \frac{\mu ^2}{G_0^2}}.`$ (18) The $`\mu ^2`$ on the r.h.s. of (18) arises from integrating over the full Fermi surface. This contribution implies a decay rate that increases substantially in dense matter. Careful algebra, however, shows that after integration over the Fermi surface, the sum of the various contributions to $`\mathrm{\Sigma }^{(0)}`$ yields $`_0=0`$. The decay rate is driven by $`\mathrm{\Sigma }^{(1)}`$ to leading order, as we now show. To facilitate the calculation of the various contributions to $`\mathrm{\Sigma }^{(1)}`$, we choose the CM kinematics with $`P_1=(q^0M/2,\stackrel{}{q})`$, $`P_2=(q^0+Q^0/2,\stackrel{}{q}+\stackrel{}{Q}/2)`$, $`P_3=(q^0+M/2,\stackrel{}{q})`$. After some algebra we obtain $`𝒯_{\mu \nu }^{A(1)}(K_1,K_2)=2{\displaystyle \frac{\stackrel{~}{e}^2}{F_T}}{\displaystyle \frac{N_f}{3}}{\displaystyle \frac{_1}{\mu ^2}}\mathrm{Tr}\left(\tau ^A\stackrel{~}{𝐐}_{}^2\right)ϵ_{\mu \nu \alpha \beta }K_1^\alpha K_2^\beta .`$ (19) Note that the $`1/3`$ results from an angular integration over the 3-momenta through the Fermi surface, while $`_1={\displaystyle \frac{d^4q}{(2\pi )^4}\frac{G^2(q)}{\mathrm{\Delta }^2(q)}}={\displaystyle \frac{i}{8}}F_T^2={\displaystyle \frac{i\mu ^2}{8\pi ^2}}`$ (20) is proportional to the square of the temporal pion decay constant . The presence of the antiparticle yields a $`\mu `$-independent contribution <sup>#5</sup><sup>#5</sup>#5The same observation applies to most of the mixing processes discussed in and found to vanish to leading order.. Note that the various contributions involving the antiparticle gap $`\overline{G}`$ in $`\mathrm{\Sigma }^{(1)}`$ contribute zero after integration. The same is true for the terms with any pair-condensate insertion on the fermion leg in between the two photon vertices. Hence, $`𝒯_{\mu \nu }^A(K_1,K_2)=i{\displaystyle \frac{\stackrel{~}{e}^2}{F_T}}{\displaystyle \frac{1}{4\pi ^2}}\mathrm{Tr}\left(\tau ^A\stackrel{~}{𝐐}_{}^2\right)ϵ_{\mu \nu \alpha \beta }K_1^\alpha K_2^\beta `$ (21) to leading order in the density. We have checked that the possible vertex corrections to (5) are subleading in the decay rate. Because of the $`F_T`$-dependence, eq. (21) suggests that the radiative decay of the generalized pion vanishes as $`1/\mu `$ in dense matter. This result, based on our explicit calculation in the CFL phase to leading order, is in agreement with the result suggested by the generalized WZW term , $$^\mu 𝐀_\mu ^3=\frac{\stackrel{~}{e}^2}{96\pi ^2}ϵ_{\mu \nu \rho \sigma }\stackrel{~}{F}^{\mu \nu }\stackrel{~}{F}^{\rho \sigma }$$ (22) with $`\stackrel{~}{F}^{\mu \nu }`$ the field strength associated to (6<sup>#6</sup><sup>#6</sup>#6The normalization of the $`\pi ^0\gamma \gamma `$ amplitude is 8 times bigger than the normalization of the corresponding effective Lagrangian or of the pertinent anomaly (22). Furthermore, $`\mathrm{Tr}\left(\tau ^3\stackrel{~}{𝐐}_{}^2\right)=1/3`$.. The finite-size of the generalized pion does not upset the geometrical normalization of the anomaly in the generalized WZW form . A direct comparison to the $`\mu =0`$ radiative decay of the usual pion in QCD, shows that (22) is off by a factor of $`1/3`$. As originally noted in , the color-flavor anomaly in the CFL phase is no longer multiple of $`N_c`$ due to the color-flavor locking in the triangle graph, hence a factor of $`1/3`$. This point is particularly explicit in the various traces in (12). 4. In the CFL phase ordinary photons are screened, and perturbative gluons are either screened or higgsed. The anomalous decay of the generalized pions occurs via the tilde-photon. We have provided a direct calculation of the radiative decay $`\stackrel{~}{\pi }^0\stackrel{~}{\gamma }\stackrel{~}{\gamma }`$ in this phase using weak-coupling arguments. To leading logarithm accuracy our result is exact and in agreement with the normalization suggested by the generalized WZW term given in . Much like in the vacuum, the radiative decay of the generalized pions is dictated by geometry in leading order, and vanishes at asymptotic densities. Clearly our analysis extends to other anomalous as well as nonanomalous processes, and should prove useful for the analysis of emission rates in dense QCD. Acknowledgments: M.A. Nowak, M. Rho and I. Zahed thank KIAS for hospitality during the completion of this work. This work was supported in part by the US DOE grant DE-FG02-88ER40388 and by the Polish Government Project (KBN) 2P03B 00814.

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